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# Mid-infrared interferometry on spectral lines: III. Ammonia and Silane around IRC +10216 and VY CMa
## 1 Introduction
Both high and low-mass stars near the end of their lives are known to emit copious amounts of material with the high mass-loss rates critically dependent on the formation of silicate or carbonaceous dust grains within a few stellar radii of the red (super-) giant photosphere. Stellar photons both impart momentum to the dust particles through absorption and scattering (which drive a wind) and heat dust close to the star to temperatures up to 1200-1500 K (e.g., Lafon & Berruyer (1991)). Thermal radiation from these hot dust grains peaks in the near-infrared (1-5 $`\mu `$m), while more distant and cooler (500 K) grains emit mostly in the mid-infrared (5-20 $`\mu `$m).
Although the physical sizes of these evolved stars are impressive (R$`{}_{}{}^{}\stackrel{>}{_{}}1`$ AU), their galactic paucity means that even the closest examples lie at distances greater than 50 pc. Thus the characteristic size scale for dusty circumstellar emission at a few stellar radii is generally a small fraction of an arcsecond, too small to be resolved using standard observing techniques from ground-based telescopes limited by telescope diffraction (in the mid-infrared) and/or atmospheric turbulence (for shorter wavelengths). However, long-baseline interferometry in the infrared can directly detect and measure the morphologies of these envelopes.
The high densities of heavy elements and the mild temperatures around evolved stars encourage the formation of myriad diatomic and polyatomic molecules in addition to dust grains. For instance, over 50 molecular species have been found around prototypical carbon star IRC +10216 (see review by Glassgold ). Attempts to understand the density, temperature, and velocity distributions as well as the formation mechanisms of these molecules have contributed to the development of the field of astrochemistry.
Using estimates of the temperature, atomic abundances, and gas density in and above the photosphere, predictions of molecular abundances can be made through detailed calculations of a network of chemical reactions. By assuming local thermodynamic equilibrium (LTE), various molecules are said to “freeze-out” in the stellar wind. The decreasing temperatures and densities cause reaction rates to fall quickly, locking the atoms in certain energetically favorable molecules. Frozen equilibria models have proven useful at explaining the abundances of many molecules observed around AGB stars (e.g., Lafont, Lucas, & Omont (1982)), however additional chemical processes are needed to explain some detected species. For instance, penetration of interstellar UV-radiation catalyzes photochemical reactions in the outer envelopes of AGB stars (e.g., Cherchneff & Glassgold (1993)), while other molecules may be produced under non-equilibrium conditions associated with shocks (Willacy & Cherchneff (1998)). Particularly relevant here is the fact that negligible amounts of silane and ammonia are predicted by most freeze-out models, yet they have been detected with some abundance in circumstellar envelopes (e.g., Keady & Ridgway (1993)); their formation is hypothesized to be catalyzed on the surfaces of dust grains in the flow. In addition, recent discoveries of brown dwarfs and extrasolar planets have catalyzed interest in the chemistry of cool stellar atmospheres (e.g., Burrows & Sharp 1999), leading to new insights into molecular formation mechanisms in these environments. Since it is often difficult or impossible to experimentally reproduce the physical conditions of circumstellar and interstellar space, molecular observations are critical to test and guide relevant theories.
This paper is the third in a series on the spatial distributions of dust and molecules in the inner envelopes of nearby red giants and supergiants. Paper I (Monnier et al. 2000a ) discussed the hardware implementation of this experiment, while Paper II (Monnier et al. 2000b ) discussed recent visibility data at 11.15 $`\mu `$m and presented the appropriate dust shell models for IRC +10216 and VY CMa. This paper (Paper III) makes use of these results as well as the first mid-infrared visibility data ever reported on spectral lines, providing critical new information on the molecular stratification around these stars.
The methods used in this work are technologically challenging and represent significant advances in infrared interferometry. By combining the high spatial resolution of long-baseline interferometry with the high spectral resolution of heterodyne spectroscopy (Paper I), this work has, for the first time, probed the brightness distribution on and off of spectral absorption features as narrow as $``$1 km s<sup>-1</sup>($`\frac{\lambda }{\mathrm{\Delta }\lambda }10^5`$). This allows the absorbing regions of polyatomic molecules to be directly measured, setting strong limits on the formation radii. Such measurements are important for determining the formation mechanism since there is presently no good theory to explain the observed high abundance of certain molecular species (e.g., NH<sub>3</sub> and SiH<sub>4</sub>).
## 2 Modeling
The observing methodology for extracting visibility data on and off of spectral lines with the Infrared Spatial Interferometer has been discussed in detail in Paper I and Monnier (1999). New dust shell models of IRC +10216 and VY CMa have been developed in Paper II. However, interpreting spectral line observations requires modeling of the distribution of circumstellar molecules as well as dust.
As for observations of continuum radiation with the ISI, it is rarely possible to directly reconstruct an image of the astrophysical source due to the sparse sampling of the Fourier plane afforded by a two-element interferometer. This limitation necessitates the use of radiative transfer models for interpreting the line data. This section discusses the assumptions used in creating models of the molecular envelopes around AGB stars, the numerical code, and the scientific goals of subsequent analyses.
### 2.1 Simple models
In Paper II, the assumptions of a spherical symmetric and uniform outflow of dust embedded in the stellar winds were adopted, approximations which have been extended to include the distribution of molecular material. As with the models of the dust shells, the molecular envelope was characterized by the inner formation radius and overall abundance factor, free parameters in fitting models to the molecular line data.
Observations have shown that the line shapes of SiH<sub>4</sub> and NH<sub>3</sub> around AGB stars typically show emission in the red wing and absorption in the blue (e.g., Goldhaber (1988)). This kind of profile is expected for molecular emission and absorption in an expanding envelope (e.g., Shu (1991)). Figure 1 shows a theoretical line profile similar to those actually observed. The bulk of the absorption is blue-shifted with respect to the star at the speed of the outflow. This is because the entire column of molecular gas along the line-of-sight connecting the observer and the star is at the same relative velocity (-v<sub>outflow</sub>), assuming uniform outflow; hence, there is a large optical depth, $`\tau _{\mathrm{line}}`$, at this frequency and self-absorption is observed. On the other hand, lines-of-sight which penetrate the molecular gas envelope at impact parameters off the observer-star axis encounter parcels of gas at a variety of relative doppler shifts. The optical depth at any particular frequency is hence generally low, allowing emission from deeper layers of hot dust to escape and to be observed without significant molecular absorption. The low-level emission observed at both blue- and red-shifted velocities arises from relaxation of IR-pumped and collisionally-excited transitions into higher energy vibrational states. Note that absorption in the line core occurs primarily because of the large optical depth at that frequency combined with the fact that the temperature decreases with increasing distance from the star. Hence, absorption under these conditions would occur even when the molecules are in local thermodynamic equilibrium (LTE) with the dust and radiation.
For the AGB stars under study here, the source of most of the mid-infrared flux is thermal emission from the dust envelope (Danchi et al. (1994)), not from the star itself. Since the blue-shifted absorption features indicate that most of the absorption is occurring along the direction towards the center of dust envelope, an interesting measurement that can be made with an interferometer is the size of the absorption region. This measurement yields direct information on the approximate location of the $`\tau _{\mathrm{line}}1`$ surface, which can be used to determine the abundance and inner radius of molecular formation. However this method is insensitive to molecules which may exist right at the inner radius of the dust, if they are in LTE. This is because there is no significant background source of radiation to absorb. These complications require the use of a sophisticated spectral line code to quantify the sensitivity of this method to molecules forming near the inner radius of the dust shell. Such a code has been used, and is described in the next section.
Previous spectral line work (Betz, McLaren & Spears (1979); McLaren & Betz (1980); Goldhaber & Betz (1984); Goldhaber (1988); Keady & Ridgway (1993)) has shown that ammonia and silane exist around AGB stars at up to 10<sup>4</sup> times the abundances predicted by equilibrium calculations of expanding outflows (for more recent work under somewhat different physical conditions, see Burrows & Sharp 1999). It has been suggested that dust catalyzes the formation of these molecules through chemical reactions on grain surfaces. If this is the case, one might expect the inner radius of molecular formation to coincide with the dust formation radius. However, atoms and molecules generally have low sticking efficiencies at high temperatures, elastically scattering off the grain surfaces (e.g., Leitch-Devlin & Williams (1985)). This fact and the lack of known chemical pathways for the formation of molecules such as ammonia and silane make direct measurements of molecular formation radii important, and this is one of the prime goals of the filterbank project of the Infrared Spatial Interferometer.
### 2.2 Radiative transfer calculations
Radiative transfer in spectral lines often is accomplished by applying the Sobolev approximation (e.g., Goldhaber (1988); Shu (1991)), which assumes that large-scale velocity gradients essentially radiatively decouple parcels of molecular gas. To accurately calculate the radiative transfer in a spectral line, the Sobolev approximation requires the natural linewidth, in this case determined by microturbulence (for cold gas) or thermal motions (for hot gas), to be much smaller than the expansion velocity. While this condition is reasonably satisfied for the narrowest lines ($`v_{\mathrm{outflow}}15\times \mathrm{\Delta }v_{\mathrm{linewidth}}`$), it is only weakly so for the broadest ones ($`v_{\mathrm{outflow}}4\times \mathrm{\Delta }v_{\mathrm{linewidth}}`$).
In order to avoid this uncertainty, all radiative transfer calculations in spectral lines were performed using a code developed by Dr. J. J. Keady (Keady (1982)), based on the method of Mihalaset al.(1975). It accurately calculates line profiles as well as frequency-dependent emission profiles, from which visibility curves can be computed on and off spectral features. This code was developed to treat the case of molecules embedded in an expanding flow, even when the absorbing molecules are comingled with continuum-emitting sources. This is likely the case for AGB stars where thermal emission by dust is believed to be “filling-in” the absorption features. The program, written in Fortran, formally solves the observer’s frame transport equation, which is necessary to correctly account for the propagation of line radiation through an expanding flow, thus avoiding the questionable Sobolev approximation. Dr. Keady has allowed his code to be used for the analyses which follow, and has assisted in porting the Fortran code to run under Solaris. This code has been extensively tested, having been used for over 15 years (Keady (1982); Keady, Hall & Ridgway (1988); Keady & Ridgway (1993); Winters et al. (1998)). While the code was already able to be used with spherical top molecules (e.g., SiH<sub>4</sub>), the appropriate partition function and statistical weights had to be programmed for NH<sub>3</sub> (a symmetric rotor). A nice summary of the mid-infrared molecular properties of both SiH<sub>4</sub> and NH<sub>3</sub> can be found in chapter 5 of Goldhaber’s thesis (1988). Readers interested in the numerical details of the radiative transfer calculation should consult Keady (1982) and Mihalaset al.(1975) for further details.
### 2.3 Assumptions
#### 2.3.1 Dust
The source function used for radiative transfer in Keady’s code assumes a single dust grain size and temperature at a given radius from the star. The code (Wolfire & Cassinelli (1986)) which was used for fitting the ISI continuum visibility data in Paper II utilized the full MRN (Mathis, Rumpl & Nordsieck (1977)) distribution of grain sizes and hence some modification must be made to adopt our previous models for use by Keady’s algorithm. The average cross-section of the MRN distribution was used as input to the Keady code, along with the size-averaged dust temperature. This necessary simplification resulted in a slight misfit to the ISI continuum visibility data, which was compensated for by empirically adjusting the overall value of the single grain-size dust opacity ($`\stackrel{<}{_{}}`$30% change).
#### 2.3.2 Temperature(s) of the gas
In addition to spherical symmetry and uniform outflow, other assumptions are also generally made for line calculations. Most important is the assumption that the occupation of various ro-vibrational states of the polyatomic molecules can be approximated by a Boltzmann distribution using separate vibrational and rotational temperatures, T<sub>vib</sub> and T<sub>rot</sub>. Ideally, the populations of the various states should be calculated using a multi-level model molecule. However, the collisional excitation constants are poorly known for SiH<sub>4</sub> and NH<sub>3</sub> and the large number of (far-IR) rotational transitions make for a difficult and uncertain result.
Goldhaber (pg.103, Goldhaber ) considered this question in some detail for IRC +10216 and his conclusions will now be summarized. For sufficiently high molecular densities, both mid-IR and far-IR (i.e., ro-vibrational and pure rotational) transitions will be optically thick, trapping line radiation. In this case, collisional excitation and de-excitation will have a strong influence on the level populations, equilibrating the vibrational and rotational temperatures with the gas (kinetic) temperature. However, the molecules under consideration here are not found at sufficiently high densities for this to apply.
At the somewhat lower densities encountered for ammonia and silane, the mid-IR (ro-vibrational) transitions become optically thin. In terms of excitation mechanism, this allows IR pumping of the vibration-rotation transitions to dominate (over collisions) within about 70 R, while rotational relaxation in the pure rotational transitions (far-IR) dominate in the outer envelope (where these transitions become optically thin). Collisional excitation and de-excitation never play a dominant role. Hence, T<sub>vib</sub> and T<sub>rot</sub> should be close to the radiation/dust temperature, T<sub>dust</sub>, suggesting that T<sub>vib</sub>=T<sub>rot</sub>=T<sub>dust</sub> is a good starting point. The radial dependence of the rotational and vibrational temperatures usually need slight adjustment during the modeling process to reproduce the relative depths of high and low excitation lines.
One interesting exception applies to the J=K states of NH<sub>3</sub>. The rotational dipole moment along the symmetry axis (z-axis) of NH<sub>3</sub> is equal to zero; hence the $`\mathrm{\Delta }K=0`$ selection rule must apply for all radiative (dipole) transitions (see Townes & Schawlow ). Therefore, a molecule rotating in the J=K state can not radiatively de-excite to a lower energy (lower J) state. Hence, the equilibrium population of J=K states is not determined by the radiation temperature but by collisions with the ambient gas (mostly H<sub>2</sub>); the rotational temperature of these states should be in equilibrium with the gas (kinetic) temperature. However, a strong (parallel-type) vibrational band at 6.1 $`\mu `$m does allow $`\mathrm{\Delta }K=\pm 1`$ transitions, thereby allowing K-ladders to come into equilibrium if sufficiently excited by the radiation field. Sufficient IR pumping in this band is likely to dominate over collisional processes out to a few hundred R, although a detailed calculation is lacking (see Goldhaber ).
#### 2.3.3 Molecular constants
The polyatomic molecules of SiH<sub>4</sub> and NH<sub>3</sub> have been modelled using physical constants, including the moments of inertia and vibrational energy levels, adopted from Keady & Ridgway (1993) and Goldhaber (1988). Readers interested in learning more about molecular spectroscopy should consult these sources; a rigorous and complete treatment of the subject can be found in Microwave Spectroscopy by Townes & Schawlow (1975).
### 2.4 Analysis method
For each star and target molecule, the analysis strategy for ISI spectral line data was straightforward and had three major steps.
1. The dust shell was modelled in order to have a reasonable approximation of the continuum source. This has been done in Paper II, where coeval continuum visibility data were used to create appropriate dust shell models of the astrophysical sources under study.
2. Molecules were assumed to form at a given radius from the star and flow out uniformly. Theoretical spectral line profiles were then calculated to fit both ISI measurements and also previous observations of spectral line depths and ratios. Lines of various excitation energies were used (when available) in order to probe the abundance, temperature, and turbulent structure of the gas. This level of analysis relied heavily on the careful modeling pre-existing in the astrophysics literature.
3. Theoretical visibilities on and off the spectral line were calculated and compared to the observed visibility ratios. The radius of molecular formation set in step 2 was varied until all available data were satisfactorily fit.
This method allows one to distinguish between gas models with high molecular abundance far from the star and those with lower molecular abundance close to the star in two different ways. While the line depth of a given transition can always be fit by varying the abundance and mean gas temperature, the relative line depths of transitions with different excitation energies can not. This probe of the average excitation temperature, when coupled with estimates of the radial temperature profile and assumption of uniform outflow, can be used to infer the location of molecular formation (e.g., Keady & Ridgway (1993)). Here in this work, we also apply a second, more direct, approach by using an interferometer. With suitable spectral and spatial resolution, one can directly detect the absorption pattern when the molecules are sufficient close to the star. This is the main thrust of this paper, where we show that molecules within 40 Rwould significantly modify the emission pattern in the line core, providing a method independent of the line ratio argument for detecting the presence of molecules close to the star.
### 2.5 Comparing model quantities to observations
After fully specifying the dust and gas characteristics, the Keady code returns a spectral line profile and monochromatic radial emission profiles. This section discusses some of the details regarding comparison of the modelling results with ISI filterbank data.
#### 2.5.1 Spectral line profiles
Because of the heterodyne detection scheme, mid-infrared radiation arriving outside of the primary beam of the interferometer was not seen. This had a small effect on the depth of the spectral line, since only a fraction of the molecular absorption occurs in the outer envelope. The effect was compensated for by multiplying the angular distribution resulting from monochromatic radial emission calculations by the effective primary beam of FWHM 3 ″(see Paper II for more discussion) and then adjusting the flux level accordingly. All spectral profiles shown in the next sections have had this correction applied.
Next, the finite bandwidth of the filterbank observations and the double-sideband (DSB) nature of the detection were accounted for. The bandpass center frequency and width were used to determine the average depth of the calculated spectral line for a corresponding single-sideband (SSB) observation. The dilution of the line depth due to the combination of the “uninteresting” additional continuum sideband encountered in practical observation was accounted for by dividing the SSB line depth by 2. This resulting DSB (diluted) line depth was then compared to observations with the filterbank.
Figure 1 shows one of the output figures from the modeling analysis suite. The bandwidth of the observation is marked off by vertical dashed lines, and line diagnostics have been calculated. The line width and average DSB line depth were also calculated for comparison with spectral line observations. Other quantities were also determined, such as the equivalent width, which were useful when comparing with lower spectral resolution data.
#### 2.5.2 Visibility ratios
The basic observable of the filterbank experiment is the visibility ratio on the spectral line compared to the dust continuum. Monochromatic radial emission profiles were calculated at 10 specific frequencies across the spectral line, sampling the stellar continuum as well as the entire absorption region. The double-sideband, radial emission profile for the finite bandwidth selected by the filterbank was determined by a flux-weighted average of the emission profiles inside the observation bandwidth and an equal bandwidth of continuum. This resultant profile was multiplied by the effective primary beam and the visibility then calculated. The visibility curve of the continuum emission was calculated as well, and the visibility ratio (on the line compared to the continuum) as a function of baseline was then determined.
Figures 1-3 illustrate this entire process, showing example output at every stage of the analysis process for a silane line around IRC +10216. All subsequent model visibility ratios have been processed in identical fashion, but these diagnostic plots are not presented for brevity. The particular gas model shown in figures 1-3 has a silane formation radius of only 10 R, small enough to produce a large change in the visibility on and off the spectral line. Since the filterbank system measures a visibility ratio, the radial emission profiles were usually of only secondary interest, and more emphasis was placed on interpretations of ratios.
To facilitate a better understanding of the visibility data, example brightness profiles appear in figure 2. The bright spike at the origin is from stellar photospheric emission seen through the dust and gas, while the secondary peak at 150 mas reveals the location of the dust shell inner radius. It is important to note that while the star itself clearly shines through the dust envelope (it is the highest surface brightness feature of the nebula), it contributes a tiny percentage of the total mid-IR emission due to the large size of the dust shell. The inset image is a two-dimensional representation of the absorption region defined by the ratio of the absorption core to the continuum profiles. One can see that absorption for this model is indeed highest near the center of the dust shell. The basic point of the filterbank experiment was to measure the size of this absorption region.
Figure 3 shows the visibility curves corresponding to the emission profiles of figure 2. The various curves are described in the figure caption, and show the effects of applying the instrumental corrections discussed in this section. In this case, a visibility ratio of about 0.92 would be expected for the range of baselines observed (2-4 m).
The analysis process discussed above took into account all the practical details of the experiment: finite bandwidth, primary beam, double-sideband detection. The next sections present the data from the filterbank experiment and detail the modeling results.
## 3 Interferometry on Spectral Lines: Results
Spectral line observations were carried out for two target sources, the carbon star IRC +10216 and the red supergiant VY CMa. General introductions to both of these stars, including discussions of previous continuum observations, can be found in Paper II. A summary of all the filterbank data can be found in figure 4; details will be discussed in subsequent sections. Recall from Paper I that the ratio of the visibility on the line compared to off the line is made from dividing the fringe amplitude ratio by the infrared power ratio. Hence, in figure 4, the “Visibility Ratios” were derived from the “Fringe (Amplitude) Ratios” and the “IR Power Ratios,” and do not represent an independent set of measurements. See tables of line frequencies in Monnier (1999, tables C.4 & C.5) for more information regarding specific molecular transitions.
### 3.1 Comparison with previous spectroscopic results
While the interferometric data presented here are the first of its kind, the depths of the spectral lines have been observed before. Figure 5 shows a comparison of all the lines observed with the double-sideband line depths observed by Goldhaber (1988). While the agreement is good in general, there is a clear signal that the line depths measured in 1998 were slightly deeper around IRC +10216 than the mid-1980s observations of Goldhaber. There are many possible explanations for this. Changes in the dust shell geometry, molecular abundance, as well as different seeing conditions during the observations themselves can cause small changes in the line depth. However, the good overall agreement is confirmation that the circumstellar environments important for forming these spectral lines around both VY CMa and IRC +10216 have not changed radically in the last decade, justifying the use of previous modeling efforts (Goldhaber (1988); Keady & Ridgway (1993)) in developing molecular gas models (see modeling step 2 in §2.4).
## 4 IRC +10216
A journal of spectral line observations for IRC +10216 can be found in table 2, while a full and concise summary of the data appears in table 3. Separate modeling for both molecules, SiH<sub>4</sub> and NH<sub>3</sub>, is presented below.
### 4.1 Silane in IRC +10216
#### 4.1.1 Previous work
The mid-infrared transitions of silane around IRC +10216 have been observed by Goldhaber & Betz (1984), Goldhaber (1988, hereafter G88), Keady & Ridgway (1993, hereafter KR93), and Holler (1999). Goldhaber and colleagues used a heterodyne spectrometer which produced the highest spectral resolution data, fully resolving the absorption line cores (spectral resolution $``$ 0.2 km s<sup>-1</sup>). KR93 employed a Fourier Transform Spectrometer at Kitt Peak with somewhat lower spectral resolution, about 3 km s<sup>-1</sup>, sufficient to determine the spectral line strengths but insufficient to fully resolve the cores. Holler used a broadband RF spectrometer coupled to the heterodyne detection system of the ISI, resulting in sub- km s<sup>-1</sup>spectral resolution (also see Isaaket al.). Because the line depths have remained largely unchanged over the last 15 years, the modeling results of these workers was assumed to hold true today. The most detailed analysis can be found in KR93, and their gas model parameters were a starting point for the analysis which follows.
KR93 found a relative abundance of silane compared to molecular hydrogen of 2.2$`\times `$10<sup>-7</sup> with a rotational temperature law, T$`{}_{\mathrm{rot}}{}^{}=2000r^{0.525}`$ (r is expressed in units of stellar radii), which falls off slightly faster than the dust temperature. The calculated column density was 2.2$`\times `$10<sup>15</sup> cm<sup>-2</sup>. In addition, T<sub>vib</sub> was taken to be about 85% of the T<sub>rot</sub> to match the weak emission observed in the higher J lines. These results assumed spherical symmetry, a uniform outflow of 14 km s<sup>-1</sup>outside of 20 R, and a microturbulent velocity of 1 km s<sup>-1</sup>. It was found that the line ratios were much better fit by truncating the silane distribution inside 40 R, and KR93 concluded that silane must be forming in the outflow at about this radius. With these assumptions, KR93 was able to satisfactorily match the line strengths of 7 silane transitions and the spectral profile of one of Goldhaber’s high spectral resolution profiles. Holler (1999) estimated the column density of SiH<sub>4</sub> from data taken in May 1999 using one ISI telescope and found it consistent with previous measurements ($``$2.7$`\times `$10<sup>15</sup> cm<sup>-2</sup>).
#### 4.1.2 Visibility observations
Table 3 reports the visibility ratios on and off of various silane absorption features in Fall 1998. A stellar recessional velocity of V<sub>LSR</sub>=-26.0 km s<sup>-1</sup>has been used to convert V<sub>LSR</sub> to expansion velocity, based on G88. In all cases, the visibility ratio has been observed to be consistent with unity, and data from all baselines (2-4 m) have been averaged together. Qualitatively, this implies that the absorption region is large compared to the spatial resolution of the interferometer baseline, $``$0.4″.
The model for the molecular envelope developed in KR93 was combined with the new dust shell model from Paper II as a starting point for this modeling work. In particular, the same power law relation for the rotational temperature was used and the vibrational temperature was set equal to 85% of the dust temperature. A microturbulent velocity of 0.6 km s<sup>-1</sup>was used to match the line widths observed by Goldhaber (1988). Next, a series of gas models was calculated using different silane formation radii: 10 R, 40 R, and 80 R. The silane abundances were then scaled to match the average line depths of the observed transitions. For each of these 3 gas models and for all 4 observed lines, the visibility ratios of the line compared to the continuum were calculated and compared to observations.
The results of these models appear in figure 6. Indeed, gas models with largest formation radii fit the line depths of the four lines the best, as was found by KR93, implying a relatively low overall rotational temperature of the absorbing molecules. In addition, the importance of the visibility ratio observations made with the ISI and filterbank can be seen. For formation radii less than about 40 R, the predicted visibility ratios are significantly below unity. Even for a formation radius of 40 R (the choice preferred by KR93 based on spectroscopy alone), the visibilities in the spectral lines are too small to be consistent with observations (although not strongly ruled out). The column density for the most favored 80 R model was found to 1.3$`\times `$10<sup>15</sup> cm<sup>-2</sup> (2.0$`\times `$10<sup>15</sup> cm<sup>-2</sup> for the 40 R case).
Since we have measurements of 4 lines of differing rotational excitation energies, the sensitivity of the visibility ratios to the rotational temperature law can be investigated. Models with identical gas density distributions were run using T<sub>rot</sub>=T<sub>vib</sub>=T<sub>dust</sub>, and the results for these calculations can be found in figure 7. This temperature law does not fit the line depths as well as the first set of models, supporting the adoption of the empirical temperature fall-off of KR93. Alternatively, the population of low-lying J transitions could be enhanced if the mass-loss rate was higher in the past; this would have the effect of increasing the density of cold silane, increasing the molecular absorption from the ground state. The visibility ratios are similar to those found in the KR93-based models above and the earlier conclusion of a large molecular formation radius is supported and now shown to be insensitive to choice of temperature profile.
#### 4.1.3 Conclusions
In short, the filterbank observations largely confirm the conclusions of KR93; namely, that silane must be forming at a large distance from the star. The visibility data directly measures the size of the absorption region to be larger than $``$0.4″and favors a formation radius even larger than that proposed by KR93. The combination of spectroscopy and interferometry convincingly demonstrates that silane is forming at a radius of $`\stackrel{>}{_{}}`$80 R. The visibility data are only able to set a lower limit on the formation radius since no visibility difference on and off the line was detected within the experimental uncertainty.
For the stellar radius of 22 mas used in the model and a distance of 135 pc, 80 R corresponds to about 240 AU from the star. It is not clear why silane would begin to form at so remote a location in the outflow. The densities of the atomic constituents of silane are proportional to r<sup>-2</sup> for a uniformly expanding envelope and would act to shut off chemical reactions which are density-dependent. If silane is formed on the surfaces of grains, then the gas phase density of silicon (and its compounds) may not be so important, but rather the amount of gas-phase silicon and H<sub>2</sub> adsorbed onto grains. Adsorption of H<sub>2</sub> on the grain surfaces would probably be small at the higher temperatures (1000-1400 K) where dust forms, and be more pronounced at larger distances form the star (for example, see calculations of Leitch-Devlin & Williams ). In these models, the temperatures of the gas and dust are about $``$290 K at 40 Rand 200 K at 80 R.
Gas-phase SiS, which forms in abundant quantities in the photosphere, has been predicted to be depleted significantly by adsorption onto grains (Glassgold & Mamon (1992)). In fact, depletion of SiS within 100 R has already been observed in both mid-infrared (Boyle et al. (1994)) and mm-wave transitions (Bieging & Nguyen-Quang-Rieu (1989); Bieging & Tafalla (1993)). We hypothesize that the large radius of SiH<sub>4</sub> creation is related to the time-scale for this depletion process and/or the adsorption of H<sub>2</sub> onto grains. It is hoped that the introduction of these new observations will stimulate renewed theoretical interest in understanding both the chemical origins and lineage behind the high abundance of silane around IRC +10216.
### 4.2 Ammonia in IRC +10216
#### 4.2.1 Previous work
The mid-infrared rovibrational transitions of ammonia around IRC +10216 have been observed by Betzet al.(1979, hereafter B79), Goldhaber (1988, G88), Keady & Ridgway (1993, KR93), and Holler (1999). Some details regarding the spectrometers employed have already been summarized in §4.1.1. Because the line depths have remained essentially unchanged over the last 15 years, the modeling results of these workers are again assumed to be valid today. As in the case for silane, the most detailed analysis can be found in KR93, and their gas model parameters acted as a starting point for the analysis which follows.
KR93 found a relative abundance of ammonia compared to molecular hydrogen of 1.7$`\times `$10<sup>-7</sup> with a rotational temperature profile, T$`{}_{\mathrm{rot}}{}^{}=2000r^{0.60}`$ (r is expressed in units of stellar radii), a faster fall-off than the profile used for silane. Even with this steeper fall-off, the predicted depth of the lowest lying aR(0,0) line was much weaker than that observed, requiring either much colder or more dense gas in the outer envelope ($`\stackrel{>}{_{}}`$400 R). The calculated column density was 2$`\times `$10<sup>15</sup> cm<sup>-2</sup>, very similar to recent results from Holler ($``$2.2$`\times `$10<sup>15</sup> cm<sup>-2</sup>, Holler (1999)). These results assumed spherical symmetry, a uniform outflow of 14 km s<sup>-1</sup>outside of 20 R , and a microturbulent velocity of 1 km s<sup>-1</sup>. Since a fit of the larger width of the high excitation aQ(6,6) line with a constant velocity flow was not possible, a more complicated velocity model was used which had the gas accelerating to terminal velocity at radii between 5 and 20 R (see figures 3a & 11 in KR93 for the velocity and density laws adopted).
Although unable to reproduce the line shapes in precise detail, the KR93 models were successful at qualitatively explaining the bulk of the spectroscopic observations available by placing the ammonia density peak around 12 R. Importantly, these models represent ideal test cases for the filterbank experiment, because the presence of ammonia absorption inside of 20 R is expected to be clearly indicated by its visibility.
#### 4.2.2 Results
Table 3 shows the ratio of the visibilities observed in and out of the aQ(3,3) NH<sub>3</sub> absorption feature in the Fall of 1998. A stellar recessional velocity of V<sub>LSR</sub>=-26.0 km s<sup>-1</sup>has been used to convert V<sub>LSR</sub> to expansion velocity, based on G88. The filterbank bandpass selected ($``$1.9 km s<sup>-1</sup>) was somewhat smaller than the full-width half-depth (FWHD) observed by G88, about 5 km s<sup>-1</sup>, and centered on the absorption core, corresponding to an outflow velocity of 14 km s<sup>-1</sup>. In this case, as for the silane observations, the visibility ratio was observed to be consistent with unity.
The modeling process began with the dust shell model from Paper II and the gas model developed in KR93. The same power law relation for the rotational temperature and a 1.0 km s<sup>-1</sup>microturbulent velocity was initially used. An ammonia density and velocity distribution similar to KR93 were used, although simplified. Specifically, the velocity was increased linearly from 0.1 to 14 km s<sup>-1</sup>between 3 and 20 R. The ammonia density profile used started at 5 R, peaked at $``$10 R, and decreased outside of 20 R consistent with a uniform outflow ($`\rho r^2`$). This corresponded to a linearly increasing ammonia shell mass until 20 R, where uniform outflow was assumed to begin. See figure 8 for a plot of the density profile used.
As was shown for silane in the last section, significant molecular absorption inside about 20 R necessarily results in substantial differences in interferometric visibilities on and off the spectral line. Figure 9 shows the spectral profile and visibility ratio for the above model of ammonia around IRC +10216, based on KR93, and also the spectral line data taken from G88. KR93 did not model this specific transition, so the agreement can not be expected to be precise; indeed, the line prediction is narrower than the profiles observed by Goldhaber. The density of ammonia has been adjusted to match the line depth observed in Fall 1998 within the bandpass of the ISI filterbank, resulting in a column density of 8.4$`\times `$10<sup>15</sup> cm<sup>-2</sup> for this model.
The model visibility ratio at a spatial frequency of $`3\times 10^5`$ radians<sup>-1</sup> is 0.926 compared to the observed ratio of 1.000$`\pm `$0.035. This model is thus ruled out at a 2-$`\sigma `$ level. This reinforces a conclusion from the modeling of silane: significant molecular absorption occurring within about 20 R for IRC +10216 results in a strong on-line and off-line difference in the interferometric visibilities, a signal not seen with the ISI filterbank experiment. The conclusions that ammonia forms at a distance of $`\stackrel{>}{_{}}`$20 Rfrom the star indicates that, as in the case of silane, ammonia molecules do not form near where the dust condenses but rather at lower temperatures and probably through surface adsorption. The dust grain temperature at 20 Ris $``$400 K in these models.
#### 4.2.3 Discussion
There are other, purely spectroscopic, observations which also suggest that the explanation put forth by KR93 for the broad aQ(6,6) line, bulk acceleration of gas within 20 R, may not be fully correct. KR93 did not attempt to fit the high spectral resolution profiles published in G88, which revealed details concerning the ammonia around IRC +10216 not reproduced by models relying on accelerating gas to broaden the absorption lines.
KR93 fit the aQ(6,6) line data originally published in B79 which was cut-off at -17 km s<sup>-1</sup>relative to the stellar velocity. Features broadened due to absorption by gas in the acceleration region would be asymmetric, with a sharp blue edge and significant red-shifted absorption by hot gas. Indeed, the model profiles of KR93 show this telltale characteristic (see figure 13 of KR93), but were unconstrained by the incomplete B79 published profile. Complete line profiles including the full blue-wing of the absorption feature were published in G88, revealing a rather symmetric line – at odds with KR93 model results.
Thus far, the linewidth data from G88 can not be reproduced in detail with this class of models, independent of the visibility data. The line widths were observed to increase with the energy of the rotational level. The 1-$`\sigma `$ widths were measured to be 0.8 km s<sup>-1</sup>, 1.7 km s<sup>-1</sup>, and 3.1 km s<sup>-1</sup>for the aR(0,0), aQ(2,2), and aQ(6,6) transitions respectively, while the relative velocity of the core remained constant to within $``$1 km s<sup>-1</sup>. While a near-star acceleration law and suitably peaked ammonia distribution can be adjusted to reasonably fit the linewidth of any single transition (e.g., figure 9), the large increase between the aQ(2,2) and aQ(6,6) line is not reproduced. Furthermore, lines broadened by accelerating gas should show significant shifts in the location of the core as a function of excitation energy.
The presence of decaying gas turbulence outside the acceleration region may explain qualitatively the linewidth behavior and the lack of a significant difference in the visibilities. The inner dust shell of IRC +10216 is already known to be quite clumpy in near-infrared images (Weigelt et al. (1998); Haniff & Buscher (1998); Tuthill, Monnier & Danchi (1998); Monnier (1999); Tuthill et al. (2000)). Furthermore, dynamical models of dust shell production in carbon stars lead to large velocity dispersion as the individual shells of dust are accelerated (Winters et al. (1995); Weigelt et al. (1998); Winters et al. (1998)). Both turbulence and velocity gradients across individual dust shells (clumps) will cause significant broadening of lines formed in these regions. Under these conditions, one would also expect the doppler velocity of the core to be nearly constant, consistent with observations, but the line widths would naturally decrease for lines forming further in the outflow as sound waves (or weak shocks) bring the outflow into hydrostatic equilibrium. The dynamical dust shell models of Winterset al.could be coupled with the radiative transfer code of Keady to quantitatively test these suppositions (as for CO in Winterset al.).
There are other potential explanations for why absorption lines by high-J states are much broader than for lower-excitation states. For instance, if the high-J states are populated by shock-heating, one would expect a higher velocity dispersion due the turbulence and velocity shifts accompanying shocks. This class of explanation, whereby the excitation mechanism (and possibly even the formation mechanism) is directly related to peculiar local conditions (e.g., shocks, high velocity dispersion), needs further consideration.
### 4.3 Summary for IRC +10216
The narrow line widths of all the silane lines and the lack of an interferometric signal place the location of silane formation at or beyond $``$80 R outside the turbulent inner envelope. In addition, the lack of interferometric signal and the broad (symmetric) lines of ammonia place their formation outside of $``$20 R, beyond the gas acceleration region, but in a turbulent (or clumpy) flow. These ideas are illustrated schematically in figure 10 and represent a significant improvement to our knowledge of the molecular stratification around IRC +10216.
## 5 VY CMa
While no silane has been observed in the oxygen-rich stellar wind of the red supergiant VY CMa, many absorption lines of ammonia have. A journal of these spectral line observations made with the ISI filterbank can be found in table 4.
#### 5.0.1 Previous work
The mid-infrared transitions of ammonia around VY CMa were first observed by McLaren & Betz (1980). However, a larger set of data with better frequency coverage was published by Goldhaber (1988, G88) and these data have been used for subsequent consideration. Details regarding the capabilities of the heterodyne spectrometer employed can be found in §4.1.1. While McLaren & Betz reported a significant change in the line profile of NH<sub>3</sub> aQ(2,2) between October 1978 and December 1979, the profile reported by G88 observed in September 1982 is nearly identical to the December 1979 data. In addition, the G88 line depth is consistent with the spectral data obtained by the filterbank in this paper (see table 5). In light of this consistency, the assumption that the G88 line profiles are still representative of the current epoch has been made.
All of the lines, aR(0,0), aQ(2,2), aQ(3,3), & aQ(6,6), observed by G88 possessed broad absorption features with 1-$`\sigma `$ widths of $``$3.8 km s<sup>-1</sup>. The cores of the lines were centered around V$`{}_{\mathrm{LSR}}{}^{}6`$ km s<sup>-1</sup>, with the high excitation line cores being slightly less blue-shifted. Assuming a stellar recessional velocity of V$`{}_{\mathrm{LSR}}{}^{}=22.3`$ km s<sup>-1</sup>(based on maser data), this corresponds to an outflow rate of $``$29 km s<sup>-1</sup>. Using these lines, G88 estimated a column density of ammonia of (6$`\pm `$4)$`\times `$10<sup>15</sup> cm<sup>-2</sup> using a simple model based on the Sobolev approximation. We have repeated these calculations using a similar gas model with the Keady code.
#### 5.0.2 Results and models
Table 5 shows the visibility ratios observed on and off of a number of NH<sub>3</sub> absorption features in Fall 1998. A stellar recessional velocity of V<sub>LSR</sub>=22.3 km s<sup>-1</sup>has been used to convert V<sub>LSR</sub> to expansion velocity, based on G88. The filterbank bandpasses selected were set to equal the full-width half-depths (FWHD) observed by G88, and were centered on the absorption cores. Data in table 5 represents a non-detection of the sP(4,3) line, which also had not been detected by any previous measurements. Just as for the IRC +10216 (see table 3), the visibility ratios were observed to be consistent with unity. In addition to the full aQ(6,6) absorption line, the red half of the line was observed separately in the hopes of observing the location of any accelerating gas. However, the SNR for this measurement ($``$10) was not high enough to justify further modeling.
A simple gas model based on the line-fit of the aR(0,0) line in G88 was developed using the new VY CMa dust shell model of Paper II. Unlike the case for IRC +10216, the linewidths for the NH<sub>3</sub> lines around VY CMa were all roughly the same. This fact and that the line core locations only showed a slight shift with excitation energy support a model in which the absorption lines all form mostly outside of the acceleration region. Therefore a constant outflow model was adopted using an expansion velocity of 29 km s<sup>-1</sup>, which matched the profiles of the aQ(2,2) and aQ(6,6) lines; the broad linewidths were created by model microturbulence of 3.5 km s<sup>-1</sup>.
The temperature of the J=K rotational states is controlled either by collisions, due to the lack of dipole transitions from these states to lower J, or by IR-pumping in the 6.1 $`\mu `$m band (see §2.3.2). During modeling, this temperature profile is normally varied using a power-law (as by KR93) to empirically fit the line depth ratios. Unfortunately there were insufficient number of lines available in this case to utilize this strategy. An alternative scheme, in which two temperature profiles were assumed to represent reasonable limits, was used here:
1. The rotational temperature was set equal to the dust temperature, T<sub>rot</sub>=T<sub>dust</sub>. This corresponded to full coupling of the internal degrees of freedom of the gas molecules to the radiation field.
2. T$`{}_{\mathrm{rot}}{}^{}=`$T$`{}_{}{}^{}r_{}^{0.60}`$, where r is expressed in units of stellar radii. This choice is similar to that found empirically for ammonia around IRC +10216 by Keady & Ridgway (1993). While VY CMa and IRC +10216 do have similar mass-loss rates (within a factor of $``$10), the differing chemistry (O-rich vs. C-rich) may affect the cooling rates, and this ad hoc temperature profile is clearly rather speculative.
Because G88 in most cases did not publish the emission components of the NH<sub>3</sub> transitions, we were not able to model this aspect of the line profile. Instead we followed KR93 (for IRC +10216) and suppressed emission in our model profiles by assuming the vibrational temperature was 85% of the dust temperature.
Figure 11 shows the results assuming T<sub>rot</sub>=T<sub>dust</sub>. The high optical depth at 11 $`\mu `$m ($`\tau 2.4`$, see Paper II) effectively shields the inner dust and molecular envelope from view. The $`\tau _{\mathrm{dust}}`$=1 surface is located at about 12 R along the central impact parameter and hence most of the observed continuum emission is occurring outside of this. This explains why none of the models with molecular formation radii between 10-80 R were decisively ruled out with the interferometry data of figure 11. However, one can see that for this assumed temperature law, the line ratios were best fit by a molecular formation radius outside of $``$40 R. Interestingly, this distance corresponds to the extent of the acceleration region deduced from proper motion studies of the H<sub>2</sub>O masers (Marvel (1996); Richards, Yates & Cohen (1998)). Figure 12 shows the results for the temperature profile of case 2, assuming T$`{}_{\mathrm{rot}}{}^{}=`$T$`{}_{}{}^{}r_{}^{0.60}`$; the same general conclusions still apply, showing that they are not highly model-dependent. It is worth noting that the equivalent widths for models with small molecular formation radii (especially for 10 R) were systematically higher than those measured by G88. While the line depths in the filterbank bandpass were fitted by increasing the ammonia density, the linewidths themselves were not. Ideally, the microturbulence parameter should be adjusted for each set of molecular formation radii to compensate for differing line broadening effects, such as from thermal broadening and the finite size of the continuum source. Such adjustments did indeed improve the quality of these fits, but significantly increased the complexity of the modeling while only weakly affecting the line and visibility ratios.
#### 5.0.3 Discussion
We conclude that NH<sub>3</sub> probably forms near the termination of the acceleration phase in the circumstellar envelope of VY CMa ($``$40 R). This hypothesis is supported by the following observations:
* Weak (but non-negligible) correlation of absorption core velocity with excitation energy suggests ammonia exists to a limited degree in the acceleration region.
* Line ratios of the aQ(2,2) and aQ(6,6) transitions were best fitted by a molecular formation radii $`\stackrel{>}{_{}}`$40 R, a result found to be insensitive to the rotational temperature profile.
* The visibility data for the high excitation aQ(6,6) line marginally supports a large molecular formation radius, but the high optical depth of the dust shell makes the interferometric data relatively insensitive to absorption in the inner circumstellar envelope.
Around IRC +10216, high excitation NH<sub>3</sub> lines were observed to be significantly broader than for the low-lying J-states. Interestingly, this behavior was not duplicated for NH<sub>3</sub> around VY CMa. Unlike the case for IRC +10216, it appears that ammonia is forming at least partly in the acceleration region around this star and hence the line formation of all the transitions is more affected by large turbulence in this region. Alternatively, the highly asymmetric inner dust shell seen in the near-infrared (Monnier et al. (1999)) and hinted at by ISI visibility data (Paper II) may be influencing the line formation characteristics. If the outflow is significantly asymmetric, linewidths will be clearly affected by the varying projected outflow velocities of absorbing gas. Continued monitoring of water masers in the inner dust shell may clarify the gas dynamics and level of gas turbulence. This information will be critical for further progress in understanding the geometry of the molecular envelope around VY CMa.
## 6 Conclusions
We have presented the first interferometric results on mid-infrared spectral lines around evolved stars. Specific results for the carbon star IRC +10216 and the red superigant VY CMa have been summarized in previous sections. In addition, a few general conclusions can be drawn:
* The formation radii of silane and ammonia are significantly beyond the dust formation zone for both evolved stars examined here.
* Since dust formation by itself does not catalyze formation of these molecules, some other physical mechanism(s), still unknown, must be at play. Probably the adsorption of gas-phase molecules (e.g., SiS or H<sub>2</sub>) onto grains set the time/length-scale for the important chemical reactions occurring on the grain surfaces.
* When coupled with spectroscopic observations, these results indicate that turbulence, or velocity dispersion, is quite high in the inner envelope and is more likely responsible for the broad linewidths of high excitation transitions than bulk acceleration.
###### Acknowledgements.
The authors would like to recognize W. Fitelson, C. Lionberger, and M. Bester for important contributions to the construction and software development of the filterbank, and K. McElroy for observing assistance. We also acknowledge productive discussions with A. Glassgold and A. Betz. JDM also would like to thank J. Keady for kindly allowing his sophisticated radiative transfer code to be used for this project, and for many stimulating discussions about the molecules of IRC +10216. This work is part of a long-standing interferometry program at U.C. Berkeley, supported by the National Science Foundation (Grant AST-9221105, AST-9321289, and AST-9731625) and by the Office of Naval Research (OCNR N00014-89-J-1583).
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# ”Particle-like” singular solutions in Einstein-Maxwell theory and in algebraic dynamics
## 1 . Introduction
As the main goal of the algebrodynamical approach , we regard the derivation of equations of physical fields (or of fundamental physical laws, in a broader sense) from a unique primary Principle of purely abstract nature, based only on the intrinsic properties of exclusive mathematical structures (groups, algebras, mappings etc.). Basic realization of this concept developed in makes use of only the differentiability conditions of functions of biquaternionic variable, i.e. of the Cauchy-Riemann (CR) equations generalized to the case of noncommutative (associative) algebras.
In the framework of the version of analysis proposed in the noncommutativity of starting algebra naturally results in the nonlinearity of generalized CR equations (GCRE), justifying the use of the latters as the dynamical equations of interacting fields.
Infact, for the algebra of biquaternions $`𝔹`$ (isomorphic to the full $`2\times 2`$ complex matrix algebra) the GCRE appear to be Lorentz invariant and carry natural 2-spinor and gauge structures. In the most examined case for which a simple geometric and physical interpretation is obvious, the conditions of $`𝔹`$-differentiability reduce to the invariant system
$$\eta (x+dx)\eta (x)d\eta =\mathrm{\Phi }(x)dX\eta (x)$$
$`(1)`$
for $``$-valued 2-spinor $`\eta (x)`$ and gauge $`\mathrm{\Phi }(x)`$ fields represented by $`2\times 1`$ and $`2\times 2`$ $``$-matrices respectively. $`()`$ in (1) denotes usual matrix multiplication (equivalent to that in $`𝔹`$), and $`dX`$ represents a $`2\times 2`$ Hermitian matrix of the increments of coordinates.
As a consequence of (1) each component of the 2-spinor $`\eta _A(x)(A,B=1,2)`$ satisfy the 4-eikonal equation . On the other hand, the compatibility conditions $`dd\eta =0`$ impose dynamical restrictions on the gauge field $`\mathrm{\Phi }(x)`$. Namely, the $`2\times 2`$ matrix $``$-valued connection 1-form
$$\mathrm{\Gamma }(x)=\mathrm{\Phi }(x)dX=\mathrm{\Gamma }^0(x)+\mathrm{\Gamma }^a(x)\sigma _a$$
$`(2)`$
($`\sigma _a,a=1,2,3`$ being the Pauli matrices) by virtue of compatibility conditions should be self-dual. Consequently, free Maxwell and Yang-Mills equations are satistied identically on the solutions of (1), for the scalar $`\mathrm{\Gamma }^0(x)`$ and the vector $`\mathrm{\Gamma }^a(x)`$ parts of the connection (2) respectively . Thus, the GCRE system (1) exhibit wonderful relations to several fundamental physical equations being for the latters in some sense generating.
In the approach regarded, an important role is performed by singular sets where the strengths of Maxwell and YM fields turn to infinity. Such sets were found to have diverse dimensions and topology. Solutions with compact structure of singular set may be then considered as particle-like, the evolution of such singularities-particles being then governed by the system (1) itself.
In this paper, we expound the process of reduction of system (1) to the equations of shear-free geodesic null congruences and, by Kerr theorem, to the solution of algebraic equation . Making use of this, we present an explicit form of the solutions to free Maxwell equations, analyze the structure and evolution of their singular set and discuss the relation of these electromagnetic fields to the Kerr-Shild metrics and, by this, to the solutions of electrovacuum Einstein-Maxwell equations. In conclusion, we consider general status and physical interpretation of ”particle-like” singular solutions in the framework of electrogravidynamics and of unified algebraic field theory proposed.
## 2 . Reduction of GCRE system and the solutions to free Maxwell equations.
Through the elimination of the gauge field $`\mathrm{\Phi }(x)`$ the system of GCRE (1) may be written in a 2-spinor form
$$\eta _A^{^{}}^{AA^{^{}}}\eta ^B^{^{}}=0$$
$`(3)`$
In the gauge $`\eta ^A^{^{}}(x)=(1,G(x))`$ the system (3) reduces to two equations for one unknown function $`G(x)`$
$$_{\overline{w}}G=G_uG,_vG=G_wG,$$
$`(4)`$
$`u,v=t\pm z,w,\overline{w}=x\pm iy`$ being the spinor coordinates. Note that as a consequence of (4) $`G(x)`$ satisfy identically both the 4-eikonal and d’Alembert wave equations . Assuming the Eqs. (4) are solved, the components of 4-potential matrix $`\mathrm{\Phi }(x)=A_0+A_a(x)\sigma _a`$ may be expressed through the function $`G(x)`$ as
$$A_w=_uG,A_v=_wG,A_u=A_{\overline{w}}=0$$
$`(5)`$
and satisfy free Maxwell equations.
Wonderfully, Eqs.(4) are completely identical to the equations of shear-free null geodesic congruences in the gauge for the spinor $`\eta (x)`$ regarded . In accord with Kerr’s theorem , we then obtain the general solution of Eqs.(4) in an implicit algebraic form
$$F(G,\overline{w}G+u,vG+w)=0,$$
$`(6)`$
$`F(G,\tau _1,\tau _2)`$ being an arbitrary golomorphic function of three complex variables including two twistor (projective) components
$$\tau _A=X_{AA^{^{}}}\eta ^A^{^{}},\tau _1=\overline{w}G+u,\tau _2=vG+w$$
$`(7)`$
Thus, a lot of solutions to free Maxwell equations may be obtained through simply examining of the algebraic Eq.(6). Singularities of related field strengths may be then found from the caustic condition
$$\frac{dF}{dG}_GF+\overline{w}_{\tau _1}F+v_{\tau _2}F=0$$
$`(8)`$
Eliminating the only unknown function $`G(x)`$ from two algebraic Eqs.(6),(8), one easily comes to the equation of singular set which determines the shape and evolution of singularities, without even taking care of explicit solving the Eq.(6) itself. The example of such procedure was presented in .
## 3 . Stationary solutions.
From the structure of twistor components (7) it’s evident that stationary solutions to Eq.(6) are exhausted by the functions $`F(G,\lambda )`$, where $`\lambda =G\tau _1\tau _2`$ doesn’t contain the time variable $`t=\frac{1}{2}(u+v)`$ at all. In accord with the results of Kerr and Wilson , for stationary solutions with compact structure of singular set the function $`F`$ should be at most quadratic in $`G`$, i.e. should have the form $`F=(G\tau _1\tau _2)+a_0G^2+a_1G+a_2,a_i`$. Linear dependence on $`G`$ immediately leads to the trivial solution with zero fields. Using 3-translations and 3-rotations, the above form may be reduced to $`F=G\tau _1\tau _22aG`$ for which we obtain from quadratic Eq.(6)
$$G(x)=\frac{x+iy}{za\pm \sqrt{(za)^2+x^2+y^2}}$$
$`(9)`$
For a real valued $`a`$ from (9) and the expression for potentials (5) we obtain the Coulomb electric field with a point singularity and a fixed value of the electric charge as a consequence of nonlinear primary system of GCRE (1). Imaginary values of $`a`$ correspond to the ring singularity with radius $`r=|a|`$ and multipole structure of EM fields with Coulomb first main term . Decomposition and separation of real parts of EM fields at a distance $`r|a|`$ gives
$$E_r\frac{e}{r^2}(1\frac{3a^2}{2r^2}(3\mathrm{cos}^2\theta 1)),E_\theta \frac{ea^2}{r^4}3\mathrm{cos}\theta \mathrm{sin}\theta ,$$
$$H_r\frac{2ea}{r^3}\mathrm{cos}\theta ,H_\theta \frac{ea}{r^3}\mathrm{sin}\theta .$$
$`(10)`$
Contrary to an ambiguous value of ring’s radius $`a`$, the dimensionless electric charge is strictly fixed (up to a sign) by field equations, so that in absolute units it may be identified with elementary charge $`e`$. Note that the solution (9) (for the case $`a=0`$) and the property of charge quantization for the GCRE system (1) have been obtained in a direct way in . (Recently there were some interesting attempts to explain electric charge quantization by topological reasons instead of dynamical considerations used here).
On the other hand, solution (9) is extremely important in the framework of GTR. Indeed, the expression $`l_\mu =\eta ^+\sigma _\mu \eta `$, where $`\eta ^T=(1,G(x))`$ is the 2-spinor related to the function $`G(x)`$, defines the principal null congruences $`l_\mu `$ of a Riemannian space-time endowed with a Kerr-Shild metric
$$g_{\mu \nu }=\eta _{\mu \nu }+H(x)l_\mu l_\nu ,$$
$`(11)`$
where $`\eta _{\mu \nu }`$ represents the metric of auxiliary Minkowski space-time. The scalar factor $`H(x)`$ should be then determined by the Einstein vacuum or electrovacuum equations and for fundamental stationary solution (9) leads to the Kerr or Kerr-Newman metrics respectively (consequently, to the Schwarzschild or Reissner-Nordström metrics for the case $`a`$=0). In our approach, it is of great importance that singularities of curvature of metric (10) are fixed by the same condition (8) as for electromagnetic field and define infact one unique particle-like object. Another wonderful fact is that for the Kerr-Newman solution of Einstein-Maxwell equations electromagnetic fields are just those defined by the GCRE system (apart from the property of quantization of charge for the latters!) and may be asymptotically presented by the Eq.(10). Moreover, these fields obey Maxwell equations both in flat space and in Riemannian space with Kerr-Shild metric (11)! Such a remarkable property of stability of electromagnetic fields under Kerr-Shild deformations of space-time geometry noticed in will be discussed elsewhere.
Making use of correspondence between the fundamental solution (9) to the GCRE system (1) and the Kerr-Newman solution of Einstein-Maxwell system, one is able to endow the solution (9) with a complete set of quantum numbers (including mass and spin). Then the gyromagnetic ratio would automatically correspond to that for Dirac particle, while the charge would be fixed in magnitude. Unfortunately, no natural reasons to ensure the quantization of mass could be seen nowadays. We’ll continue the discussion below.
### 3.1 . Nonstationary solutions.
Let us consider the general quadratic form of the function $`F(G,\tau _1,\tau _2)`$. When the terms bilinear in $`\tau _1`$, $`\tau _2`$ are absent, such functions (under the restriction on singular set to be compact) correspond to the boosted or rotated Kerr solution . On the other hand, the function $`F=\tau _1\tau _2b^2G`$ has been considered in detail in . The explicit expression for $`G`$ in this case is
$$G=\frac{2uw}{\sigma ^2+\rho ^2+b^2\pm \sqrt{\mathrm{\Delta }}},\mathrm{\Delta }(\sigma ^2+\rho ^2+b^2)^24\sigma ^2\rho ^2,$$
$`(12)`$
where $`\sigma ^2=uv=t^2z^2,\rho ^2=w\overline{w}=x^2+y^2`$. EM fields correspondent to (12) are
$$E_\rho =\frac{8b^2\rho z}{\mathrm{\Delta }^{3/2}},E_z=\pm \frac{4b^2}{\mathrm{\Delta }^{3/2}}(t^2z^2+\rho ^2+b^2),H_\phi =\frac{8b^2\rho t}{\mathrm{\Delta }^{3/2}}.$$
$`(13)`$
For real $`b`$ the fields (13) are identical to the well-known Born solution for two point-like charged ”particles” performing uniformly accelerated counter-motion. The value of electric charge for each particle does not depend on $`b`$ being fixed and equal to the charge of fundamental solution (9). For the case of imaginary $`b`$ one has the singularity of rather exotic toroidal structure, defined by the equation $`z^2+(\rho \pm b)^2=t^2`$ (see for details). In general case of complex-valued $`b`$ singular set manifests itself as the two rings of fixed radii performing again the oncoming hyperbolic motion along $`z`$-axis.
It may be proved that (up to the transformations of Poincare group) the axisymmetric solutions to the Eq.(6) (and to the GCRE system (1) respectively!) generated by quadratic function $`F`$ are exhausted by the Kerr-like solution (9) and the nonstationary bisingular solution (12) together with (toroidal or double ring-like) modifications of the latter.
It seems, however, that the solutions to GCRE with compact singularity and non-axial symmetries may be of interest too. Here we present an example of such solution which may be obtained from the generating function $`F=\tau _1\tau _2a^2G^2`$. Resolving the equation $`F=0`$, one comes to the following expression
$$G=\frac{2uw}{\pm \sqrt{\mathrm{\Delta }}+uv+w\overline{w}},$$
$`(14)`$
$`\mathrm{\Delta }`$ where $`\mathrm{\Delta }=(t^2x^2y^2z^2)^24a^2(t+z)(x+iy)`$. The singular set for this solution is defined by the condition $`\mathrm{\Delta }=0`$ and for $`t=0`$ has the form of flat figure ”8” curve (Fig.a) ). The time evolution of this singularity is illustrated by Fig.b).
EM fields related to the solution (14) and being represented by the complex combination $`\stackrel{}{}=\stackrel{}{E}i\stackrel{}{H}`$ are:
$$_+_1+i_2=\frac{2a^2w^2}{\mathrm{\Delta }^{3/2}},_{}_1i_2=\frac{2a^2u^2}{\mathrm{\Delta }^{3/2}},_3=\frac{2a^2uw}{\mathrm{\Delta }^{3/2}}.$$
$`(15)`$
For each finite moment of time they decrease rapidly (as $`r^4`$) with the distance from the centre of singularity. The fields are neutral (with total charge being equal to zero) and null ($`\stackrel{}{E}^2\stackrel{}{H}^2=0,\stackrel{}{E}\stackrel{}{H}=0`$).
In conclusion, we present a peculiar solution with noncompact singularity which serves as the analogue of electromagnetic wave in GRCE dynamics. For the solutions of wave-like type the generating function $`F`$ should depend only on one twistor component, say, $`\tau _1`$. Then, for the equation $`F(G,\overline{w}G+u)=0`$ the initial distribution of $`G(u)`$ may be arbitrary fixed at $`\overline{w}=0`$, i.e. at the $`Z`$-axis. Choosing for the latter the monochromatic dependence and resolving the equation $`GA\mathrm{exp}i\mathrm{\Omega }(\overline{w}G+u)=0`$, where the parameters $`A,\mathrm{\Omega }`$ are assumed to be real and positive, we find
$$G=iW(iA\mathrm{\Omega }\overline{w}\mathrm{exp}i\mathrm{\Omega }u)/\mathrm{\Omega }\overline{w},$$
$`(16)`$
$`W`$ being the principal branch of the so called Lambert function which is the solution of the equation $`W(z)\mathrm{exp}W(z)=z`$.
The structure of singular set is simply derived after that and appears to be a neutral helix of radius $`1/\mathrm{\Omega }Ae`$ and of lead $`2\pi /\mathrm{\Omega }`$ propagating along $`Z`$\- direction with the speed of light. Electromagnetic fields are mutually orphogonal and transversal while polarization depends on the distance from the axis. In the direction perpendicular to the axis the fields fall at large distancies as $`1/r`$. As before, the fields are globally defined only up to a sign.
Shear-free null geodesic congruences $`l_\mu `$ and the Kerr-Shild metrics (12) may be associated with the nonstationary solutions above-presented up to the scalar factor $`H(x)`$. At present it’s not clear if the latter may be choosed so that the Einstein-Maxwell electrovacuum system would be satisfied. By this, an interesting representation for the shear-free congruences (through consideration of null cone emanated by the source moving along some curve in complex space) as well as the condition of stability for electromagnetic fields under the Kerr-Schild deformations of space-time geometry may be of great use.
## 4 . General status of ”particle-like” singular solutions
Well-known are the numerous problems arising in GTR and in quantum field theory in respect to the singularities of solutions of field equations (violation of causality , divergences etc.). On the other hand, just the naked singularity of Kerr-Newman solution (which appears instead of black hole solution in the case of a large angular moment) manifests itself many remarkable properties related to that of elementary particles. Accordingly, several attempts to construct the model of electron on the base of Kerr-type solutions (KTS) have been undertaken .
However, they all dealt with the problem of physically suitable source for KTS to be found which is tigthly related to the well-known twovaluedness of Kerr-like geometry and electromagnetic fields in particular. Infact, the introduction of source becomes admissible only after the cut of space which restore the global uniqueness of the $``$-valued functions representing the fields of KTS.
Unfortunately, the surface of cut is quite ambigious: it may be either the disk spanning the Kerr singular ring or the oblate spheroid covering the singular ring on which the Kerr-Newman metric turns surprisingly into the Minkowski one. Consequently, one may think of the source of KTS as of the ”rotating relativistic disk”, of the ”bubble of flat geometry” within the external Kerr-Newman space-time etc. Thus, we are to conclude that there are no grounds to speak about the ”source” of KTS solutions at all since the twovaluedness is the unavoidable feature of their internal mathematical structure.
To illustrate the above statement, let us consider a simplier case of the singular ”particle-like” solutions to free Maxwell equations in flat space-time presented in this paper. Note that all of them (apart from the Coulomb and Born solutions with point-like singularities) are of the same two-valued structure being in each point defined up to a sign. Certainly, by no $`\delta `$\- function distribution of charge and current along the singular curve one can reproduce the field distribution in the whole space.
On the other hand, such solutions are locally well defined and may be analytically continuated from the region of regularity so that the full structure of singular set is established in a unique way. One cannot in any way change either the shape and topology of the singularity or its time evolution (the latter property being the most important from physical point of view). Suppose we really hope to describe the interactions and transmutations of elementary particles by means of the solutions regarded (which are rather to be multisingular for real physical process). Then we’ll proceed in well-defined and unique predictions in spite of partial indefinitness of EM fields and escape any divergences at all!
Moreover, one may think of such solutions as of the only possibility to explain the ”spin 1/2” structure at a purely classical level and with transparent picture of space-time dynamics being preserved. This is still more true in the framework of the algebrodynamical approach we develop, at least for two reasons. The first one is that in respect to the internal structure of GCRE system (1), gauge (electromagnetic plus Yang-Mills) fields stand there hand by hand with the 2-spinor structure so that the latter appears naturally together with Maxwell equations. The second reason is that, apart from the right value of gyromagnetic ratio, the value of electric charge is automatically fixed by the field equations themselves.
Of course, the stationary KTS as well as bisingular and ”figure 8” solutions presented here can say nothing about real dynamics of an ansamble of compact ”particle-like” singularities. Even the problem of interaction of two Kerr-Newman objects is far from solution. In the framework of algebrodynamics the overdetermined structure of GCRE impose restrictions even on the initial distribution of the fields so that the scattering problem should be fully reformulated.
Historically, the solution of Maxwell equations with ring singularity have been obtained by Appel in 1887 and revived in the works of Newman and Burinskii . General study of singular solution to Maxwell equations have been undertaken by Bateman . Nowadays the concept of naked singularities of KTS as the model for elementary particles is successfully developed, say, in the works of Clement .
To conclude, we argue that hostile attitude of physicists to singularities of field equations could be quite unjustified. There exist no restrictions of principal character for the compact multisingular solutions to describe the interactions of particle-like objects in a self-consistent way. Then their transmutations could be treated as perestroikas of singularities in terms of catastrophe theory. This programme should be implemented independently both in the framework of Einstein-Maxwell dynamics and of the algebraic dynamics based in particular on the GCRE system (1).
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# 1 Introduction
## 1 Introduction
Noncommutative structures in string theory seem worth studying not only because they arise as a nonperturbative effect as in D-branes , Matrix theory , or IIB matrix models , but because they would reveal essential features of string theory as a unique theory known to date describing various extended objects (including strings themselves) in a consistent manner. However, recent progress in the study of these structures suggests that the extendedness of a string in the space-time cannot be described by a mere specification of noncommutativity relation such as
$$[X^\mu ,X^\nu ]=i\theta ^{\mu \nu }$$
(1.1)
but requires more specification concerning microscopic details. In fact, quantum field theory on a noncommutative space defined by the above commutation relation with $`\theta ^{0i}0`$, where $`i`$ denotes one of the spatial directions, is found to behave as a theory consisting of extended objects like dipoles . Consequently, from the viewpoint of pointlike quanta, the theory exhibits acausal behavior, nonlocality and nonunitarity . On the other hand, it is observed in that open strings living on the same noncommutative space do not show such pathologies and exhibit causal time-delays, reflecting their extendedness. Furthermore, a space/time noncommutative field theory cannot be obtained as the low-energy limit of string theory . Thus the open string exhibits another origin of its consistent extendedness other than noncommutativity described by (1.1). From this point of view, it is important to examine causal time-delay behavior of the open string theory in further details and check its consistency. Another importance is in that it seems related to the space-time uncertainty principle . This is as it should be, because the space-time uncertainty principle basically prescribes how a string extends in space-time in a way consistent with (perturbative) unitarity and analyticity. Therefore, studying the causal time-delay behavior may shed light on identifying central features in string theory dynamics.
In the present paper, we shall be exploring causal time-delay that show up in high-energy scattering of open strings at fixed order in string perturbation theory. Our motivation for the study comes from various sides. The high-energy scattering may enable us to explore the short-distance structure of string theory. The reason why we go beyond the tree-level analysis done in is that in the high-energy scattering the contribution from higher genus to the amplitude is known to dominate the lower order ones . Therefore, it is important to examine the causal time-delay at higher orders and see if it is pronounced similarly as compared to the result of . Another motivation is that at high energy the space-time picture of the scattering becomes transparent because a certain class of classical trajectories become dominant . This effect permits examining whether the interpretation of time-delays given in persists to hold at higher genus as well.
The organization of this paper is as follows. In the next section, we give a brief recapitulation of the high-energy open string scattering at an arbitrary order in perturbation theory . In section 3 we proceed to evaluate the causal time-delay at each order in perturbation theory and find that it is reduced gradually at higher orders. In section 4, we discuss space-time interpretation of our result using various features of the saddle point configuration . Our analysis shows that the result is consistent with the interpretation in based on the property that the string grows in the longitudinal direction with its energy. In the final section, we make brief comments on relation with the space-time uncertainty principle and comparison with the case of space/time noncommutative field theory.
## 2 High-Energy Scattering of Open Strings
In this section, for our foregoing analysis, we recapitulate relevant results concerning the high-energy scattering of open strings . Consider, at genus $`G`$, four-point scattering amplitude in closed string theory given by:
$$𝒜_c^{(G)}(p_1,\mathrm{},p_4)=g^{2G+2}\frac{𝒟g_{\alpha \beta }}{𝒩}𝒟X^\mu \mathrm{exp}\left(\frac{1}{4\pi \alpha ^{}}d^2\xi \sqrt{g}g^{\alpha \beta }_\alpha X^\mu _\beta X_\mu \right)\mathrm{\Pi }_{i=1}^4V_i(p_i),$$
(2.1)
where $`g`$ is the closed string coupling constant, $`𝒩`$ is the volume of the group of the diffeomorphisms and Weyl rescalings, and $`V_i(p_i)`$ are the vertex operators. In the case of the scattering of four tachyons, the Gaussian integration over $`X^\mu `$ yields
$$𝒜_c^{(G)}(p_i)=g^{2G+2}__G[dm]\mathrm{\Pi }_id^2\xi _i\sqrt{g(\xi _i)}\mathrm{\Omega }_c(m,\xi _i)\mathrm{exp}\left(\alpha ^{}\underset{i<j}{}p_ip_jG_m(\xi _i,\xi _j)\right),$$
where $`m`$ are a set of coordinates for the moduli space $`_G`$ of genus $`G`$ Riemann surfaces with four punctures, $`\mathrm{\Omega }_c(m,\xi _i)`$ is the standard measure on $`_G`$, and $`G_m`$ is the scalar Green function on the genus-$`G`$ Riemann surface. In the kinematical limit all $`p_ip_j`$’s become large, one can perform the integration via saddle-point method and get
$$𝒜_c^{(G)}(p_i)g^{2G+2}\mathrm{\Omega }_c(\widehat{m},\widehat{\xi }_i)(_c^{\prime \prime })^{\frac{1}{2}}\mathrm{exp}\left(\alpha ^{}_c(p_i,\widehat{\xi }_i,\widehat{m}_i)\right),$$
where
$$_c(p_i,\xi _i,m)\underset{i<j}{}p_ip_jG_m(\xi _i,\xi _j)$$
is the electrostatic energy of an analog system of charges $`p_i`$ at $`\xi _i`$ on a Riemann surface with moduli $`m`$. $`\widehat{\xi }_i`$ and $`\widehat{m}`$ denote the values at the saddle point where $`(_c^{\prime \prime })^{1/2}`$, the determinant of the second derivative of $`_c`$, is evaluated.
Following the strategy of , high-energy scattering amplitudes of open strings may be obtained from those of closed strings by utilizing the reflection principle: doubling an open Riemann surface with boundaries yields a closed Riemann surface of higher genus. It turns out that, in the oriented open string theory, all relevant quantities associated with the original open Riemann surface can be extracted from the doubled, closed Riemann surface by imposing invariance under the reflection with respect to an appropriate symmetry plane . Since (2.1) can be easily generalized to Riemann surfaces with boundaries, the above argument shows that, in the high energy scattering of four open strings as well, the amplitude are approximated by its saddle point expression at the $`G`$-th order in open string perturbation theory
$$𝒜_o^{(G)}(p_i)g^{G+1}\mathrm{\Omega }_o(\widehat{m},\widehat{\xi }_i)(_o^{\prime \prime })^{\frac{1}{2}}\mathrm{exp}\left(\alpha ^{}_o(p_i,\widehat{\xi }_i,\widehat{m}_i)\right),$$
where $`\mathrm{\Omega }_o`$ denotes the measure on the moduli space of open Riemann surfaces with boundaries,
$$_o(p_i,\xi _i,m)2\underset{i<j}{}p_ip_jG_m(\xi _i,\xi _j),$$
is the analog electrostatic energy on the open Riemann surface, while $`G_m`$ denotes the scalar Green function for the doubled, closed Riemann surface of genus $`G`$. In fact, it can be shown that every saddle-point configuration of the oriented open string can be constructed from the associated saddle-point configuration of the closed string via reflection principle . The closed string saddle-point configuration at the $`G`$-th order was obtained in as the $`(G+1)`$-sheeted Riemann surface of the form
$$y^{G+1}=\mathrm{\Pi }_{i=1}^4(za_i)^{L_i},$$
(2.2)
where $`L_i`$’s are relatively prime to $`G+1`$, $`_iL_i=0`$ (mod $`G+1`$)Note that this condition prevents from considering all possible Riemann surfaces. Just for brevity, we will not treat the case of $`L_i0`$ (mod $`G+1`$). and the branch points $`a_i`$ are separated by $`1/(G+1)`$ times the period. Therefore one can construct the corresponding open string saddle-points and also evaluate the electrostatic energy $`_o^{(G)}`$ on them. It is given by
$$_o^{(G)}=\frac{_o^{(0)}}{G+1},$$
where
$$_o^{(0)}=s\mathrm{ln}|s|+t\mathrm{ln}|t|+u\mathrm{ln}|u|,$$
is the electrostatic energy on the disk, and $`s`$, $`t`$, $`u`$ are the Mandelstam variables: $`s=(p_1+p_2)^2`$, $`t=(p_2+p_3)^2`$, $`u=(p_1+p_3)^2`$. Thus, the high-energy scattering amplitude of four open tachyons at the $`G`$-th order in perturbation theory is approximated by
$$𝒜_o^{(G)}g^{G+1}\mathrm{exp}\left(\alpha ^{}\frac{s\mathrm{ln}|s|+t\mathrm{ln}|t|+u\mathrm{ln}|u|}{G+1}\right).$$
(2.3)
For later convenience, let us consider the nearly backward scattering $`u0`$ and compare the amplitude Eq.(2.3) at $`G=0`$ with the Veneziano amplitude in the same kinematical region. The Veneziano amplitude is given by
$$𝒜_{st}g\frac{\mathrm{\Gamma }(\alpha ^{}s)\mathrm{\Gamma }(\alpha ^{}t)}{\mathrm{\Gamma }(1+\alpha ^{}u)}+(tu)+(su),$$
(2.4)
up to kinematical factors. For small $`u`$, the first term can be expanded as
$$𝒜_{st}g\frac{\pi }{s\mathrm{sin}(\pi \alpha ^{}s)}(1\alpha ^{}u\mathrm{ln}(\alpha ^{}s)+\mathrm{}),$$
(2.5)
where we use the identity
$$z\mathrm{\Gamma }(z)\mathrm{\Gamma }(z)=\frac{\pi }{\mathrm{sin}\pi z}.$$
It should be noted that as long as $`u`$ is finite and non-zero, the second term in (2.5) is dominant in the high energy region, while setting $`u`$ to zero, only the first term in (2.5) would survive. This manifests the fact that the high-energy limit $`s\mathrm{}`$ does not commute with the backward scattering limit $`u0`$. Anther way to see this is that the high-energy scattering amplitude (2.3) becomes trivial if one sets $`u=0`$, a result which does not agree with (2.5). In fact, the above analysis cannot be applied to the exactly forward ($`t=0`$) or exactly backward ($`u=0`$) scattering, as, in such cases, the saddle point configurations approach boundaries of the moduli space $`_G`$, where two of the points $`\xi _i`$ at which vertex operators are inserted become coincident.
To proceed further and examine causal time-delay, we will need to have at hand asymptotic scattering amplitudes which are valid even at higher perturbative order. Therefore, in the following, we shall consider the high-energy, fixed angle scattering with a moderate scattering angle $`\varphi 0,\pi `$, where we can rely on the Gross-Mañes amplitude (2.3). Note that this is in contrast to the case in where the exactly backward scattering is considered and accordingly the scattering amplitude becomes singular in the sense that all poles contribute to it.
In concluding this section, it is worth pointing out a universal structure of the Gross-Mañes amplitude (2.3). At the $`G`$-th order in string perturbation theory, the $`G`$-dependence in (2.3) prompts to perform the following rescaling:
$$\alpha ^{}\widehat{\alpha }^{}\frac{\alpha ^{}}{G+1}.$$
(2.6)
Defining rescaled kinematic invariants
$`\widehat{s}\alpha ^{}`$ $``$ $`s\widehat{\alpha }^{}={\displaystyle \frac{s}{G+1}}\alpha ^{},`$
$`\widehat{t}\alpha ^{}`$ $``$ $`t\widehat{\alpha }^{}={\displaystyle \frac{t}{G+1}}\alpha ^{},`$
it follows that
$$𝒜_o^{(G)}(s,t)g^G𝒜_o^{(0)}(\widehat{s},\widehat{t}),$$
(2.7)
viz. the $`G`$-th order high-energy scattering amplitude can be reinterpreted as the tree-level amplitude at reduced values of kinematical invariants. As we will see later, universal features of the saddle point configurations yield a general and reasonable physical interpretation to the scaling (2.6).
## 3 Causal Time-Delay at Higher Genus
In this section, we would like to examine causal time-delay in the higher-order, high-energy scattering of four massless open string states. In order to do this, we first note that the saddle-point trajectory is universal and independent of the quantum numbers of scattered particles as the contribution from the vertex operators in (2.1) is at most the polynomials of momenta . Thus, up to kinematical prefactors, the tachyon amplitude (2.3) should be applicable to the massless case.
Let us consider the scattering of two massless open string states $`1+23+4`$, closely following the analysis of . For simplicity, we assume that the scattering takes place in a plane in a flat space-time and choose the center-of-mass system for the transverse momenta. As suggested in , one of the conceivable ways at present to extract the extendedness of strings from a scattering amplitude is that we first treat initial and final string states just as particle states by constructing their wave-packets with respect to the center-of-mass coordinates, and then observe the uncertainties of the interaction region, on which the extendedness would be reflected.
Following this method, we choose the wave-packet for each massless open string states as
$$\mathrm{\Phi }_i(x_i,p_i)=𝑑𝐤_if_i(𝐤_i𝐩_i)\mathrm{exp}\left(i(𝐤_i𝐱_i|𝐤_i|t_i)\right),$$
(3.1)
where $`f_i(𝐤)`$ is any function with a peak at $`𝐤=0`$. Then the S-matrix element is given by
$$3,4|S|1,2=\left(\mathrm{\Pi }_{i=1}^4𝑑𝐤_i\right)f_3^{}f_4^{}f_1f_2\delta (\underset{i=1}{\overset{4}{}}k_i)𝒜(s,t).$$
(3.2)
As proposed in , the uncertainty of the interaction region is measured by examining the response of the S-matrix under appropriate shifts of the particle trajectories in space-time. If we make shifts $`t_it_i+\mathrm{\Delta }t_i`$, then the wave-packet (3.1) becomes
$$\mathrm{\Phi }_i(x_i,p_i;\mathrm{\Delta }t_i)=𝑑𝐤_if_i(𝐤_i𝐩_i)\mathrm{exp}\left(i(𝐤_i𝐱_i|𝐤_i|t_i)\right)\mathrm{exp}(i|𝐤_i|\mathrm{\Delta }t_i).$$
(3.3)
For simplicity, set $`\mathrm{\Delta }t_1=\mathrm{\Delta }t_2=\mathrm{\Delta }t_3=\mathrm{\Delta }t_4=\mathrm{\Delta }t/2`$ and $`|𝐤_i|=E`$ for all $`i`$, the integrand in (3.2) would acquire an additional phase factor $`\mathrm{exp}(2iE\mathrm{\Delta }t)`$ due to this shift. Then the time-delay $`\mathrm{\Delta }T`$ (uncertainty in time) can be estimated as the decay width of (3.2) with respect to $`|\mathrm{\Delta }t|`$ under the insertion of this additional phase factor. As argued in , as $`\mathrm{\Delta }t`$ increases, the decay becomes appreciable when the variation of $`\mathrm{ln}𝒜(E)`$ is exceeded by the variation of the additional phase $`2E\mathrm{\Delta }t`$. Therefore, we can set
$$\mathrm{\Delta }T|\mathrm{\Delta }t|\frac{1}{2}\left|\frac{}{E}\mathrm{ln}𝒜(E)\right|.$$
(3.4)
Substituting the Gross-Mañes amplitude (2.3) into this equation and setting $`E=|𝐩|=p_0`$, we immediately obtain
$$\mathrm{\Delta }T\frac{\alpha ^{}p_0}{G+1}f(\varphi ),$$
(3.5)
where $`f(\varphi )`$ is defined as
$$f(\varphi )\mathrm{sin}^2\frac{\varphi }{2}\mathrm{ln}\mathrm{sin}^2\frac{\varphi }{2}\mathrm{cos}^2\frac{\varphi }{2}\mathrm{ln}\mathrm{cos}^2\frac{\varphi }{2},$$
(3.6)
thus fixed in our kinematical region. Although we cannot determine the sign of the time uncertainty by this method, we conclude that this is causal, because at tree level it is argued in that the string scattering is indeed causal, and even at higher genus, the scattering process remains almost the same as that at tree level due to universal features of the Gross-Mañes saddle point, as we will discuss in the next section. The reduced causal time-delay (3.5) is the main result in this section. It displays the fact that the causal time-delay at $`G`$-th order in string perturbation theory is reduced by $`1/(G+1)`$ compared to that at tree level. Here we note that, in (3.4), a numerical constant we have omitted should be independent of the genus, because it originates only from kinematics of the scattering process. Rather, it may be possible to take (3.4) as one of definitions of $`\mathrm{\Delta }T`$ in the framework here.
It should be emphasized that the time-delay at each order in perturbation theory is purely theoretical: what we will observe here must be a time-delay associated with the full amplitude. Nevertheless we believe our reduced time delay is important because, as we discussed later, it provides reasonable physical pictures and important implications of high-energy behavior of strings in the framework of perturbative string theory.
## 4 Reduction of Time-Delay: Space-Time Interpretation
In this section, we would like to draw certain space-time interpretation concerning the higher-genus reduction of the causal time-delay (3.5) by collecting various features the saddle-point exhibits.
Let us begin with the saddle-point trajectory of the closed string as originally found in . It is given by
$$X^\mu (z)=\frac{i}{G+1}\underset{i=1}{\overset{4}{}}\alpha ^{}p_i^\mu \mathrm{ln}|za_i|+𝒪(1/s),$$
(4.1)
where $`z`$ denotes the point on the $`(G+1)`$-sheeted Riemann surface (2.2). It is worth noting that, from (4.1), even at higher order, the shape of the saddle point trajectory remains the same as that at tree level, except that the overall space-time scale is reduced by a factor $`1/(G+1)`$. By examining the behavior at the vicinity of the branch points $`a_i`$, where the four vertex operators are inserted, one can easily get the following space-time picture of the saddle point trajectory : each of the two incoming closed strings winds around a closed curve $`(G+1)`$ times, then they interact and propagate as $`(G+1)`$ many intermediate short closed strings.To be more precise, there exist the cases in which we have only one intermediate string. They occur when $`L_i+L_j0`$ (mod $`G+1`$) in (2.2). However, they does not make a significant difference in later discussions. See also the footnote below Eq. (2.2). Subsequently, the $`G+1`$ short strings then rejoin together and finally produce two separate $`(G+1)`$-times wound outgoing closed strings. This picture is consistent with the fact that, as shown in (4.1), the string trajectory at $`G`$-th order is scaled down by a factor $`1/(G+1)`$ compared to that at tree level.
As recapitulated in the previous section, any oriented open string diagram can be obtained from corresponding closed string diagram by cutting each closed string into two open strings and keep one side of them. Therefore, high-energy scattering amplitude and space-time picture for open string can be obtained, via reflection principle, from those for closed string. From the space-time picture of the open string trajectory deduced this way, it follows immediately that, for the open strings as well, the space-time trajectory of the saddle point configuration at higher order in perturbation theory is exactly the same as that at tree level, except now that the space-time size is reduced by a factor $`1/(G+1)`$.
In interpreting the causal time-delay of open string scattering at tree level , key observations have been that the extension of the open string along the longitudinal direction grows linearly with energy, $`Lp_0\alpha ^{}`$, and that this leads to causal time-delay whose magnitude is proportional to $`L`$ and hence to $`p_0`$.<sup>§</sup><sup>§</sup>§Here, one assumes implicitly that the longitudinal direction of a string could be identified definitely. However, it is not evident that this is always possible especially at high energy where a string would oscillate frequently. Presumably our scattering should be argued in the infinite momentum frame which is the most natural in high-energy scattering. From this point of view, the fact that saddle-point trajectories of the open string are all universal enables us to understand intuitively reduction of the causal time-delay (3.5) at higher orders. At $`G`$-th order in perturbation theory, size of each open string along the longitudinal direction (scattering axis) is effectively scaled down by $`1/(G+1)`$ compared to that at tree level. Consequently, when scattered off each other, the open strings at $`G`$-th order would display the causal time-delay only by $`1/(G+1)`$ times that at tree level. This intuitive explanation is in complete agreement with our earlier result (3.5). Stated differently, as the incoming strings are folded or wound around $`(G+1)`$ times, their tension is effectively increased by a factor $`(G+1)`$ as in (2.6).As the string is folded over, one might have anticipated that the tension is reduced instead of being increased. However, this argument applies only to nonrelativistic string such as polymer. For relativistic string, the tension remains always equal to the mass density. Correspondingly, the causal time-delay is reduced as given in (3.5). Turning around the argument, from high-energy scattering in perturbation theory, we have confirmed that the causal time-delay exhibited therein is consistent with the interpretation that the longitudinal size of the string grows linearly with energy, $`Lp_0\alpha ^{}`$ (see (4.1)). This seems also consistent with the space-time uncertainty principle , on which we will discuss more later.
In concluding this section, we summarize the space-time trajectory of high-energy open string scattering as follows: according to (4.1) and the reflection principle, each of the two incoming open strings is in $`(G+1)`$-times multiply wound or folded configuration. As the two open string endpoints interact, at the scattering point, the wound or folded strings rearrange themselves into a configuration consisting of at most two short open strings and $`G/2`$ many singly wound, short closed strings. Characteristic size of these short strings is that of incoming open strings, viz. $`1/(G+1)`$-times smaller than the size at tree level. After interaction, the little strings rejoin and split off into two outgoing open strings, each of them are again $`(G+1)`$-times multiply wound or folded. In general, these outgoing open strings would have different lengths from the incoming ones. A cartoon view of the space-time trajectory of the scattering process is depicted in Figure 1.
## 5 Discussions
### 5.1 Comparison with Space/Time Noncommutative Field Theory
Let us compare our results with those expected in the space/time noncommutative field theory. As shown in , at tree level, the noncommutative field theory exhibits both advanced and retarded scattering. Even though the theory is generically non-unitary , one might put it aside and inquire of acausality at higher orders in perturbation theory: are the advanced/retarded effects present at higher orders in perturbation theory and, if so, are the effects increased or reduced as compared to those at tree level? It is well known that, in the maximal noncommutativity limit, $`\theta ^{\mu \nu }\mathrm{}`$, the scattering amplitudes of noncommutative field theory are dominated by planar loop graphs. This is so because all nonplanar loop graphs are completely suppressed due to the maximal noncommutativity. In planar graphs, there is no $`\theta `$-dependence apart from the overall Moyal phase associated with external momenta. These overall Moyal phase is the same for all planar graphs no matter how higher order they come from in perturbation theory. As such, at any order in perturbation theory, one does not expect any reduction of the retarded or advanced scattering behavior in the maximally noncommutative field theory. This should be contrasted to the fact that, in open string theory, universal structure of the saddle points has played an essential role in the reduction of the causal time delay. Indeed, exponential fall off and universal structure of the high-energy scattering amplitude (2.3) are characteristic features of open string theory, which are not present in quantum field theories with a finite number of field degrees of freedom.
### 5.2 Relation to Space-time Uncertainty Relation
Let us make comparison of what we have found with the space-time uncertainty relation:
$$\mathrm{\Delta }X\mathrm{\Delta }T>\alpha ^{}.$$
(5.1)
First it is worth noting that in our case this inequality is far from being saturated at lower order in perturbation theory. One might naively estimate that, in the present context,
$$\mathrm{\Delta }X\frac{p_0\alpha ^{}}{G+1}\mathrm{and}\mathrm{\Delta }T\frac{1}{E}\frac{1}{p_0},$$
such that the space-time uncertainty relation (5.1) is not obeyed. However, $`\mathrm{\Delta }T`$ should include not only the quantum fluctuation uncertainty but also the causal time-delay effect and, at high energy $`p_0^2\alpha ^{}1`$, the latter is the dominant effect:
$$\mathrm{\Delta }T\frac{p_0\alpha ^{}}{G+1}.$$
(5.2)
Thus, at a fixed order in perturbation theory, the left-hand side of (5.1) is much bigger than $`\alpha ^{}`$.We would like to thank T. Yoneya for pointing out this effect. Even including time-delay effect, the fact that both $`\mathrm{\Delta }X`$ and $`\mathrm{\Delta }T`$ become smaller at higher order as estimated above implies that the space-time uncertainty relation (5.1) would be saturated at the order $`G_{\mathrm{max}}\sqrt{p_0^2\alpha ^{}}`$. Does this mean that the space-time uncertainty relation is violated beyond that order in perturbation theory? We believe not so.
A possible resolution is that the causal time-delay may be enhanced further at higher-order in perturbation theory. Indeed, in case of high-energy closed string scattering, via Borel-resummation of fixed order saddle-point results , it has been already argued that high-energy asymptotic behavior could be significantly modified at nonperturbative level. After Borel-resummation, the scattering amplitude turns out to behaves as
$$𝒜_{resum}(s)e^\sqrt{s}$$
in the kinematical range $`\mathrm{log}(1/g^2)s1/g^{4/3}`$ and exhibits much bigger amplitude than any fixed-order perturbative behavior, $`e^{s/(G+1)}`$. Thus, once Borel-resummed, it might be that causal time-delay is considerably different from the fixed-order estimate (5.2). It would even be the case that the inequality in the space-time uncertainty relation (5.1) is saturated nonperturbatively in the high-energy regime, where the symmetry of the string theory is also believed to be enhanced enormously. In , the time-delay is indeed discussed based on (5.2). However, it should be noted that (5.2) was derived by the Borel resummation and as such it is one of possible guesses for nonperturbative high-energy amplitude. On the other hand, our time-delay is based on the amplitude at the saddle point which is certainly valid in the high-energy scattering with a fixed moderate angle. In this sense, we believe that although our time-delay is theoretical in itself, its reduced nature is important and that it indeed suggests the nonperturbative saturation of the space-time uncertainty relation. This is because if it were to grow with the order of the perturbation theory, it is impossible to expect such saturation. One can also find another importance of our time-delay in the following heuristic argument: suppose it is also possible to take advantage of the saddle point method in the summation over $`G`$ as discussed in , we find that the dominant order is given by $`G\sqrt{\alpha ^{}s/\mathrm{log}g}`$, and then we get $`\mathrm{\Delta }X\mathrm{\Delta }T\sqrt{\alpha ^{}}`$ up to some logarithmic corrections by substituting this into our time-delay. Therefore, this argument using our time-delay also implies that the space-time uncertainty relation is almost saturated at nonperturbative level.
Another possible resolution is that yet another new sort of nonperturbative effects begin to appear at the order $`G_{\mathrm{max}}`$. For example, as the strings are boosted to infinite momenta, $`p_0p_{}N/R`$ ( $`R`$ is an appropriate infrared cutoff scale) and hence $`G_{\mathrm{max}}N`$. This is the same order as the light-cone momenta $`p_+`$ and hence the width of the light-cone string diagram. We have seen, in previous sections, that the saddle-point trajectory is such that the string folds or winds around $`G`$-times. The fact that there is a maximal folding or winding and that each elementary folding or winding carries $`𝒪(1)`$ unit of $`p_+`$ momentum indicate that sub-structure of the open string begins to show up. Indeed, the effect is strikingly reminiscent of the matrix string theory, where a generic string configuration is built out of minimal length strings, each carrying precisely one unit of $`p_+`$. Using matrix string theory, it has even been found that D-brane pair-production becomes an important effect when the fragmented strings produced during the high-energy scattering are all of minimal length, viz. $`GG_{\mathrm{max}}`$. Whether this striking similarity is a mere coincidence or not deserves further study, especially, in light of the fact that the D-brane pair production at high-energy scattering has been found for closed strings and hence cannot be a feature unique to open strings only.
### Acknowledgement
We would like to thank S. de Haro, Y. Kazama, H.B. Nielsen, G. Veneziano and T. Yoneya for invaluable discussions and helpful comments on high-energy scattering in string theory. SJR acknowledges warm hospitality of Theory Division at CERN, where the work is completed.
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# On self-similar singular solutions of the complex Ginzburg-Landau equation.
## 1 Introduction
In this paper we study singular solutions to the initial value problem for the complex Ginzburg-Landau equation (CGL)
$`i{\displaystyle \frac{u}{t}}+(1iϵ)\mathrm{\Delta }u+(1+i\delta )|u|^{2\sigma }u`$ $`=`$ $`f,\text{in }^d\times (0,T),`$ (1)
$`u(x,0)`$ $`=`$ $`u_0(x).`$ (2)
where $`u=u(x,t)`$ is a complex-valued function defined in $`^d\times (0,T)`$ and satisfying suitable decay conditions as $`|x|\mathrm{}`$, the parameters $`ϵ`$, $`\delta `$, $`\sigma `$ are non-negative real numbers, $`u_0`$ is a given initial condition, which is assumed to be smooth, with suitable decay as $`|x|\mathrm{}`$, and $`f=f(x,t)`$ is a given function, also assumed to be smooth, with suitable decay as $`|x|\mathrm{}`$.
We are mostly interested in the case $`ϵ>0`$, $`\delta 0`$, and $`2/d<\sigma <2/d+1/2`$. However, some of our results are new even for the case $`ϵ0`$, $`\delta 0`$, $`\sigma >2/d`$, i.e., they also include the super-critical non-linear Schrödinger equation. Under the assumption $`ϵ>0`$, $`0<\sigma <2/d+1/2`$ it is possible to prove that for each smooth $`u_0`$ and $`f`$, with an appropriate decay at infinity, the problem (1)-(2) has a suitable weak solution, see and . Such a solution is regular away from a closed set $`𝒮^d\times (0,T)`$ with $`𝒫^{d2/\sigma }(𝒮)=0`$, where $`𝒫^\alpha `$ denotes the parabolic $`\alpha `$-dimensional Hausdorff measure, see . It has been an open problem whether the singular set $`𝒮`$ can be non-empty.
We present a very strong evidence, which is based on a combination of rigorous analysis and numerical computations, that singularities may indeed exist. We shall see that there appears to exist a countable set of different types of singular solutions. Among these solutions we identify (numerically) those which are stable under a suitable notion of stability defined in Section 4. It turns out that while most of the solutions are unstable, in certain cases there may exist more than one type of stable singularities.
The case of the non-linear Schrödinger equation (NLS) (i.e., $`ϵ=\delta =0`$) and $`f=0`$ has been studied by numerous authors, see for example or the monograph . In particular, Zacharov () conjectured the existence of self-similar singularities of the form
$$u(x,t)=\left(2\kappa (Tt)\right)^{\frac{1}{2}(\frac{1}{\sigma }+i\frac{\omega }{\kappa })}Q\left((2\kappa (Tt))^{1/2}|x|\right),$$
(3)
where $`Q(\xi )`$ is a complex valued function defined on $`(0,\mathrm{})`$, with asymptotic behavior
$$Q(\xi )\xi ^{1/\sigma i\omega /\kappa },\text{as}\xi \mathrm{}.$$
While no rigorous proof of this conjecture seems to be available, there is an overwhelming evidence based on numerical and formal analytical calculations supporting the existence of such singularities (see, e.g., ). We also refer the reader to for a rigorous result supporting the conjecture.
In this paper we will argue that these singularities persist also for a certain range of $`ϵ>0`$ and $`\delta >0`$. In fact, we shall find many new self-similar singularities even for the case $`ϵ=\delta =0`$ (and $`f=0`$). Using the self-similar singular solutions of the form (3) one can easily construct singular solutions of (1) and (2) with compactly supported, smooth $`u_0`$ and compactly supported, smooth $`f`$.
From the form (3) of self-similar solutions one obtains the following boundary value problem for the function $`Q`$:
$`(1iϵ)(Q^{\prime \prime }+{\displaystyle \frac{d1}{\xi }}Q^{})+i\kappa \xi Q^{}+i{\displaystyle \frac{\kappa }{\sigma }}Q\omega Q+(1+i\delta )|Q|^{2\sigma }Q=0,`$ (4)
$`Q^{}(0)=0,\text{and}`$ (5)
$`Q(\xi )\xi ^{1/\sigma i\omega /\kappa }\text{as}\xi \mathrm{}.`$ (6)
In general, this problem does not have a solution for arbitrary values of parameters $`\kappa `$, $`\omega `$ and hence the unknowns in (4)-(6) are $`Q`$, $`\kappa `$ and $`\omega `$.
When $`ϵ=0`$, see, for example, for a discussion of (6). For $`ϵ>0`$ the condition $`Q(\xi )\xi ^{1/\sigma i\omega /\kappa }`$ at infinity is dictated by the partial regularity result mentioned earlier, since $`u`$ defined by (3) must be regular at almost all points of the form $`(x,T)`$, $`x0`$. The presented results are based on a detailed analysis of the boundary value problem (4)-(6). The analysis employs both analytical and numerical techniques and is naturally divided into the following two steps:
First, we rigorously prove that for $`\xi _11`$ and for each set of parameter values $`ϵ0`$, $`\delta 0`$ and $`2/d<\sigma <2/d+1/2`$ there exists a two-dimensional manifold of solutions to (4) on $`(\xi _1,\mathrm{})`$ with the correct asymptotic behavior at infinity.
Second, we solve numerically the boundary value problem on $`(0,\xi _1)`$ with an appropriate (approximate) boundary condition for $`Q`$ at the boundary point $`\xi =\xi _1`$. In this step the choice of $`\xi _1`$ is based on numerical evidence of convergence, and not on the rigorous estimates obtained in (i), which contain constants we did not try to evaluate exactly.
We briefly summarize main results of our calculations. We note a (scaling) symmetry in the problem (4)-(6): If $`\lambda >0`$, $`\theta `$ and $`(Q(\xi ),\kappa ,\omega )`$ is a solution of (4)-(6) then $`(\lambda ^{1/\sigma +i\theta }Q(\lambda \xi ),\lambda ^2\kappa ,\lambda ^2\omega )`$ is also a solution. We chose a representative in each of these families of solutions by imposing suitable normalization conditions. We mostly work with the normalization $`\omega =1`$ and $`\mathrm{Im}Q(0)=0`$, $`Q(0)>0`$. Such solutions will be called normalized solutions and they are uniquely determined by two parameters $`(\kappa ,\mu )`$, where $`\mu =Q(0)`$. Some numerical results are better presented in a different parameterization in which $`Q(0)=1`$ is fixed, $`(\kappa ,\omega )`$ are used as parameters. We shall state the use of the latter normalization explicitly whenever it is used.
We observe that for fixed values of $`d1`$, $`\sigma >2/d`$ and $`ϵ=\delta =0`$ there exists a countable family of normalized solutions $`(Q_j(\xi ),\kappa _j)=(Q_j(\xi ),\kappa _j,1)`$, $`j=1,2,\mathrm{}`$ of (4)-(6). The $`j`$th profile $`|Q_j|`$, when extended to $`(\mathrm{},\mathrm{})`$ as an even function, has exactly $`j`$ local maxima, and somewhat resembles the profile of the $`j`$th state of an elementary quantum mechanical oscillator. A precise quantum mechanical interpretation of the solutions is more complicated, and is related to the so-called resonances or quasi-stationary states. The first solution of the family has been known, see or , for example. We are not aware of any mentioning of the other solutions in the literature. All the solutions persist if $`ϵ`$ and $`\delta `$ are perturbed to (small) strictly positive values. We now describe the behavior of these perturbed solutions for $`\delta =0`$. We let $`\mu _j=Q_j(0)`$. As mentioned above, the solution $`(Q_j,\kappa _j)`$ is determined by $`(\kappa _j,\mu _j)`$. We observe that a branch of solutions parameterized by $`(\kappa _j(ϵ),\mu _j(ϵ))`$ emanates from each point $`(\kappa _j,\mu _j)`$. We plot the curves $`(ϵ,\kappa _j(ϵ))`$, $`j=1,\mathrm{},5`$, in Figure 1 for $`d=1`$, $`\sigma =2.3`$ and in Figure 6 for $`d=3`$, $`\sigma =1`$. In the other figures we plot the profiles $`|Q_j^ϵ(\xi )|`$ of the corresponding solutions $`Q_j^ϵ(\xi )`$ at certain points along each branch.
An interesting feature observed in the behavior of the branches, is the existence of a turning point on each branch at $`ϵ=ϵ_j^{}`$. We see from the graphs that as $`ϵ`$ returns to zero along the branch, $`\kappa _j(ϵ)`$ tends to zero, suggesting that the solution $`Q_j^ϵ`$ converges to a (radial) solution of the equation
$$\mathrm{\Delta }QQ+|Q|^{2\sigma }Q=0,\text{in }^d,$$
(7)
satisfying $`Q(x)0`$ as $`|x|\mathrm{}`$.
Our computations presented in Section 3 clarify the structure of the diagram which turns out to be slightly more complicated than the picture suggested above. We conjecture the following: If $`d>1`$, the solutions $`Q_j^ϵ`$ corresponding to branch $`(\kappa _j(ϵ),\mu _j(ϵ))`$ with an odd index $`j=2k1`$ converge (as $`ϵ`$ and $`\kappa _j(ϵ)`$ approach zero) to the $`k`$-th (normalized) radial solution of (7). We recall that the first of these solutions is usually called the ground state, and that for $`d=1`$ there are no other solutions of (7) satisfying the appropriate boundary conditions. (See, for example, for more details.) If $`d=1`$ and also for $`j=2k`$ in the case $`d>1`$, as $`ϵ`$ and $`\kappa _j(ϵ)`$ approach zero, the profiles $`Q_j^ϵ`$ separate into $`j`$ approximate ground states which move away from each other. In particular, for $`j=2k`$ the profiles converge locally uniformly to zero. When $`d=1`$ and $`j=2k1`$, the profiles converge locally uniformly to the ground state.
We tabulate results of our numerical calculations in Table 1 and Table 2.
We also looked at branches of solutions when $`\delta `$ is related to $`ϵ`$ by $`\delta =rϵ`$, with $`r>0`$ of order $`10^1`$ and $`10^0`$. The behavior was similar to the case $`\delta =0`$, with the turning point $`ϵ^{}`$ getting closer to zero as $`r`$ increased, as one might heuristically expect.
Questions related to stability of the singularities are addressed in Section 4. Our calculations indicate that for the non-linear Schrödinger equation all the new singularities we found are unstable, and the singularity corresponding to $`(\kappa _1,\mu _1)`$ is stable. The situation is more complicated for $`ϵ>0`$, see Section 4 for details.
Our interest in singular solutions to CGL stems from analogies between (1) in the case $`d=3`$, $`\sigma =1`$, $`ϵ,\delta >0`$ and the three-dimensional Navier-Stokes equation (NSE). The two equations have the same scaling properties and the same energy identity. Moreover, the existence and partial regularity theory of weak solutions for NSE and CGL are similar (with CGL being technically easier), see . The analogy between NSE and CGL may be rather superficial and may break down at any deeper level. However, at the same time there are no known properties of solutions to the Navier-Stokes equation which would prevent the same scenario as presented here for CGL.
The formula for the Navier-Stokes equation corresponding to (3) would be
$$u(x,t)(2\kappa (Tt))^{1/2}U((2\kappa (Tt))^{1/2}x,\tau ),$$
(8)
where $`\tau =\frac{1}{2\kappa }\mathrm{ln}\frac{T}{Tt}`$, and $`U`$ is a suitable divergence-free vector field periodic in $`\tau `$ with suitable decay in the self-similar variable $`y=(2\kappa (Tt))^{1/2}x`$.
The case, $`U/\tau 0`$, was already considered by Leray . It was proved in and in greater generality in that NSE does not admit non-trivial solutions of the form (8) with $`U`$ independent of $`\tau `$. The problem is open for $`U`$ periodic in $`\tau `$.
We finish the introduction with the following speculation. Most of the singularities we have found are unstable, hence it is unlikely they would be observed in direct numerical simulations of the initial value problem (1)-(2) or in physical experiments that are modeled by CGL. Could it perhaps be the case that NSE does admit singular solutions (say of the form (8)), but all of them are unstable and therefore more or less impossible to be detected in direct numerical simulations or physical experiments ?
This intriguing scenario was once suggested to one of the authors by Sergiu Klainerman during a lunch-break conversation at a conference in Southern California.
## 2 Analysis of the profile equation at infinity
In this section we study solutions of (4) in the interval $`(\xi _1,\mathrm{})`$ satisfying the condition $`Q(\xi )\xi ^{1/\sigma i\omega /\kappa }`$ as $`\xi \mathrm{}`$. Heuristically one expects that the behavior of such solutions is mainly governed by the linear part of (4):
$$(1iϵ)u^{\prime \prime }+(1iϵ)\frac{d1}{\xi }u^{}+i\kappa \xi u^{}+\frac{i\kappa }{\sigma }u\omega u=0$$
(9)
Equation (9) is equivalent to Kummer’s equation, also known as the confluent hypergeometric equation. The solutions of this equation are well-understood, see for example , and one can hence get a more or less complete picture of the behavior of solutions to (9). We shall use analytical tools from the theory of confluent hyper-geometric equations to describe solutions of the full equation (4).
A canonical form of Kummer’s equation is
$$z\frac{\mathrm{d}^2w}{\mathrm{d}z^2}+(bz)\frac{\mathrm{d}w}{\mathrm{d}z}aw=0,$$
(10)
and the equation (9) is transformed into this form by letting
$$z=\frac{i\kappa }{(1iϵ)}\frac{\xi ^2}{2},a=\frac{1}{2}\left(\frac{1}{\sigma }+\frac{i\omega }{\kappa }\right),b=\frac{d}{2}.$$
(11)
There is voluminous literature on this equation and properties of special functions (confluent hyper-geometric functions) which appear as its solutions. We recall some of the properties and, for the convenience of the reader, we also sketch how to derive them. For more details about confluent hyper-geometric functions we refer the reader to , , .
A classical formula for a solution of (10) is given by
$$U(a,b,z)=\frac{1}{\mathrm{\Gamma }(a)}_0^{\mathrm{}}e^{tz}t^{a1}(1+t)^{ba1}𝑑t,$$
(12)
where $`\mathrm{\Gamma }`$ is Euler’s gamma function. The integral is clearly well-defined for $`\mathrm{Re}a>0`$ and $`\mathrm{Re}z>0`$. The factor $`1/\mathrm{\Gamma }(a)`$ is not essential for our analysis, nevertheless, we include it to keep our notation in agreement with the standard one. The role of this factor is to normalize the leading term in the asymptotic $`U(a,b,z)=z^a(1+\mathrm{O}(z^1))`$ as $`z\mathrm{}`$. It is easy to check by direct calculation that the function $`U`$ given by (12) solves the equation (10). We have
$$\frac{d^k}{dz^k}U(a,b,z)=\frac{1}{\mathrm{\Gamma }(a)}_0^{\mathrm{}}(t)^ke^{tz}t^{a1}(1+t)^{ba1}𝑑t,$$
and after substituting to (10) we see from a simple integration by parts that the equation is satisfied.
By letting $`zt=s`$ in (12) we obtain (for $`\mathrm{Re}z>0`$, $`\mathrm{Re}a>0`$)
$$U(a,b,z)=\frac{1}{\mathrm{\Gamma }(a)}z^a_0^{\mathrm{}}e^ss^{a1}\left(1+\frac{s}{z}\right)^{ba1}𝑑s.$$
(13)
The above expression is used to extend the definition of $`U`$ for $`\mathrm{Re}a>0`$ and $`z`$ with $`\pi <\mathrm{arg}z<\pi `$, since the integral is convergent and an analytic function of $`z`$ under these assumptions. By formally expanding the term $`(1+s/z)^{ba1}`$ and integrating the resulting (formal) series term-by-term, we obtain the following asymptotic expansion
$$U(a,b,z)=z^a\left(\underset{k=0}{\overset{n1}{}}\frac{(a)_k(1+ab)_k}{k!}(z)^k+\mathrm{O}(|z|^n)\right),$$
(14)
where $`(a)_0=1`$, $`(a)_k=a(a+1)\mathrm{}(a+k1)`$. The formal calculations of the asymptotic expansion (14) can be easily justified rigorously by splitting the integral in (13) as $`_0^{\mathrm{}}=_0^{s_1}+_{s_1}^{\mathrm{}}`$ with $`s_1=|z|^{1/2}`$, for example.
One can check easily by a direct calculation that the function
$$V(a,b,z)=e^zU(ba,b,z)$$
(15)
is another solution of Kummer’s equation and that the functions $`U`$ and $`V`$ are linearly independent. The Wronskian of $`U`$ and $`V`$ is given by
$$U\frac{\mathrm{d}V}{\mathrm{d}z}V\frac{\mathrm{d}U}{\mathrm{d}z}=e^{\pm i\pi (ba)}z^be^z,$$
(16)
where the sign $`+`$ is for $`\mathrm{Im}z>0`$ and $``$ in the opposite case. (Formula (16) is easily derived from the fact that the Wronskian satisfies the differential equation $`y^{}+(b/z1)y=0`$ and from the asymptotic expansion (14).)
We need to know the behavior of $`U`$ and $`V`$ in the region $`\pi /2\mathrm{arg}z<0`$. Taking into account the definition of $`V`$, we see that it is sufficient to control $`U`$ in the region $`\pi /2\mathrm{arg}z<\pi `$. Formula (13) is suitable for analysis in this region if $`\mathrm{Re}(ba)>0`$ since the integral in (13) is the uniformly absolutely convergent whenever $`z`$ approaches a point in $`(\mathrm{},0)`$ from the upper half-plane. In our applications the condition $`\mathrm{Re}(ba)>0`$ is always satisfied and therefore (13) is sufficient for our analysis.
We now have sufficient information about the solutions of (10), and hence also (9), to be able to proceed with the analysis of the inhomogeneous equation
$$(1iϵ)u^{\prime \prime }+(1iϵ)\frac{d1}{\xi }u^{}+i\kappa \xi u^{}+\frac{i\kappa }{\sigma }u\omega u=f(\xi ),$$
(17)
for $`\xi (\xi _1,\mathrm{})`$. We assume that the function $`f`$ is decaying sufficiently fast as $`\xi \mathrm{}`$. We are interested in solutions of (17) which have the asymptotics $`u\xi ^{1/\sigma i\omega /\kappa }`$ as $`\xi \mathrm{}`$. We denote $`P`$, $`E`$ two linearly independent solutions of (4)
$`P(\xi )P(\kappa ,\omega ,ϵ;\xi )=U(a,b,z),`$
$`E(\xi )E(\kappa ,\omega ,ϵ;\xi )=V(a,b,z),`$
$`\text{where }z={\displaystyle \frac{i\kappa }{1iϵ}}{\displaystyle \frac{\xi ^2}{2}},a={\displaystyle \frac{1}{2}}\left({\displaystyle \frac{1}{\sigma }}+i{\displaystyle \frac{\omega }{\kappa }}\right),b={\displaystyle \frac{d}{2}}.`$
The parameters $`d`$ and $`\sigma `$ are held fixed in the perturbation analysis described below, therefore we do not indicate the dependence of $`P`$ and $`E`$ on them. The Wronskian $`W=PE^{}P^{}E`$ is easily computed from (16).
$$W(\xi )=W(\omega ,\kappa ,ϵ;\xi )=\frac{i\kappa }{1iϵ}e^{\pm \pi i(ba)}\xi z^be^z.$$
(18)
The next step is to use the standard variation of constant to obtain solutions to (17) satisfying $`u(\xi )\xi ^{1/\sigma i\omega /\kappa }`$ as $`\xi \mathrm{}`$. We write the solution in the form
$`u(\xi )`$ $`=`$ $`c_1(\xi )P(\xi )+c_2(\xi )E(\xi ),\text{with}`$ (19)
$`u^{}(\xi )`$ $`=`$ $`c_1(\xi )P^{}(\xi )+c_2(\xi )E^{}(\xi ),`$ (20)
and $`c_2(\xi )0`$ “sufficiently fast” as $`\xi \mathrm{}`$. We obtain
$$c_1^{}(\xi )=f(\xi )\frac{E(\xi )}{(1iϵ)W(\xi )},c_2^{}(\xi )=f(\xi )\frac{P(\xi )}{(1iϵ)W(\xi )},$$
(21)
which together with the condition $`c_2(\xi )0`$ at the infinity gives a formal expression for the solution
$$u(\xi )=\gamma P(\xi )+_{\xi _1}^{\mathrm{}}K(\xi ,\eta )f(\eta )𝑑\eta ,$$
(22)
where $`\gamma `$ is a constant and
$$K(\xi ,\eta )=\{\begin{array}{cc}\frac{1}{(1iϵ)}P(\xi )E(\eta )W^1(\eta )\hfill & \text{for }\xi _1<\eta \xi \hfill \\ \frac{1}{(1iϵ)}E(\xi )P(\eta )W^1(\eta )\hfill & \text{for }\xi \eta \hfill \end{array}.$$
The strategy for finding solutions of the non-linear problem (4) in $`(\xi _1,\mathrm{})`$ with the required decay $`u(\xi )\xi ^{1/\sigma i\omega /\kappa }`$ at infinity should now be clear. Heuristically we expect that the manifold of such solutions will be a deformation of the one-dimensional complex subspace $`\{\gamma P|\gamma \}`$, at least in a neighborhood of the origin. The deformation is from $`\gamma P`$ to the fixed point of the operator $`T`$
$$uTu(\xi )=\gamma P(\xi )_{\xi _1}^{\mathrm{}}(1+i\delta )K(\xi ,\eta )|u(\eta )|^{2\sigma }u(\eta )𝑑\eta .$$
The outlined strategy can be successfully carried out by using the properties of Kummer’s functions recalled above. The fixed point theorem can be applied in the Banach space
$$𝒳_\vartheta =\{uC([\xi _1,\mathrm{}))|\underset{\xi \xi _1}{sup}|\xi |^{1/\sigma \vartheta }|u(\xi )|<\mathrm{}\}$$
equipped with the norm
$$u_\vartheta =\underset{\xi \xi _1}{sup}|\xi |^{1/\sigma \vartheta }|u(\xi )|.$$
This approach is obviously standard. However, there are some subtle points in the situation studied here due to oscillatory behavior of the function $`E`$, see the Appendix.
The main result of this section, which can be derived in a fully rigorous way from the above analysis is the following theorem. (A complete proof of the theorem is presented in the Appendix.)
###### Theorem 1.
Assume $`1d3`$, $`2/d<\sigma <2/d+1/2`$, $`0<\kappa _1<\kappa _2`$, $`\omega _1<\omega _2`$, $`0<ϵ_1`$, $`\delta _1<\delta _2`$.
There exists $`\rho _1>0`$ such that for each $`\xi _11`$ and each
$$(\beta ,\kappa ,\omega ,ϵ,\delta )\mathrm{\Lambda }_{\rho _1}\{\beta ||\beta |\rho _1\}\times [\kappa _1,\kappa _2]\times [\omega _1,\omega _2]\times [0,ϵ_1]\times [\delta _1,\delta _2]$$
the boundary value problem
$`(1iϵ)(Q^{\prime \prime }+{\displaystyle \frac{d1}{\xi }}Q^{})+i\kappa \xi Q^{}+i{\displaystyle \frac{\kappa }{\sigma }}Q\omega Q+(1+i\delta )|Q|^{2\sigma }Q=0,`$ (23)
$`Q(\xi _1)=\beta ,\text{and}`$ (24)
$`Q(\xi )\xi ^{1/\sigma i\omega /\kappa }\text{as}\xi \mathrm{}.`$ (25)
considered in $`(\xi _1,\mathrm{})`$ has a solution
$$F(\xi )=F(\beta ,\kappa ,\omega ,ϵ,\delta ,\xi _1;\xi ).$$
Moreover, $`F`$ can be constructed in such a way that the following conditions are satisfied
The mapping from $`\mathrm{\Lambda }_{\rho _1}`$ to $`𝒳_\vartheta `$ defined by
$$(\beta ,\kappa ,\omega ,ϵ,\delta )F(\beta ,\kappa ,\omega ,ϵ,\delta ,\xi _1;)$$
is $`C^1`$ up to the boundary for each $`\vartheta >0`$.
The complex-valued function defined by
$$(\beta ,\kappa ,\omega ,ϵ,\delta )\frac{F}{\xi }(\beta ,\kappa ,\omega ,ϵ,\delta ,\xi _1,\xi _1)$$
is $`C^1`$ (up to the boundary) in $`\mathrm{\Lambda }_{\rho _1}`$.
$`F`$ and its derivatives $`F^{(k)}`$ have the following asymptotic expansions
$`F(\xi )=\xi ^{1/\sigma i\omega /\kappa }\left({\displaystyle \underset{l=0}{\overset{n}{}}}a_l\xi ^{2l}+\mathrm{O}(\xi ^{2(n+1)})\right),`$ (26)
$`F^{(k)}(\xi )={\displaystyle \frac{^k}{\xi ^k}}\left(\xi ^{1/\sigma i\omega /\kappa }{\displaystyle \underset{l=0}{\overset{n}{}}}a_l\xi ^{2l}\right)+\mathrm{O}(\xi ^{1/\sigma 2(n+1)k}).`$ (27)
We have $`F(0,\kappa ,\omega ,ϵ,\delta ,\xi _1,\xi )=0`$ and
$$\frac{F}{\beta }(0,\kappa ,\omega ,ϵ,\delta ,\xi _1,\xi )=\frac{P(\kappa ,\omega ,ϵ,\xi )}{P(\kappa ,\omega ,ϵ,\xi _1)}.$$
Proof: See the Appendix
###### Remark 2.1.
For $`ϵ>0`$ the function $`u(x,t)`$ given by (3) has to be regular at all points $`(x,t)`$ with $`x0`$ by the partial regularity theorem proved in . Therefore any solution of (4)-(6) must admit an asymptotic expansion of the form stated in (iii) of Theorem 1. We note that the convergence of the series $`a_l\xi ^{2l}`$ in the asymptotic expansion of $`Q`$ is equivalent to the analycity of $`u`$ (in $`t`$) at the points $`(x,T)`$, $`x0`$. The asymptotic expansion of the profiles $`Q`$ does not converge and therefore $`u`$ is not analytic in $`t`$ at any point $`(x,T)`$.
Theorem 1 together with elementary perturbation arguments can be used to show that if some non-degeneracy conditions are satisfied, then every solution of (4)-(6) for $`ϵ=\delta =0`$ will persist (with a slight deformation) for small $`ϵ>0`$ and $`\delta >0`$. We will briefly describe this standard procedure for the convenience of the reader.
First we consider solutions on the finite interval $`(0,\xi _1]`$ to the initial value problem
$`(1iϵ)(Q^{\prime \prime }+{\displaystyle \frac{d1}{\xi }}Q^{})+i\kappa \xi Q^{}+i{\displaystyle \frac{\kappa }{\sigma }}Q\omega Q+(1+i\delta )|Q|^{2\sigma }Q=0,`$ (28)
$`Q(0)=\mu ,Q^{}(0)=0.`$ (29)
We denote the solution (if it exists) by $`G(\xi )=G(\mu ,\kappa ,\omega ,ϵ,\delta ;\xi )`$. Clearly the set of parameters $`(\mu ,\kappa ,\omega ,ϵ,\delta )`$ for which $`G`$ is well-defined is open. We define
$$\beta (\mu ,\kappa ,\omega ,ϵ,\delta )=G(\mu ,\kappa ,\omega ,ϵ,\delta ;\xi _1).$$
(30)
Assume that the boundary-value problem (4)-(6) has a solution which satisfies $`Q(\xi _1)`$$`=\mu `$. With a slight abuse of notation, let us denote such a solution by $`Q(\mu ,\kappa ,\omega ,ϵ,\delta ;\xi )`$. Clearly $`Q(\mu ,\kappa ,\omega ,ϵ,\delta ;\xi )`$ is defined only on a submanifold of the parameter space, but this will not be important in what follows. As we have seen in the introduction, we have
$$\lambda ^{1/\sigma +i\theta }Q(\mu ,\kappa ,\omega ,ϵ,\delta ,\lambda \xi )=Q(\lambda ^{1/\sigma +i\theta }\mu ,\lambda ^2\kappa ,\lambda ^2\omega ,ϵ,\delta ;\xi ),$$
for all $`\lambda >0`$ and $`\theta [0,2\pi )`$. Therefore we can work with normalized solutions, i. e. we assume that $`\omega =1`$ and that $`Q(0)`$ is real and non-negative. Assume $`\mu _0>0,\kappa _0>0`$ and suppose $`Q(\mu _0,\kappa _0,1,0,0;\xi )`$ exists.
We set
$$g(\mu ,\kappa ,ϵ,\delta )=\frac{G}{\xi }(\mu ,\kappa ,1,ϵ,\delta ;\xi _1)\frac{F}{\xi }(\beta (\mu ,\kappa ,1,ϵ,\delta ),\kappa ,1,ϵ,\delta ,\xi _1,\xi _1).$$
(31)
where $`\beta (\mu ,\kappa ,\omega ,ϵ,\delta )`$ is defined by (30). By our assumptions and by Theorem 1 the mapping $`g`$ is well defined and continuously differentiable in a set of the form $`(\mu _1,\mu _2)\times (\omega _1,\omega _2)\times [0,ϵ_1)\times (\delta _1,\delta _2)`$ containing the point $`(\mu _0,\kappa _0,0,0)`$. Since $`g(\mu _0,\kappa _0,0,0)=0`$, we see that the equation
$$g(\mu ,\kappa ,ϵ,\delta )=0$$
has solutions for small $`ϵ>0,\delta >0`$ if the following non-degeneracy condition is satisfied:
$$det\left(\begin{array}{cc}g_\mu ^1& g_\kappa ^1\\ g_\mu ^2& g_\kappa ^2\end{array}\right)0\text{at}(\mu _0,\kappa _0,0,0),$$
(32)
where $`g^1=\mathrm{Re}g`$, $`g^2=\mathrm{Im}g`$ and subscripts denote partial derivatives with respect to the corresponding variables.
Thus a non-trivial solution of (4)-(6) for $`ϵ=\delta =0`$ also gives a solution of (4)-(6) for $`ϵ,\delta >0`$ if (32) is satisfied. Based on our numerical calculations described in the next section, we conjecture that (32) is satisfied for $`\sigma >2/d`$.
## 3 Numerical results
Theorem 1 allows us to rewrite the boundary value problem (4)-(6) as a boundary value problem on a finite interval $`(0,\xi _1)`$ in the following way.
$`(1iϵ)(Q^{\prime \prime }+{\displaystyle \frac{d1}{\xi }}Q^{})+i\kappa \xi Q^{}+i{\displaystyle \frac{\kappa }{\sigma }}QQ+(1+i\delta )|Q|^{2\sigma }Q=0,`$ (33)
$`Q^{}(0)=0,`$ (34)
$`Q(\xi _1)=\beta ,`$ (35)
$`Q^{}(\xi _1)={\displaystyle \frac{F}{\xi }}(\beta ,\kappa ,1,ϵ,\delta ,\xi _1;\xi _1),`$ (36)
where the unknown quantities are $`Q`$, $`\beta `$ and $`\kappa `$. Of course, the problem (33)-(36) is equivalent to the equation $`g=0`$ in the previous section. If we approximate $`F`$ by the first few terms of its asymptotic expansion, the problem (33)-(36) can be solved numerically. In our numerical computations we investigated the dependence on $`\xi _1`$ and on the number of terms of the asymptotic expansion of $`F`$. It turned out that $`\xi _130`$ and the first term of the asymptotic expansion already worked very well. However, many of our computations were done with the first two terms of the asymptotic expansion (26). In the case $`ϵ=\delta =0`$, i.e. NLS, the value of $`(\mu _1,\kappa _1)`$ was computed in using completely different approach. The values presented in that paper are in an excellent agreement with our computations, see below.
The first term of the asymptotic expansion of $`F`$ is $`F\beta (\xi /\xi _1)^{1/\sigma i\omega /\kappa }`$. Using this approximation, we obtain from (33)-(36)
$`(1iϵ)(Q^{\prime \prime }+{\displaystyle \frac{d1}{\xi }}Q^{})+i\kappa \xi Q^{}+i{\displaystyle \frac{\kappa }{\sigma }}QQ+(1+i\delta )|Q|^{2\sigma }Q=0,`$ (37)
$`Q^{}(0)=0,`$ (38)
$`\xi _1Q^{}(\xi _1)+\left({\displaystyle \frac{1}{\sigma }}+i{\displaystyle \frac{1}{\kappa }}\right)Q(\xi _1)=0.`$ (39)
Higher order approximations can be derived in a similar way. Note that in the formulation of (37) we already fixed the normalization $`\omega =1`$, so that the unknowns are $`Q`$ and $`\kappa `$. The boundary condition (39) is also closely related to the boundary condition used in , for simulations based on solving time dependent problem in the PDE (1)-(2).
There are essentially two approaches to the numerical solution of the boundary-value problem (37)-(39). One can use collocation methods to approximate the boundary-value problem and then to apply Newton’s method to the discretization of the non-linear operator which defines the equation (37). Implementation of this strategy requires further changes in the formulation since the linearization of the non-linear operator always has zero in its spectrum due to the $`S^1`$-equivariance of the equation. Therefore a shooting method was easier to implement and it also proved to be sufficiently accurate. Because of well-known sensitivity of shooting methods to problem parameters we performed computations in different normalizations: with $`Q(0)=1`$ fixed and parameters $`(\kappa ,\omega )`$ as the unknowns as well as with $`\omega =1`$ fixed and parameters $`(\kappa ,\mu )`$ as the unknowns. Moreover we compared both backward and forward shooting methods on the interval $`(0,\xi _1)`$. All computed solutions turned out to be in a very good agreement.
Here we describe only the shooting method using the normalization $`\omega =1`$, in which we calculate $`\xi _1Q^{}(\xi _1)+(1/\sigma +i/\kappa )Q(\xi _1)`$ as a function of the parameters $`\mu =Q(0),\kappa ,ϵ`$, and $`\delta `$, i.e., we solve the equation
$$f(\mu ,\kappa ,ϵ,\delta )\xi _1Q^{}(\xi _1)+\left(1/\sigma +i/\kappa \right)Q(\xi _1)=0.$$
(40)
The integration of the underlying ODE must be done with sufficient accuracy. We compared various ODE solvers. A variable-order, variable-step Adams method as implemented, for example, in the NAG library proved to be sufficiently accurate in most of the calculations. To locate the initial values for each branch we inspected the two dimensional subspace of the parameter space given by $`(\mu ,\kappa ,0,0)`$ for $`d=1`$ and $`\sigma =2.3`$ and computed the degree of the function $`f`$ restricted to that subspace along various curves. The other solutions were then calculated by continuation. This was done with the help of bifurcation analysis package developed as a part of . The implementation of the path-following procedure with a Newton corrector step can be done efficiently as the linearization along a solution is evaluated at the same integration step as the solution.
As we described in the introduction, it appears that the equation $`f(\mu ,\kappa ,0,0)`$ has countably many solutions $`(\kappa _j,\mu _j)`$ in the region $`\kappa >0`$, $`\mu >0`$. The corresponding profiles $`|Q_j|`$, when extended to $`(\mathrm{},\mathrm{})`$ as even functions, have exactly $`j`$ local maxima.
###### Remark 3.1.
As mentioned in the beginning of this section the solutions $`(\kappa _j,\mu _j)`$ may slightly depend on the value $`\xi _1`$ and the approximation of $`F`$. For solutions described here we tested the dependence of the results on $`\xi _1`$ for $`\xi _1[20,100]`$, and also on the approximation of $`F`$ by taking either one or two terms in the asymptotics expansion of $`F`$. We also directly compared our numerical solutions in intervals of the form $`(\eta _1,\xi _1)`$ with the explicit formulae given by one or two terms of the asymptotic expansion of the solution in $`(\eta _1,\xi _1)`$ for various values of $`\eta _1`$. All these tests indicated a good convergence of our approximations. Based on these tests, we estimate that the error in the values of the “roots” $`(\kappa _j,\mu _j)`$ is of the order $`10^3`$ or better.
Case I ($`d=1`$, $`\sigma =2.3`$, $`\delta =0`$): Results for the one-dimensional case are tabulated in Table 1.
The continuation of solutions parameterized by $`ϵ`$ is depicted in Figure 1. To give the reader a good idea about the form of the solutions we plot profiles $`|Q(\xi )|`$ at a few points on each branch (at the point $`ϵ=0`$ (solution to NLS), at a point on the upper part of the branch and at another point on the lower part of the branch). We used both normalizations, $`Q(0)=1`$ in Figure 23, and $`\omega =1`$ in Figure 45. Normalization $`Q(0)=1`$ is convenient for computing solutions along the odd branches as these solutions exhibit a maximum at the origin which is also present in the solution for $`ϵ0`$ on the lower part of the branch.
We now describe the behavior of solutions along the branches. We observe that as $`ϵ`$ returns to zero after passing through the turning point, we have $`\kappa _j(ϵ)0`$.
For $`j=1`$ the solution $`Q_1^ϵ(\xi )`$ approaches a specific solution, usually called the ground state, of the equation
$$u^{\prime \prime }+u=|u|^{2\sigma }u.$$
(41)
For $`j>1`$ the profile $`|Q_j^ϵ(\xi )|`$ seems to separate into $`j`$ copies of the ground-state solution which are moving to infinity. For $`j`$ even all of them “escape” to infinity, while for $`j`$ odd one will stay at the origin and the rest will move to the infinity as $`ϵ0`$ and $`\kappa _j(ϵ)0`$.
Case II ($`d=3`$, $`\sigma =1`$, $`\delta =0`$): Results for the three-dimensional case are tabulated in Table 2.
Some data for NLS ($`ϵ=\delta =0`$) and the “basic solution” (corresponding to the beginning of our first branch) are available in the literature and can be used to estimate accuracy of our calculations. One can see a very good agreement of our solution for $`j=1`$ with values $`\kappa =\mathtt{0.917}`$ and $`\mu Q(0)=\mathtt{1.885}`$ obtained in from simulations that used the dynamical rescaling method applied to the initial value problem (1)-(2).
The behavior of solutions $`Q_j^ϵ`$ when $`ϵ`$ returns to zero is qualitatively similar to the one-dimensional case, but there are also some new interesting features.
First, the maximal possible $`ϵ_j^{}`$ is not attained on the first branch ($`j=1`$). Second, the behavior as $`ϵ0`$ and $`\kappa _j(ϵ)0`$ is different for even and odd values $`j`$. For $`j=2k1`$ (odd) the solutions $`Q_j^ϵ(\xi )`$ converge (along the $`j`$th-branch) to the $`k`$-th radial solution of the problem
$`\mathrm{\Delta }u+u`$ $`=`$ $`|u|^{2\sigma }u\text{in}^3,`$ (42)
$`u`$ $``$ $`0\text{as}|x|\mathrm{}.`$ (43)
For $`j=2k`$ (even), one observes a similar behavior as in the one-dimensional case.
###### Remark 3.2.
From the regularity theory for CGL (see ) one expects that the value of $`ϵ`$ cannot cross a certain finite value $`\stackrel{~}{ϵ}`$, as we follow solutions along any branch. What we see in the calculations is a turning point on each branch. In Table 1 and Table 2 we list for each of the calculated branches the parameter $`ϵ^{}`$ which is the maximal value of the perturbation parameter $`ϵ`$ reached along the branch.
## 4 Stability of singular solutions
A natural question regarding the stability of the self-similar singularities constructed in the previous section is for example the following: What will be the behavior of solutions to (1) and (2) if $`u_0`$ is a slight perturbation of a solution $`Q`$ of the profile equation (4)-(6) and $`f`$ is small? (One can think for example of taking $`u_0=\phi Q+u_1`$, where $`\phi `$ is a compactly supported smooth cut-off function which is identically one in a large ball, $`Q`$ is a non-trivial solution of (4)-(6), $`u_1`$ is small and compactly supported, and at the same time taking $`f`$ which is small and compactly supported in $`x`$.)
Here we present an approach to this problem which uses the method of dynamical rescaling (see for example ), which seems to be natural in this context. We emphasize that we will not obtain fully rigorous analytical results which would completely answer the question raised above. Our goal is to present some preliminary calculations which seem to be adequate to the issue of stability in the context of numerical simulations.
We briefly recall the main idea of the method of dynamical rescaling. We consider a solution $`u:^d\times (0,T)`$ of the complex Ginzburg-Landau equation
$$i\frac{u}{t}+(1iϵ)\mathrm{\Delta }u+(1+i\delta )|u|^{2\sigma }u=0.$$
(44)
Motivated by the scaling invariance of solutions
$$u(x,t)\lambda ^{1/\sigma }u(\lambda x,\lambda ^2t)$$
of (44), we write the solution $`u`$ as
$$u(x,t)=L^{1/\sigma }(t)v(L^1(t)x,\tau ),$$
where $`L(t)>0`$ is to be chosen and $`d\tau =L^2(t)dt`$. From (44) we obtain
$$i\frac{v}{\tau }+i\kappa (\tau )\left(\xi \frac{v}{\xi }+\frac{1}{\sigma }v\right)+(1iϵ)\mathrm{\Delta }v+(1+i\delta )|v|^{2\sigma }v=0,$$
(45)
where $`\kappa (\tau )=L^1(\tau )dL(\tau )/d\tau =L(t)dL(t)/dt`$.
We now chose $`L(t)`$ so that, roughly speaking, the typical length-scale over which $`v`$ oscillates is $`1`$. If $`u`$ develops a singularity at $`t=T`$ (and it is regular in $`(0,T_1)`$ for each $`T_1<T`$), then $`L(t)>0`$ in $`(0,T)`$ and $`L(t)0`$ as $`tT`$.
One way to control oscillations of $`v`$ is to impose a condition $`𝒥(v(.,\tau ))=1`$, where $`𝒥`$ is a suitable functional controlling regularity of solutions to (45). One is then led to the following system
$`i{\displaystyle \frac{v}{\tau }}+i\kappa (\tau )\left(\xi {\displaystyle \frac{v}{\xi }}+{\displaystyle \frac{1}{\sigma }}v\right)+(1iϵ)\mathrm{\Delta }v+(1+i\delta )|v|^{2\sigma }v=0,`$ (46)
$`𝒥(v(.,\tau ))=1,`$ (47)
where the unknowns are $`v`$ and $`\kappa (\tau )`$. Various choices of the functional $`𝒥`$ have been used for numerical calculations. For example, in the functional
$$𝒥_0(v)=_^d|v|^2𝑑x$$
was used. Although it does not directly follow from the known regularity theory that $`𝒥_0`$ controls the regularity of solutions to (45) in the case $`ϵ=\delta =0`$ this choice turned out to work satisfactorily in the numerical computations.
In the following analysis we assume that the functional $`𝒥`$ is invariant under the action of the symmetry group $`S^1`$, i.e., $`𝒥(e^{i\theta }v)=𝒥(v)`$ for all $`\theta [0,2\pi )`$. Self-similar singularities of the form (3) correspond to solutions of (46) of the form $`v(\xi ,\tau )=e^{i\omega \tau }Q(\xi )`$, i.e., to $`S^1`$-orbits of solutions to the problem (4)-(6). We will consider the linearized stability of these orbits, which seems to be the simplest natural notion of stability in our context. Suppose $`(v_0(\xi ,\tau ),\kappa _0(\tau ))`$ is a solution of (46) such that $`v_0(\xi ,\tau )=e^{i\omega _0\tau }Q(\xi )`$ and $`\kappa _0(\tau )=\kappa _0=\mathrm{const}`$. We consider a perturbed solution $`(v(\xi ,\tau ),\kappa (\tau ))`$ in the form
$`v(\xi ,\tau )=e^{i\omega _0\tau }(Q(\xi )+w(\xi ,\tau ))`$
$`\kappa (\tau )=\kappa _0+\varkappa (\tau ),`$
where $`w`$ and $`\varkappa `$ are infinitesimally small. A simple calculation gives
$`{\displaystyle \frac{w}{\tau }}=w+\varkappa (\tau )\left(\xi Q^{}+{\displaystyle \frac{1}{\sigma }}Q\right),`$ (48)
$`𝒥^{}(Q),w=0,`$ (49)
where the operator $``$ is defined as the linearization of the equation (45) along the solution $`(v_0,\kappa _0)`$, i.e.,
$$w=(i+ϵ)\mathrm{\Delta }w\kappa _0\left(\xi \frac{w}{\xi }+\frac{1}{\sigma }w\right)i\omega _0w+i(1+\sigma )|Q|^{2\sigma }w+i\sigma |Q|^{2\sigma 2}Q^2\overline{w}$$
and $`𝒥^{}`$ denotes the derivative of $`𝒥`$. (Note that the operator $``$ is not complex linear and therefore it is natural to carry out the analysis in the real representation.)
There are two eigenfunctions of $``$ that can be formally derived directly from the invariance of the profile equation (4) under the scaling symmetries
$`(Q(\xi ),\kappa ,\omega )`$ $``$ $`(e^{i\theta }Q(\xi ),\kappa ,\omega ),`$ (50)
$`(Q(\xi ),\kappa ,\omega )`$ $``$ $`(\lambda ^{1/\sigma +i\omega /\kappa }Q(\lambda \xi ),\lambda ^2\kappa ,\lambda ^2\omega ).`$ (51)
The invariance under the $`S^1`$-symmetry (50) implies existence of an eigenfunction $`Y_1=iQ`$ with the eigenvalue zero, while the symmetry under (51) leads (formally) to an eigenfunction $`Y_2=\xi Q^{}+(1/\sigma +i\omega _0/\kappa _0)Q`$ with the eigenvalue $`2\kappa _0`$. We note that $`Y_2(\xi )=\mathrm{O}(\xi ^{1/\sigma 2})`$ as $`\xi \mathrm{}`$, due to (26) and (27). We have
$`Y_1=0,`$ (52)
$`Y_2=2\kappa _0Y_2.`$ (53)
We rewrite (48) as
$`{\displaystyle \frac{w}{\tau }}=w+\varkappa (\tau )\left(Y_2{\displaystyle \frac{\omega _0}{\kappa _0}}Y_1\right)`$ (54)
$`𝒥^{}(Q),w=0,`$ (55)
where the unknowns are $`w(\xi ,\tau )`$ and $`\varkappa (\tau )`$.
Proceeding further with our formal reasoning, we view (54) as a dynamical system with one linear constraint in a suitable linear space $`𝒲`$ of functions on $`^d`$. In this formal analysis we will not try to specify $`𝒲`$. For a rigorous analysis it would be natural to try to find a suitable Banach space containing all smooth compactly supported functions, together with the functions $`Y_1`$ and $`Y_2`$. Suppose that the space $`𝒲`$ can be decomposed into a direct sum $`𝒲=Y_1Y_2𝒵`$ where $`𝒵𝒲`$ is invariant under $``$. In the finite-dimensional situation, a sufficient condition for this would be that the eigenvalues $`0`$ and $`2\kappa _0`$ are simple. Assuming that such a decomposition exists, we write
$$w(\xi ,\tau )=w_1(\tau )Y_1(\xi )+w_2(\tau )Y_2(\xi )+Z(\xi ,\tau ),$$
(56)
where $`w_1`$, $`w_2`$ are scalar functions and $`Z(.,\tau )𝒵`$.
After substituting (56) into (54) and using $`𝒥^{}(Q),Y_1=0`$ (a consequence of the invariance of $`𝒥`$), we obtain
$`{\displaystyle \frac{Z}{\tau }}`$ $`=`$ $`Z`$ (57)
$`{\displaystyle \frac{w_1}{\tau }}`$ $`=`$ $`\varkappa (\tau ){\displaystyle \frac{\omega _0}{\kappa _0}}`$ (58)
$`{\displaystyle \frac{w_2}{\tau }}`$ $`=`$ $`2\kappa _0w_2+\varkappa (\tau )`$ (59)
$`w_2`$ $`=`$ $`{\displaystyle \frac{𝒥^{}(Q),Z}{𝒥^{}(Q),Y_2}}`$ (60)
where we assume that $`𝒥^{}(Q),Y_20`$. This is a completely natural assumption in the context of (46). From (57) - (60) one can see that in a finite-dimensional situation and under the assumptions stated above, the following conditions would be equivalent:
If $`(w,\varkappa )`$ is a solution of (48), then $`w(\xi ,\tau )`$ approaches exponentially $`aY_1(\xi )`$ for some $`a`$ and $`\varkappa (\tau )`$ approaches exponentially zero.
The spectrum of $`|_𝒵`$ belongs to the set $`\{z|\mathrm{Re}z<0\}`$.
In a finite-dimensional situation and under the assumption that the eigenvalues $`0`$ and $`2\kappa _0`$ are simple, condition (ii) would be equivalent to the condition that all the spectrum of $``$ except $`0`$ and $`2\kappa _0`$ lies in $`\{z|\mathrm{Re}z<0\}`$.
The condition (i) is exactly the linearized orbital stability of the orbit $`(e^{i\omega _0\tau }Q,\kappa _0)`$ for the system (46), (48).
It seems to be a non-trivial problem to put the above formal analysis on a rigorous basis in the infinite-dimensional setting. However, the formal analysis strongly suggests that in a finite dimensional situation which arises in numerical approximations of (46) the stability of the solutions $`(e^{i\omega _0\tau }Q(\xi ),\kappa _0)`$ should be governed by the spectrum of an appropriate approximation of the operator $``$.
We remark that our stability analysis is independent of $`𝒥`$, except for the natural assumption $`𝒥^{}(Q),Y_20`$.
We numerically calculated the approximation of eigenvalues for a discrete approximation of $``$ in the space of radial functions along the branches of solutions parameterized by $`ϵ`$. The original system (46) was truncated to a finite interval $`(0,\xi _1)`$ by imposing the time-dependent boundary condition at $`\xi =\xi _1`$
$$\frac{v}{\tau }(\xi _1,\tau )+\kappa (\tau )\xi _1\frac{v}{\xi }(\xi _1,\tau )+\frac{\kappa (\tau )}{\sigma }v(\xi _1,\tau )=0,$$
(61)
which is a time-dependent equivalent of the boundary condition derived from the asymptotic expansion in Section 2. The condition (61) was also used in the numerical simulations in and . The computations were carried out for $`\xi _1=30`$.
The accuracy of our numerical approximation can be checked indirectly by comparing the predicted eigenvalues $`0`$ and $`2\kappa _0`$ with the corresponding eigenvalues we obtained from our numerical calculations. We saw a very good agreement, in most cases the error was of the order $`10^4`$. Both $`0`$ and $`2\kappa _0`$ appeared simple, except in some natural degenerate cases when other eigenvalues were crossing them as we moved along branches.
Case I ($`d=1`$, $`\sigma =2.3`$, $`\delta =0`$): The calculations confirmed what one intuitively expects: The solutions on the upper part of the branch $`j=1`$ are stable and all other solutions are unstable. In Figure 1 we used a solid line for the stable parts of the curves and a dashed line for the unstable parts.
Case II ($`d=3`$, $`\sigma =1`$, $`\delta =0`$): The situation is similar with one notable exception. In our computations we detected stable solutions also on the lower part of the branch $`j=2`$. Accuracy of our approximation did not allow us to decide whether all solutions on the lower part of the branch $`j=2`$ for $`0<ϵ<ϵ^{}`$ are stable, since as $`ϵ`$ approached zero we observed some eigenvalues very close to the imaginary axis. In Figure 6 we used a solid line for the stable parts of the curves and a dashed line for the unstable parts.
We emphasize again that the calculations we carried out only deal with stability in the space of radial functions. Based on computations in , it appears that the solutions on the upper part of the branch $`j=1`$ are also stable with respect to perturbations that break the radial symmetry.
It is clear that further work is required to fully clarify issues concerning the stability of solutions described in this paper.
## 5 Appendix
In this section we give a full proof of Theorem 1.
We fix $`r_1>0`$ and $`\vartheta >0`$. Values of these parameters will be chosen later. We denote by $`𝒳_\vartheta `$ the Banach space of continuous functions $`u:[\xi _1,\mathrm{})`$ for which the norm
$$u_\vartheta =\underset{\xi \xi _1}{sup}|\xi |^{1/\sigma \vartheta }|u(\xi )|$$
is finite. In this proof we use the following notation
$`\stackrel{~}{\mathrm{\Lambda }}=[\kappa _1,\kappa _2]\times [\omega _1,\omega _2]\times [0,ϵ_1]\times [\delta _1,\delta _2]`$
$`\mathrm{\Lambda }=\{\gamma ||\gamma |r_1\}\times \stackrel{~}{\mathrm{\Lambda }}.`$
A point in $`\stackrel{~}{\mathrm{\Lambda }}`$ is denoted by $`\stackrel{~}{\lambda }=(\lambda _1,\lambda _2,\lambda _3,\lambda _4)(\kappa ,\omega ,ϵ,\delta )`$ and similarly $`\lambda \mathrm{\Lambda }`$ represents $`(\lambda _0,\lambda _1,\lambda _2,\lambda _3,\lambda _4)(\gamma ,\kappa ,\omega ,ϵ,\delta )`$. The functions $`P`$, $`E`$, $`W`$, and $`K`$ introduced in Section 2 are written as $`P=P(\stackrel{~}{\lambda },\xi )`$, $`E=E(\stackrel{~}{\lambda },\xi )`$, $`W=W(\stackrel{~}{\lambda },\xi ,\eta )`$, $`K=K(\stackrel{~}{\lambda },\xi )`$. With a slight abuse of notation we sometimes write also $`P=P(\lambda ,\xi )`$ even if the function does not depend on $`\lambda _0`$.
###### Lemma 2.
There exists $`C>0`$ such that for $`\stackrel{~}{\lambda }\stackrel{~}{\mathrm{\Lambda }}`$, $`\xi ,\eta \xi _0`$ we have
$$|K(\stackrel{~}{\lambda },\xi ,\eta )|\{\begin{array}{cc}C\xi ^{1/\sigma }\eta ^{1/\sigma 1}\hfill & \text{ for }\xi _0\eta \xi ,\hfill \\ C\xi ^{d+1/\sigma }\eta ^{11/\sigma +d}\hfill & \text{ for }\xi _0\xi \eta .\hfill \end{array}$$
(62)
Proof: This statement is an easy consequence of the definition of $`K`$ and of the asymptotic expansion for $`U`$ as $`z\mathrm{}`$ in the sector $`\pi /2\mathrm{arg}z\pi `$. An important point in connection with (14) is that the constant in the remainder $`\mathrm{O}(|z|^n)`$ can be taken same when the parameter $`a`$ runs through a compact subset of $`\{z|\mathrm{Re}z>0\}`$. (This can be seen for example from (13).)$`\mathrm{}`$
###### Lemma 3.
Assume $`(2\sigma +1)\vartheta <2+2/\sigma d`$. Then the formula
$$T(\lambda ,u)(\xi )=\gamma P(\stackrel{~}{\lambda },\xi )_{\xi _1}^{\mathrm{}}(1+i\delta )K(\stackrel{~}{\lambda },\xi ,\eta )|u(\eta )|^{2\sigma }u(\eta )𝑑\eta $$
(63)
defines a continuous mapping $`T:\mathrm{\Lambda }\times 𝒳_\vartheta 𝒳_\vartheta `$.
Moreover, there exists $`C>0`$ such that
$$T(\lambda ,u)_\vartheta C|\gamma |\xi _1^\vartheta +C\xi _1^{2+2\sigma \vartheta }u_\vartheta ^{2\sigma +1}$$
(64)
and
$$T(\lambda ,u)T(\lambda ,v)_\vartheta C\xi _1^{2+2\sigma \vartheta }uv_\vartheta \left(u_\vartheta ^{2\sigma }+v_\vartheta ^{2\sigma }\right)$$
(65)
for all $`\lambda \mathrm{\Lambda }`$ and all $`u,v𝒳_\vartheta `$.
Proof: The convergence of the integrals in (63) and (64) follows easily from Lemma 2. To get the estimate (65) we use Lemma 2 together with the elementary inequality
$$\left||z_1|^{2\sigma }z_1|z_2|^{2\sigma }z_2\right|M|z_1z_2|\left(|z_1|^{2\sigma }+|z_2|^{2\sigma }\right),$$
which holds for all $`z_1,z_2`$ if $`M>0`$ is properly chosen.
The continuity of the mapping $`T`$ can be proved by suitably splitting the integral over $`(\xi _1,\mathrm{})`$ into two integrals, one over $`(\xi _1,\xi _2)`$ and the other one over $`(\xi _2,\mathrm{})`$. We can then estimate the integral over $`(\xi _2,\mathrm{})`$ by using Lemma 2. $`\mathrm{}`$
The estimates (64) and (65) show that $`T(\lambda ,.)`$ is a contraction of the ball $`_\rho =\{u𝒳_\vartheta |u_\vartheta \rho \}`$ into itself if $`\rho `$ is such that
$`Cr_1\xi _1^\vartheta +C\xi _1^{2+2\sigma \vartheta }\rho ^{2\sigma +1}\rho ,`$ (66)
$`2C\rho ^{2\sigma }\xi _1^{2+2\sigma \vartheta }<1.`$ (67)
###### Proposition 4.
Assume $`(2\sigma +1)\vartheta <2+2/\sigma d`$. If (66) and (67) are satisfied, then for each $`\lambda \mathrm{\Lambda }`$ the mapping $`T(\lambda ,.)`$ has a unique fixed point $`u(\lambda )=u(\lambda ,\xi )`$ in $`_\rho =\{u𝒳_\vartheta |u_\vartheta \rho \}`$.
Moreover, the mapping $`\lambda u(\lambda )`$ is a continuous mapping from $`\mathrm{\Lambda }`$ to $`𝒳_\vartheta `$.
Proof: The proposition follows directly from Lemma 3 and the Banach Fixed Point Theorem. $`\mathrm{}`$
Next we need to establish that the solutions $`u(\lambda ,\xi )`$ have the required regularity properties at the infinity.
###### Lemma 5.
The function $`u(\lambda )=u(\lambda ,\xi )`$ from Proposition 4 exhibits the following behavior as $`\xi \mathrm{}`$:
$$u(\lambda ,\xi )=\xi ^{1/\sigma i\omega /\kappa }\left(\underset{j=0}{\overset{n}{}}a_j(\lambda )\xi ^{2j}+\mathrm{O}(\xi ^{2n2})\right),$$
(68)
where $`a_j(\lambda )`$ are continuous functions of $`\lambda `$ and the constant in $`\mathrm{O}(\xi ^{2n2})`$ is independent of $`\lambda \mathrm{\Lambda }`$.
Moreover, the expansion (68) can be differentiated in the following sense:
$$\frac{^ku}{\xi ^k}(\lambda ,\xi )=\frac{^k}{\xi ^k}\left(\xi ^{1/\sigma i\omega /\kappa }\underset{j=0}{\overset{n}{}}a_j(\lambda )\xi ^{2j}\right)+\mathrm{O}(\xi ^{1/\sigma 2n2k}),$$
(69)
where the constant in $`\mathrm{O}(\xi ^{1/\sigma 2n2k})`$ is independent of $`\lambda \mathrm{\Lambda }`$.
Proof: The proof is similar to bootsrapping arguments used in the regularity theory. From (63) and the assumption $`(2\sigma +1)\vartheta <2+2/\sigma d`$ we get immediately
$$u(\lambda ,\xi )=\mathrm{O}(\xi ^{1/\sigma })\text{as}\xi \mathrm{}.$$
(70)
We now rewrite the equation $`u=Tu`$ as
$$u(\xi )=(\gamma +\gamma _1)P(\xi )\frac{1+i\delta }{1iϵ}P(\xi )_\xi ^{\mathrm{}}EW^1|u|^{2\sigma }u+\frac{1+i\delta }{1iϵ}E(\xi )_\xi ^{\mathrm{}}PW^1|u|^{2\sigma }u,$$
(71)
where
$$\gamma _1=\frac{1+i\delta }{1iϵ}_{\xi _1}^{\mathrm{}}EW^1|u|^{2\sigma }u.$$
Using (70) and (71), we obtain (for a suitable $`a_0`$)
$$u(\xi )=\xi ^{1/\sigma i\omega /\kappa }\left(a_0+\mathrm{O}(\xi ^2)\right)$$
(72)
We now repeat the procedure and use (71) with (72) instead of (70). After integrating by parts in integrals of the form
$$_\xi ^{\mathrm{}}\eta ^\alpha e^{\nu \eta ^2/2}=\frac{1}{\nu }_\xi ^{\mathrm{}}\eta ^{\alpha 1}\frac{}{\eta }e^{\nu \eta ^2/2}$$
which come up in the calculation, we get easily
$$u(\xi )=\xi ^{1/\sigma i\omega /\kappa }\left(a_0+a_1\xi ^2+\mathrm{O}(\xi ^4)\right).$$
Repeating this procedure, we get (68) by induction. The expansion (69) can be obtained in a similar way, after we differentiate (71). $`\mathrm{}`$
###### Lemma 6.
The mapping $`\lambda u(\lambda )`$ from Proposition 4 is continuously differentiable as a mapping from $`\mathrm{\Lambda }`$ to $`𝒳_\vartheta `$
Proof: We define $`\varphi :\mathrm{\Lambda }\times 𝒳_\vartheta 𝒳_\vartheta `$ as $`\varphi (\lambda ,v)=vT(\lambda ,v)`$ and we denote $`=\{(\lambda ,u(\lambda )),\lambda \mathrm{\Lambda }\}\mathrm{\Lambda }\times 𝒳_\vartheta `$. By using similar estimates as those used in the proof of Lemma 3, one can easily see that $`\varphi `$ is differentiable with respect to $`v`$ and that the (partial) derivative $`\mathrm{D}_v\varphi `$ is continuous as a map from $`\mathrm{\Lambda }\times 𝒳_\vartheta `$ into the space $`(𝒳_\vartheta ,𝒳_\vartheta )`$ of bounded linear operators on $`𝒳_\vartheta `$. Moreover, a calculation similar to the one leading to (65) gives
$$I\mathrm{D}_v\varphi 2C\xi _1^{2+2\sigma \vartheta }\rho ^{2\sigma }$$
(73)
in $`\mathrm{\Lambda }\times _\rho `$. Therefore, since we assume (67), $`\mathrm{D}_v\varphi `$ is invertible for all $`(\lambda ,v)\mathrm{\Lambda }\times _\rho `$. By Lemma 7 below, $`\varphi `$ is also differentiable with respect to $`\lambda `$ at each point of $``$ and the partial derivative $`\mathrm{D}_\lambda \varphi (\lambda ,v)`$ is continuous in $``$. Now we can conclude with the same arguments as in the Implicit Function Theorem to show that $`u`$ is differentiable and
$$\frac{u}{\lambda }(\lambda )=\left[\mathrm{D}_v\varphi (\lambda ,u(\lambda ))\right]^1\left[\mathrm{D}_\lambda \varphi (\lambda ,u(\lambda ))\right].$$
(74)
$`\mathrm{}`$
Lemma 6 also implies the differentiability of the function $`\lambda \frac{u}{\xi }(\lambda ,\xi _1)`$, since the partial derivative of $`u`$ at $`\xi _1`$ can be expressed in terms of $`u`$ due to the fact that $`u`$ solves a differential equation. For example, one can take a smooth function $`\phi :[\xi _1,\mathrm{})`$ with $`\phi (\xi _1)=1`$ and $`\phi (\xi )0`$, for $`\xi >\xi _1+1`$ and write $`u(\xi _1)=_{\xi _1}^{\mathrm{}}(u\phi )^{\prime \prime }𝑑\xi `$. Then we express $`u^{\prime \prime }`$ from the differential equation for the profile (4) and integrate by parts to obtain
$$\frac{u}{\xi }(\xi _1)=c_1u(\xi _1)+_{\xi _1}^{\mathrm{}}(u(\xi )\psi (\xi )+|u(\xi )|^{2\sigma }u(\xi )\stackrel{~}{\psi }(\xi ))𝑑\xi ,$$
where $`\psi `$, $`\stackrel{~}{\psi }`$ are supported on a finite interval and $`c_1`$ is a suitable constant. $`\mathrm{}`$
###### Lemma 7.
The mapping $`\varphi :\mathrm{\Lambda }\times 𝒳_\vartheta 𝒳_\vartheta `$ defined as $`\varphi (\lambda ,v)=vT(\lambda ,v)`$ is differentiable with respect to $`\lambda `$ at each point $`=\{(\lambda ,u(\lambda )),\lambda \mathrm{\Lambda }\}\mathrm{\Lambda }\times 𝒳_\vartheta `$. Moreover, the partial derivative $`\mathrm{D}_\lambda \varphi `$ is continuous in $``$.
Proof: Using (11) we note that
$$\frac{U}{a}(a,b,z)=z^a\mathrm{log}z\left(1+\mathrm{O}(|z|^1)\right),$$
(75)
as $`z\mathrm{}`$ in the sector $`\pi /2\mathrm{arg}z<\pi `$ and $`a`$ belongs to a compact subset of $`\{\zeta |\mathrm{Re}\zeta >0\}`$. Let us carry out the proof for the partial derivative $`\varphi /\lambda _1`$, for example. The other partial derivatives can be handled in a similar way. (If fact, for $`\lambda _0`$,$`\lambda _2`$ and $`\lambda _4`$ the proof is much easier.) When $`ϵ>0`$ one sees easily (using the rapid decay of $`W^1(\eta )`$ in that case) that $`\varphi `$ is differentiable in $`\lambda _1`$ and that
$$\frac{\varphi }{\lambda _1}(\lambda ,v)(\xi )=\gamma \frac{P}{\lambda _1}(\lambda ,\xi )+_{\xi _1}^{\mathrm{}}\frac{K}{\lambda _1}(\lambda ,\xi ,\eta )|v(\eta )|^{2\sigma }v(\eta )𝑑\eta .$$
(76)
It follows from (75) that, for any $`\vartheta >0`$, the first term on the right-hand side is continuous as a mapping from $`\mathrm{\Lambda }`$ to $`𝒳_\vartheta `$. For any $`ϵ_2>0`$, the integral on the right-hand side has the required continuity properties in the region $`ϵ_2ϵϵ_1`$. due to the rapid decay of $`W^1(\eta )`$. The problem is to obtain estimates which are uniform for $`0<ϵ<ϵ_1`$. A closer inspection of the integral on the right-hand side of (76) reveals that the only terms for which the required continuity is not obvious come from differentiating the exponential $`e^{i\kappa \eta ^2/2(1iϵ)}`$ that appears in $`W^1(\eta )`$ This term can be handled with using integration by parts as follows. Recalling that $`\lambda _1=\kappa `$, we must estimate an integral of the form
$$_\xi ^{\mathrm{}}P(\lambda ,\eta )\eta ^{d1}\frac{}{\kappa }\left(e^{\frac{i\kappa }{1iϵ}\frac{\eta ^2}{2}}\right)|v(\eta )|^{2\sigma }v(\eta )𝑑\eta .$$
(77)
To estimate this integral we write
$$\frac{}{\kappa }\left(e^{\frac{i\kappa }{1iϵ}\frac{\eta ^2}{2}}\right)=\frac{\eta }{2\kappa }\frac{}{\eta }\left(e^{\frac{i\kappa }{1iϵ}\frac{\eta ^2}{2}}\right)$$
and integrate by parts. This eliminates the power $`\eta ^2`$ obtained after differentiating in (77), and the required estimate follows. $`\mathrm{}`$
To complete the proof of Theorem 1 we need to show that the coordinate $`\gamma =\lambda _0`$ can be replaced by $`\beta =u(\lambda ,\xi _1)`$, if $`r_1`$ and $`\rho `$ are chosen properly. We have
$$\frac{u}{\lambda _0}(\lambda ,\xi _1)=P(\lambda ,\xi _1)+_{\xi _1}^{\mathrm{}}K(\stackrel{~}{\lambda },\xi _1,\eta )\frac{}{\lambda _0}|u|^{2\sigma }u(\lambda ,\eta )𝑑\eta .$$
(78)
(Since $`\lambda _0`$ is complex, we interpret this, with a slight abuse of notation, as an equation between real $`2\times 2`$ matrices. Another possibility would be to interpret it literally and do a similar calculation for $`/\overline{\lambda }_0`$.) We carry out the differentiation in the integral on the right-hand side and use (74), which gives
$$\frac{u}{\lambda _0}(\lambda )=\left[\mathrm{D}_v\varphi (\lambda ,u(\lambda ))\right]^1P.$$
(79)
We note from (73) that by a suitable choice of the constant $`\rho `$ (to be specified below) we can achieve that
$$[\mathrm{D}_v\varphi (\lambda ,u(\lambda ))]^12,$$
(80)
in $`\mathrm{\Lambda }`$ where the norm is taken in the space of linear operators on the space $`𝒳_0`$. Using (78) - (80), we see easily, by a similar calculation as in the proof of (65), that by a suitable choice of $`\rho `$ we can achieve that $`\frac{u}{\lambda _0}(\lambda ,\xi _1)/P(\lambda ,\xi _1)`$ is close to the identity uniformly in $`\mathrm{\Lambda }`$. A suitable choice for the constants $`\rho `$ and $`r_1`$ is, for example, as follows
$$\rho =\alpha \xi _1^{1/\sigma \vartheta },r_1=\alpha ^{2\sigma +1}\xi _1^{1/\sigma },$$
where $`\alpha `$ is sufficiently small. The proof of Theorem 1 can now be easily completed. $`\mathrm{}`$
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# Some linear Jacobi structures on vector bundles
## 1 Jacobi manifolds and Lie algebroids
Let $`M`$ be a differentiable manifold of dimension $`n`$. We will denote by $`C^{\mathrm{}}(M,\text{I}\text{R})`$ the algebra of $`C^{\mathrm{}}`$ real-valued functions on $`M`$, by $`\mathrm{\Omega }^1(M)`$ the space of 1-forms, by $`X(M)`$ the Lie algebra of vector fields and by $`[,]`$ the Lie bracket of vector fields.
A Jacobi structure on $`M`$ is a pair $`(\mathrm{\Lambda },E)`$, where $`\mathrm{\Lambda }`$ is a 2-vector and $`E`$ is a vector field on $`M`$ satisfying the following properties:
$$[\mathrm{\Lambda },\mathrm{\Lambda }]_{SN}=2E\mathrm{\Lambda },[E,\mathrm{\Lambda }]_{SN}=0.$$
(1)
Here $`[,]_{SN}`$ denotes the Schouten-Nijenhuis bracket (). The manifold $`M`$ endowed with a Jacobi structure is called a Jacobi manifold. A bracket of functions (the Jacobi bracket) is defined by $`\{\overline{f},\overline{g}\}=\mathrm{\Lambda }(d\overline{f},d\overline{g})+\overline{f}E(\overline{g})\overline{g}E(\overline{f})`$, for all $`\overline{f},\overline{g}C^{\mathrm{}}(M,\text{I}\text{R})`$. Note that
$$\{\overline{f},\overline{g}\overline{h}\}=\overline{g}\{\overline{f},\overline{h}\}+\overline{h}\{\overline{f},\overline{g}\}\overline{g}\overline{h}\{\overline{f},1\}.$$
(2)
In fact, the space $`C^{\mathrm{}}(M,\text{I}\text{R})`$ endowed with the Jacobi bracket is a local Lie algebra in the sense of Kirillov (see ). Conversely, a structure of local Lie algebra on $`C^{\mathrm{}}(M,\text{I}\text{R})`$ defines a Jacobi structure on $`M`$ (see ). If the vector field $`E`$ identically vanishes then $`(M,\mathrm{\Lambda })`$ is a Poisson manifold. Jacobi and Poisson manifolds were introduced by Lichnerowicz () (see also ).
A Lie algebroid structure on a differentiable vector bundle $`\pi :AM`$ is a pair that consists of a Lie algebra structure $`[[,]]`$ on the space $`\mathrm{\Gamma }(A)`$ of the global cross sections of $`\pi :AM`$ and a homomorphism of vector bundles $`\rho :ATM`$, the anchor map, such that if we also denote by $`\rho :\mathrm{\Gamma }(A)X(M)`$ the homomorphism of $`C^{\mathrm{}}(M,\text{I}\text{R})`$-modules induced by the anchor map then: (i) $`\rho :(\mathrm{\Gamma }(A),[[,]])(X(M),[,])`$ is a Lie algebra homomorphism and (ii) for all $`\overline{f}C^{\mathrm{}}(M,\text{I}\text{R})`$ and for all $`\mu ,\eta \mathrm{\Gamma }(A)`$, one has $`[[\mu ,\overline{f}\eta ]]=\overline{f}[[\mu ,\eta ]]+(\rho (\mu )(\overline{f}))\eta `$ (see ).
If $`(A,[[,]],\rho )`$ is a Lie algebroid over $`M`$, one can introduce the Lie algebroid cohomology complex with trivial coefficients (for the explicit definition of this complex we remit to ). The space of 1-cochains is $`\mathrm{\Gamma }(A^{})`$, where $`A^{}`$ is the dual bundle to $`A`$, and a 1-cochain $`\varphi \mathrm{\Gamma }(A^{})`$ is a 1-cocycle if and only if
$$\varphi [[\mu ,\eta ]]=\rho (\mu )(\varphi (\eta ))\rho (\eta )(\varphi (\mu )),\text{ for all }\mu ,\eta \mathrm{\Gamma }(A).$$
(3)
A Jacobi manifold $`(M,\mathrm{\Lambda },E)`$ has an associated Lie algebroid $`(T^{}M\times \text{I}\text{R},[[,]]_{(\mathrm{\Lambda },E)},\mathrm{\#}_{(\mathrm{\Lambda },E)})`$, where $`T^{}M`$ is the cotangent bundle of $`M`$ and $`[[,]]_{(\mathrm{\Lambda },E)}`$, $`\mathrm{\#}_{(\mathrm{\Lambda },E)}`$ are defined by
$$\begin{array}{ccc}[[(\alpha ,\overline{f}),(\beta ,\overline{g})]]_{(\mathrm{\Lambda },E)}& =\hfill & (_{\mathrm{\#}_\mathrm{\Lambda }(\alpha )}\beta _{\mathrm{\#}_\mathrm{\Lambda }(\beta )}\alpha d(\mathrm{\Lambda }(\alpha ,\beta ))+\overline{f}_E\beta \overline{g}_E\alpha i_E(\alpha \beta ),\hfill \\ & & \mathrm{\Lambda }(\beta ,\alpha )+\mathrm{\#}_\mathrm{\Lambda }(\alpha )(\overline{g})\mathrm{\#}_\mathrm{\Lambda }(\beta )(\overline{f})+\overline{f}E(\overline{g})\overline{g}E(\overline{f})),\hfill \\ & & \\ \mathrm{\#}_{(\mathrm{\Lambda },E)}(\alpha ,\overline{f})& =\hfill & \mathrm{\#}_\mathrm{\Lambda }(\alpha )+\overline{f}E,\hfill \end{array}$$
(4)
for $`(\alpha ,\overline{f}),(\beta ,\overline{g})\mathrm{\Omega }^1(M)\times C^{\mathrm{}}(M,\text{I}\text{R})`$, $``$ being the Lie derivative operator and $`\mathrm{\#}_\mathrm{\Lambda }:\mathrm{\Omega }^1(M)X(M)`$ the mapping given by $`\beta (\mathrm{\#}_\mathrm{\Lambda }(\alpha ))=\mathrm{\Lambda }(\alpha ,\beta )`$ (see ).
In the particular case when $`(M,\mathrm{\Lambda })`$ is a Poisson manifold we recover, by projection, the Lie algebroid $`(T^{}M,[[,]]_\mathrm{\Lambda },`$ $`\mathrm{\#}_\mathrm{\Lambda })`$, where $`[[,]]_\mathrm{\Lambda }`$ is the bracket of 1-forms defined by (see ):
$$[[,]]_\mathrm{\Lambda }:\mathrm{\Omega }^1(M)\times \mathrm{\Omega }^1(M)\mathrm{\Omega }^1(M),[[\alpha ,\beta ]]_\mathrm{\Lambda }=_{\mathrm{\#}_\mathrm{\Lambda }(\alpha )}\beta _{\mathrm{\#}_\mathrm{\Lambda }(\beta )}\alpha d(\mathrm{\Lambda }(\alpha ,\beta )).$$
## 2 Some linear Jacobi structures on vector bundles
Let $`\pi :AM`$ be a vector bundle and $`A^{}`$ the dual bundle to $`A`$. Suppose that $`\pi ^{}:A^{}M`$ is the canonical projection. If $`\mu \mathrm{\Gamma }(A)`$ and $`\overline{f}C^{\mathrm{}}(M,\text{I}\text{R})`$ then $`\mu `$ determines a linear function on $`A^{}`$ which we will denote by $`\stackrel{~}{\mu }`$ and $`f=\overline{f}\pi ^{}`$ is a $`C^{\mathrm{}}`$ real-valued function on $`A^{}`$ which is basic.
Now, assume that $`(A,[[,]],\rho )`$ is a Lie algebroid over $`M`$. Then $`A^{}`$ admits a Poisson structure $`\mathrm{\Lambda }_A^{}`$ such that the Poisson bracket of linear functions is again linear (see ). The local expression of $`\mathrm{\Lambda }_A^{}`$ is given as follows. Let $`U`$ be an open coordinate neighbourhood of $`M`$ with coordinates $`(x^1,\mathrm{},x^m)`$ and $`\{e_i\}_{i=1,\mathrm{},n}`$ a local basis of sections of $`\pi :AM`$ in $`U`$. Then, $`(\pi ^{})^1(U)`$ is an open coordinate neighbourhood of $`A^{}`$ with coordinates $`(x^i,\mu _j)`$ such that $`\mu _j=\stackrel{~}{e}_j`$, for all $`j`$. In these coordinates the structure functions and the components of the anchor map are
$`[[e_i,e_j]]={\displaystyle \underset{k=1}{\overset{n}{}}}c_{ij}^ke_k,\rho (e_i)={\displaystyle \underset{l=1}{\overset{m}{}}}\rho _i^l{\displaystyle \frac{}{x^l}},i,j\{1,\mathrm{},n\},`$ (5)
with $`c_{ij}^k,\rho _i^lC^{\mathrm{}}(U,\text{I}\text{R})`$, and the Poisson structure $`\mathrm{\Lambda }_A^{}`$ is given by
$$\mathrm{\Lambda }_A^{}=\underset{i<j}{}\underset{k}{}c_{ij}^k\mu _k\frac{}{\mu _i}\frac{}{\mu _j}+\underset{i,l}{}\rho _i^l\frac{}{\mu _i}\frac{}{x^l}.$$
(6)
Next, we will show an extension of the above results for the Jacobi case.
We will denote by $`\mathrm{\Delta }`$ the Liouville vector field of $`A^{}`$ and by $`\varphi ^vX(A^{})`$ the vertical lift of $`\varphi \mathrm{\Gamma }(A^{})`$. Note that if $`(x^i,\mu _j)`$ are fibred coordinates on $`A^{}`$ as above and $`\varphi ={\displaystyle \underset{i=1}{\overset{n}{}}}\varphi _ie^i`$, with $`\varphi _iC^{\mathrm{}}(U,\text{I}\text{R})`$ and $`\{e^i\}`$ the dual basis of $`\{e_i\}`$, then
$$\mathrm{\Delta }=\underset{i=1}{\overset{n}{}}\mu _i\frac{}{\mu _i},\varphi ^v=\underset{i=1}{\overset{n}{}}\varphi _i\frac{}{\mu _i}.$$
(7)
Thus, using (1), (3), (5), (6) and (7), we deduce
###### Theorem 1
Let $`(A,[[,]],\rho )`$ be a Lie algebroid over $`M`$ and $`\varphi \mathrm{\Gamma }(A^{})`$ a 1-cocycle. Then, there is a unique Jacobi structure $`(\mathrm{\Lambda }_{(A^{},\varphi )},E_{(A^{},\varphi )})`$ on $`A^{}`$ with Jacobi bracket $`\{,\}_{(A^{},\varphi )}`$ satisfying
$$\{\stackrel{~}{\mu },\stackrel{~}{\eta }\}_{(A^{},\varphi )}=\stackrel{~}{[[\mu ,\eta ]]},\{\stackrel{~}{\mu },\overline{f}\pi ^{}\}_{(A^{},\varphi )}=(\rho (\mu )(\overline{f})+\varphi (\mu )\overline{f})\pi ^{},\{\overline{f}\pi ^{},\overline{g}\pi ^{}\}_{(A^{},\varphi )}=0,$$
for $`\mu ,\eta \mathrm{\Gamma }(A)`$ and $`\overline{f},\overline{g}C^{\mathrm{}}(M,\text{I}\text{R})`$. The Jacobi structure is given by
$$\mathrm{\Lambda }_{(A^{},\varphi )}=\mathrm{\Lambda }_A^{}+\mathrm{\Delta }\varphi ^v,E_{(A^{},\varphi )}=\varphi ^v.$$
Now, we will prove a converse of Theorem 1.
###### Theorem 2
Let $`\pi :AM`$ be a vector bundle over $`M`$ and let $`(\mathrm{\Lambda },E)`$ be a Jacobi structure on the dual bundle $`A^{}`$ satisfying:
* The Jacobi bracket of linear functions is again linear.
* The Jacobi bracket of a linear function and the constant function 1 is a basic function.
Then, there is a Lie algebroid structure on $`\pi :AM`$ and a 1-cocycle $`\varphi \mathrm{\Gamma }(A^{})`$ such that $`\mathrm{\Lambda }=\mathrm{\Lambda }_{(A^{},\varphi )}`$ and $`E=E_{(A^{},\varphi )}`$.
Proof: Denote by $`\{,\}`$ the Jacobi bracket on $`A^{}`$ induced by the Jacobi structure $`(\mathrm{\Lambda },E)`$ and suppose that $`\mu ,\eta \mathrm{\Gamma }(A)`$ and that $`\overline{f},\overline{g}C^{\mathrm{}}(M,\text{I}\text{R})`$. If $`\pi ^{}:A^{}M`$ is the canonical projection, the function $`\{(\overline{f}\pi ^{})\stackrel{~}{\mu },1\}=\{\stackrel{~}{\overline{f}\mu },1\}`$ is basic. Thus, from (2) and (C2), we have that
$$\{\overline{f}\pi ^{},1\}=0.$$
(8)
On the other hand, the function $`\{\stackrel{~}{\mu },(\overline{f}\pi ^{})\stackrel{~}{\eta }\}=\{\stackrel{~}{\mu },\stackrel{~}{\overline{f}\eta }\}`$ is linear. Therefore, from (2), (C1) and (C2), we obtain that the function $`\{\stackrel{~}{\mu },\overline{f}\pi ^{}\}`$ is basic. Consequently, the Jacobi bracket of a linear function and a basic function is a basic function. In particular, $`\{\overline{f}\pi ^{},(\overline{g}\pi ^{})\stackrel{~}{\mu }\}=\{\overline{f}\pi ^{},\stackrel{~}{\overline{g}\mu }\}`$ is basic. This implies that (see (2) and (8))
$$\{\overline{f}\pi ^{},\overline{g}\pi ^{}\}=0.$$
(9)
Now, we define the section $`[[\mu ,\eta ]]`$ of the vector bundle $`\pi :AM`$ and the $`C^{\mathrm{}}`$ real-valued functions on $`M`$, $`\varphi (\mu )`$ and $`\rho (\mu )(\overline{f})`$, which are characterized by the following relations
$`\stackrel{~}{[[\mu ,\eta ]]}=\{\stackrel{~}{\mu },\stackrel{~}{\eta }\},\varphi (\mu )\pi ^{}=\{\stackrel{~}{\mu },1\},\rho (\mu )(\overline{f})\pi ^{}=\{\stackrel{~}{\mu },\overline{f}\pi ^{}\}(\overline{f}\pi ^{})\{\stackrel{~}{\mu },1\}.`$ (10)
From (2), (8), (9) and (10), we deduce that $`\varphi `$ can be considered as a $`C^{\mathrm{}}(M,\text{I}\text{R})`$-linear map $`\varphi :\mathrm{\Gamma }(A)C^{\mathrm{}}(M,\text{I}\text{R})`$ (that is, $`\varphi \mathrm{\Gamma }(A^{})`$) and that $`\rho `$ can be considered as a $`C^{\mathrm{}}(M,\text{I}\text{R})`$-linear map $`\rho :\mathrm{\Gamma }(A)X(M)`$. Moreover, using (2), (3), (10) and the fact that $`\{,\}`$ is the Jacobi bracket of a Jacobi structure (see Section 1), it follows that the triple $`(A,[[,]],\rho )`$ is a Lie algebroid over $`M`$ and that $`\varphi \mathrm{\Gamma }(A^{})`$ is a 1-cocycle. Finally, by (9), (10) and Theorem 1, we conclude that $`(\mathrm{\Lambda },E)=(\mathrm{\Lambda }_{(A^{},\varphi )},E_{(A^{},\varphi )}).`$ $`QED`$
###### Remark 1
That condition (C1) does not necessarily imply condition (C2) is illustrated by the following simple example. Let $`M`$ be a single point and $`A^{}=\text{I}\text{R}^2`$ endowed with the Jacobi structure $`(\mathrm{\Lambda },E)`$, where $`\mathrm{\Lambda }=xy\frac{}{x}\frac{}{y}`$ and $`E=x\frac{}{x}`$. It is easy to prove that the Jacobi bracket satisfies (C1) but not (C2).
Let $`M`$ be a differentiable manifold and $`\pi :AM`$ a vector bundle. Denote by $`\stackrel{~}{𝒜}`$ and $`𝒥`$ the following sets. $`\stackrel{~}{𝒜}`$ is the set of the pairs $`(([[,]],\rho ),\varphi )`$, where $`([[,]],\rho )`$ is a Lie algebroid structure on $`\pi :AM`$ and $`\varphi \mathrm{\Gamma }(A^{})`$ is a 1-cocycle. $`𝒥`$ is the set of the Jacobi structures $`(\mathrm{\Lambda },E)`$ on $`A^{}`$ which satisfy the conditions (C1) and (C2) (see Theorem 2).
Then, using Theorems 1 and 2, we obtain
###### Theorem 3
The mapping $`\mathrm{\Psi }:\stackrel{~}{𝒜}𝒥`$ between the sets $`\stackrel{~}{𝒜}`$ and $`𝒥`$ given by
$$\mathrm{\Psi }(([[,]],\rho ),\varphi )=(\mathrm{\Lambda }_{(A^{},\varphi )},E_{(A^{},\varphi )})$$
is a bijection.
Note that $`\mathrm{\Psi }(𝒜)=𝒫`$, where $`𝒫`$ is the subset of the Jacobi structures of $`𝒥`$ which are Poisson and $`𝒜`$ is the subset of $`\stackrel{~}{𝒜}`$ of the pairs of the form $`(([[,]],\rho ),0)`$, that is, $`𝒜`$ is the set of the Lie algebroid structures on $`\pi :AM`$. Therefore, from Theorem 3, we deduce a well known result (see ): the mapping $`\mathrm{\Psi }`$ induces a bijection between the sets $`𝒜`$ and $`𝒫`$.
## 3 Examples and applications
In this section we will present some examples and applications of the results obtained in Section 2.
1.- Let $`(g,[,])`$ be a real Lie algebra of dimension $`n`$. Then, $`g`$ is a Lie algebroid over a point. The resultant Poisson structure $`\mathrm{\Lambda }_g^{}`$ on $`g^{}`$ is the well known Lie-Poisson structure (see (6)). Thus, if $`\varphi g^{}`$ is a 1-cocycle then, using Theorem 1, we deduce that the pair $`(\mathrm{\Lambda }_{(g^{},\varphi )},E_{(g^{},\varphi )})`$ is a Jacobi structure on $`g^{}`$, where
$$\mathrm{\Lambda }_{(g^{},\varphi )}=\mathrm{\Lambda }_g^{}+RC_\varphi ,E_{(g^{},\varphi )}=C_\varphi ,$$
$`R`$ is the radial vector field on $`g^{}`$ and $`C_\varphi `$ is the constant vector field on $`g^{}`$ induced by $`\varphi g^{}`$.
2.- Let $`(TM,[,],Id)`$ be the trivial Lie algebroid. In this case, the Poisson structure $`\mathrm{\Lambda }_{T^{}M}`$ on $`T^{}M`$ is the canonical symplectic structure. Therefore, if $`\varphi `$ is a closed 1-form on $`M`$, then the pair
$$\mathrm{\Lambda }_{(T^{}M,\varphi )}=\mathrm{\Lambda }_{T^{}M}+\mathrm{\Delta }\varphi ^v,E_{(T^{}M,\varphi )}=\varphi ^v,$$
is a Jacobi structure on $`T^{}M`$. Furthermore, we can prove that the map $`\mathrm{\#}_{\mathrm{\Lambda }_{(T^{}M,\varphi )}}:\mathrm{\Omega }^1(T^{}M)X(T^{}M)`$ is an isomorphism and consequently, using the results of (see also ), it follows that $`(\mathrm{\Lambda }_{(T^{}M,\varphi )},E_{(T^{}M,\varphi )})`$ is a locally conformal symplectic structure.
3.- Let $`(M,\mathrm{\Lambda })`$ be a Poisson manifold and $`(T^{}M,[[,]]_\mathrm{\Lambda },\mathrm{\#}_\mathrm{\Lambda })`$ the associated cotangent Lie algebroid (see Section 1). The induced Poisson structure on $`TM`$ is the complete lift $`\mathrm{\Lambda }^c`$ to $`TM`$ of $`\mathrm{\Lambda }`$ (see ). Thus, if $`XX(M)=\mathrm{\Gamma }(TM)`$ is a 1-cocycle, that is, $`X`$ is a Poisson infinitesimal automorphism ($`_X\mathrm{\Lambda }=0`$), we deduce that
$$\mathrm{\Lambda }_{(TM,X)}=\mathrm{\Lambda }^c+\mathrm{\Delta }X^v,E_{(TM,X)}=X^v,$$
is a Jacobi structure on $`TM`$.
4.- The triple $`(TM\times \text{I}\text{R},\text{[}\text{ , }\text{]},\pi )`$ is a Lie algebroid over $`M`$, where $`\pi :TM\times \text{I}\text{R}TM`$ is the canonical projection over the first factor and \[ , \] is the bracket given by
$$\text{[}(X,\overline{f}),(Y,\overline{g})\text{]}=([X,Y],X(\overline{g})Y(\overline{f})),\text{ for }(X,\overline{f}),(Y,\overline{g})X(M)\times C^{\mathrm{}}(M,\text{I}\text{R}).$$
(11)
In this case, the Poisson structure $`\mathrm{\Lambda }_{T^{}M\times \text{I}\text{R}}`$ on $`T^{}M\times \text{I}\text{R}`$ is just the canonical cosymplectic structure of $`T^{}M\times \text{I}\text{R}`$, that is, $`\mathrm{\Lambda }_{T^{}M\times \text{I}\text{R}}=\mathrm{\Lambda }_{T^{}M}`$. Now, it is easy to prove that $`\varphi =(0,1)\mathrm{\Omega }^1(M)\times C^{\mathrm{}}(M,\text{I}\text{R})=\mathrm{\Gamma }(T^{}M\times \text{I}\text{R})`$ is a 1-cocycle (see (3) and (11)). Moreover, using Theorem 1, we have that the Jacobi structure $`(\mathrm{\Lambda }_{(T^{}M\times \text{I}\text{R},\varphi )},E_{(T^{}M\times \text{I}\text{R},\varphi )})`$ on $`T^{}M\times \text{I}\text{R}`$ is the one defined by the canonical contact 1-form $`\eta _M`$. We recall that $`\eta _M`$ is the 1-form on $`T^{}M\times \text{I}\text{R}`$ given by $`\eta _M=dt+\lambda _M`$, $`\lambda _M`$ being the Liouville 1-form of $`T^{}M`$ (see ).
5.- Let $`(M,\mathrm{\Lambda },E)`$ be a Jacobi manifold and $`(T^{}M\times \text{I}\text{R},[[,]]_{(\mathrm{\Lambda },E)},\mathrm{\#}_{(\mathrm{\Lambda },E)})`$ the associated Lie algebroid (see Section 1). From (1), (3) and (4), it follows that $`\varphi =(E,0)X(M)\times C^{\mathrm{}}(M,\text{I}\text{R})=\mathrm{\Gamma }(TM\times \text{I}\text{R})`$ is a 1-cocycle. On the other hand, a long computation, using (4), (6), (7) and Theorem 1, shows that
$$\mathrm{\Lambda }_{(TM\times \text{I}\text{R},\varphi )}=\mathrm{\Lambda }^c+\frac{}{t}E^ct\left(\mathrm{\Lambda }^v+\frac{}{t}E^v\right),E_{(TM\times \text{I}\text{R},\varphi )}=E^v,$$
where $`\mathrm{\Lambda }^c`$ (resp. $`\mathrm{\Lambda }^v`$) is the complete (resp. vertical) lift to $`TM`$ of $`\mathrm{\Lambda }`$ and $`E^c`$ (resp. $`E^v`$) is the complete (resp. vertical) lift to $`TM`$ of $`E`$. We remark that in the authors characterize the conformal infinitesimal automorphisms of $`(M,\mathrm{\Lambda },E)`$ as Legendre-Lagrangian submanifolds of the Jacobi manifold $`(TM\times \text{I}\text{R},\mathrm{\Lambda }_{(TM\times \text{I}\text{R},\varphi )},E_{(TM\times \text{I}\text{R},\varphi )})`$.
6.- Let $`(A,[[,]],\rho )`$ be a Lie algebroid over $`M`$ and $`\varphi \mathrm{\Gamma }(A^{})`$ a 1-cocycle. Denote by $`\widehat{\mathrm{\Lambda }}_{A^{}\times \text{I}\text{R}}`$ the Poissonization of the Jacobi structure $`(\mathrm{\Lambda }_{(A^{},\varphi )},E_{(A^{},\varphi )})`$, that is, $`\widehat{\mathrm{\Lambda }}_{A^{}\times \text{I}\text{R}}`$ is the Poisson structure on $`\widehat{A}^{}=A^{}\times \text{I}\text{R}`$ given by (see )
$$\widehat{\mathrm{\Lambda }}_{A^{}\times \text{I}\text{R}}=e^t\left(\mathrm{\Lambda }_{(A^{},\varphi )}+\frac{}{t}E_{(A^{},\varphi )}\right).$$
(12)
$`\widehat{A}^{}`$ is the total space of a vector bundle over $`M\times \text{I}\text{R}`$ and, from (12), we obtain that the Poisson bracket of two linear functions on $`\widehat{A}^{}`$ is again linear. This implies that the dual vector bundle $`\widehat{A}=A\times \text{I}\text{R}M\times \text{I}\text{R}`$ admits a Lie algebroid structure $`([[,]]\widehat{},\widehat{\rho })`$. Note that the space $`\mathrm{\Gamma }(\widehat{A})`$ can be identified with the set of time-dependent sections of $`AM`$. Under this identification, we deduce that (see (10) and (12))
$$[[\widehat{\mu },\widehat{\eta }]]\widehat{}=e^t\left([[\widehat{\mu },\widehat{\eta }]]+\varphi (\widehat{\mu })(\frac{d\widehat{\eta }}{dt}\widehat{\eta })\varphi (\widehat{\eta })(\frac{d\widehat{\mu }}{dt}\widehat{\mu })\right),\widehat{\rho }(\widehat{\mu })=e^t\left(\rho (\widehat{\mu })+\varphi (\widehat{\mu })\frac{}{t}\right),$$
for all $`\widehat{\mu },\widehat{\eta }\mathrm{\Gamma }(\widehat{A})`$, where $`{\displaystyle \frac{d\widehat{\mu }}{dt}}`$ (resp. $`{\displaystyle \frac{d\widehat{\eta }}{dt}}`$) is the derivative of $`\widehat{\mu }`$ (resp. $`\widehat{\eta }`$) with respect to the time. Note that if $`t\text{I}\text{R}`$ then the sections $`\widehat{\mu }`$ and $`\widehat{\eta }`$ induce, in a natural way, two sections $`\widehat{\mu }_t`$ and $`\widehat{\eta }_t`$ of $`AM`$ and that $`[[\widehat{\mu },\widehat{\eta }]]`$ is the time-dependent section of $`AM`$ given by $`[[\widehat{\mu },\widehat{\eta }]](x,t)=[[\widehat{\mu }_t,\widehat{\eta }_t]](x)`$, for all $`(x,t)M\times \text{I}\text{R}`$.
The construction of the Lie algebroid $`(\widehat{A},[[,]]\widehat{},\widehat{\rho })`$ from the Lie algebroid $`(A,[[,]],\rho )`$ and the cocycle $`\varphi `$ plays an important role in .
Acknowledgments. Research partially supported by DGICYT grant PB97-1487 (Spain). D. Iglesias wishes to thank the Spanish Ministerio de Educación y Cultura for a FPU grant.
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# 1 Introduction
## 1 Introduction
Neutrinoless double beta (0$`\nu `$2$`\beta `$) decay is forbidden in the Standard Model (SM) since it violates lepton number ($`L`$) conservation. However many extensions of the SM incorporate $`L`$ violating interactions and thus could lead to 0$`\nu `$2$`\beta `$ decay . Currently, besides the conventional left-handed neutrino ($`\nu `$) exchange mechanism, there are many other possibilities to trigger this process : right-handed $`\nu `$ exchange in left-right symmetric models; exchange of squarks, sneutrinos, etc. via supersymmetric (SUSY) interactions; exchange of leptoquarks in models with leptoquarks; exchange of excited Majorana neutrinos in models with composite heavy neutrinos, and so on. In that sense 0$`\nu `$2$`\beta `$ decay has a great conceptual importance due to the strong statement obtained in a gauge theory of the weak interaction that a non-vanishing 0$`\nu `$2$`\beta `$ decay rate requires neutrinos to be massive Majorana particles, independently of which mechanism induces it . Therefore, at present 0$`\nu `$2$`\beta `$ decay is considered as a powerful method to test new physical effects beyond the SM, while absence of this process – established at the present level of sensitivity – would yield strong restrictions on parameters of manifold extensions of the SM and narrow the wide choice of the theoretical models. At the same time 0$`\nu `$2$`\beta `$ decay is very important in light of the current status of neutrino physics (see ). Indeed, the solar neutrino problem (in particular lack of <sup>7</sup>Be neutrinos) and the deficit of the atmospheric muon neutrino flux could be explained by means of neutrino oscillations, which require nonzero neutrino masses. Also indication for $`\nu _\mu `$/$`\nu _e`$ oscillations was found by the LSND collaboration . Oscillation experiments are only sensitive to neutrino mass difference, while measuring the 0$`\nu `$2$`\beta `$ decay rate can give the absolute value of the $`\nu `$ mass scale, and hence provide a crucial test of neutrino mass models.
Despite the numerous efforts to observe 0$`\nu `$2$`\beta `$ decay beginning from 1948 up to the present this process still remains unobserved. The highest half-life limits were set in direct experiments with several nuclides: $`T_{1/2}`$(0$`\nu )10^{22}`$ yr for <sup>82</sup>Se , <sup>100</sup>Mo ; $`T_{1/2}`$(0$`\nu )10^{23}`$ yr for <sup>116</sup>Cd , <sup>130</sup>Te and <sup>136</sup>Xe ; and T<sub>1/2</sub>(0$`\nu )10^{25}`$ yr for <sup>76</sup>Ge .
The present theoretical and experimental status of 2$`\beta `$ decay investigations makes it necessary to extend the number of candidate nuclides studied at a sensitivity comparable to or better than that for <sup>76</sup>Ge (neutrino mass limit $`m_\nu `$ $`0.10.5`$ eV). With this aim we consider in the present paper the use of the super-low background liquid scintillation detector – the BOREXINO Counting Test Facility (CTF) – for high sensitivity 2$`\beta `$ decay research.
## 2 The CTF and choice of candidate nuclei
The full description of the CTF and its performance have been published elsewhere . Here we recall the main features of this apparatus, which are important in the following.
The CTF is installed in the Gran Sasso Underground Laboratory and consists of an external cylindrical water tank ($``$11$`\times `$10 m; $``$1000 t of water) serving as passive shielding for 4.8 m<sup>3</sup> of liquid scintillator contained in an inner spherical vessel of $``$2.1 m. High purity water is supplied by the BOREXINO water plant, which provides its radio-purity level of $`10^{14}`$ g/g (U, Th), $`10^{10}`$ g/g (K natural) and $`<`$5 $`\mu `$Bq/l for <sup>222</sup>Rn .
The liquid scintillator is a binary solution of 1.5 g/l of PPO in pseudocumene. The fluorescence peak emission is at 365 nm and the yield of emitted photons is $``$10<sup>4</sup> per MeV of energy deposited. The attenuation length is larger than $`5`$ m above 380 nm . The principal scintillator decay time was measured as $``$3.5 ns in a small volume, while as 4.5–5.0 ns with a source placed in the center of the CTF. The scintillator is purified by recirculating it from the inner vessel through a Rn stripping tower, a water extraction unit, a Si-Gel column extraction unit, and a vacuum distillation unit. It ensures that <sup>232</sup>Th and <sup>238</sup>U contaminations in the liquid scintillator are less than $`(2`$$`5)10^{16}`$ g/g.
The inner vessel for the liquid scintillator is made of nylon film, 500 $`\mu `$m thick, with excellent optical clarity at 350-500 nm, which allows collection of scintillation light with the help of 100 phototubes (PMT) fixed to a 7 m diameter support structure inside the water tank. The PMTs are 8” Thorn EMI 9351 tubes made of low radioactivity Schott 8246 glass, and characterized by high quantum efficiency (26% at 420 nm), limited transit time spread ($`\sigma `$ = $`1`$ ns), good pulse height resolution for single photoelectron pulses (Peak/Valley = $`2.5`$), low dark noise rate (0.5 kHz), low after pulse probability (2.5%), and a gain of 10<sup>7</sup>. The PMTs are fitted with light concentrators 57 cm long and 50 cm diameter aperture. They provide 20% optical coverage. The number of photoelectrons per MeV measured experimentally is (300 $`\pm `$ 30)/MeV on average. An upgrade of the CTF was realized in 1999 when an additional nylon barrier against radon convection and a muon veto system were installed.
For each event the charge and timing (precision of 1 ns) of hit PMTs are recorded. Each channel is doubled by an auxiliary channel to record all other events coming within a time window of 8 ms after the first event. For longer delay the computer clock is used (accuracy of $``$0.1 s). Event parameters measured in the CTF include:
\- the total charge collected by the PMTs during $`0`$$`500`$ ns (event energy);
\- the tail charge ($`48`$$`548`$ ns) used for pulse shape discrimination;
\- PMT timing to reconstruct the event in space (resolution of 10–15 cm);
\- the time elapsed between sequential events, used to tag time-correlated events.
Due to all these measures the CTF is the best super-low background scintillator of large volume at present. Indeed, the total background rate in the energy region 250 – 800 keV (so called ”solar neutrino energy window”) is about 0.3 counts/yr$``$keV$``$kg and is dominated by external background from Rn in the shielding water ($``$30 mBq/m<sup>3</sup> in the region surrounding the inner vessel), while the internal background is less than 0.01 counts/yr$``$keV$``$kg. Therefore one can conclude that the CTF is the ideal apparatus for super-low background 2$`\beta `$ decay research.
For the choice of 2$`\beta `$ candidate nuclei let us express the $`0\nu 2\beta `$ decay probability (neglecting right-handed contributions) as follows :
$$\left(T_{1/2}^{0\nu }\right)^1=G_{mm}^{0\nu }\left|ME\right|^2m_\nu ^2,$$
(1)
where $`\left|ME\right|`$ is the nuclear matrix element of the $`0\nu 2\beta `$ decay, and $`G_{mm}^{0\nu }(Z,Q_{\beta \beta })`$ is the phase space factor. Ignoring for the moment the $`\left|ME\right|`$ calculation , it is evident from Eq. (1) that for the sensitivity of the $`2\beta `$ decay study with a particular candidate the most important parameter is the available energy release ($`Q_{\beta \beta }`$). First, because the phase space integral (hence, $`0\nu 2\beta `$ decay rate) strongly depends on $`Q_{\beta \beta }`$ value (roughly as $`Q_{\beta \beta }^5`$). Second, the larger the $`2\beta `$ decay energy, the simpler – from an experimental point of view – to overcome background problems. The crucial value is 2614 keV which is the energy of the most dangerous $`\gamma `$’s from <sup>208</sup>Tl decay (<sup>232</sup>Th family). Among 35 candidates there are only six nuclei with $`Q_{\beta \beta }`$ larger than 2.6 MeV : <sup>48</sup>Ca ($`Q_{\beta \beta }`$ $`=4272`$ keV, natural abundance $`\delta =0.187\%`$), <sup>82</sup>Se ($`Q_{\beta \beta }`$ $`=2995`$ keV, $`\delta =8.73\%`$), <sup>96</sup>Zr ($`Q_{\beta \beta }`$ $`=3350`$ keV, $`\delta =2.80\%`$), <sup>100</sup>Mo ($`Q_{\beta \beta }`$ $`=3034`$ keV, $`\delta =9.63\%`$), <sup>116</sup>Cd ($`Q_{\beta \beta }`$ $`=2805`$ keV, $`\delta =7.49\%`$), and <sup>150</sup>Nd ($`Q_{\beta \beta }`$ $`=3367`$ keV, $`\delta =5.64\%`$). The values of the phase space integral $`G_{mm}^{0\nu }`$ for these candidates are (in units 10<sup>-14</sup> yr): 6.4 (<sup>48</sup>Ca), 2.8 (<sup>82</sup>Se), 5.7 (<sup>96</sup>Zr), 4.6 (<sup>100</sup>Mo), 4.9 (<sup>116</sup>Cd), and $``$20 (<sup>150</sup>Nd) . In comparison, $`G_{mm}^{0\nu }`$ for <sup>76</sup>Ge is equal $``$0.6 (in the same units) because of the lower $`2\beta `$ decay energy ($`Q_{\beta \beta }`$ $`=2039`$ keV). From this list <sup>100</sup>Mo and <sup>116</sup>Cd were chosen as candidates for $`2\beta `$ decay study with the CTF in the first phase. The main reason for the choice of <sup>100</sup>Mo – in addition to its high $`Q_{\beta \beta }`$ value – is the fact that the Institute for Nuclear Research (Kiev) possesses $``$1 kg of <sup>100</sup>Mo enriched to $``$99%. The <sup>116</sup>Cd was chosen because during the last decade the INR (Kiev) has developed and performed $`2\beta `$ decay experiments with this nuclide , which can be considered as a pilot step for the proposed project with the CTF.
## 3 CAMEO-I experiment with <sup>100</sup>Mo in the CTF
There are two different classes of $`2\beta `$ decay experiments: (a) an ”active” source experiment, in which the detector (containing $`2\beta `$ candidate nuclei) serves as source and detector simultaneously; (b) an experiment with ”passive” source which is introduced in the detector system . The sensitivity of any $`2\beta `$ decay apparatus is determined first, by the available source strengths (mass of the source), and second, by the detector background. Another factor very essential for determining the sensitivity is the energy resolution of the detector. Indeed, for a detector with poor energy resolution the events from the high energy tail of the $`2\nu 2\beta `$ decay distribution run into the energy window of the $`0\nu 2\beta `$ decay peak, and therefore generate background which cannot be discriminated from the $`0\nu 2\beta `$ decay signal, even in principle. All of the decay features are similar: the same two particles are emitted simultaneously from one point of the source, in the same energy region and with identical angular distribution. However, the better the energy resolution – the smaller the $`2\nu `$ tail becomes within the $`0\nu `$ interval, and thus the irreducible background becomes lower. Hence, we can conclude that the ultimate sensitivity to detect $`0\nu 2\beta `$ decay is really limited by the energy resolution of the detector, which is the most crucial parameter for any kind of set up for $`2\beta `$ decay study.
For the second class of experiments the sensitivity is also restricted by the trade-off between source strengths and detection efficiency. The number of $`2\beta `$ decay candidate nuclei can be enlarged by increasing the source thickness, which at the same time will lead to lower detection efficiency caused by absorption of electrons in the source and transformation of the measured $`2\beta `$ decay spectra (broadening of the peak and shifting it to lower energies).
These experimental considerations are illustrated in Fig. 1, where results of model experiments to study $`2\beta `$ decay of <sup>100</sup>Mo are presented. The following assumptions were accepted for simulation: mass of <sup>100</sup>Mo source is 1 kg ($``$6$`10^{24}`$ nuclei of <sup>100</sup>Mo); measuring time is 5 years; half-life of <sup>100</sup>Mo two neutrino $`2\beta `$ decay $`T_{1/2}`$(2$`\nu 2\beta )`$ $`=`$ $`10^{19}`$ yr (e.g. see ref. ), while for 0$`\nu `$ mode $`T_{1/2}`$(0$`\nu 2\beta )`$ = $`10^{24}`$ yr. The simulations were performed with the GEANT3.21 package and event generator DECAY4 , which describes the initial kinematics of the events in $`\alpha `$, $`\beta `$, and 2$`\beta `$ decay (how many particles are emitted, their types, energies, directions and times of emission). The initial $`2\beta `$ decay spectra are shown in Fig. 1a and Fig. 1b for different vertical scales. These spectra simulate $`2\beta `$ decay of <sup>100</sup>Mo nuclei placed in an ideal detector (”active” source technique) with 100% efficiency for $`2\beta `$ decay events and with zero background (energy resolution and energy threshold of 10 keV are presumed). For the next step the <sup>100</sup>Mo source was introduced in the same detector as a foil (”passive” source technique). The simulated spectra are depicted in Fig. 1c (thickness of <sup>100</sup>Mo foil is 15 mg/cm<sup>2</sup>) and Fig. 1d (60 mg/cm<sup>2</sup>). Finally the energy resolution of the detector ($`FWHM`$) was taken into account and results are shown in Fig. 1e ($`FWHM`$ = 4% at 3 MeV) and Fig. 1f ($`FWHM`$ = 8.8% at 3 MeV). It should be stressed that Fig. 1 represents the results of an ideal experiment, while in any real study the available results can only be worse by reason of the actual background, higher energy threshold and lower detection efficiency, etc. In fact, this is very strong statement because it allows to set the sensitivity limit for any real apparatus. For instance, it is evident from Fig. 1f that 0$`\nu 2\beta `$ decay of <sup>100</sup>Mo with half-life $`T_{1/2}`$ = $`10^{24}`$ yr will hardly be observed by using the ”passive” source technique with the following characteristics: (i) product of detection efficiency by number of <sup>100</sup>Mo nuclei $``$6$`10^{24}`$ (e.g., 1 kg of <sup>100</sup>Mo at 100% efficiency, or 10 kg of <sup>100</sup>Mo at 10% efficiency); (ii) <sup>100</sup>Mo source thickness of 60 mg/cm<sup>2</sup>; (iii) $`FWHM`$ = 8.8% at 3 MeV. At the same time, it is obvious from Fig. 1e that such a goal can be reached with similar apparatus with $`FWHM`$ = 4% at 3 MeV, and with a 15 mg/cm<sup>2</sup> source.
The last solution requires an increase of source area, which is usually limited by the available dimensions of the low background detectors used for $`2\beta `$ decay search. However, the unique features of the CTF – large sensitive volume and super-low background rate of the detector – permit an advanced $`2\beta `$ decay study of <sup>100</sup>Mo with the help of large square ($``$7 m<sup>2</sup>) and thin <sup>100</sup>Mo foil placed inside the liquid scintillator.
The <sup>100</sup>Mo source in the CTF is a complex system placed in the inner vessel with liquid scintillator. It can be represented by three mutually perpendicular and crossing flat disks with diameter of 180 cm whose centers are aligned with the center of inner vessel of the CTF. Each disk is composed of three layers: inner <sup>100</sup>Mo source<sup>1</sup><sup>1</sup>1In fact, in the plane of the disk the <sup>100</sup>Mo source itself consists of four sectors with spacing between them of 12 cm, which helps in spatial reconstruction when events occur near the crossing of the disks. (thickness of $``$15 mg/cm<sup>2</sup>) placed between two plastic scintillators of 1 mm thickness. The inner side of each plastic is coated with thin Al foil serving as light reflector, while the whole source ”sandwich” must be encapsulated by thin transparent film (made of teflon or syndiotactic polypropylene) to avoid plastic dissolution in liquid scintillator. Plastic detectors can have a much longer decay constant (f.e. as Bicron plastic BC-444 with $`\tau 260`$ ns and light output $``$40% of anthracene) with respect to the liquid scintillator, thus their pulses can be discriminated easily from the liquid scintillator signals. The plastics tag electrons emitted from the <sup>100</sup>Mo source, reducing the background of this complex detector system significantly. The energy loss measured by the plastics are added to the electron energy deposit in the liquid scintillator to obtain an accurate value of the electron energy and to improve the energy resolution of the whole detector.
### 3.1 Light collection, energy and spatial resolution
The energy resolution of the scintillation detector depends mainly on the quality of scintillator itself, the fraction of light collected by PMTs, uniformity of light collection, quantum efficiency and noise of the PMT photocathodes, stability and noise of the electronics. The excellent liquid scintillator used in the CTF yields about 10<sup>4</sup> emitted photons per MeV of energy deposited. In the present CTF design the actual optical coverage is 20% and the number of photoelectrons (p.e.) per MeV measured experimentally is (300 $`\pm `$ 30)/MeV on average. Thus with fourfold increase of light collection it would yield $``$1200 p.e. per 1 MeV or $``$3600 p.e. for 3 MeV.
To increase the actual optical coverage in the CTF, the PMTs can be mounted closer to the center of the detector. For instance, if 200 PMTs are fixed at diameter 5 m (and correspondingly the light concentrators’ entrances at diameter 4 m), or 96 PMTs are fixed at diameter 3.8 m, the optical coverage is equal $``$80%. We consider below the last configuration because it is the worst case for background contribution from the PMTs<sup>2</sup><sup>2</sup>2Special R&D is in progress now to find optimal solution for the required 80% optical coverage in the CTF.. Since the whole volume of the scintillator is divided by <sup>100</sup>Mo sources into 8 sectors, the PMTs are split into 8 groups of 12 PMTs each, so that one sector is viewed by one PMT group. Within the single sector (three mutually perpendicular reflector plates) scintillation photons would undergo less than 1.5 reflections on average before reaching the light concentrator aperture. The Monte Carlo simulation of the light propagation in such a geometry were performed with the help of GEANT3.21 program . The emission spectrum and angular distribution of scintillation photons were added to GEANT code. The simulation finds that 3 MeV energy deposit would yield $``$3700 photoelectrons allowing a measurement of the neutrinoless $`2\beta `$ decay peak of <sup>100</sup>Mo ($`Q_{\beta \beta }`$ $`=3034`$ keV) with an energy resolution $`FWHM`$ = 4%.
This goal can be reached if the non-uniformity of light collection is corrected by using accurate spatial information about each event; hence, the spatial reconstruction ability of the CTF has to be enhanced also. The results of the Monte Carlo simulation prove such a possibility and show that spatial resolution of $``$5–6 cm can be obtained with the upgraded CTF<sup>3</sup><sup>3</sup>3The value obtained for the spatial resolution should be considered as indicative because in this preliminary phase of study a simplified model for light propagation in the CTF liquid scintillator has been used.. Primarily this is due to better light collection (increased by a factor of four). Secondly, it is owing to the spatial reconstruction method based on the comparison of pulse amplitudes from the different PMTs (within one group of 12 PMTs), and at the same time due to analysis of the time structure of each pulse (which can include direct and reflected light).
### 3.2 Background simulation
The model of the CAMEO-I set up with <sup>100</sup>Mo (described above) is used for the calculations. We distinguish here between so called ”$`\beta `$” layers of the liquid scintillator, 15 cm thick<sup>4</sup><sup>4</sup>4These ”$`\beta `$” layers are separated from the total volume of the liquid scintillator by using the spatial information from the CTF. The thickness of 15 cm is chosen to guarantee the proper spatial reconstruction accuracy and the full absorption of the electrons emitted in the 2$`\beta `$ decay of <sup>100</sup>Mo. on both sides of the complex <sup>100</sup>Mo source, and the rest of the liquid scintillator volume serving as an active shield for these main inner layers. In such a detector system the following energies are measured:
i) $`E_1^{pl}`$ and $`E_2^{pl}`$ are the energy losses in the first and second plastic;
ii) $`E_1^\beta `$ and $`E_2^\beta `$ are the energy deposits in the first and second ”$`\beta `$” layer;
iii) $`E^{ls}`$ is the energy loss in the liquid scintillator active shield.
The energy threshold values of the detectors are set as $`E_{thr}^{pl}`$= 15 keV for the plastic and $`E_{thr}^{ls}`$ = $`E_{thr}^\beta `$ = 10 keV for the liquid scintillator. The energy resolution is $`FWHM^{pl}=\sqrt{10.8E^{pl}}`$ for the plastic and $`FWHM^\beta =\sqrt{4.8E^\beta }`$ for the liquid scintillator ($`FWHM`$, $`E^{pl}`$ and $`E^\beta `$ are in keV). The latter corresponds to the value of 4% at 3 MeV. The following cuts are used in the simulation in order to recognize the double $`\beta `$ decay events:
i) $`E_1^{pl}`$ or $`E_2^{pl}`$ $`E_{thr}^{pl}`$;
ii) $`E_1^{pl}+E_2^{pl}`$ $``$ 300 keV;
iii) if $`E_i^\beta `$ $`E_{thr}^\beta `$, the corresponding $`E_i^{pl}`$ must be $``$ $`E_{thr}^{pl}`$, necessarily;
iv) $`E^{ls}`$ $`E_{thr}^{ls}`$, i. e. there is no signal in the liquid scintillator active shield.
The simulation of the background and decays of radioactive nuclides in the installation were performed with the help of GEANT3.21 and the event generator DECAY4 .
#### 3.2.1 Two neutrino 2$`\beta `$ decays of <sup>100</sup>Mo
The half-life of 2$`\nu `$2$`\beta `$ decay of <sup>100</sup>Mo has been already measured as $``$10<sup>19</sup> yr (e.g. ref. ), hence the corresponding activity of a 1 kg <sup>100</sup>Mo source equals $``$13.2 mBq. The response functions of the CAMEO-I set up for $`2\nu 2\beta `$ decay of <sup>100</sup>Mo with $`T_{1/2}`$ $`=`$ 10<sup>19</sup> yr as well as for $`0\nu 2\beta `$ decay with $`T_{1/2}`$ $`=`$ 10<sup>24</sup> yr (for comparison) were simulated as described above, and results are depicted in Fig. 2a. The calculated values of efficiency for the neutrinoless channel are 80% (within energy window 2.8 – 3.15 MeV), 74% (2.85 – 3.15 MeV), and 63.5% (2.9 – 3.15 MeV). Background due to $`2\nu 2\beta `$ decay distribution are 19.7 counts (2.8 – 3.15 MeV), 6.1 counts (2.85 – 3.15 MeV), and 1.3 counts (2.9 – 3.15 MeV) for 5 years measuring period.
#### 3.2.2 Radioactive contamination of the <sup>100</sup>Mo source
The $``$1 kg sample of metallic molybdenum enriched in <sup>100</sup>Mo to $``$99% – which has to be applied in the present project – was already used in the quest for $`2\beta `$ decay of <sup>100</sup>Mo to the excited states of <sup>100</sup>Ru . In that experiment the radioactive impurities of the <sup>100</sup>Mo source by <sup>40</sup>K and nuclides from <sup>232</sup>Th and <sup>238</sup>U chains were measured. Only <sup>208</sup>Tl (measured activity is $``$ 0.5 mBq/kg) and <sup>214</sup>Bi (measured activity is (12$`\pm `$3) mBq/kg) – due to their high energy release – could generate the background in the $`0\nu 2\beta `$ decay energy window for <sup>100</sup>Mo. The background problem associated with <sup>100</sup>Mo radioactive contamination was carefully investigated by the NEMO collaboration, which is going to begin $`2\beta `$ decay measurements with $``$10 kg of <sup>100</sup>Mo . It was found that for the NEMO-3 detector the maximum acceptable internal activities of <sup>100</sup>Mo are 0.3 mBq/kg for <sup>214</sup>Bi and 0.02 mBq/kg for <sup>208</sup>Tl . The intensive R&D were performed by the NEMO collaboration with an aim to show that these severe requirements can be reached by using presently available physical and chemical methods for <sup>100</sup>Mo purification . On this basis, in our calculation the <sup>100</sup>Mo contamination criterion for <sup>214</sup>Bi has been taken as 0.3 mBq/kg, and for <sup>208</sup>Tl as 0.1 mBq/kg, which is only 5 times better than the actual activity limit in our <sup>100</sup>Mo sample (0.5 mBq/kg)<sup>5</sup><sup>5</sup>5We have accepted the less severe and more realistic criterion for <sup>208</sup>Tl, because it was shown that chemical purification of Mo is very successful concerning <sup>226</sup>Ra chain impurities, while for <sup>208</sup>Tl the procedure is very difficult and less effective ..
The results of simulation, performed as described above, are as following: i) <sup>214</sup>Bi contribution to background within the energy interval 2.9 – 3.15 MeV is 6.5 counts/yr$``$kg; ii) <sup>208</sup>Tl contribution is equal to 0.06 counts/yr$``$kg. The mentioned impurities can be really dangerous for the experiment. However, there exists a possibility to reduce these background substantially by using information on the arrival time of each event for analysis and selection of some decay chains in Th and U families . With this aim, let us consider the <sup>226</sup>Ra chain containing <sup>214</sup>Bi: <sup>222</sup>Rn ($`T_{1/2}`$ = 3.82 d; $`Q_\alpha `$ = $`5.59`$ MeV) $``$ <sup>218</sup>Po (3.10 m; $`Q_\alpha `$ = $`6.11`$ MeV) $``$ <sup>214</sup>Pb (26.8 m; $`Q_\beta `$ = $`1.02`$ MeV) $``$ <sup>214</sup>Bi (19.9 m; $`Q_\beta `$ = $`3.27`$ MeV) $``$ <sup>214</sup>Po (164.3 $`\mu `$s; $`Q_\alpha `$ = $`7.83`$ MeV) $``$ <sup>210</sup>Pb. The great advantage of the CAMEO-I experiment is the very thin <sup>100</sup>Mo source ($``$15 mg/cm<sup>2</sup>), which allows detection of most of $`\alpha `$ and $`\beta `$ particles emitted before or after <sup>214</sup>Bi decay, and tags the latter with the help of time analysis of the measured events. Indeed, our calculation gives the following values of the detection efficiencies: $`\epsilon _1`$ = 55% for <sup>214</sup>Po ($`\alpha `$ particles); $`\epsilon _2`$ = 80% for <sup>214</sup>Pb ($`\beta `$); $`\epsilon _3`$ = 37% for <sup>218</sup>Po ($`\alpha `$); $`\epsilon _4`$ = 32% for <sup>222</sup>Rn ($`\alpha `$). The probability to detect at least one of these decays (<sup>214</sup>Po or <sup>214</sup>Pb or <sup>218</sup>Po or <sup>222</sup>Rn) can be expressed as:
$`\epsilon =1(1\epsilon _1)(1\epsilon _2)(1\epsilon _3)(1\epsilon _4).`$
By substituting in this formula the calculated efficiency values, it yields $`\epsilon `$ = 96.1%, which means that only $``$4% of the <sup>214</sup>Bi decays would not be tagged, i.e. <sup>214</sup>Bi contribution to background can be reduced by a factor of 25 (to the value of $``$0.26 counts/yr$``$kg). The expected total $`\alpha `$ decay rate from <sup>238</sup>U and <sup>232</sup>Th families in the entire <sup>100</sup>Mo source is $``$300 decays/day, however for an area of 10$`\times 10`$ cm it is only 0.4 decays/day, which allows use of the chains with half-life of 26.8 and 19.9 minutes for time analysis. The simulated background spectrum from the internal <sup>100</sup>Mo source contamination by the <sup>214</sup>Bi and <sup>208</sup>Tl is presented in Fig. 2a, where the total internal background rate in the energy interval 2.9 – 3.15 MeV is 0.3 counts/yr$``$kg or $``$1.5 counts for 5 years measuring period.
#### 3.2.3 Cosmogenic activities in <sup>100</sup>Mo source
To estimate the cosmogenic activity produced in the <sup>100</sup>Mo foil, we used the program COSMO which calculates the production of all radionuclides with half-lives in the range from 25 days to 5 million years by nucleon-induced reactions in a given target. This program takes into account the variation of spallation, evaporation, fission and peripheral reaction cross sections with nucleon energy, target and product charge and mass numbers, as well as the energy spectrum of cosmic ray nucleons near the Earth’s surface.
For the CAMEO-I project cosmogenic activities were calculated for <sup>100</sup>Mo source enriched in <sup>100</sup>Mo to 98.5% (other Mo isotopes: <sup>98</sup>Mo – 0.7%, <sup>97</sup>Mo – 0.1%, <sup>96</sup>Mo – 0.2%, <sup>95</sup>Mo – 0.2%, <sup>94</sup>Mo – 0.1%, <sup>92</sup>Mo – 0.2%). It was assumed 5 years exposure period and deactivation time of about one year in the underground laboratory. The calculation shows that among several nuclides with $`T_{1/2}25`$ d produced in <sup>100</sup>Mo source only two can give some background in the energy window of the <sup>100</sup>Mo neutrinoless $`2\beta `$ decay. These are <sup>88</sup>Y ($`Q_{EC}`$=3.62 MeV; $`T_{1/2}`$=107 d) and <sup>60</sup>Co ($`Q_\beta `$=2.82 MeV; $`T_{1/2}`$=5.3 yr). Fortunately their activities are very low ($``$190 decays/yr for <sup>88</sup>Y and $``$50 decays/yr for <sup>60</sup>Co), thus the estimated background in the energy region of 2.7 – 3.2 MeV is practically negligible: $``$ 0.02 counts/yr$``$kg from <sup>88</sup>Y activity and $``$ 0.005 counts/yr$``$kg from <sup>60</sup>Co.
#### 3.2.4 External background
There are several origins of the external background in the CAMEO-I experiment, for example, neutrons and $`\gamma `$ quanta from natural environmental radioactivity (mainly from concrete walls of the Gran Sasso Underground Laboratory), contamination of PMTs by <sup>40</sup>K and nuclides from U and Th families, Rn impurities in the shielding water, cosmic muons ($`\mu `$ showers and muon induced neutrons, inelastic scattering and capture of muons), etc. From all of them only $`\gamma `$ quanta caused by PMT contamination and by Rn impurities in the shielding water were simulated in the present work, while others were simply estimated as negligible on the basis of the results of ref. , where such origins and contributions for the GENIUS project were investigated carefully.
The radioactivity values of the EMI 9351 PMT accepted for the simulation are: 0.194 Bq/PMT (<sup>208</sup>Tl); 1.383 Bq/PMT (<sup>214</sup>Bi); and 191 Bq/PMT (<sup>40</sup>K) . Also possible <sup>222</sup>Rn activity $``$30 mBq/m<sup>3</sup> in the shielding water (in the region surrounding the inner vessel) is taken into account. The model of the CAMEO-I detector system described above was used in the calculations, but with two differences: i) ”$`\beta `$” layers are considered as liquid scintillator blocks with dimensions 10$`\times `$10$`\times `$10 cm<sup>3</sup>; ii) the threshold of the liquid scintillator active shield is increased to 30 keV. The simulation performed under these assumptions gives the following background rate in the $`0\nu 2\beta `$ decay energy interval (2.9 – 3.15 MeV): i) 0.32 counts/yr$``$kg due to <sup>214</sup>Bi in PMT; ii) practically zero rates from <sup>208</sup>Tl in the PMTs and <sup>222</sup>Rn in the shielding water. The total simulated background contributions due to <sup>214</sup>Bi and <sup>208</sup>Tl contamination of the PMTs is shown in Fig. 2a also. After 5 years it yields $``$1.6 counts in the energy window 2.9 – 3.15 MeV. Summarizing all background sources for 5 years of measurements, one can obtain the total number of $``$4.4 counts in the energy range 2.9 – 3.15 MeV.
### 3.3 Sensitivity of the CAMEO-I experiment with <sup>100</sup>Mo
The sensitivity of the proposed experiment can be expressed in the term of a lower half-life limit for the $`0\nu 2\beta `$ decay of <sup>100</sup>Mo as following:
$$T_{1/2}\mathrm{ln}2N\eta t/S,$$
(2)
where $`N`$ is the number of <sup>100</sup>Mo nuclei ($``$6$``$10<sup>24</sup> in our case); $`t`$ is the measuring time (5 years); $`\eta `$ is the detection efficiency (63.5%); and $`S`$ is the number of effect’s events, which can be excluded with a given confidence level on the basis of measured data. Thus for the five years CAMEO-I experiment $`T_{1/2}`$ $``$ $`(`$13/$`S`$)$``$10<sup>24</sup> yr. Taking into account the expected background of 4.4 counts, we can accept 3–5 events as the value for $`S`$ (depending on the method of estimating $`S`$ ) which gives $`T_{1/2}`$ $``$ (3–5)$``$10<sup>24</sup> yr and, in accordance with , the limit on the neutrino mass $`<m_\nu >0.5`$ eV. On the other hand, it is evident from Fig. 2a that neutrinoless $`2\beta `$ decay of <sup>100</sup>Mo with half-life $`T_{1/2}`$ = 10<sup>24</sup> yr can certainly be registered: the signal (13 counts) to background (4.4 counts) ratio is approximately 3:1.
Similar limits $`T_{1/2}(0\nu 2\beta )`$ $``$ (3–5)$``$10<sup>24</sup> yr can be obtained by the CAMEO-I set up with other nuclides, <sup>82</sup>Se ($`Q_{\beta \beta }`$ $`=2996`$ keV), <sup>96</sup>Zr ($`Q_{\beta \beta }`$ $`=3350`$ keV), <sup>116</sup>Cd ($`Q_{\beta \beta }`$ $`=2804`$ keV), and <sup>150</sup>Nd ($`Q_{\beta \beta }`$ $`=3368`$ keV)<sup>6</sup><sup>6</sup>6We do not include <sup>48</sup>Ca ($`Q_{\beta \beta }`$ $`=4272`$ keV) in that list because of its very low natural abundance ($`0.187\%`$), and hence extremely high cost of a one kg <sup>48</sup>Ca source.. Due to its reasonable cost <sup>116</sup>Cd is the preferable second candidate after <sup>100</sup>Mo. Note, however, that a half-life limit of $``$5$``$10<sup>24</sup> yr for <sup>150</sup>Nd would lead – on the basis of the nuclear matrix elements calculation – to a restriction on the neutrino mass $`<m_\nu >0.08`$ eV.
## 4 High sensitivity $`2\beta `$ decay study of <sup>116</sup>Cd with the CTF
The most sensitive $`0\nu 2\beta `$ results are obtained by using an ”active” source technique . We recall the highest limits $`T_{1/2}^{0\nu }`$ $``$ (1–2)$`10^{25}`$ yr established for <sup>76</sup>Ge with the help of high purity (HP) enriched <sup>76</sup>Ge detectors , and bounds $`T_{1/2}^{0\nu }`$ $``$ $``$1$`0^{23}`$ yr set for <sup>136</sup>Xe with a high pressure Xe TPC , <sup>130</sup>Te with TeO<sub>2</sub> low temperature bolometers , and <sup>116</sup>Cd with <sup>116</sup>CdWO<sub>4</sub> scintillators .
Continuing this line, we propose to advance the experiment with <sup>116</sup>CdWO<sub>4</sub> to the sensitivity level of $``$10<sup>26</sup> yr by exploiting the advantages of the CTF. The idea is to place $``$65 kg of enriched <sup>116</sup>CdWO<sub>4</sub> crystals in the liquid scintillator of the CTF, which would be used as a light guide and anticoincidence shield for the main <sup>116</sup>CdWO<sub>4</sub> detectors (CAMEO-II project). To prove the feasibility of this task we are considering in the following discussion a pilot <sup>116</sup>Cd experiment, and then the design concept of the present proposal, as well as problems concerning the light collection, energy and spatial resolution, background sources and sensitivity estimates of the CAMEO-II project with <sup>116</sup>Cd.
### 4.1 The pilot <sup>116</sup>Cd study
Here we briefly recall the main results of <sup>116</sup>Cd research performed during the last decade by the INR (Kiev)<sup>7</sup><sup>7</sup>7From 1998 this experiment was carried out by the Kiev-Firenze collaboration . in the Solotvina Underground Laboratory (in a salt mine 430 m underground ), and published elsewhere . The cadmium tungstate crystal scintillators, enriched in <sup>116</sup>Cd to 83%, were grown for research . The light output of this scintillator is relatively large: $``$40% of NaI(Tl) . The refractive index is 2.3. The fluorescence peak emission is at 480 nm with principal decay time of $``$14 $`\mu `$s . The density of CdWO<sub>4</sub> crystal is 7.9 g/cm<sup>3</sup>, and the material is non-hygroscopic and chemically inert.
In the first phase of the study only one <sup>116</sup>CdWO<sub>4</sub> crystal (15.2 cm<sup>3</sup>) was placed inside a veto plastic scintillator and viewed by a PMT through a light-guide 51 cm long. Outer passive shielding consisted of high-purity copper (5 cm), lead (23 cm) and polyethylene (16 cm). The background rate in the energy range 2.7–2.9 MeV ($`Q_{2\beta }`$=2805 keV) was $``$0.6 counts/yr$``$kg$``$keV . With 19175 h statistics the half-life limit for 0$`\nu `$2$`\beta `$ decay of <sup>116</sup>Cd was set as $`T_{1/2}`$(0$`\nu `$) $``$ 3.2$``$1$`0^{22}`$ yr (90% C.L.) , while for neutrinoless 2$`\beta `$ decay with emission of one (M1) or two (M2) Majorons as $`T_{1/2}`$(0$`\nu `$M1) $``$ 1.2$`10^{21}`$ yr and $`T_{1/2}`$(0$`\nu `$M2) $``$ 2.6$`10^{20}`$ yr (90% C.L.) .
In 1998 a new set up with four <sup>116</sup>CdWO<sub>4</sub> crystals (total mass 339 g) was mounted in the Solotvina Laboratory. The enriched detectors are viewed by a special low background 5” EMI tube (with RbCs photocathode) through one light-guide ($``$10$`\times `$55 cm), which is composed of two glued parts: quartz 25 cm long and plastic scintillator (Bicron BC-412) 30 cm long. The main detectors are surrounded by an active shield made of 15 natural CdWO<sub>4</sub> crystals of large volume (total mass 20.6 kg). These are viewed by a low background PMT (FEU-125) through an active plastic light-guide ($``$17$`\times `$49 cm). The whole CdWO<sub>4</sub> array is situated within an additional active shield made of plastic scintillator 40$`\times `$40$`\times `$95 cm, thus, together with active light-guides, complete 4$`\pi `$ active shielding of the main <sup>116</sup>CdWO<sub>4</sub> detectors is provided. The outer passive shield consists of high-purity copper ($`3`$$`6`$ cm), lead ($`22.5`$$`30`$ cm) and polyethylene (16 cm). The set up is isolated carefully against penetration of air which could be contaminated by radon.
The multichannel event-by-event data acquisition is based on two IBM personal computers (PC) and a CAMAC crate. For each event the following information is stored on the hard disk of the first PC: the amplitude (energy), arrival time and additional tags. The second computer records the pulse shape of the <sup>116</sup>CdWO<sub>4</sub> scintillators in the energy range $`0.25`$$`5`$ MeV. It is based on a fast 12 bit ADC (AD9022) and is connected to the PC by a parallel digital I/O board (PC-DIO-24 from National Instruments) . Two PC-DIO-24 boards are used to link both computers and establish – with the help of proper software – a one-to-one correspondence between the pulse shape data recorded by the second computer and the information stored in the first PC.
The energy scale and resolution of the main detector – four enriched crystals taken as a whole – were measured with different sources (<sup>22</sup>Na, <sup>40</sup>K, <sup>60</sup>Co, <sup>137</sup>Cs, <sup>207</sup>Bi, <sup>226</sup>Ra, <sup>232</sup>Th and <sup>241</sup>Am) as $`FWHM(`$keV) = $`44+\sqrt{23.4E+2773}`$, where energy $`E`$ is in keV. In particular, the energy resolution is equal to 11.5% at 1064 keV and 8.0% at 2615 keV. Also the relative light yield and energy resolution for $`\alpha `$ particles were determined as following: $`\alpha /\beta =0.12+1.1`$1$`0^5E_\alpha `$ and $`FWHM_\alpha (`$keV) = $`0.033E_\alpha `$ ($`E_\alpha `$ is in keV). The routine calibration was carried out weekly with a <sup>207</sup>Bi source and once per two weeks with <sup>232</sup>Th. The dead time was monitored permanently by an LED optically connected to the main PMT (actual value $``$4.2$`\%`$).
The background spectrum measured during 4629 h in the new set up with four <sup>116</sup>CdWO<sub>4</sub> crystals is given in Fig. 3, where old data obtained with one <sup>116</sup>CdWO<sub>4</sub> crystal of 121 g are also shown for comparison. The background is lower in the whole energy range, except for the $`\beta `$ spectrum of <sup>113</sup>Cd ($`Q_\beta `$ = $`316`$ keV) <sup>8</sup><sup>8</sup>8Abundance of <sup>113</sup>Cd in enriched <sup>116</sup>CdWO<sub>4</sub> crystals is $``$2% .. In the energy region $`2.5`$$`3.2`$ MeV (location of expected 0$`\nu `$2$`\beta `$ peak) the background rate is 0.03 counts/yr$``$kg$``$keV, twenty times lower than before. It was achieved first, due to new passive and active shielding, and secondly, as a result of the time-amplitude and pulse-shape analysis of the data.
As an example of the time-amplitude technique we consider here in detail the analysis of the following sequence of $`\alpha `$ decays from the <sup>232</sup>Th family: <sup>220</sup>Rn ($`Q_\alpha `$ = $`6.40`$ MeV, $`T_{1/2}`$ = $`55.6`$ s) $``$ <sup>216</sup>Po ($`Q_\alpha `$ = $`6.91`$ MeV, $`T_{1/2}`$ = $`0.145`$ s) $``$ <sup>212</sup>Pb. The electron equivalent energy for <sup>220</sup>Rn $`\alpha `$ particles in <sup>116</sup>CdWO<sub>4</sub> is $``$1.2 MeV, thus all events in the energy region $`0.7`$$`1.8`$ MeV were used as triggers. Then all signals following the triggers in the time interval $`10`$$`1000`$ ms (94.5% of <sup>216</sup>Po decays) were selected. The spectra of the <sup>220</sup>Rn and <sup>216</sup>Po $`\alpha `$ decays obtained in this way from data – as well as the distribution of the time intervals between the first and second events – are in an excellent agreement with those expected from $`\alpha `$ particles of <sup>220</sup>Rn and <sup>216</sup>Po . Using these results and taking into account the efficiency of the time-amplitude analysis and the number of accidental coincidences (3 pairs from 218 selected), the activity of <sup>228</sup>Th (<sup>232</sup>Th family) inside <sup>116</sup>CdWO<sub>4</sub> crystals was determined to be 38(3) $`\mu `$Bq/kg .
The same technique was applied to the sequence of $`\alpha `$ decays from the <sup>235</sup>U family, and yields 5.5(14) $`\mu `$Bq/kg for <sup>227</sup>Ac impurity in the crystals.
The pulse shape (PS) of <sup>116</sup>CdWO<sub>4</sub> events ($`0.25`$$`5`$ MeV) is digitized by a 12-bit ADC and stored in 2048 channels with 50 ns/channel width. Due to different shapes of the scintillation signal for various kinds of sources the PS technique based on the optimal digital filter was developed, and clear discrimination between $`\gamma `$ rays and $`\alpha `$ particles was achieved . In the energy region 4.5– 6 MeV for $`\alpha `$ particles (or 0.8–1.2 MeV for $`\gamma `$ quanta) numerical characteristics of the shape (shape indicator, $`SI`$) are as following: $`SI_\gamma `$= 21.3 $`\pm `$ 2.0, and $`SI_\alpha `$= 32.5 $`\pm `$ 2.9. The PS selection technique ensures the possibility to discriminate ”illegal” events: double pulses, $`\alpha `$ events, etc., and thus suppress the background. For instance, PS selection of the background events, whose $`SI`$ lie in the interval $`SI_\gamma +2.4\sigma _\gamma <`$ $`SI`$ $`<`$ $`SI_\alpha +2.4\sigma _\alpha `$ ($``$90% of $`\alpha `$ events) yields the total $`\alpha `$ activity of <sup>116</sup>CdWO<sub>4</sub> crystals as 1.4(3) mBq/kg. The last value can be adjusted with activities determined by the time-amplitude analysis under the usual assumption that secular radioactive equilibriums in some chains of the <sup>232</sup>Th and <sup>238</sup>U families (e.g. <sup>230</sup>Th $``$ <sup>226</sup>Ra chain) are broken.
Since $`SI`$ characterizes the full signal, it is also useful to examine the pulse front edge. It was found that at least 99% of ”pure” $`\gamma `$ events (measured with calibration <sup>232</sup>Th source) satisfy the following restriction on pulse rise time : $`\mathrm{\Delta }t`$($`\mu `$s) $``$ 1.24 – 0.5$`E_\gamma `$ \+ 0.078$`E_\gamma ^2`$, where $`E_\gamma `$ is in MeV. Hence, the background pulses which do not pass this filter, were excluded from the residual $`\beta /\gamma `$ spectrum.
The results of PS analysis of the data are presented in Fig. 4. The initial (without PS selection) spectrum of the <sup>116</sup>CdWO<sub>4</sub> scintillators in the energy region $`1.2`$$`4`$ MeV – collected during 4629 h in anticoincidence with the active shield – is depicted in Fig. 4a, while the spectrum after PS selection of the $`\beta /\gamma `$ events, whose $`SI`$ lie in the interval $`SI_\gamma 3.0\sigma _\gamma `$ $`SISI_\gamma +2.4\sigma _\gamma `$ and $`\mathrm{\Delta }t`$($`\mu `$s) $``$ 1.24 – 0.5$`E_\gamma `$ \+ 0.078$`E_\gamma ^2`$ (containing 98% of $`\beta /\gamma `$ events), is shown in Fig. 4b. From these figures the background reduction due to pulse-shape analysis is evident. Furthermore, Fig. 4c represents the difference between spectra in Fig. 4a and 4b. These events, at least for energy above 2 MeV, can be produced by <sup>228</sup>Th activity from the intrinsic contamination of the <sup>116</sup>CdWO<sub>4</sub> crystals (measured by the time-amplitude analysis as described above). Indeed, two decays in the fast chain <sup>212</sup>Bi ($`Q_\beta `$= $`2.25`$ MeV) $``$ <sup>212</sup>Po ($`Q_\alpha `$= $`8.95`$ MeV, $`T_{1/2}`$= $`0.3`$ $`\mu `$s) $``$ <sup>208</sup>Pb cannot be time resolved in the CdWO<sub>4</sub> scintillator (with an exponential decay time $``$15 $`\mu `$s , ) and will result in one event. To determine the residual activity of <sup>228</sup>Th in the crystals, the response function of the <sup>116</sup>CdWO<sub>4</sub> detectors for the <sup>212</sup>Bi $``$ <sup>212</sup>Po $``$ <sup>208</sup>Pb chain was simulated with GEANT3.21 and event generator DECAY4. The simulated function is shown in Fig. 4c, from which one can see that the high energy part of the experimental spectrum is well reproduced ($`\chi ^2`$ = $`1.3`$) by the expected response for <sup>212</sup>Bi $`^{212}`$Po $`^{208}`$Pb decays<sup>9</sup><sup>9</sup>9The rest of the spectrum below 1.9 MeV (Fig. 4c) can be explained as the high energy tail of the PS selected $`\alpha `$ particles.. Corresponding activity of <sup>228</sup>Th inside the <sup>116</sup>CdWO<sub>4</sub> crystals, deduced from the fit in the 1.9–3.7 MeV energy region, is 37(4) $`\mu `$Bq/kg, that is in a good agreement with the value determined by the time-amplitude analysis of the chain <sup>220</sup>Rn $``$ <sup>216</sup>Po $``$ <sup>212</sup>Pb 38(3) $`\mu `$Bq/kg.
To estimate the half-life limits for different neutrinoless $`2\beta `$ decay mode, the simple background model was used. In fact, in the $`0\nu 2\beta `$ decay energy region only three background components (presented in fig. 3) are important: (i) external $`\gamma `$ background from U/Th contamination of the PMTs; (ii) the tail of the $`2\nu 2\beta `$ decay spectrum; and (iii) the internal background distribution expected from the <sup>212</sup>Bi $``$ <sup>212</sup>Po $``$ <sup>208</sup>Pb decay (<sup>228</sup>Th chain). The limits for the neutrinoless mode of 2$`\beta `$ decay are set as $`T_{1/2}`$ $`0.7(2.5`$)$``$1$`0^{23}`$ yr at 90%(68%) C.L. (for transition to the ground state of <sup>116</sup>Sn), while for decays to the first 2$`{}_{}{}^{+}{}_{1}{}^{}`$ and second 0$`{}_{}{}^{+}{}_{1}{}^{}`$ excited levels of <sup>116</sup>Sn as $`T_{1/2}`$ $`1.3(4.8`$)$``$1$`0^{22}`$ yr and $`0.7(2.4`$)$``$1$`0^{22}`$ yr at 90%(68%) C.L., accordingly. For $`0\nu `$ decay with emission of one or two Majorons, the limits are: $`T_{1/2}`$(0$`\nu `$M1) $`3.7(5.9`$)$``$1$`0^{21}`$ yr and $`T_{1/2}`$(0$`\nu `$M2) $`5.9(9.4`$)$``$1$`0^{20}`$ yr at 90%(68%) C.L. Also the half-life of <sup>116</sup>Cd two neutrino $`2\beta `$ decay is determined as $`T_{1/2}(2\nu 2\beta )=2.6\pm 0.1`$(stat)$`{}_{}{}^{+0.7}{}_{0.4}{}^{}`$(syst)$`10^{19}`$ yr .
For illustration the high energy part of the experimental spectrum of the <sup>116</sup>CdWO<sub>4</sub> detectors measured during 4629 h (histogram) is shown in Fig. 5 together with the fit by the 2$`\nu `$2$`\beta `$ contribution ($`T_{1/2}`$=2.6$``$1$`0^{19}`$ yr). The smooth curves $`0\nu 2\beta `$M1 and $`0\nu 2\beta `$M2 are 90% C.L. exclusion distributions of M1 and M2 decays of <sup>116</sup>Cd with $`T_{1/2}`$=$`3.`$7$``$1$`0^{21}`$ yr and $`T_{1/2}`$= $`5.`$9$``$1$`0^{20}`$ yr, respectively. In the insert, the peak from $`0\nu 2\beta `$ decay with $`T_{1/2}`$(0$`\nu )`$ = $`1.`$0$``$10<sup>22</sup> yr and 90% C.L. exclusion (solid histogram) with $`T_{1/2}`$(0$`\nu )`$ = $`7.`$0$``$1$`0^{22}`$ yr are depicted for comparison.
The following restrictions on the neutrino mass (using calculations ) and neutrino-Majoron coupling constant (on the basis of calculation ) are derived from the experimental results obtained: $`m_\nu 2.6(1.4)`$ eV and $`g_M12(9.5`$)$``$1$`0^5`$ at 90%(68%) C.L. . It is expected that after $``$5 years of measurements a neutrino mass limit of $`m_\nu 1.2`$ eV would be found. However further advance of this limit to the sub-eV neutrino mass domain is impossible without substantial sensitivity enhancement, which can be reached with a greater number of <sup>116</sup>CdWO<sub>4</sub> detectors ($``$65 kg) placed in the liquid scintillator of the CTF.
### 4.2 Design concept of the CAMEO-II project with <sup>116</sup>Cd
In the preliminary design concept of the CAMEO-II experiment, 24 enriched <sup>116</sup>CdWO<sub>4</sub> crystals of large volume ($``$350 cm<sup>3</sup>) are located in the liquid scintillator of the CTF and fixed at 0.4 m distance from the CTF center, thus homogeneously spread out on a sphere with diameter 0.8 m. With a mass of 2.7 kg per crystal ($``$7$`\times `$9 cm) the total <sup>116</sup>Cd mass is 20 kg ($``$10<sup>26</sup> nuclei of <sup>116</sup>Cd). It is proposed that 200 PMTs with light concentrators are fixed at diameter 5 m providing the total optical coverage of 80% (see footnote 3).
The light output of CdWO<sub>4</sub> scintillator is (35–40)% of NaI(Tl) which yields $``$1.5$``$10<sup>4</sup> emitted photons per MeV of energy deposited. With a total light collection of $``$80% and PMT quantum efficiency of $``$25% energy resolution of several % at 1 MeV can be obtained.
To justify this value a GEANT Monte Carlo simulation of the light propagation in this geometry was performed, which gave $``$4000 photoelectrons for 2.8 MeV energy deposit. Therefore, with total optical coverage 80% the neutrinoless $`2\beta `$ decay peak of <sup>116</sup>Cd ($`Q_{\beta \beta }`$ $`=`$ 2805 keV) can be measured with energy resolution $`FWHM`$ = 4%. The principal feasibility to obtain such an energy resolution with CdWO<sub>4</sub> crystal situated in a liquid has been successfully demonstrated by measurements with the help of a simple device. A cylindrical CdWO<sub>4</sub> crystal (40 mm in diameter and 30 mm in height) was fixed in the centre of a teflon container with inner diameter 70 mm. This was coupled on opposite sides with two PMTs Philips XP2412, so that the distance from each flat surface of crystal to the corresponding PMT’s photocathode is 30 mm, while the gap between the side surface of the crystal and the inner surface of the container is 15 mm. The container was filled up with the pure and transparent paraffin oil (refractive index $``$1.5). Two PMTs work in coincidence and results of measurements with <sup>207</sup>Bi source are depicted in Fig. 6, where the spectrum obtained with the standard detector arrangement (CdWO<sub>4</sub> crystal wrapped by teflon diffuser and directly coupled to the PMT’s photocathode with optical glue) is also shown for comparison. As evident from Fig. 6, a substantial ($``$42%) increase in light collection was obtained with CdWO<sub>4</sub> in the liquid. It resulted in improvement of the detector energy resolution in the whole energy region 50–3000 keV (see Fig. 7 where the spectra measured with <sup>137</sup>Cs, <sup>60</sup>Co and <sup>232</sup>Th sources are presented). It should be stressed that $`FWHM`$ values (7.4% at 662 keV; 5.8% at 1064 keV; 5.4% at 1173 keV and 4.3% at 2615 keV) are similar to those for NaI(Tl) crystals and have never been reached before with CdWO<sub>4</sub> crystal scintillators .
Moreover, a strong dependence of the light collected by each individual PMT versus position of the emitting source in the crystal was found. Such a dependence can be explained by the large difference of the refraction indexes of CdWO<sub>4</sub> crystal ($`n=2.3`$) and liquid scintillator ($`n^{}=1.51`$), which leads to light redistribution between reflection and refraction processes due to changes of the source’s position.
General formulae for the angular distribution of the light emitted in the crystal and propagating in the liquid scintillator are quite cumbersome, and below we give expressions for some simple cases, when a CdWO<sub>4</sub> crystal with equal diameter and height ($`d=h=2a`$) is placed in the center of the CTF detector and a light source is positioned on the crystal axis. Assuming a ratio of refraction indexes $`n/n^{}`$ = $`\sqrt{2}`$ (which is close to the real value) and neglecting light absorption, the equation for the particular case with light source in the center of crystal is of the form:
$$\frac{dW}{d\mathrm{cos}\theta }=\frac{1}{2\sqrt{2}}[\frac{\left|\mathrm{cos}\theta \right|}{\sqrt{1+\mathrm{cos}^2\theta }}+1],$$
(3)
where $`\theta `$ is the angle between the axis of the crystal (z-axis in our coordinate system) and direction of the photon in the liquid scintillator. This function is depicted in Fig. 8a together with the distribution for another case with the light source on the bottom of crystal (on the axis again):
$$\frac{dW}{d\mathrm{cos}\theta }=\frac{1}{2\sqrt{2}}[\frac{\left|\mathrm{cos}\theta \right|}{\sqrt{1+\mathrm{cos}^2\theta }}+2\mathrm{\Theta }(\theta )],$$
(4)
where $`\mathrm{\Theta }(\theta )`$ is the unit step function depending on angle $`\theta `$, as following:
$`\mathrm{\Theta }(\theta )=\{\begin{array}{c}1,0\theta \pi /2\\ 0,\pi /2<\theta \pi \end{array}.`$
The angular distribution for a more general case, when the light source is positioned on an arbitrary point of the axis with coordinate $`z`$ ($`aza`$), is shown in Fig. 8b. When the location of the source is shifted from $`z`$ to $`z+\mathrm{}z`$, the values of angles $`\theta _z`$ and $`\pi \theta _z`$, for which $`dW/d\mathrm{cos}\theta `$ changes sharply, are also changed: $`\mathrm{cos}\theta _z=\sqrt{2(a\left|z\right|)^2/(a^2+(a\left|z\right|)^2)}`$. The difference $`dW(z+\mathrm{}z)/d\mathrm{cos}\theta dW(z)/d\mathrm{cos}\theta `$ represented in Fig. 8c behaves as two peaks. The area of both peaks is equal to $`\mathrm{}S=a^2\mathrm{}z/\left[a^2+(a\left|z\right|)^2\right]^{3/2}`$. Assuming that $`N`$ is the full number of photoelectrons detected by all PMTs, and $`\mathrm{}N`$ is the difference of the number of photoelectrons due to the source shift $`\mathrm{}z`$, we find $`\mathrm{}N=Na^2\mathrm{}z/\left[a^2+(a\left|z\right|)^2\right]^{3/2}`$. From the last expression one gets the formula for the spatial resolution in the CdWO<sub>4</sub> crystal by supposing that difference $`\mathrm{}N`$, which can be registered with 95% C.L., is equal $`\mathrm{}N=2\sqrt{N}`$:
$$\mathrm{}z=\frac{2a}{\sqrt{N}}\left\{1+\left[1\frac{\left|z\right|}{a}\right]^2\right\}^{3/2}.$$
(5)
Substituting in Eq. (5) the crystal’s dimensions $`a=4`$ cm, and the total number of photoelectrons $`N=4000`$ yields spatial resolution of $`4`$ mm in the center of the crystal ($`z=0`$) and $`1.5`$ mm near the top or bottom of the crystal ($`z=\pm a`$). Our GEANT Monte Carlo simulation proves these values and shows that, with a cylindrical CdWO<sub>4</sub> crystal ($`7\times 9`$ cm) viewed by 200 PMTs, spatial resolution of 1–5 mm can be reached depending on the event’s location and the energy deposited in the crystal (see, however, footnote 4). The simulated distributions of the photoelectron number among PMTs due to scintillations in CdWO<sub>4</sub> crystal are depicted in Fig. 9. The distance between the source’s positions on the crystal axis is equal to $`\mathrm{}z`$ = 2.5 mm (Fig. 9a), and $`\mathrm{}z`$ = 1.5 cm (Fig. 9b).
These interesting features of light collection from <sup>116</sup>CdWO<sub>4</sub> in the CTF would allow a reduction in the contribution from high energy $`\gamma `$ quanta (e.g. <sup>208</sup>Tl) to background in the energy region of interest. Besides, the non-uniformity of light collection can be accurately corrected by using spatial information about each event in the CdWO<sub>4</sub> crystal, hence, helping to reach the required energy resolution of the detectors.
### 4.3 Background simulation
The background simulations for the CAMEO-II experiment with <sup>116</sup>Cd were performed with the GEANT code and event generator DECAY4, and by using the model described above.
#### 4.3.1 Two neutrino 2$`\beta `$ decays of <sup>116</sup>Cd
The half-life of two neutrino 2$`\beta `$ decay of <sup>116</sup>Cd has been measured as $``$2.7$``$10<sup>19</sup> yr . The response functions of the CAMEO-II set up for $`2\nu 2\beta `$ decay of <sup>116</sup>Cd with $`T_{1/2}`$ $`=`$ 2.7$``$10<sup>19</sup> yr, as well as for $`0\nu 2\beta `$ decay with $`T_{1/2}`$ $`=`$ 10<sup>25</sup> yr were simulated, and results are depicted in Fig. 2b. The calculated values of efficiency for the neutrinoless channel are 86.1% (for energy window 2.7 – 2.9 MeV), and 75.3% (2.75 – 2.9 MeV). Background in the corresponding energy interval from $`2\nu 2\beta `$ decay distribution is 2.3 counts/yr (2.7 – 2.9 MeV) or 0.29 counts/yr (2.75 – 2.9 MeV).
#### 4.3.2 Radioactive contamination of <sup>116</sup>CdWO<sub>4</sub> crystals
The very low levels of radioactive impurities by <sup>40</sup>K and nuclides from natural radioactive chains of <sup>232</sup>Th and <sup>238</sup>U in the enriched and natural CdWO<sub>4</sub> crystals were demonstrated by the INR (Kiev) experiment . On this basis the contamination criterion for <sup>214</sup>Bi and <sup>208</sup>Tl has been accepted in our calculation as $``$10 $`\mu `$Bq/kg, which is equal to the actual activity value or limit determined for different samples of CdWO<sub>4</sub> crystals .
The calculated background contribution from the sum of <sup>208</sup>Tl and <sup>214</sup>Bi activities is $``$2000 counts/yr in the energy interval 2.7 – 2.9 MeV. However, applying the time-amplitude analysis with spatial resolution and pulse-shape discrimination technique developed this background rate can be reduced to $``$0.2 counts/yr or $``$1.0 counts for 5 years measuring period.
#### 4.3.3 Cosmogenic activities in <sup>116</sup>CdWO<sub>4</sub>
For the CAMEO-II project cosmogenic activities in the <sup>116</sup>CdWO<sub>4</sub> detectors were calculated by the program COSMO . A 1 month exposure period on the Earth’s surface was assumed and a deactivation time of about three years in the underground laboratory. Only two nuclides produce background in the energy window of the <sup>116</sup>Cd neutrinoless $`2\beta `$ decay. These are <sup>110m</sup>Ag ($`Q_\beta `$=3.0 MeV; $`T_{1/2}`$=250 d) and <sup>106</sup>Ru ($`Q_\beta `$ 40 keV; $`T_{1/2}`$=374 d) $``$ <sup>106</sup>Rh ($`Q_\beta `$=3.5 MeV; $`T_{1/2}`$=30 s). Fortunately <sup>106</sup>Ru activity is low and the time-amplitude analysis can be applied ($`T_{1/2}`$=30 s), its estimated background is practically negligible: $``$0.1 counts/yr in the energy region 2.7 – 2.9 MeV, and 0.05 counts/yr (2.75 – 2.9 MeV). The background from <sup>110m</sup>Ag is quite large: $``$23 (or $``$20) counts/yr for the energy interval 2.7 – 2.9 MeV (2.75 – 2.9 MeV). However its contribution can be reduced significantly by using spatial information because <sup>110m</sup>Ag decays are accompanied by cascades of $`\gamma `$ quanta with energies $``$ 600 keV, which would be absorbed in spatially separated parts of the detector giving an anticoincidence signature. Simulation under the assumption that the <sup>116</sup>CdWO<sub>4</sub> crystal consists of small independent detectors with $`h=d=1.2`$ cm, yields the residual background rates $``$0.3 (or 0.2) counts/yr in the corresponding energy region 2.7 – 2.9 MeV (2.75 – 2.9 MeV). The simulated spectrum from the cosmogenic activity of <sup>110m</sup>Ag is depicted in Fig. 2b.
#### 4.3.4 External background
As in the previous case with <sup>100</sup>Mo from the various sources of external background only $`\gamma `$ quanta due to PMT contamination and from Rn impurities in the shielding water were simulated, while others were simply estimated as negligible on the basis of the results of ref. . The radioactivity values of the EMI 9351 PMT accepted for the simulation are: 0.194 Bq/PMT for <sup>208</sup>Tl; 1.383 Bq/PMT for <sup>214</sup>Bi; and 191 Bq/PMT for <sup>40</sup>K . Also possible <sup>222</sup>Rn activity in the shielding water (in the region surrounding the inner vessel) at the level of $``$30 mBq/m<sup>3</sup> was taken into account. The simulation performed under these assumptions finds that the only important contribution to the background in the vicinity of the $`0\nu 2\beta `$ decay energy is <sup>208</sup>Tl activity from the PMTs. The calculated values are $``$0.8 and 0.05 counts/yr in the energy interval 2.7 – 2.9 MeV (2.75 – 2.9 MeV). However, with the help of spatial information available for each event occurring inside the <sup>116</sup>CdWO<sub>4</sub> crystal, these contributions can be reduced further to the level of $``$0.08 (or 0.005) counts/yr in the energy region 2.7 – 2.9 MeV (2.75 – 2.9 MeV). The simulated background contribution from <sup>208</sup>Tl contamination of the PMTs is shown in Fig. 2b. Summarizing all background sources gives $``$3 counts/yr (0.6 counts/yr) in the energy interval 2.7 – 2.9 MeV (2.75 – 2.9 MeV).
### 4.4 Sensitivity of the <sup>116</sup>Cd experiment
As earlier we will estimate the sensitivity of the CAMEO-II experiment with the help of Eq. (2). Taking into account the number of <sup>116</sup>Cd nuclei ($``$10<sup>26</sup>), measuring time of 5–8 years, detection efficiency of 75%, and with expected background of 3–4 counts, one can obtain a half-life limit $`T_{1/2}`$(0$`\nu 2\beta `$) $``$ 10<sup>26</sup> yr. On the other hand, it is evident from Fig. 2b that neutrinoless $`2\beta `$ decay of <sup>116</sup>Cd with half-life of $``$10<sup>25</sup> yr would be clearly registered.
It should be stressed that such a level of sensitivity for 0$`\nu 2\beta `$ decay cannot be reached in the presently running $`2\beta `$ decay experiments (perhaps only with <sup>76</sup>Ge), as well as for approved projects, like NEMO-3 and CUORICINO , which are under construction now. It was shown above that the sensitivity of the NEMO-3 set up is limited at the level of $``$4$``$10<sup>24</sup> yr by the detection efficiency and energy resolution (see Fig. 1f). The CUORICINO project is designed to study $`2\beta `$ decay of <sup>130</sup>Te with the help of 60 low temperature bolometers made of TeO<sub>2</sub> crystals (750 g mass each). Another aim of CUORICINO is to be a pilot step for a future CUORE project (not approved yet), which would consists of one thousand TeO<sub>2</sub> bolometers with total mass of 750 kg . Despite the excellent energy resolution of these detectors ($``$5 keV at 2.5 MeV) the main disadvantage of the cryogenic technique is its complexity, which requires the use of a lot of different construction materials in the apparatus. This fact, together with the lower $`2\beta `$ decay energy of <sup>130</sup>Te ($`Q_{\beta \beta }`$=2528 keV), makes it quite difficult to reach the same super-low level of background as obtained in experiments with semiconductor and scintillation detectors. For example, one can compare the background rate of 0.6 counts/yr$``$kg$``$keV at 2.5 MeV from the current Milano experiment with twenty TeO<sub>2</sub> bolometers with the value of 0.03 counts/yr$``$kg$``$keV in the energy region $`2.5`$$`3.2`$ MeV, which was reached in the Kiev-Firenze experiment with <sup>116</sup>CdWO<sub>4</sub> crystal scintillators . In that sense the CAMEO-II experiment has a great fundamental advantage because signaling from the <sup>116</sup>CdWO<sub>4</sub> crystals to the PMTs (placed far away from crystals) is provided by light propagating in the super-low background medium of liquid scintillator, whereas cryogenic or semiconductor detectors must be connected with receiving modules by cables. These additional materials (wires, insulators, etc.), whose radioactive contamination are much larger in comparison with TeO<sub>2</sub> crystals, Ge detectors or liquid scintillators, must be introduced in the neighborhood of the main detectors, giving rise to additional background<sup>10</sup><sup>10</sup>10There are two origins of such a background: i) radioactive contamination of the materials introduced; ii) external background penetrating through the slots in the detector shielding required for the connecting cables..
Another drawback of cryogenic detectors is their low reliability. At the same time, the CAMEO-II experiment with <sup>116</sup>CdWO<sub>4</sub> crystals is simple and reliable, and therefore can run for decades without problems and with very low maintenance cost<sup>11</sup><sup>11</sup>11It should be noted that the <sup>116</sup>CdWO<sub>4</sub> crystals produced for the CAMEO-II experiment can also be used as cryogenic detectors with high energy resolution . In the event of a positive effect seen by CAMEO-II these crystals could be measured by the CUORE apparatus; in some sense both projects are complementary..
Moreover, the CAMEO-II project can be advanced farther by exploiting one ton of <sup>116</sup>CdWO<sub>4</sub> detectors ($``$1.5$`1`$0<sup>27</sup> nuclei of <sup>116</sup>Cd) and the BOREXINO apparatus (CAMEO-III). With this aim 370 enriched <sup>116</sup>CdWO<sub>4</sub> crystals (2.7 kg mass of each) would be placed at a diameter 3.2 m in the BOREXINO liquid scintillator. The simulated response functions of such a detector system for $`2\nu 2\beta `$ decay of <sup>116</sup>Cd with $`T_{1/2}`$ $`=`$ 2.7$``$10<sup>19</sup> yr, as well as for $`0\nu 2\beta `$ decay with $`T_{1/2}`$ $`=`$ 10<sup>26</sup> yr considering a 10-year measuring period are depicted in Fig. 2c. Because background in BOREXINO should be even lower than in the CTF, the sensitivity of CAMEO-III for neutrinoless $`2\beta `$ decay of <sup>116</sup>Cd is estimated as $`T_{1/2}`$ $``$ 10<sup>27</sup> yr, while $`0\nu 2\beta `$ decay with half-life of $``$10<sup>26</sup> yr can be detected. This level of sensitivity can be compared only with that of the GENIUS project , which is under discussion now and intends to operate one ton of Ge (enriched in <sup>76</sup>Ge) semiconductor detectors placed in a tank ($`12\times 12`$ m) with extremely high-purity liquid nitrogen (required demands on its radioactive contamination are $`10^{15}`$ g/g for <sup>40</sup>K, <sup>238</sup>U, <sup>232</sup>Th, and 0.05 mBq/m<sup>3</sup> for <sup>222</sup>Rn ) serving as cooling medium and shielding for the detectors simultaneously. Let us estimate the sensitivity of GENIUS in the same way as that for CAMEO. One ton of Ge detectors with enrichment $``$86% (as in the current <sup>76</sup>Ge experiments ) would provide $``$7$``$10<sup>27</sup> nuclei of <sup>76</sup>Ge, thus in the optimistic case of zero background a sensitivity of $`T_{1/2}^{0\nu }`$ $``$ 5$``$10<sup>27</sup> yr can be reached there.
By the aid of Eq. (1) one can obtain an expression for the neutrino mass bound derived from the experimental half-life limit for 0$`\nu 2\beta `$ decay as $`limm_\nu =\left\{limT_{1/2}^{0\nu }G_{mm}^{0\nu }\left|ME\right|^2\right\}^{1/2}.`$ Because of the lower $`2\beta `$ decay energy the phase-space integral $`G_{mm}^{0\nu }`$ for 0$`\nu 2\beta `$ decay of <sup>76</sup>Ge is about ten times lower than for <sup>116</sup>Cd. Hence, it is evident from the last equation that neglecting the complicated problem of nuclear matrix elements calculation, the CAMEO-III experiment will bring at least the same restriction on the neutrino mass as the GENIUS project. Indeed, on the basis of the CAMEO half-life limit $`T_{1/2}^{0\nu }`$ $``$ 10<sup>27</sup> yr and using calculations one can derive a limit on the neutrino mass of $`0.02`$ eV at 90% C.L., which is practically equal to the value $`0.01`$ eV claimed as the main goal of GENIUS . At the same time it is obvious that the technical tasks, whose solutions are required for the realization of these super–high sensitive projects (GENIUS, CUORE, and CAMEO) are simpler for CAMEO. In fact, the super-low background apparatus needed for the last experiment is already running (this is the CTF) or under construction (BOREXINO), while for the CUORE and GENIUS proposals such unique apparati should be designed and constructed.
## 5 Conclusions
1. The unique features of the CTF and BOREXINO (super-low background and large sensitive volume) are used to develop a realistic, competitive, and efficient program for high sensitivity 2$`\beta `$ decay research (CAMEO project). This program includes three natural steps, and each of them would bring substantial physical results:
CAMEO-I. With a passive 1 kg source made of <sup>100</sup>Mo (<sup>116</sup>Cd, <sup>82</sup>Se, <sup>150</sup>Nd) and located in the liquid scintillator of the CTF, the sensitivity (in terms of the half-life limit for $`0\nu 2\beta `$ decay) is (3–5)$``$10<sup>24</sup> yr. It corresponds to a bound on the neutrino mass $`m_\nu `$ 0.1–0.3 eV, which is similar to or better than those of running (<sup>76</sup>Ge), and future NEMO-3 (<sup>100</sup>Mo) and CUORICINO (<sup>130</sup>Te) experiments.
CAMEO-II. With 24 enriched <sup>116</sup>CdWO<sub>4</sub> crystal scintillators (total mass of 65 kg) placed as ”active” detectors in the liquid scintillator of the CTF the sensitivity would be $``$10<sup>26</sup> yr. Such a half-life limit could be obtained only by future CUORE (<sup>130</sup>Te) and GENIUS (<sup>76</sup>Ge) projects. Pilot <sup>116</sup>Cd research, performed by the INR (Kiev) during the last decade, as well as Monte Carlo simulation show the feasibility of CAMEO-II, which will yield a limit on the neutrino mass $`m_\nu `$ 0.05–0.07 eV.
CAMEO-III. By exploiting one ton of <sup>116</sup>CdWO<sub>4</sub> detectors (370 enriched <sup>116</sup>CdWO<sub>4</sub> crystals) introduced in the BOREXINO liquid scintillator, the half-life limit can be advanced to the level of $``$10<sup>27</sup> yr, corresponding to a neutrino mass bound of $`0.02`$ eV.
2. In contrast to other projects CAMEO has three principal advantages:
i) Practical realization of the CAMEO project is simpler due to the use of already existing super-low background CTF or (presently under construction) BOREXINO apparatus;
ii) Signaling from the <sup>116</sup>CdWO<sub>4</sub> crystals to PMT (placed far away) is provided by light propagating in the high-purity medium of liquid scintillator – this allows practically zero background to be reached in the energy region of the 0$`\nu 2\beta `$ decay peak;
iii) Extreme simplicity of the technique used for 2$`\beta `$ decay study leads to high reliability and low maintenance costs for the CAMEO experiment, which therefore can run permanently and stably for decades.
3. Fulfillment of the CAMEO program would be a real breakthrough in the field of $`2\beta `$ decay investigation, and will bring outstanding results for particle physics, cosmology and astrophysics. Discovery of neutrinoless $`2\beta `$ decay will clearly and unambiguously manifest new physical effects beyond the Standard Model. In the event of a null result, the limits obtained by the CAMEO experiments would yield strong restrictions on parameters of manifold extensions of the SM (neutrino mass and models; dark matter and solar neutrinos; right-handed contributions to weak interactions; leptoquark masses; bounds for parameter space of SUSY models; neutrino-Majoron coupling constant; composite heavy neutrinos; Lorentz invariance, etc.), which will help to advance basic theory and our understanding of the origin and evolution of the Universe.
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# A tensor form of the Dirac equation
## Introduction
The Dirac equation for an electron can be written in several different but equivalent forms. In this paper we consider two forms of the Dirac equation. Namely Hestenes’ form of the Dirac Equation (HDE) and the Tensor form of the Dirac Equation (TDE). The aim of this paper is to prove the following theorem.
Theorem 1.The Dirac equation for an electron (2) (see ) can be written in a form of linear tensor equation (58).
An equation is said to be written in a form of tensor equation if all values in it are tensors and all operations in it take tensors to tensors.
A column of four complex valued functions (a bispinor) represents the wave function of the electron in the Dirac equation. This wave function has unusual (compared to tensors) transformation properties under Lorentz changes of coordinates. These properties are investigated in the theory of spinors (see ,,,,).
There were attempts to find tensor equations equivalent to the Dirac equation ,,. The resulting equations were nonlinear. This fact leads to difficulties with the superposition principle, etc. The tensor equation under consideration in our paper is linear.
In the first part of paper we consider the Clifford algebra $`𝒞\mathrm{}(1,3)`$ and in the second part we consider the exterior algebra of Minkowski space $`\mathrm{\Lambda }()`$. Elements of $`\mathrm{\Lambda }()`$ are covariant antisymmetric tensors and elements of $`𝒞\mathrm{}(1,3)`$ are not tensors.
The tensor form of the Dirac equation (58) under consideration has three sources. The first source is the Ivanenko-Landau-Kähler equation (65) (see ,). This is a tensor equation and a wave function of the electron is represented in it by a nonhomogeneous covariant antisymmetric tensor field with 16 complex valued components (four times more than in the Dirac equation).
The second source is Hestenes’ form of the Dirac equation (46) (see ,). The wave function of the electron is represented in it by a real even element of the Clifford algebra $`𝒞\mathrm{}^{\mathrm{even}}(1,3)`$ and has 8 real valued components. This equation is equivalent to the Dirac equation (see the proof in ). HDE contains nontensor values namely elements of the Clifford algebra $`𝒞\mathrm{}(1,3)`$. Hence it is not a tensor equation.
The third source is the constructions developed in , ,, .
Also we want to mention an interesting approach to the construction of the 2D Dirac tensor equation suggested by D.Vassiliev .
## 1 Part I.
### 1.1 The Dirac equation for an electron.
Let $``$ be Minkowski space with the metric tensor
$$g=g_{\mu \nu }=g^{\mu \nu }=\mathrm{diag}(1,1,1,1),$$
(1)
$`x^\mu `$ coordinates, and $`_\mu =/x^\mu `$. Greek indices run over $`(0,1,2,3)`$. Summation convention over repeating indices is assumed. Our system of units is such that speed of light $`c`$, Plank’s constant $`\mathrm{}`$, and the positron charge have the value $`1`$.
The Dirac equation for an electron has the form
$$\gamma ^\mu (_\mu \psi +ia_\mu \psi )+im\psi =0,$$
(2)
where $`\psi =(\psi _1\psi _2\psi _3\psi _4)^T`$ is a column of four complex valued functions of $`x=(x^0,x^1,x^2,x^3)`$ ($`\psi `$ is the wave function of an electron), $`a_\mu =a_\mu (x)`$ is a real valued covector of electromagnetic potential, $`m>0`$ is a real constant (the electron mass), and $`\gamma ^\mu `$ are the Dirac $`\gamma `$-matrices
$`\gamma ^0`$ $`=`$ $`\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& 1\end{array}\right),\gamma ^1=\left(\begin{array}{cccc}0& 0& 0& 1\\ 0& 0& 1& 0\\ 0& 1& 0& 0\\ 1& 0& 0& 0\end{array}\right),`$ (3)
$`\gamma ^2`$ $`=`$ $`\left(\begin{array}{cccc}0& 0& 0& i\\ 0& 0& i& 0\\ 0& i& 0& 0\\ i& 0& 0& 0\end{array}\right),\gamma ^3=\left(\begin{array}{cccc}0& 0& 1& 0\\ 0& 0& 0& 1\\ 1& 0& 0& 0\\ 0& 1& 0& 0\end{array}\right).`$
The equation (2) is invariant under the gauge transformation
$$\psi \psi \mathrm{exp}(i\lambda ),a_\mu a_\mu _\mu \lambda ,$$
(4)
where $`\lambda =\lambda (x)`$ is a smooth real valued function.
The matrices (3) satisfy the relations
$$\gamma ^\mu \gamma ^\nu +\gamma ^\nu \gamma ^\mu =2\eta ^{\mu \nu }\mathrm{𝟏},$$
(5)
where $`\mathrm{𝟏}`$ is the $`4\times 4`$-identity matrix and
$$\eta ^{\mu \nu }=\mathrm{diag}(1,1,1,1).$$
Note that if $`4\times 4`$-matrix $`T`$ is invertible, then the matrices $`\stackrel{´}{\gamma }^\mu =T^1\gamma ^\mu T`$ also satisfy (5).
Let us denote
$$\gamma ^{\mu _1\mathrm{}\mu _k}=\gamma ^{\mu _1}\mathrm{}\gamma ^{\mu _k},\text{for}0\mu _1<\mathrm{}<\mu _k3.$$
(6)
Then the 16 matrices
$$\mathrm{𝟏},\gamma ^0,\gamma ^1,\gamma ^2,\gamma ^3,\gamma ^{01},\gamma ^{02},\gamma ^{03},\gamma ^{12},\gamma ^{13},\gamma ^{23},\gamma ^{012},\gamma ^{013},\gamma ^{023},\gamma ^{123},\gamma ^{0123}$$
(7)
form a basis of the matrix algebra $`(4,𝒞)`$ ($`𝒞`$ is the field of complex numbers and $``$ is the field of real numbers).
Remark. We assume that the matrix $`\eta ^{\mu \nu }`$ is independent of $`g`$ from (1). And more than this, the values $`g^{\mu \nu }`$ and $`\eta ^{\mu \nu }`$ have different mathematical meaning. Namely metric tensor $`g^{\mu \nu }`$ is a geometrical object which is attribute of Minkowski space. But the matrix $`\eta ^{\mu \nu }`$ is an algebraic object which contains structure constants of an underlying algebra (as we see in a moment this algebra is the Clifford algebra). Accidentally, the matrices $`\eta ^{\mu \nu }`$ and $`g^{\mu \nu }`$ have the same diagonal form $`\mathrm{diag}(1,1,1,1)`$.
### 1.2 The Clifford algebra and the spinor group.
Let $``$ be a real 16-dimensional vector space with basis elements enumerated by ordered multi-indices
$$\mathrm{},\mathrm{}^0,\mathrm{}^1,\mathrm{}^2,\mathrm{}^3,\mathrm{}^{01},\mathrm{}^{02},\mathrm{}^{03},\mathrm{}^{12},\mathrm{}^{13},\mathrm{}^{23},\mathrm{}^{012},\mathrm{}^{013},\mathrm{}^{023},\mathrm{}^{123},\mathrm{}^{0123}.$$
(8)
Suppose the multiplication of elements of $``$ is defined by the following rules:
$``$ is an associative algebra (with the unity element $`\mathrm{}`$) w.r.t. this multiplication;
$`\mathrm{}^\mu \mathrm{}^\nu +\mathrm{}^\nu \mathrm{}^\mu =2\eta ^{\mu \nu }\mathrm{}`$;
$`\mathrm{}^{\mu _1}\mathrm{}\mathrm{}^{\mu _k}=\mathrm{}^{\mu _1\mathrm{}\mu _k},\text{for}0\mu _1<\mathrm{}<\mu _k3`$.
Then this algebra $``$ is called the (real) Clifford algebra<sup>1</sup><sup>1</sup>1The Clifford algebra was invented in 1878 by the English mathematician W.K.Clifford and is denoted by $`𝒞\mathrm{}(1,3)`$, where the numbers $`1`$ and $`3`$ determine the signature of the matrix $`\eta ^{\mu \nu }`$. The complex Clifford algebra is denoted by $`𝒞\mathrm{}_𝒞(1,3)`$.
Elements of $`𝒞\mathrm{}(1,3)`$ of the form
$$U=\underset{\mu _1<\mathrm{}<\mu _k}{}u_{\mu _1\mathrm{}\mu _k}\mathrm{}^{\mu _1\mathrm{}\mu _k}$$
(9)
are said to be elements of rank $`k`$. For every $`k=0,1,2,3,4`$ the set of elements of rank $`k`$ is a subspace $`𝒞\mathrm{}^k(1,3)`$ of $`𝒞\mathrm{}(1,3)`$ and
$$𝒞\mathrm{}(1,3)=𝒞\mathrm{}^0(1,3)\mathrm{}𝒞\mathrm{}^4(1,3)=𝒞\mathrm{}^{\mathrm{even}}(1,3)𝒞\mathrm{}^{\mathrm{odd}}(1,3),$$
where
$`𝒞\mathrm{}^{\mathrm{even}}(1,3)`$ $`=`$ $`𝒞\mathrm{}^0(1,3)𝒞\mathrm{}^2(1,3)𝒞\mathrm{}^4(1,3),`$
$`𝒞\mathrm{}^{\mathrm{odd}}(1,3)`$ $`=`$ $`𝒞\mathrm{}^1(1,3)𝒞\mathrm{}^3(1,3).`$
The dimensions of the spaces $`𝒞\mathrm{}^k(1,3),k=0,1,2,3,4`$ are equal to $`1,4,6,4,1`$ respectively and the dimensions of $`𝒞\mathrm{}^{\mathrm{even}}(1,3)`$ and $`𝒞\mathrm{}^{\mathrm{odd}}(1,3)`$ are equal to $`8`$. Four elements of $`𝒞\mathrm{}(1,3)`$ are called generators of the Clifford algebra if any element of $`𝒞\mathrm{}(1,3)`$ can be represented as a linear combination of products of these generators. Four generators $`L^\mu 𝒞\mathrm{}(1,3)`$ are said to be primary generators if $`L^\mu L^\nu +L^\nu L^\mu =2\eta ^{\mu \nu }\mathrm{}`$.
Hence, the basis elements $`\mathrm{}^\mu `$ are primary generators of the Clifford algebra $`𝒞\mathrm{}(1,3)`$.
Let us define the trace of a Clifford algebra element as a linear operation $`\mathrm{Tr}:𝒞\mathrm{}`$ or $`\mathrm{Tr}:𝒞\mathrm{}_𝒞𝒞`$ such that $`\mathrm{Tr}(\mathrm{})=1`$ and $`\mathrm{Tr}(\mathrm{}^{\mu _1\mathrm{}\mu _k})=0,k=1,2,3,4`$. The reader can easily prove that
$$\mathrm{Tr}(UVVU)=0,\mathrm{Tr}(V^1UV)=\mathrm{Tr}U,U,V𝒞\mathrm{}_𝒞.$$
For
$$U=\underset{\mu _1<\mathrm{}<\mu _k}{}u_{\mu _1\mathrm{}\mu _k}\mathrm{}^{\mu _1\mathrm{}\mu _k}𝒞\mathrm{}_𝒞^k(1,3)$$
we may define an involution $`:𝒞\mathrm{}_𝒞^k𝒞\mathrm{}_𝒞^k`$, $`k=0,\mathrm{},4`$ by
$$U^{}=\underset{\mu _1<\mathrm{}<\mu _k}{}\overline{u}_{\mu _1\mathrm{}\mu _k}\mathrm{}^{\mu _k}\mathrm{}\mathrm{}^{\mu _1},$$
(10)
where the bar means complex conjugation. For $`U𝒞\mathrm{}^k(1,3)`$ we have
$$U^{}=(1)^{\frac{k(k1)}{2}}U.$$
It is readily seen that
$$U^{}=U,(UV)^{}=V^{}U^{},U,V𝒞\mathrm{}(1,3).$$
(11)
Let us define the group with respect to multiplication
$$\mathrm{Spin}(1,3)=\{S𝒞\mathrm{}^{\mathrm{even}}(1,3):S^{}S=\mathrm{}\},$$
which is called the spinor group. For any $`F𝒞\mathrm{}^{\mathrm{even}}(1,3)`$ consider the linear operator $`G_F:𝒞\mathrm{}(1,3)𝒞\mathrm{}(1,3)`$ such that
$$G_F(U)=F^{}UF.$$
(12)
In the sequel we use the following well known propositions (see proofs in ).
If $`T𝒞\mathrm{}^{\mathrm{even}}(1,3)`$ or $`T𝒞\mathrm{}^{\mathrm{odd}}(1,3)`$, then
$`G_T`$ $`:`$ $`𝒞\mathrm{}^k(1,3)𝒞\mathrm{}^k(1,3),k=1,2,3;`$
$`G_T`$ $`:`$ $`𝒞\mathrm{}^0(1,3)𝒞\mathrm{}^4(1,3)𝒞\mathrm{}^0(1,3)𝒞\mathrm{}^4(1,3).`$
If $`S\mathrm{Spin}(1,3)`$, then
$$G_S:𝒞\mathrm{}^k(1,3)𝒞\mathrm{}^k(1,3),k=0,1,2,3,4.$$
and
$$S^{}\mathrm{}^\nu S=p_\mu ^\nu \mathrm{}^\mu ,\text{for}S\mathrm{Spin}(1,3),$$
(13)
where the matrix $`P=p_\mu ^\nu `$ satisfies the relations
$$P^TgP=g,\mathrm{det}P=1,p_0^0>0.$$
(14)
Therefore if we transform the coordinate system
$$\stackrel{~}{x}^\nu =p_\mu ^\nu x^\mu ,$$
(15)
with the aid of this matrix $`P=P(S)`$, then we get the proper orthochronous Lorentz transformation from the group $`\mathrm{SO}^+(1,3)`$. Conversely, if some matrix $`P`$ specifies a transformation (15) from the group $`\mathrm{SO}^+(1,3)`$, then there exist two elements $`\pm S\mathrm{Spin}(1,3)`$ such that the formula (13) is satisfied (in other words $`\mathrm{Spin}(1,3)`$ is a double covering of $`\mathrm{SO}^+(1,3)`$). We say that the transformation of coordinates (15),(14) from the group $`\mathrm{SO}^+(1,3)`$is associated with the element $`S\mathrm{Spin}(1,3)`$ if (13) holds.
### 1.3 Secondary generators of the Clifford algebra.
Denote $`\mathrm{}^5=\mathrm{}^{0123}=\mathrm{}^0\mathrm{}^1\mathrm{}^2\mathrm{}^3`$. Then $`(\mathrm{}^5)^2=\mathrm{}`$ and $`\mathrm{}^5`$ commutes with all even elements and anticommutes with all odd elements of $`𝒞\mathrm{}(1,3)`$.
Definition 2.If elements $`H𝒞\mathrm{}^1(1,3)`$ and $`I,K𝒞\mathrm{}^2(1,3)`$ satisfy the relations
$$H^2=\mathrm{},I^2=K^2=\mathrm{},[H,I]=[H,K]=0,\{I,K\}=IK+KI=0,$$
(16)
then the elements $`H,\mathrm{}^5,I,K`$ are said to be secondary generators of the Clifford algebra $`𝒞\mathrm{}(1,3)`$.
In particular, if we take
$$\stackrel{´}{H}=\mathrm{}^0,\stackrel{´}{I}=\mathrm{}^{12},\stackrel{´}{K}=\mathrm{}^{13},$$
(17)
then these elements satisfy (16) and hence, the elements $`\stackrel{´}{H},\mathrm{}^5,\stackrel{´}{I},\stackrel{´}{K}`$ are secondary generators of $`𝒞\mathrm{}(1,3)`$.
If $`H,\mathrm{}^5,I,K`$ are secondary generators of $`𝒞\mathrm{}(1,3)`$, then the 16 elements
$$\mathrm{},H,I,K,HI,HK,IK,HIK,\mathrm{}^5,\mathrm{}^5H,\mathrm{}^5I,\mathrm{}^5K,\mathrm{}^5HI,\mathrm{}^5HK,\mathrm{}^5IK,\mathrm{}^5HIK$$
are the basis elements of $`𝒞\mathrm{}(1,3)`$ (linear independent) and the trace $`\mathrm{Tr}`$ of every element of this basis, except $`\mathrm{}`$, is equal to zero.
Let $`H,\mathrm{}^5,I,K`$ be secondary generators of $`𝒞\mathrm{}(1,3)`$. The first pair $`H,\mathrm{}^5`$ is such that
$$H^2=\mathrm{},(\mathrm{}^5)^2=\mathrm{},\{H,\mathrm{}^5\}=0.$$
(18)
Thus the elements $`H,\mathrm{}^5`$ are generators of the Clifford algebra $`𝒞\mathrm{}(1,1)`$.
The second pair $`I,K`$ is such that
$$I^2=K^2=\mathrm{},\{I,K\}=0.$$
(19)
Therefore the elements $`I,K`$ are generators of the Clifford algebra $`𝒞\mathrm{}(0,2)`$ (which is isomorphic to the algebra of quaternions). Furthermore, the elements $`H,\mathrm{}^5`$ are commute with the elements $`I,K`$
$$[H,I]=[H,K]=[\mathrm{}^5,I]=[\mathrm{}^5,K]=0.$$
(20)
Consequently the Clifford algebra $`𝒞\mathrm{}(1,3)`$ is isomorphic to the direct product
$$𝒞\mathrm{}(1,3)𝒞\mathrm{}(1,1)𝒞\mathrm{}(0,2).$$
This relation leads to the well known fact that $`𝒞\mathrm{}(1,3)`$ can be represented by the algebra $`(2,)`$ of $`2\times 2`$ matrices with quaternion elements.
We omit proofs of the following three propositions.
Proposition 3. Suppose elements $`H𝒞\mathrm{}^{\mathrm{odd}}(1,3)`$ and $`I𝒞\mathrm{}^{\mathrm{even}}(1,3)`$ satisfy the relations
$$H^2=\mathrm{},I^2=\mathrm{},[H,I]=0.$$
(21)
Then there exists an invertible element $`T𝒞\mathrm{}^{\mathrm{even}}(1,3)`$ such that
$$T^1HT=\mathrm{}^0,T^1IT=\mathrm{}^{12}.$$
Proposition 4. Suppose elements $`H𝒞\mathrm{}^1(1,3)`$ and $`I𝒞\mathrm{}^2(1,3)`$ satisfy the relations (21). Then there exists an element $`S\mathrm{Spin}(1,3)`$ such that
$$S^{}HS=\mathrm{}^0,S^{}IS=\mathrm{}^{12}.$$
Proposition 5. Suppose $`H,\mathrm{}^5,I,K`$ are secondary generators of $`𝒞\mathrm{}(1,3)`$; then there exists a unique element $`S\mathrm{Spin}(1,3)`$ such that
$$S^{}HS=\mathrm{}^0,S^{}IS=\mathrm{}^{12},S^{}KS=\mathrm{}^{13}.$$
### 1.4 Idempotents, left ideals, and matrix representations of $`𝒞\mathrm{}(1,3)`$.
An element $`\tau `$ of an algebra $`𝒜`$ is said to be an idempotent if $`\tau ^2=\tau `$. A subspace $`𝒜`$ is called a left ideal of the algebra $`𝒜`$ if $`au`$ for all $`a𝒜,u`$. Every idempotent $`\tau 𝒜`$ generates the left ideal
$$(\tau )=\{a\tau :a𝒜\}.$$
Consider the idempotent
$$t=\frac{1}{4}(\mathrm{}+H)(\mathrm{}iI)$$
(22)
and the left ideal $`(t)`$ of $`𝒞\mathrm{}_𝒞(1,3)`$. The complex dimension of $`(t)`$ is equal to $`4`$. Let us denote
$$t_k=F_kt,k=1,2,3,4,$$
(23)
where
$$F_1=\mathrm{},F_2=K,F_3=I\mathrm{}^5,F_4=KI\mathrm{}^5.$$
(24)
Elements $`t_k(t),k=1,2,3,4`$ are linear independent and form a basis of $`(t)`$. It is easy to check that $`t_kt_1=t_k`$ and $`t_kt_n=0`$ for $`n1`$.
Let us define an operation of Hermitian conjugation
$$U^{}:=HU^{}H,U𝒞\mathrm{}_𝒞(1,3)$$
(25)
such that $`(UV)^{}=V^{}U^{}`$, $`U^{}=U`$.
Now we may introduce a scalar product of elements of the left ideal $`(t)`$
$$(U,V):=4\mathrm{Tr}(UV^{}),U,V(t).$$
(26)
This scalar product converts the left ideal $`(t)`$ into the four dimensional unitary space.
Theorem .The basis elements $`t_k=t^k`$, $`k=1,2,3,4`$ of $`(t)`$ are mutually orthogonal w.r.t. the scalar product (26)
$$(t_k,t^n)=\delta _k^n,k,n=1,2,3,4$$
where $`\delta _k^k=1`$ and $`\delta _k^n=0`$ for $`kn`$.
Proof is by direct calculation.
In what follows we use the formulas
$$\mathrm{\Psi }=(\mathrm{\Psi },t^k)t_k\text{for}\mathrm{\Psi }(t)$$
(27)
and
$$(KU,V)=(U,K^{}V)\text{for}U,V(t),K𝒞\mathrm{}_𝒞(1,3).$$
Now we introduce the following concept. We claim that the set of secondary generators $`H,\mathrm{}^5,I,K`$ uniquely defines a matrix representation for $`𝒞\mathrm{}(1,3)`$. Indeed, for $`U𝒞\mathrm{}(1,3)`$ the products $`Ut_k`$ belong to $`(t)`$ and can be represented as linear combinations of basis elements $`t_n`$ with certain coefficients $`\gamma (U)_k^n`$
$$Ut_k=\gamma (U)_k^nt_n,k=1,2,3,4.$$
(28)
Thus, the matrix $`\gamma (U)`$ with the elements $`\gamma (U)_k^n`$ (an upper index enumerates lines and a lower index enumerate columns of a matrix) is associated with the element $`U𝒞\mathrm{}(1,3)`$. In particular, the matrices $`\gamma ^\mu =\gamma (\mathrm{}^\mu )`$ are defined by the formulas
$$\mathrm{}^\mu t_k=\gamma (\mathrm{}^\mu )_k^nt_n,k=1,\mathrm{},4;\mu =0,\mathrm{},3.$$
(29)
Considering scalar products of the left and right hand sides of (28),(29) by $`t^n`$ and using mutual orthogonality of $`t_k`$, we get
$$\gamma (U)_k^n=(Ut_k,t^n)$$
(30)
and, in particular
$$\gamma _k^{n\mu }=\gamma (\mathrm{}^\mu )_k^n=(\mathrm{}^\mu t_k,t^n).$$
(31)
It can be shown that
$$\gamma (UV)=\gamma (U)\gamma (V),\gamma (\mathrm{})=\mathrm{𝟏},\gamma (\alpha U)=\alpha \gamma (U),\alpha 𝒞.$$
Therefore the map $`\gamma :𝒞\mathrm{}(4,𝒞)`$ defined by (28) is a matrix representation of the Clifford algebra $`𝒞\mathrm{}`$.
If we take secondary generators of the form (17), then the matrices $`\gamma ^\mu `$ from (31) are equal to the matrices from (3).
Let $`S`$ be an element of the group $`\mathrm{Spin}(1,3)`$ and $`\gamma (S)`$, $`\gamma (S^{})`$ be the matrix representation of $`S`$, $`S^{}`$ given by (28). From the formula (13) we have
$$S^{}\mathrm{}^\mu S=p_\nu ^\mu \mathrm{}^\nu ,S\mathrm{}^\nu S^{}=q_\mu ^\nu \mathrm{}^\mu ,$$
(32)
where
$$p_\nu ^\mu q_\lambda ^\nu =\delta _\lambda ^\mu ,q_\nu ^\mu p_\lambda ^\nu =\delta _\lambda ^\mu .$$
For the secondary generators $`H,\mathrm{}^5,I,K`$ consider the transformation
$$(H,\mathrm{}^5,I,K)(S^{}HS,\mathrm{}^5,S^{}IS,S^{}KS),$$
which leads to the transformation of the left ideal $`(t)(S^{}tS)`$ and the basis elements
$$t_k\stackrel{´}{t_k}=S^{}t_kS.$$
Now we may define a new matrix representation of the Clifford algebra $`\stackrel{´}{\gamma }:𝒞\mathrm{}(4,𝒞)`$ with the aid of the formula
$$U\stackrel{´}{t}_k=\stackrel{´}{\gamma }(U)_k^n\stackrel{´}{t}_n.$$
(33)
Theorem .For every $`U𝒞\mathrm{}`$ the matrix $`\gamma (U)`$ defined with the aid of (30) connected with the matrix representation $`\stackrel{´}{\gamma }(U)`$ by the formula
$$\stackrel{´}{\gamma }(U)=\gamma (S)\gamma (U)\gamma (S^{}).$$
Proof. It is sufficient to prove this theorem for the primary generators $`\mathrm{}^\mu `$. We have
$$\mathrm{}^\mu t_k=\gamma (\mathrm{}^\mu )_k^nt_n.$$
Multiplying both sides of this relation from the left by $`S^{}`$ and from the right by $`S`$, we obtain
$$(S^{}\mathrm{}^\mu S)(S^{}t_kS)=\gamma (\mathrm{}^\mu )_k^n(S^{}t_nS).$$
Substituting $`S^{}t_kS=\stackrel{´}{t}_k`$ and $`S^{}\mathrm{}^\mu S=p_\nu ^\mu \mathrm{}^\nu `$ from (32), we get
$$p_\nu ^\mu \mathrm{}^\nu \stackrel{´}{t}_k=\gamma (\mathrm{}^\mu )_k^n\stackrel{´}{t}_n.$$
Multiplying both sides by $`q_\mu ^\lambda `$ from (32) and summing over $`\mu `$, we obtain
$`q_\mu ^\lambda p_\nu ^\mu \mathrm{}^\nu \stackrel{´}{t}_k=\mathrm{}^\lambda \stackrel{´}{t}_k`$ $`=`$
$`=`$ $`\gamma (q_\mu ^\lambda \mathrm{}^\mu )_k^n\stackrel{´}{t}_n=\gamma (S\mathrm{}^\lambda S^{})_k^n\stackrel{´}{t}_n`$
$`=`$ $`(\gamma (S)\gamma (\mathrm{}^\lambda )\gamma (S^{}))_k^n\stackrel{´}{t}_n.`$
This completes the proof.
Let us note that $`\stackrel{´}{\gamma }(S)=\gamma (S)`$ and $`\stackrel{´}{\gamma }(S^{})=\gamma (S^{})`$
### 1.5 A one-to-one correspondence between $`(t)`$ and $`𝒞\mathrm{}_𝒞(1,3)`$.
The dimension of the linear space $`𝒞\mathrm{}^{\mathrm{even}}(1,3)`$ is equal to $`8`$. The left ideal $`(t)`$ has the complex dimension $`4`$ and, thus, the real dimension $`8`$. We may consider the map $`𝒞\mathrm{}^{\mathrm{even}}(1,3)(t)`$ given by the formula
$$\mathrm{\Psi }\mathrm{\Psi }t=\varphi ^kt_k,$$
where $`\mathrm{\Psi }𝒞\mathrm{}^{\mathrm{even}}(1,3)`$ and $`\varphi ^k=(\mathrm{\Psi }t,t^k)`$. Let us prove that this map gives the one-to-one correspondence between the even Clifford algebra $`𝒞\mathrm{}^{\mathrm{even}}(1,3)`$ and the left ideal $`(t)`$.
Theorem 3.Suppose $`\mathrm{\Phi }=\varphi ^kt_k(t)`$. Then the equation for $`\mathrm{\Omega }𝒞\mathrm{}^{\mathrm{even}}(1,3)`$
$$\mathrm{\Omega }t=\mathrm{\Phi }$$
(34)
has a unique solution
$$\mathrm{\Omega }=F_k(\alpha ^k\mathrm{}+\beta ^kI),$$
(35)
where $`\varphi ^k=\alpha ^k+i\beta ^k`$ and $`F_k`$ are defined in (24).
Proof. Multiplying both sides of (35) from the right by $`t`$ and using the relation $`It=it`$, we see that the formula (35) really gives the solution of (34). Now we must prove that the homogeneous equation $`\mathrm{\Omega }t=0`$ has only trivial solution $`\mathrm{\Omega }=0`$. Firstly we give the proof for the secondary generators from (17). Suppose $`\mathrm{\Omega }`$ is of the form
$$\mathrm{\Omega }=\omega \mathrm{}+\underset{\mu <\nu }{}\omega _{\mu \nu }\mathrm{}^{\mu \nu }+\omega _5\mathrm{}^5.$$
Then the element $`\mathrm{\Omega }t(t)`$ can be expanded in the basis $`t_k`$
$$\mathrm{\Omega }t=(\omega i\omega _{12})t_1+(\omega _{13}i\omega _{23})t_2+(\omega _{03}+i\omega _5)t_3+(\omega _{01}i\omega _{02})t_4.$$
Therefore from $`\mathrm{\Omega }t=0`$ we get $`\mathrm{\Omega }=0`$. For the case of general secondary generators we must represent $`\mathrm{\Omega }`$ as the linear combination of products of the generators $`\mathrm{}^5,I,K`$ and arguing as above. So the solution (35) of (34) is unique. This completes the proof.
Theorem .If $`\mathrm{\Psi }𝒞\mathrm{}^{\mathrm{even}}(1,3)`$ and $`S\mathrm{Spin}(1,3)`$, then the relation
$$\mathrm{\Psi }t=\psi ^kt_k,$$
(36)
where $`\psi ^k=(\mathrm{\Psi }t,t^k)𝒞`$, is invariant under the following transformation:
$`t_k`$ $``$ $`\stackrel{´}{t}_k=S^{}t_kS,`$
$`\mathrm{\Psi }`$ $``$ $`\stackrel{´}{\mathrm{\Psi }}=\mathrm{\Psi }S,`$ (37)
$`\psi ^k`$ $``$ $`\stackrel{´}{\psi }^k=\gamma (S)_l^k\psi ^l=(St_l,t^k)\psi ^l.`$
Proof. We have
$`(\stackrel{´}{\mathrm{\Psi }}\stackrel{´}{t}\stackrel{´}{\psi }^k\stackrel{´}{t}_k)S^{}`$ $`=`$ $`\mathrm{\Psi }t(St_l,t^k)\psi ^lS^{}t_k=`$
$`=`$ $`\mathrm{\Psi }t\psi ^l(St_l,t^k)(S^{}t_k,t^n)t_n=`$
$`=`$ $`\mathrm{\Psi }t\psi ^nt_n.`$
Here we use the formulas
$$S^{}t_k=(S^{}t_k,t^n)t_n$$
and
$$(St_l,t^k)(S^{}t_k,t^n)=\gamma (S)_l^k\gamma (S^{})_k^n=\delta _l^n.$$
This completes the proof.
Consider the correspondence $`\mathrm{}^\mu \gamma ^\mu `$ between primary generators of the Clifford algebra and Dirac’s $`\gamma `$-matrices. Let us extend this correspondence to the map $`\gamma :𝒞\mathrm{}(1,3)(4,𝒞)`$ or $`\gamma :𝒞\mathrm{}_𝒞(1,3)(4,𝒞)`$ in such a way that any element of the Clifford algebra
$$U=u\mathrm{}+u_\mu \mathrm{}^\mu +\underset{\mu _1<\mu _2}{}u_{\mu _1\mu _2}\mathrm{}^{\mu _1\mu _2}+\underset{\mu _1<\mu _2<\mu _3}{}u_{\mu _1\mu _2\mu _3}\mathrm{}^{\mu _1\mu _2\mu _3}+u_{0123}\mathrm{}^{0123}$$
(38)
corresponds to the matrix $`𝐔(4,𝒞)`$
$$𝐔=u\mathrm{𝟏}+u_\mu \gamma ^\mu +\underset{\mu _1<\mu _2}{}u_{\mu _1\mu _2}\gamma ^{\mu _1\mu _2}+\underset{\mu _1<\mu _2<\mu _3}{}u_{\mu _1\mu _2\mu _3}\gamma ^{\mu _1\mu _2\mu _3}+u_{0123}\gamma ^{0123}$$
(39)
And let $`\mathrm{}:(4,𝒞)𝒞\mathrm{}_𝒞(1,3)`$ or $`\mathrm{}:(4,𝒞)𝒞\mathrm{}(1,3)`$ be the inverse map such that any matrix (39) corresponds to the element (38) of the (complex) Clifford algebra.
### 1.6 The covariance of the Dirac equation.
Let us consider the transformation of the Dirac equation (2) under a change of coordinates $`(x)(\stackrel{~}{x})`$ from the group $`\mathrm{SO}^+(1,3)`$. This change of coordinates (15),(14) is associated with an element $`S\mathrm{Spin}(1,3)`$ in accordance with (13). By $`𝐒=\gamma (S)`$ denote the matrix representation of the element $`S`$. Then the matrix $`𝐒`$ satisfies
$$𝐒^{}\gamma ^\nu 𝐒=p_\mu ^\nu \gamma ^\mu .$$
(40)
The covectors $`_\mu `$ and $`a_\mu `$ are transformed under the change of coordinates (15) as
$$_\mu =p_\mu ^\nu \stackrel{~}{}_\nu ,a_\mu =p_\mu ^\nu \stackrel{~}{a}_\nu ,$$
(41)
where $`\stackrel{~}{}_\nu =/\stackrel{~}{x}^\nu `$ and $`\stackrel{~}{a}_\nu `$ are components of the covector $`a_\mu `$ in coordinates $`(\stackrel{~}{x})`$. By substituting (41),(40) into (2), we obtain
$`\gamma ^\mu (_\mu \psi +ia_\mu \psi )+im\psi `$ $`=`$ $`p_\mu ^\nu \gamma ^\mu (\stackrel{~}{}_\nu \psi +i\stackrel{~}{a}_\nu \psi )+im\psi =`$
$`=`$ $`𝐒^{}\gamma ^\nu 𝐒(\stackrel{~}{}_\nu \psi +i\stackrel{~}{a}_\nu \psi )+im\psi =`$
$`=`$ $`𝐒^{}(\gamma ^\nu (\stackrel{~}{}_\nu (𝐒\psi )+i\stackrel{~}{a}_\nu (𝐒\psi ))+im(𝐒\psi ))`$
Hence, if a column of four complex valued functions $`\psi =\psi (x)`$ in the coordinates $`(x)`$ satisfies the equation (2), then the column $`\stackrel{~}{\psi }=𝐒\psi (x(\stackrel{~}{x}))`$ in coordinates $`(\stackrel{~}{x})`$ satisfies the equation
$$\gamma ^\nu (\stackrel{~}{}_\nu \stackrel{~}{\psi }+i\stackrel{~}{a}_\nu \stackrel{~}{\psi })+im\stackrel{~}{\psi }=0,$$
(42)
which has the same form as (2).
Definition 1.A column of four complex valued functions $`\psi `$ is called a bispinor if $`\psi `$ transforms under the change of coordinates (15),(14),(13) as $`\psi \stackrel{~}{\psi }=𝐒\psi (x(\stackrel{~}{x}))`$, where $`𝐒=\gamma (S)`$.
### 1.7 Algebraic bispinors and the Dirac equation.
Consider the equation
$$\mathrm{}^\mu (_\mu \rho +ia_\mu \rho )+im\rho =0,$$
(43)
where $`\rho =\rho (x)𝒞\mathrm{}_𝒞(1,3)`$. The element $`\rho `$ has 16 complex components, i.e., four times more than the bispinor. Multiplying the equation (43) from the right by $`t`$, we obtain that the element $`\theta =\rho t(t)`$ satisfying the same equation
$$\mathrm{}^\mu (_\mu \theta +ia_\mu \theta )+im\theta =0.$$
(44)
Theorem 2.An element $`\theta =\psi _kt^k(t)`$ satisfies the equation (44) iff the column $`\psi =(\psi _1\psi _2\psi _3\psi _4)^T`$ satisfies the Dirac equation (2).
Proof. Necessity. Substituting $`\theta =\psi _kt^k`$ into (44) and using (28), we get
$$\mathrm{}^\mu t^k(_\mu \psi _k+ia_\mu \psi _k)+im\psi _lt^l=(\gamma _l^{\mu k}(_\mu \psi _k+ia_\mu \psi _k)+im\psi _l)t^l=0.$$
(45)
Taking into account the linear independence of $`t^k`$, we obtain the equations
$$\gamma _l^{\mu k}(_\mu \psi _k+ia_\mu \psi _k)+im\psi _l=0,l=1,2,3,4,$$
which is evidently equivalent to the Dirac equation (2).
Arguing as above but in inverse order, we prove sufficiency. This completes the proof.
We say that for $`t`$ from (22) the left ideal $`(t)𝒞\mathrm{}_𝒞(1,3)`$ is the spinor space. Elements of the spinor space are called algebraic bispinors. The formula $`\theta =\psi _kt^k`$ gives the relation between the algebraic bispinor $`\theta (t)`$ and the bispinor $`\psi =(\psi _1\psi _2\psi _3\psi _4)^T`$.
### 1.8 Hestenes’ form of the Dirac equation.
Let $`H,\mathrm{}^5,I,K`$ be secondary generators of $`𝒞\mathrm{}(1,3)`$ and be independent of $`x`$. Consider the equation for $`\mathrm{\Psi }=\mathrm{\Psi }(x)𝒞\mathrm{}^{\mathrm{even}}(1,3)`$
$$\mathrm{}^\mu (_\mu \mathrm{\Psi }+a_\mu \mathrm{\Psi }I)+m\mathrm{\Psi }HI=0,$$
(46)
which was invented by D.Hestenes in 1966. This equation is called Hestenes’ form of the Dirac equation (HDE). Let us show that the equation (46) is equivalent to the equation (44) and consequently to the Dirac equation (2). To prove this fact we need the following theorem.
Theorem 4.An element $`\mathrm{\Psi }=\mathrm{\Psi }(x)𝒞\mathrm{}^{\mathrm{even}}(1,3)`$ satisfies the equation (46) iff the element $`\theta =\mathrm{\Psi }t(t)`$ satisfies the equation (44).
Proof. Necessity. Suppose $`\mathrm{\Psi }𝒞\mathrm{}^{\mathrm{even}}(1,3)`$ satisfies the equation (46). Let us multiply (46) from the right by the idempotent $`t`$ and use the relations
$$Ht=t,It=it.$$
Then we get the equation (44) for $`\theta =\mathrm{\Psi }t(t)`$
$$\mathrm{}^\mu (_\mu \mathrm{\Psi }t+a_\mu \mathrm{\Psi }It)+m\mathrm{\Psi }HIt=\mathrm{}^\mu (_\mu \mathrm{\Psi }t+ia_\mu \mathrm{\Psi }t)+im\mathrm{\Psi }t=0.$$
Sufficiency. Suppose the element $`\theta (t)`$ satisfies the equation (44). Let us multiply this equation from the right by $`t`$ and use the relations
$$t=tH,it=tI,\theta t=\theta .$$
Then we obtain
$$\mathrm{}^\mu (_\mu \theta +a_\mu \theta I)+m\theta HI=0.$$
(47)
Now using Theorem 3 we may take $`\mathrm{\Psi }𝒞\mathrm{}^{\mathrm{even}}(1,3)`$ as the solution of the equation $`\mathrm{\Psi }t=\theta `$. We claim that this $`\mathrm{\Psi }`$ satisfies the equation (46). In fact, substituting $`\theta =\mathrm{\Psi }t`$ into the equation (47) and multiplying from the left by $`H`$, we get
$$0=H(\mathrm{}^\mu (_\mu \mathrm{\Psi }+a_\mu \mathrm{\Psi }I)+m\mathrm{\Psi }HI)t\mathrm{\Omega }t,$$
where $`\mathrm{\Omega }𝒞\mathrm{}^{\mathrm{even}}(1,3)`$. By Theorem 3 the equation $`\mathrm{\Omega }t=0`$ has only the trivial solution $`\mathrm{\Omega }=0`$. Hence $`\mathrm{\Psi }`$ satisfies (46). The theorem is proved.
Thus, we have proved that HDE (46) for $`\mathrm{\Psi }𝒞\mathrm{}^{\mathrm{even}}(1,3)`$ is equivalent to the Dirac equation (2) for $`\psi =(\psi _1\psi _2\psi _3\psi _4)^T`$, $`\psi _k=\alpha _k+i\beta _k`$ and the relation between these solutions is
$$\mathrm{\Psi }=F^k(\alpha _k\mathrm{}+\beta _kI),$$
(48)
where the summation over $`k=1,2,3,4`$ is assumed.
Consider a change of coordinates (15),(14),(13) from the group $`\mathrm{SO}^+(1,3)`$associated with the element $`S\mathrm{Spin}(1,3)`$. Arguing as in section 3, we see that HDE transforms under this change of coordinates as follows
$$\mathrm{}^\mu (_\mu \mathrm{\Psi }+a_\mu \mathrm{\Psi }I)+m\mathrm{\Psi }HI=S^{}(\mathrm{}^\nu (\stackrel{~}{}_\nu (S\mathrm{\Psi })+\stackrel{~}{a}_\nu (S\mathrm{\Psi })I)+m(S\mathrm{\Psi })HI).$$
Hence, if $`\mathrm{\Psi }=\mathrm{\Psi }(x)𝒞\mathrm{}^{\mathrm{even}}(1,3)`$ in coordinates $`(x)`$ satisfies HDE (46), then the element $`\stackrel{~}{\mathrm{\Psi }}=S\mathrm{\Psi }(x(\stackrel{~}{x}))`$ in coordinates $`(\stackrel{~}{x})`$ satisfies the equation
$$\mathrm{}^\nu (\stackrel{~}{}_\nu \stackrel{~}{\mathrm{\Psi }}+\stackrel{~}{a}_\nu \stackrel{~}{\mathrm{\Psi }}I)+m\stackrel{~}{\mathrm{\Psi }}HI=0,$$
which has the same form as (46). As was shown, the relation in coordinates $`(x)`$ between a solution $`\psi `$ of the Dirac equation and the solution $`\mathrm{\Psi }`$ of HDE is given by the formula
$$\mathrm{\Psi }t=\psi _kt^k,\psi =(\psi _1\psi _2\psi _3\psi _4)^T.$$
And the relation in coordinates $`(\stackrel{~}{x})`$ between the solution $`\stackrel{~}{\psi }=𝐒\psi `$ of the Dirac equation and the solution $`\stackrel{~}{\mathrm{\Psi }}=S\mathrm{\Psi }`$ of HDE is given by the formula
$$\stackrel{~}{\mathrm{\Psi }}t=\stackrel{~}{\psi }_lt^l,\stackrel{~}{\psi }=(\stackrel{~}{\psi }_1\stackrel{~}{\psi }_2\stackrel{~}{\psi }_3\stackrel{~}{\psi }_4)^T.$$
Indeed, using the formula $`St=\gamma (S)_l^kt^l`$, where $`\gamma (S)_l^k`$ are elements of the matrix $`𝐒`$, we obtain
$$\stackrel{~}{\mathrm{\Psi }}t=S\mathrm{\Psi }t=\psi _kSt^k=\psi _k\gamma (S)_l^kt^l=\stackrel{~}{\psi }_lt^l.$$
### 1.9 The Grassmann-Clifford bialgebra.
Suppose that for elements of $`𝒞\mathrm{}(1,3)`$ the exterior multiplication (denoted by $``$) is defined by the following rules:
$`𝒞\mathrm{}(1,3)`$ is an associative algebra (with the unity element $`\mathrm{}`$) with respect to exterior multiplication;
$`\mathrm{}^\mu \mathrm{}^\nu =\mathrm{}^\nu \mathrm{}^\mu `$, $`\mu ,\nu =0,1,2,3`$;
$`\mathrm{}^{\mu _1}\mathrm{}\mathrm{}^{\mu _k}=\mathrm{}^{\mu _1\mathrm{}\mu _k}`$ for $`0\mu _1<\mathrm{}<\mu _k3`$.
The resulting algebra (equipped with the Clifford multiplication and with the exterior multiplication) is called the Grassmann-Clifford bialgebra and is denoted by $`\mathrm{\Lambda }(1,3)`$. The complex valued Grassmann-Clifford bialgebra is denoted by $`\mathrm{\Lambda }_𝒞(1,3)`$. Any element $`U\mathrm{\Lambda }(1,3)`$ can be expanded in the basis as in (38). The coefficients $`u_{\mu _1\mathrm{}\mu _k}`$ in (38) are enumerated by ordered multi-indices. Let us take the coefficients that are antisymmetric w.r.t. all indices
$$u_{\mu _1\mathrm{}\mu _k}=u_{[\mu _1\mathrm{}\mu _k]},$$
where square brackets denote the operation of alternation (with the division by $`k!`$). Then elements of the form
$$\underset{\mu _1<\mathrm{}<\mu _k}{}u_{\mu _1\mathrm{}\mu _k}\mathrm{}^{\mu _1\mathrm{}\mu _k}=\frac{1}{k!}u_{\nu _1\mathrm{}\nu _k}\mathrm{}^{\nu _1}\mathrm{}\mathrm{}^{\nu _k}=\frac{1}{k!}u_{\nu _1\mathrm{}\nu _k}\mathrm{}^{\nu _1}\mathrm{}\mathrm{}^{\nu _k}$$
(49)
are said to be elements of rank $`k`$ and belong to $`\mathrm{\Lambda }^k(1,3)`$, where $`\mathrm{\Lambda }^k(1,3)`$, $`\mathrm{\Lambda }^{\mathrm{even}}(1,3)`$, $`\mathrm{\Lambda }^{\mathrm{odd}}(1,3)`$ are the same as $`𝒞\mathrm{}^k(1,3)`$, $`𝒞\mathrm{}^{\mathrm{even}}(1,3)`$, $`𝒞\mathrm{}^{\mathrm{odd}}(1,3)`$. If $`U\mathrm{\Lambda }^r(1,3),V\mathrm{\Lambda }^s(1,3)`$ then
$$UV=(1)^{rs}VU\mathrm{\Lambda }^{r+s}(1,3).$$
(50)
## 2 Part II.
### 2.1 The exterior algebra of Minkowski space.
Let $``$ be Minkowski space with the metric tensor (1), with coordinates $`x^\mu `$, with basis coordinate vectors $`e_\mu `$, and with basis covectors $`e^\mu =g^{\mu \nu }e_\nu `$. Consider a covariant antisymmetric tensor field of rank $`0k3`$ on $``$ with components
$$u_{\mu _1\mathrm{}\mu _k}=u_{[\mu _1\mathrm{}\mu _k]}(x).$$
It is suitable to write this field with the aid of the expression
$$\frac{1}{k!}u_{\nu _1\mathrm{}\nu _k}e^{\nu _1}\mathrm{}e^{\nu _k},$$
(51)
where the expression $`e^{\nu _1}\mathrm{}e^{\nu _k}`$ in the fixed coordinates $`(x)`$ can be considered as an element of the Grassmann algebra. Under a linear change of coordinates
$$x^\mu =\frac{x^\mu }{\stackrel{~}{x}^\nu }\stackrel{~}{x}^\nu =q_\nu ^\mu \stackrel{~}{x}^\nu $$
(52)
the transformation rule for the expression $`e^{\nu _1}\mathrm{}e^{\nu _k}`$ corresponds to the transformation rule for basis covectors $`e^\nu `$. That is,
$$e^\nu =q_\mu ^\nu \stackrel{~}{e}^\mu ,e^{\nu _1}\mathrm{}e^{\nu _k}=q_{\mu _1}^{\nu _1}\mathrm{}q_{\mu _k}^{\nu _k}\stackrel{~}{e}^{\mu _1}\mathrm{}\stackrel{~}{e}^{\mu _k}.$$
(53)
Therefore under the change of coordinates (52) the expression (51) is invariant
$$\frac{1}{k!}u_{\nu _1\mathrm{}\nu _k}e^{\nu _1}\mathrm{}e^{\nu _k}=\frac{1}{k!}\stackrel{~}{u}_{\mu _1\mathrm{}\mu _k}\stackrel{~}{e}^{\mu _1}\mathrm{}\stackrel{~}{e}^{\mu _k},$$
where
$$\stackrel{~}{u}_{\mu _1\mathrm{}\mu _k}=p_{\nu _1}^{\lambda _1}\mathrm{}p_{\nu _k}^{\lambda _k}u_{\lambda _1\mathrm{}\lambda _k}$$
are components of the tensor field $`u_{\nu _1\mathrm{}\nu _k}`$ in the coordinates $`(\stackrel{~}{x})`$ and
$$p_\nu ^\lambda =\frac{\stackrel{~}{x}^\lambda }{x^\nu },p_\nu ^\lambda q_\mu ^\nu =\delta _\mu ^\lambda (\delta _\lambda ^\lambda =1,\delta _\mu ^\lambda =0\text{for}\mu \lambda ).$$
The expressions (51) are called exterior forms of rank $`k`$ or $`k`$-forms. The set of all $`k`$-forms is denoted by $`\mathrm{\Lambda }^k()`$. The formal sum of $`k`$-forms
$$\underset{k=0}{\overset{4}{}}\frac{1}{k!}u_{\nu _1\mathrm{}\nu _k}e^{\nu _1}\mathrm{}e^{\nu _k}$$
(54)
are said to be (nonhomogeneous) exterior form. The set of all exterior forms is denoted by $`\mathrm{\Lambda }()`$ and
$$\mathrm{\Lambda }()=\mathrm{\Lambda }^0()\mathrm{}\mathrm{\Lambda }^4()=\mathrm{\Lambda }^{\mathrm{even}}()\mathrm{\Lambda }^{\mathrm{odd}}().$$
It is well known that the exterior product of exterior forms is an exterior form.
Consider the Hodge star operator $`:\mathrm{\Lambda }^k()\mathrm{\Lambda }^{4k}()`$. If $`U\mathrm{\Lambda }^k()`$ has the form (51), then
$$U=\frac{1}{k!(4k)!}\epsilon _{\mu _1\mathrm{}\mu _4}u^{\mu _1\mathrm{}\mu _k}e^{\mu _{k+1}}\mathrm{}e^{\mu _4},$$
where
$$u^{\mu _1\mathrm{}\mu _k}=g^{\mu _1\nu _1}\mathrm{}g^{\mu _k\nu _k}u_{\nu _1\mathrm{}\nu _k}$$
and $`\epsilon _{\mu _1\mathrm{}\mu _4}`$ is the sign of the permutation $`(\mu _1\mathrm{}\mu _4)`$. $`U`$ is an exterior form (a covariant antisymmetric tensor) w.r.t. any change of coordinates with a positive Jacobian.
Remark. In this paper we consider changes of coordinates only from the group $`\mathrm{SO}^+(1,3)`$(a Jacobian is equal to 1) and do not distinguish tensors and pseudotensors.
It is easy to prove that for any $`U\mathrm{\Lambda }^k()`$
$$(U)=(1)^{k+1}U.$$
(55)
Further on we consider the bilinear operator $`\mathrm{Com}:\mathrm{\Lambda }^2()\times \mathrm{\Lambda }^2()\mathrm{\Lambda }^2()`$ such that for basis 2-forms
$$\mathrm{Com}(e^{\mu _1}e^{\mu _2},e^{\nu _1}e^{\nu _2})=2g^{\mu _1\nu _1}e^{\mu _2}e^{\nu _2}2g^{\mu _2\nu _2}e^{\mu _1}e^{\nu _1}+2g^{\mu _1\nu _2}e^{\mu _2}e^{\nu _1}+2g^{\mu _2\nu _1}e^{\mu _1}e^{\nu _2}$$
Evidently
$$\mathrm{Com}(U,V)=\mathrm{Com}(V,U),U,V\mathrm{\Lambda }^2().$$
Now we define the Clifford multiplication of exterior forms with the aid of the following formulas:
$`\stackrel{0}{U}\stackrel{k}{V}`$ $`=`$ $`\stackrel{k}{V}\stackrel{0}{U}=\stackrel{0}{U}\stackrel{k}{V}=\stackrel{k}{V}\stackrel{0}{U},`$
$`\stackrel{1}{U}\stackrel{k}{V}`$ $`=`$ $`\stackrel{1}{U}\stackrel{k}{V}(\stackrel{1}{U}\stackrel{k}{V}),`$
$`\stackrel{k}{U}\stackrel{1}{V}`$ $`=`$ $`\stackrel{k}{U}\stackrel{1}{V}+(\stackrel{k}{U}\stackrel{1}{V}),`$
$`\stackrel{2}{U}\stackrel{2}{V}`$ $`=`$ $`\stackrel{2}{U}\stackrel{2}{V}+(\stackrel{2}{U}\stackrel{2}{V})+{\displaystyle \frac{1}{2}}\mathrm{Com}(\stackrel{2}{U},\stackrel{2}{V}),`$
$`\stackrel{2}{U}\stackrel{3}{V}`$ $`=`$ $`\stackrel{2}{U}\stackrel{3}{V}(\stackrel{2}{U}\stackrel{3}{V}),`$
$`\stackrel{2}{U}\stackrel{4}{V}`$ $`=`$ $`\stackrel{2}{U}\stackrel{4}{V},`$
$`\stackrel{3}{U}\stackrel{2}{V}`$ $`=`$ $`\stackrel{3}{U}\stackrel{2}{V}(\stackrel{3}{U}\stackrel{2}{V}),`$
$`\stackrel{3}{U}\stackrel{3}{V}`$ $`=`$ $`\stackrel{3}{U}\stackrel{3}{V}+(\stackrel{3}{U}\stackrel{3}{V}),`$
$`\stackrel{3}{U}\stackrel{4}{V}`$ $`=`$ $`\stackrel{3}{U}\stackrel{4}{V},`$
$`\stackrel{4}{U}\stackrel{2}{V}`$ $`=`$ $`\stackrel{4}{U}\stackrel{2}{V},`$
$`\stackrel{4}{U}\stackrel{3}{V}`$ $`=`$ $`\stackrel{4}{U}\stackrel{3}{V},`$
$`\stackrel{4}{U}\stackrel{4}{V}`$ $`=`$ $`\stackrel{4}{U}\stackrel{4}{V},`$
where ranks of exterior forms are denoted as $`\stackrel{k}{U}\mathrm{\Lambda }^k()`$ and $`k=0,1,2,3,4`$.
From this definition we may get some properties of the Clifford multiplication of exterior forms.
If $`U,V\mathrm{\Lambda }()`$, then $`UV\mathrm{\Lambda }()`$.
The axioms of associativity and distributivity are satisfied for the Clifford multiplication.
$`e^\mu e^\nu =e^\mu e^\nu +g^{\mu \nu }e,e^\mu e^\nu +e^\nu e^\mu =2g^{\mu \nu }e`$.
$`e^{\mu _1}\mathrm{}e^{\mu _k}=e^{\mu _1}\mathrm{}e^{\mu _k}=e^{\mu _1\mathrm{}\mu _k}`$ for $`0\mu _1<\mathrm{}<\mu _k3`$.
If $`U,V\mathrm{\Lambda }^2()`$, then $`\mathrm{Com}(U,V)=UVVU`$.
Taking into account these properties of Clifford multiplication, we may conclude that Propositions 1 – 5 of Part I initially formulated for elements of $`𝒞\mathrm{}(1,3)`$ are also valid for elements of $`\mathrm{\Lambda }()`$.
In the sequel, we use the group
$$\mathrm{Spin}()=\{S\mathrm{\Lambda }^{\mathrm{even}}():S^{}S=e\}.$$
(56)
### 2.2 Operators $`d,\delta ,\mathrm{{\rm Y}},\mathrm{\Delta }`$.
First consider the operator
$$dV=e^\mu _\mu V,V\mathrm{\Lambda }()$$
such that
$`d:\mathrm{\Lambda }^k()\mathrm{\Lambda }^{k+1}()`$;
$`d^2=0`$;
$`d(UV)=dUV+(1)^kUdV`$ for $`U\mathrm{\Lambda }^k(),V\mathrm{\Lambda }()`$.
Secondly, consider the operator $`\delta `$
$$\delta U=dU\text{for}U\mathrm{\Lambda }()$$
such that
$`\delta :\mathrm{\Lambda }^k()\mathrm{\Lambda }^{k1}()`$;
$`\delta ^2=0`$.
Thirdly, consider the operator (Upsilon)
$$\mathrm{{\rm Y}}=d\delta $$
such that
$`\mathrm{{\rm Y}}:\mathrm{\Lambda }^k()\mathrm{\Lambda }^{k+1}()\mathrm{\Lambda }^{k1}()`$;
$`\mathrm{{\rm Y}}U=e^\mu _\mu U`$.
The second property in this list follows from the definition of Clifford multiplication
$$\stackrel{1}{U}\stackrel{k}{V}=\stackrel{1}{U}\stackrel{k}{V}(\stackrel{1}{U}\stackrel{k}{V})$$
if we formally substitute $`\stackrel{1}{U}=e^\mu _\mu `$.
Finally, consider the Beltrami-Laplace operator
$$\mathrm{\Delta }=\mathrm{{\rm Y}}^2=(d\delta )^2=(d\delta +\delta d)=g^{\mu \nu }_\mu _\nu $$
such that
$`\mathrm{\Delta }:\mathrm{\Lambda }^k()\mathrm{\Lambda }^k()`$;
$`\mathrm{\Delta }`$ commutes with the operators $`d,\delta ,\mathrm{{\rm Y}},`$.
### 2.3 A tensor form of the Dirac equation.
Let $`H\mathrm{\Lambda }^1(),I\mathrm{\Lambda }^2()`$ be two independent of $`x`$ exterior forms such that
$$H^2=e,I^2=e,[H,I]=0.$$
(57)
Now we consider the equation
$$\mathrm{{\rm Y}}\mathrm{\Phi }+A\mathrm{\Phi }I+m\mathrm{\Phi }HI=0,$$
(58)
where $`\mathrm{\Phi }=\mathrm{\Phi }(x)\mathrm{\Lambda }^{\mathrm{even}}()`$, $`A=a_\mu (x)e^\mu \mathrm{\Lambda }^1()`$, and $`m0`$ is a real constant. All the values $`(\mathrm{\Phi },A,I,H)`$ in (58) are exterior forms (covariant antisymmetric tensors). Two operations are used in (58). Namely the differential operator $`\mathrm{{\rm Y}}=d\delta =e^\mu _\mu `$ and the Clifford multiplication of exterior forms. Both operations take exterior forms to exterior forms. In other words, (58) is a tensor equation. We say that the equation (58) is the Tensor form of the Dirac Equation (TDE). The TDE is invariant under the following global (independent of $`x`$) transformation
$`\mathrm{\Phi }`$ $``$ $`\mathrm{\Phi }S,`$
$`H`$ $``$ $`S^{}HS,`$ (59)
$`I`$ $``$ $`S^{}IS,`$
where $`S\mathrm{Spin}()`$, $`_\mu S=0`$.
In a fixed coordinate system $`(x)`$ the TDE is equivalent to HDE (46) and the connection between (58) and (46) is given by the formula
$$(\mathrm{\Phi })_{e^\mu \mathrm{}^\mu }=\mathrm{\Psi }.$$
Let us remind that HDE (46) is equivalent to the Dirac equation (2) and the connection between them is given by the formula (48).
Remark. Taking into account the theorem 5, it is clear that we may use in the TDE $`H\mathrm{\Lambda }^{\mathrm{odd}}(),I\mathrm{\Lambda }^{\mathrm{even}}()`$ which satisfies (57) instead of $`H\mathrm{\Lambda }^1(),I\mathrm{\Lambda }^2()`$. In this case the equation (58) is invariant under the global transformation
$$\mathrm{\Phi }\mathrm{\Phi }T,HT^1HT,IT^1IT,$$
where $`T\mathrm{\Lambda }^{\mathrm{even}}()`$ is an invertible element.
Under a linear change of coordinates $`(x)(\stackrel{~}{x})`$ all exterior forms are invariants. Therefore in coordinates $`(\stackrel{~}{x})`$ the TDE has the form
$$\stackrel{~}{\mathrm{{\rm Y}}}\stackrel{~}{\mathrm{\Phi }}+\stackrel{~}{A}\stackrel{~}{\mathrm{\Phi }}\stackrel{~}{I}+m\stackrel{~}{\mathrm{\Phi }}\stackrel{~}{H}\stackrel{~}{I}=0,$$
where $`\stackrel{~}{\mathrm{{\rm Y}}}=\stackrel{~}{e}^\mu /\stackrel{~}{x}^\mu =\mathrm{{\rm Y}}`$, $`\stackrel{~}{\mathrm{\Phi }}=\mathrm{\Phi }`$, $`\stackrel{~}{A}=A`$, $`\stackrel{~}{H}=H`$, $`\stackrel{~}{I}=I`$ and $`\stackrel{~}{\mathrm{\Phi }},\stackrel{~}{A},\stackrel{~}{H},\stackrel{~}{I}`$ are the exterior forms written in coordinates $`(\stackrel{~}{x})`$.
Let us define the trace of an exterior form as a linear operation $`\mathrm{Tr}:\mathrm{\Lambda }()`$ or $`\mathrm{Tr}:\mathrm{\Lambda }_𝒞()𝒞`$ such that $`\mathrm{Tr}(e)=1`$ and $`\mathrm{Tr}(e^{\mu _1\mathrm{}\mu _k})=0,k=1,2,3,4`$. The reader can easily prove that
$$\mathrm{Tr}(UVVU)=0,\mathrm{Tr}(V^1UV)=\mathrm{Tr}U,U,V\mathrm{\Lambda }().$$
Now we can find the conservative current for the TDE. For this let us denote
$$C=\mathrm{\Phi }^{}(\mathrm{{\rm Y}}\mathrm{\Phi }+A\mathrm{\Phi }I+m\mathrm{\Phi }HI),\overline{\mathrm{\Phi }}=H\mathrm{\Phi }^{}.$$
Then
$$HC=\overline{\mathrm{\Phi }}(e^\mu _\mu \mathrm{\Phi }+A\mathrm{\Phi }I+m\mathrm{\Phi }HI),HC^{}=(_\mu \overline{\mathrm{\Phi }}e^\mu I\overline{\mathrm{\Phi }}AmIH\overline{\mathrm{\Phi }})\mathrm{\Phi }.$$
Using the formula $`\mathrm{Tr}(UVVU)=0`$, we get
$$\mathrm{Tr}(I\overline{\mathrm{\Phi }}A\mathrm{\Phi }\overline{\mathrm{\Phi }}A\mathrm{\Phi }I)=0,\mathrm{Tr}(HI\overline{\mathrm{\Phi }}\mathrm{\Phi }\overline{\mathrm{\Phi }}\mathrm{\Phi }HI)=0.$$
(60)
Suppose $`\mathrm{\Phi }\mathrm{\Lambda }^{\mathrm{even}}()`$ is a solution of the TDE, then, with the aid of (60), we obtain
$$0=\mathrm{Tr}(H(C+C^{}))=\mathrm{Tr}(\overline{\mathrm{\Phi }}e^\mu _\mu \mathrm{\Phi }+_\mu \overline{\mathrm{\Phi }}e^\mu \mathrm{\Phi })=\mathrm{Tr}(_\mu (\overline{\mathrm{\Phi }}e^\mu \mathrm{\Phi }))=_\mu j^\mu ,$$
where
$$j^\mu =\mathrm{Tr}(\overline{\mathrm{\Phi }}e^\mu \mathrm{\Phi }).$$
Therefore the vector $`j^\mu `$ is the conservative current. If we take the 1-form
$$J=g_{\mu \nu }j^\nu e^\mu =j_\mu e^\mu =\mathrm{\Phi }\overline{\mathrm{\Phi }},$$
then the divergence $`_\mu j^\mu =0`$ can be rewritten in the form
$$\delta J=0.$$
Finally let us define the Lagrangian from which the TDE can be derived
$$Lagr_1=\mathrm{Tr}(HCI).$$
Adding the term that describes the free field $`A`$ to $`Lagr_1`$, we obtain
$$Lagr=Lagr_1+\mathrm{Tr}(F^2),$$
(61)
where $`F=dA`$ is a 2-form, $`\mathrm{Tr}(F^2)=\frac{1}{2}f^{\mu \nu }f_{\mu \nu }`$, $`f_{\mu \nu }=_\mu a_\nu _\nu a_\mu `$. Hence the Lagrangian $`Lagr`$ depends on the following exterior forms: $`\mathrm{\Phi }\mathrm{\Lambda }^{\mathrm{even}}()`$, $`\overline{\mathrm{\Phi }}\mathrm{\Lambda }^{\mathrm{odd}}()`$, $`A\mathrm{\Lambda }^1()`$. Using the variational principle we suppose that the exterior forms $`\mathrm{\Phi }`$ and $`\overline{\mathrm{\Phi }}`$ are independent and as the variational variables we take the 8 functions which are the coefficients of the exterior form $`\overline{\mathrm{\Phi }}`$ and the 4 functions which are the coefficients of the exterior form $`A`$. The Lagrange-Euler equations with respect to these variables give us the system of equations, which can be written in the form
$`(d\delta )\mathrm{\Phi }+A\mathrm{\Phi }I+m\mathrm{\Phi }HI=0,`$
$`dA=F,`$ (62)
$`\delta F=J,`$
where $`J=\mathrm{\Phi }\overline{\mathrm{\Phi }}=\mathrm{\Phi }H\mathrm{\Phi }^{}`$. This system of equation can also be written in the form
$`e^\mu (_\mu \mathrm{\Phi }+a_\mu \mathrm{\Phi }I)+m\mathrm{\Phi }HI=0,`$
$`_\mu a_\nu _\nu a_\mu =f_{\mu \nu },`$ (63)
$`_\mu f^{\mu \nu }=j^\nu ,`$
where $`j^\nu =\mathrm{Tr}(\overline{\mathrm{\Phi }}e^\nu \mathrm{\Phi })`$, $`f^{\mu \nu }=g^{\mu \lambda }g^{\nu ϵ}f_{\lambda ϵ}`$.
The Lagrangian (61) and the systems of equations (62),(63) are invariant under the gauge transformation with the symmetry group $`\mathrm{U}(1)`$
$$\mathrm{\Phi }\mathrm{\Phi }\mathrm{exp}(\lambda I),AAd\lambda ,$$
(64)
where $`\lambda =\lambda (x)\mathrm{\Lambda }^0()`$ and $`\mathrm{exp}(\lambda I)=e\mathrm{cos}\lambda +I\mathrm{sin}\lambda `$. And they are also invariant under the global (independent of $`x`$) transformation (59) with the symmetry group $`\mathrm{Spin}()`$. The gauge invariance (64) expresses the interaction of the electron (fermion) with an electromagnetic field. The global invariance (59) leads to unusual (compared to tensors) transformation properties of bispinors under Lorentz changes of coordinates.
The Dirac equation can be generalized to the (pseudo) Riemannian space $`𝒱`$ using a special technique known as the tetrad formalism. We suppose that the TDE gives another possibility to describe the electron in the presence of gravity. Here we must take into account the fact that in Riemannian space there is no invariance of equations under global transformations of the form (59). Consequently we must use a gauge (local) transformation instead of the global transformation (59). This leads to a new gauge field with the $`\mathrm{Spin}(𝒱)`$ symmetry group, which we interpret as the gravitational field. For details of such an approach see .
### 2.4 Other tensor equations.
Consider the Ivanenko-Landau-Kähler equation ,
$$\mathrm{{\rm Y}}\rho +iA\rho +im\rho =0,$$
(65)
where $`\rho =\rho (x)\mathrm{\Lambda }_𝒞()`$ and $`A=A(x)\mathrm{\Lambda }^1()`$. We arrive at the TDE (58) by multiplying (65) from the right by the idempotent
$$t=\frac{1}{4}(e+H)(eiI)$$
and using the relations
$$t=tH,it=tI.$$
Similarly, multiplying (65) by
$$t=\frac{1}{2}(e+H)$$
and using the relation $`t=tH`$, we arrive at the following equation:
$$\mathrm{{\rm Y}}\eta +iA\eta +im\eta H=0,$$
(66)
where $`\eta \mathrm{\Lambda }_𝒞^{\mathrm{even}}()`$. In the same way, multiplying (65) by
$$t=\frac{1}{2}(eie^5),(e^5=e^{0123}=e^0e^1e^2e^3=e^0e^1e^2e^3),$$
and using the relation $`it=te^5`$ we arrive at the equation
$$\mathrm{{\rm Y}}\omega +A\omega e^5+m\omega e^5=0,$$
(67)
where $`\omega \mathrm{\Lambda }()`$. Evidently, (65),(66),(67) are tensor equations.
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# Perron-Frobenius theory for positive maps on trace ideals
## 1. Introduction
In the theory of quantum information the transmission through noisy channels plays an important role. Usually it is described by what physicists either call a quantum operation (see e.g. ) or a stochastic map (, see also ) or a super-operator (see e.g. ) and what mathematicians call a completely positive (trace preserving) map. In mathematics positive maps were first studied by Kadison in the context of $`C^{}`$-algebras and completely positive maps by Stinespring . In quantum physics it was first introduced by Haag and Kastler, who called it an operation . They were then studied in more detail in . contains an extensive discussion, how this concept naturally arises in quantum physics.
Because they preserve positivity such maps naturally fit into a context where it makes sense to ask about a formulation of a corresponding Perron-Frobenius theory.
As the name indicates, the earliest results relating positivity of a finite dimensional linear map to the non-degeneracy of an eigenvalue and the positivity of the resulting (Perron-Frobenius) eigenvector is due to Perron and Frobenius . The first extension to the infinite dimensional was given by Jentzsch . Since then there has been an extensive development (see e.g. ).
The interest in quantum physics and quantum field theory originates from the observation, that very often the ground state of a quantum system with Hamilton operator $`H0`$ is non-degenerate and nowhere vanishing. Indeed, under suitable circumstances such a ground state can be viewed as the Perron-Frobenius vector for $`\mathrm{exp}tH,t>0`$, an observation first made by Glimm and Jaffe (see also for an overview and with references to other articles).
The Perron-Frobenius theory and related topics for positive linear maps on linear spaces of operators was first studied in , where in the first article the space of operators was the $`C^{}`$-algebra of all linear operators on a finite dimensional Hilbert space. In the second and third articles the discussion was extended to the infinite dimensional case and in the fourth article the analysis was carried out in the context of a von Neumann algebra.
The present article can be viewed as an extension of the discussion in to the infinite dimensional case. Indeed, we will consider positive maps on the Schatten classes $`𝒥_p,\mathrm{\hspace{0.33em}1}p\mathrm{}`$ and provide a Perron-Frobenius theory on such spaces. As in the usual context the notions of ergodicity and positivity improving will play an essential role. Since density matrices are elements of $`𝒥_1`$ our main focus will be on the case $`p=1`$. In particular we will show that any compact, ergodic, completely positive and trace preserving map in $`𝒥_1`$ has a unique density matrix invariant under this map (i.e. the eigenvalue equals 1) and which therefore can be viewed as the Perron-Frobenius vector for this map.
The article is organized as follows. In Section 2 we will briefly review the concepts needed for our discussion. In Section 3 we will prove a Perron-Frobenius theorem (non-degeneracy of a certain eigenvalue), provided such an eigenvalue exists, a typical condition needed in infinite dimensional contexts. Section 4 will provide sufficient conditions for such an eigenvalue to exist. In Section 5 we will give examples. Also we will prove a converse in the sense that to any density matrix $`>0`$ there are completely positive maps for which this density matrix is the Perron-Frobenius eigenvector.
## 2. Positive and Completely Positive Maps
We start by recalling and establishing some definitions and facts related to trace ideals and to the concept of positive and completely positive maps ( for more details on trace ideals see e.g. and on positive and completely positive maps e.g. ). Let $``$ be any complex, separable Hilbert space with scalar product denoted by $`,`$. By $`()`$ we denote the $`C^{}`$-algebra of all bounded operators on $``$ equipped with the norm $`||||`$. By $`𝕀()`$ we denote the identity map on $``$.
For any $`A()`$ we set $`|A|=(A^{}A)^{1/2}()`$ where denotes the adjoint. Also we let $`\mathrm{Tr}`$ be the trace operation on $``$. For $`1p<\mathrm{}`$ let $`𝒥_p=𝒥_p()`$ be the Schatten class consisting of all elements $`A()`$ such that $`|A|^p`$ is trace class. $`𝒥_p`$ is equipped with the norm $`A_p=(\mathrm{Tr}(|A|^p))^{1/p}`$ making each $`𝒥_p`$ a complete Banach space. One has $`A^{}𝒥_p`$ if $`A𝒥_p`$, i.e. the space $`𝒥_p`$ is self-adjoint in $`()`$, such that $`A^{}_p=A_p`$. Also we set $`𝒥_{\mathrm{}}=()`$ again with the same norm, i.e. $`||||_{\mathrm{}}=||||`$. We will frequently make use of the trivial identity $`A_p=|A|_p`$ valid for all $`A𝒥_p`$ and all $`1p\mathrm{}`$. Via the trace $`𝒥_q`$ is the dual of $`𝒥_p`$ with $`1/p+1/q=1`$ for all $`1p<\mathrm{}`$, i. e. each continuous linear functional on $`𝒥_p`$ is given by an $`A𝒥_q`$ in the the form $`\mathrm{Tr}(A^{}B),B𝒥_p`$. In particular the spaces $`𝒥_p`$ are reflexive for all $`1<p<\mathrm{}`$. We prefer to use the adjoint in describing linear functionals since for $`p=2`$, in which case $`𝒥_2`$ is the space of Hilbert-Schmidt operators on $``$, the scalar product making it a Hilbert space is just given by $`A,B_2=\mathrm{Tr}(A^{}B)`$. Also one has the Hölder inequality $`|\mathrm{Tr}(A^{}B)|A_qB_p`$. Let $`𝒞=𝒞()=𝒞_{\mathrm{}}`$ be the closed positive cone of all elements $`A0`$ in $`()`$. As usual we write $`A0`$ and $`A>0`$ if the relations $`\phi ,A\phi \mathrm{\hspace{0.17em}0}`$ and $`\phi ,A\phi >\mathrm{\hspace{0.17em}0}`$ respectively hold for all $`0\phi `$. Correspondingly we write $`A>B`$ (or $`B<A`$) and $`AB`$ (or $`BA`$) if $`AB>0`$ and $`AB0`$ respectively. Obviously $`AB>0`$ or $`A>B0`$ implies $`A>0`$. As is common, $`A`$ is said to be positive definite if $`A>0`$. Set $`𝒞_p=𝒞_p()=𝒞𝒥_p`$, a closed set in $`𝒥_p`$. Also $`A0`$ in $`𝒥_p`$ if and only if $`\mathrm{Tr}(BA)0`$ for all $`0B𝒥_q`$. Correspondingly $`A>0`$ in $`𝒥_p`$ if and only if $`\mathrm{Tr}(BA)>0`$ for all $`0B0`$ in $`𝒥_q(1/p+1/q=1,\mathrm{\hspace{0.17em}1}p\mathrm{})`$. In this sense the cones $`𝒞_p`$ and $`𝒞_q`$ are dual to each other. The closed set $`𝒞_{1,1}=\{A𝒞_1|\mathrm{Tr}A=A_1=1\}`$ in $`𝒥_1`$ is the set of all density matrices in $``$.
By definition a positive map $`\varphi `$ in $`𝒥_p`$ is a linear map from $`𝒥_1`$ into itself, which leaves $`𝒞_p`$ invariant. $`\varphi `$ is called $`n`$-positive if the induced map $`\varphi _n=\varphi 𝕀_n`$ in $`𝒥_p(_n)`$ also leaves the corresponding cone $`𝒞_p(_n)`$ of non-negative elements invariant. Here $`_n`$ is any Hilbert space of dimension $`1n<\mathrm{}`$. Obviously if $`\varphi `$ is $`n`$-positive, then it is also $`n^{}`$-positive for any $`1n^{}n`$ as is $`\lambda \varphi ,\lambda >0`$. If $`\varphi `$ and $`\varphi ^{}`$ are both $`n`$-positive on $`𝒥_p`$ then so is their composite $`\varphi \varphi ^{}`$ and their sum $`\varphi +\varphi ^{}`$. So the $`n`$-positive maps in $`𝒥_p`$ form a cone in the linear space of all linear maps in $`𝒥_p`$.
If $`\varphi `$ is $`n`$-positive for all $`n`$, then $`\varphi `$ is called completely positive. Thus the map $`\varphi _\alpha :A\alpha A\alpha ^{}`$ for any $`\alpha ()`$ is completely positive on $`𝒥_p`$ for all $`1p\mathrm{}`$ as is any finite linear combination $`\varphi _{\underset{¯}{\alpha }}=_i\varphi _{\alpha _i}`$ with $`\alpha _i()`$. If $`\alpha `$ has an inverse $`\alpha ^1`$ in $`()`$ then $`\varphi _\alpha ^1=\varphi _{\alpha ^1}`$. Also $`\varphi _\alpha \varphi _\alpha ^{}=\varphi _{\alpha \alpha ^{}}`$ and $`\varphi _{\lambda \alpha }=|\lambda |^2\varphi _\alpha `$ for $`\alpha ,\alpha ^{}()`$ and $`\lambda `$. As it turns out at least for $`p=1`$ the induced linear map $`\varphi 𝕀_{^{}}`$ from $`𝒥_{p=1}(^{})`$ then also leaves the corresponding cone invariant for any separable Hilbert space $`^{}`$ (see e.g. ). Since density matrices are in $`𝒥_1`$, the case $`p=1`$ is of most interest in quantum physics. Then there are good physical reasons to consider completely positive maps rather than only positive maps (see e.g, ).
Although the following lemma is well known in similar contexts, we still will provide the short proof.
###### Lemma 2.1.
The following relation holds if the map $`\varphi `$ in $`𝒥_p(1p\mathrm{})`$ is positive.
(1)
$$\varphi (A)^{}=\varphi (A^{}),A𝒥_p.$$
###### Proof.
We first consider the case when $`A`$ is self-adjoint, i.e. $`A^{}=A`$. Write $`A=A_+A_{}`$ with $`A_\pm 𝒞_p,A_+A_{}=A_{}A_+=0`$ and $`|A|=A_++A_{}`$. Here $`\pm A_\pm `$ are the positive and negative parts of $`A`$, obtainable from the spectral representation of $`A`$ or more explicitly as $`A_\pm =1/2(|A|\pm A)`$. Obviously $`A_\pm _pA_p`$. Since $`\varphi (A_\pm )0`$, $`\varphi (A)=\varphi (A_+)\varphi (A_{})`$ is self-adjoint. For arbitrary $`A()`$ write $`A=\mathrm{}A+i\mathrm{}A()`$ with $`\mathrm{}A=1/2(A+A^{}),\mathrm{}A=1/2i(AA^{})`$, such that $`\mathrm{}A,\mathrm{}A()`$ are self-adjoint with $`\mathrm{}A,\mathrm{}A𝒥_p`$ whenever $`A𝒥_p`$. More precisely we have the a priori bound $`\mathrm{}A_pA_p,\mathrm{}A_pA_p`$. This gives the decomposition
$$A=(\mathrm{}A)_+(\mathrm{}A)_{}+i((\mathrm{}A)_+(\mathrm{}A)_{})$$
and hence by the linearity of $`\varphi `$
(2)
$$\varphi (A)=\varphi ((\mathrm{}A)_+)\varphi ((\mathrm{}A)_{})+i\varphi ((\mathrm{}A)_+)i\varphi ((\mathrm{}A)_{})$$
with $`\varphi (A)_p\varphi ((\mathrm{}A)_+)_p+\varphi ((\mathrm{}A)_{})_p+\varphi ((\mathrm{}A)_+)_p+\varphi ((\mathrm{}A)_{})_p`$. In particular (2) shows that indeed (1) holds. ∎
We also note the following. The relation $`|A|A|A|`$ for $`A=A^{}𝒥_p`$ implies $`\varphi (|A|)\varphi (A)\varphi (|A|)`$, whenever the map $`\varphi `$ in $`𝒥_p`$ is positive. This in turn gives
(3)
$$|\varphi (A)|\varphi (|A|)$$
valid for any $`A=A^{}𝒥_p`$ and any positive map $`\varphi `$ in $`𝒥_p`$. We do not know whether (or when) this relation continues to hold when the condition $`A=A^{}`$ is dropped (see, however, (23) below for a weaker result, which will suffice for our purposes). Observe by comparison that for any $`n\times n`$ matrix $`S`$ with non-negative entries, the classical context for the Perron-Frobenius theorem, one has $`|S\underset{¯}{z}|S|\underset{¯}{z}|`$ for $`\underset{¯}{z}^n`$. Here $`|\underset{¯}{z}|`$ is the vector whose components are the absolute values of the corresponding components of $`\underset{¯}{z}`$. Also for real vectors $`\underset{¯}{x}=\{x_i\}`$ and $`\underset{¯}{y}=\{y_i\}`$ by definition $`\underset{¯}{x}\underset{¯}{y}`$ if and only if $`x_iy_i`$ holds for all $`i`$.
Furthermore one does not necessarily have $`\varphi (A)_\pm =\varphi (A_\pm )`$ for any $`A=A^{}𝒥_p`$. However, since $`A_{}AA_+`$ implies $`\varphi (A_{})\varphi (A)\varphi (A_+)`$, the inequalities $`\varphi (A)_\pm \varphi (A_\pm )`$ hold.
###### Lemma 2.2.
Any positive map $`\varphi `$ in $`𝒥_p(1p\mathrm{})`$ is continuous, i.e. it satisfies $`\varphi _p<\mathrm{}`$.
Here we denote by $`\varphi _p`$ the norm of any continuous linear map $`\varphi `$ in $`𝒥_p`$, such that
(4)
$$\varphi _p=\underset{A_p1}{sup}\varphi (A)_p.$$
In particular
(5)
$$\varphi _1=\underset{A𝒞_{1,1}}{sup}\mathrm{Tr}(\varphi (A)),$$
if $`\varphi `$ is completely positive on $`𝒥_1`$ (see e.g. ).
###### Proof.
We adapt a standard proof (see e.g. , p. 19) used in the context of positive maps on non-unital $`C^{}`$-algebras. We first claim that it suffices to prove boundedness on $`𝒞_p`$. Indeed, this follows from the decomposition (2) and the related bounds. Assume that $`\varphi `$ is not bounded on $`𝒞_p`$. Then there are $`A_n𝒞_p`$ with $`A_n_p1`$ and $`\varphi (A_n)_pn^3`$. Let $`A=_n1/n^2A_n𝒞_p`$, such that $`01/n^2A_nA`$ holds for all $`n`$.
We need the fact that
(6)
$$B_pB^{}_p$$
holds for any $`0BB^{}𝒞_p`$. The case $`p=\mathrm{}`$ is trivial. When $`1p\mathrm{}`$ take $`\phi _n`$ to be a complete orthonormal basis in $``$ diagonalizing $`B`$. Then we have
$`B_p^p`$ $`=`$ $`{\displaystyle \underset{n}{}}\phi _n,B\phi _n^p{\displaystyle \underset{n}{}}\phi _n,B^{}\phi _n^p`$
$``$ $`{\displaystyle \underset{n}{}}\phi _n,B^^p\phi _n=B^{}_p^p,`$
proving (6). Here we have used the estimate
(7)
$$|\phi ,A\phi |^p\phi ,|A|^p\phi ,$$
valid for any $`A=A^{}()`$, any normalized $`\phi `$ and any $`1p`$ (see e.g. , p.21). Since $`\varphi `$ is positive (6) gives $`n1/n^2\varphi (A_n)_p\varphi (A)_p`$, which is a contradiction. ∎
For any continuous linear map $`\varphi `$ in $`𝒥_p,\mathrm{\hspace{0.17em}1}p<\mathrm{}`$ let $`\varphi ^{}`$ be the “adjoint” continuous linear map from $`𝒥_q,\mathrm{\hspace{0.17em}1}/q+1/p=1`$ into itself given by the relation $`\mathrm{Tr}((\varphi ^{}(A))^{}B)=\mathrm{Tr}(A^{}\varphi (B))`$ for all $`A𝒥_q`$ and $`B𝒥_p`$. In particular $`\varphi ^{}=\varphi `$ holds whenever $`1<p<\mathrm{}`$. Also $`\varphi _{\underset{¯}{\alpha }}^{}=\varphi _{\underset{¯}{\alpha }^{}}`$ with $`\underset{¯}{\alpha }^{}=(\alpha _1^{},\mathrm{}\alpha _i^{},\mathrm{})`$ by the cyclicity of the trace.
By definition a continuous linear map $`\varphi `$ in $`𝒥_1`$ is called trace preserving if $`\mathrm{Tr}(\varphi (A))=\mathrm{Tr}(A)`$ holds for all $`A𝒥_1`$ and this in turn is equivalent to the relation $`\varphi ^{}(𝕀)=𝕀`$. Obviously trace preserving maps leave the set $`𝒞_{1,1}`$ of density matrices invariant. If $`\varphi `$ in $`𝒥_1`$ is completely positive and trace invariant, then by (5) $`\varphi _1=1`$.
All completely positive maps on $`𝒥_1`$ have the Kraus representation , a consequence of a theorem of Stinespring (see also for a proof in the finite dimensional case): Given such a $`\varphi `$ there is an at most denumerable set of elements $`\underset{¯}{\alpha }=\{\alpha _i\}_i`$ in $`()`$ satisfying
(8)
$$\underset{iK}{}\alpha _i^{}\alpha _i\varphi _1𝕀$$
for any finite subset $`K`$ such that $`\varphi =\varphi _{\underset{¯}{\alpha }}`$, again with
(9)
$$\varphi _{\underset{¯}{\alpha }}(A)=\underset{i}{}\alpha _iA\alpha _i^{},$$
which now may be an infinite sum. This representation is not unique. Conversely each such $`\varphi _{\underset{¯}{\alpha }}`$ is completely positive. More precisely $`\varphi _{\underset{¯}{\alpha }}`$ is defined as follows. Let $`\underset{¯}{\alpha }_N=(\alpha _1,\mathrm{}..,\alpha _N)`$. Then the $`\varphi _{\underset{¯}{\alpha }_N}`$ form a Cauchy sequence with respect to the norm $`||||_1`$ and $`\varphi _{\underset{¯}{\alpha }}`$ is defined as the limit (see e.g. ). If $`_i\alpha _i^{}\alpha _i=𝕀`$, such that $`\varphi _1=1`$, then $`\varphi _{\underset{¯}{\alpha }}`$ is trace preserving in $`𝒥_1`$. In the finite dimensional case $`(dim<\mathrm{})`$ the representation may always be chosen such that the index $`i`$ runs through a finite set of the order at most $`(dim)^2`$.
An important and interesting feature of completely positive maps in the context of quantum physics is that they not necessarily map pure states, i.e. one-dimensional orthogonal projections, into pure states.
Again we have the following representation for the adjoint of $`\varphi _{\underset{¯}{\alpha }}`$
(10)
$$\varphi _{\underset{¯}{\alpha }}^{}(A)=\underset{i}{}\alpha _i^{}A\alpha _i,$$
i.e. $`\varphi _{\underset{¯}{\alpha }}^{}=\varphi _{\underset{¯}{\alpha }^{}}`$. So $`\varphi _{\underset{¯}{\alpha }}`$ is trace preserving if and only if $`\varphi _{\underset{¯}{\alpha }}^{}(𝕀)=𝕀`$ and then $`\varphi _{\underset{¯}{\alpha }}_1=1`$. Observe that to any completely positive $`\varphi `$ with $`\varphi _11`$ or equivalently $`\varphi ^{}(𝕀)𝕀`$ ( this condition is again natural in the context of quantum physics, see e.g. ) we may in a natural way associate a trace preserving, completely positive map $`\widehat{\varphi }`$ given as $`\widehat{\varphi }(A)=\varphi (A)+\stackrel{~}{\alpha }A\stackrel{~}{\alpha }`$ where $`\stackrel{~}{\alpha }=(𝕀\varphi ^{}(𝕀))^{1/2}0`$.
An example to which we shall return below is when the $`\alpha _i`$’s are orthogonal projection operators which are pairwise orthogonal, i.e. satisfy
(11)
$$\alpha _i\alpha _j=\delta _{ij}\alpha _i=\delta _{ij}\alpha _i^{}.$$
Then $`\varphi _{\underset{¯}{\alpha }}^{}(𝕀)=𝕀`$ if and only if the $`\alpha _i`$ provide a decomposition of unity, i.e. if $`_i\alpha _i=𝕀`$ holds. As already remarked $`\varphi _{\underset{¯}{\alpha }}`$ is then also trace preserving. More generally if the $`\alpha _i`$ are all selfadjoint with $`_i\alpha _i^2=𝕀`$, then $`\varphi _{\underset{¯}{\alpha }}`$ is trace preserving.
We remark that there are some situations (see e.g. and which actually was the motivation for the present discussion) where one starts with completely positive maps of the following form. Let $`(\mathrm{\Omega },\mu )`$ be a measure space, i.e. $`\mu `$ is countably additive. Suppose we are given a measurable map $`\omega \alpha (\omega )`$ from $`\mathrm{\Omega }`$ into $`()`$ such that
(12)
$$\alpha (\omega )^2𝑑\mu (\omega )<\mathrm{}.$$
We then set
(13)
$$\varphi (A)=\alpha (\omega )A\alpha (\omega )^{}𝑑\mu (\omega ),$$
such that its adjoint takes the form
(14)
$$\varphi ^{}(A)=\alpha (\omega )^{}A\alpha (\omega )𝑑\mu (\omega ),$$
The condition corresponding to (8) is given as
(15)
$$\varphi ^{}(𝕀)=\alpha (\omega )^{}\alpha (\omega )𝑑\mu (\omega )c𝕀$$
for some $`0<c<\mathrm{}`$. To cast $`\varphi `$ in the form (9) such that (15) gives (8), let $`\chi _i`$ be any orthonormal basis in the separable Hilbert space $`L^2(\mathrm{\Omega },d\mu )`$ and set
$$\alpha _i=\overline{\chi _i}(\omega )\alpha (\omega )𝑑\mu (\omega ).$$
As an example one may choose $`\mathrm{\Omega }`$ to be a compact group, $`\mu `$ its Haar measure and $`\omega \alpha (\omega )`$ a unitary representation. Then the image of $`\varphi `$ consists of all trace class operators which commute with all $`\alpha (\omega )`$. Also $`\varphi `$ is trace preserving and an idempotent, i.e. $`\varphi ^2=\varphi `$. Moreover $`\varphi `$ is a conditional expectation when $`dim<\mathrm{}`$, i.e. $`\varphi (A\varphi (B))=\varphi (A)\varphi (B)=\varphi (\varphi (A)B)`$ holds for each $`A,B`$.
## 3. The Perron-Frobenius Theorem
Before we turn to a discussion of the Perron-Frobenius theorem we make some general remarks on $`\sigma _p(\varphi )`$ for the map $`\varphi `$ in $`𝒥_p`$, in particular when $`\varphi `$ is positive. Here $`\sigma _p(\varphi )`$ is the spectrum of $`\varphi `$, i.e. the set of all $`\lambda `$ for which $`(\lambda \varphi )`$ does not have a bounded inverse in $`𝒥_p`$. Let $`r_p(\varphi )`$ be the spectral radius of any bounded linear map $`\varphi `$ in $`𝒥_p`$ , i.e.
(16)
$$r_p(\varphi )=\underset{n\mathrm{}}{lim}\varphi ^n_p^{1/n}\varphi _p.$$
By a celebrated general result of Gelfand , one has $`sup|\sigma _p(\varphi )|=r_p(\varphi )`$. If the map $`\varphi `$ in $`𝒥_1`$ is trace preserving and completely positive then $`r_1(\varphi )=\varphi _1=1`$. Indeed, $`\varphi ^n`$ is then also completely positive and trace preserving such that by (5) $`\varphi ^n_1=1`$ holds for all $`n`$.
Lemma 2.1 allows us to draw the following conclusions on $`\sigma _p(\varphi )`$ for positive maps $`\varphi `$ in $`𝒥_p`$. First we observe that in general $`\sigma _p(\varphi )`$ is not contained in the real axis. In fact, in the finite dimensional case $`\sigma _p(\varphi )`$ if and only if $`\varphi ^{}=\varphi `$ and in the infinite dimensional case the same statement is valid when $`p=2`$. Let $`\rho _p(\varphi )`$ be the resolvent set of the map $`\varphi `$ in $`𝒥_p`$, i.e. the complement in $``$ of $`\sigma _p(\varphi )`$. If $`\lambda \rho _p(\varphi )`$, then to each $`A𝒥_p`$ there is a unique $`A^{}𝒥_p`$, such that $`(\lambda \varphi )(A^{})=A`$. Taking the adjoint and using (1) gives $`(\overline{\lambda }\varphi )(A^{^{}})=A^{}`$. Since $`A`$ can be chosen arbitrary also $`A^{}`$ can be chosen arbitrary giving a unique $`A^{^{}}`$, this shows that both sets $`\rho _p(\varphi )`$ and $`\sigma _p(\varphi )`$ lie symmetric with respect to the real axis. Let $`\lambda \sigma (\varphi )`$ be a real eigenvalue with eigenvector $`A`$: $`\varphi (A)=\lambda A`$. So by (1) $`A^{}`$ is also an eigenvector with the same eigenvalue. With the decomposition $`A=A_1+iA_2`$ we see that both self-adjoint operators $`A_1`$ and $`A_2`$ are eigenvectors, so any eigenspace for a real eigenvalue is spanned by self-adjoint elements. If in addition the map $`\varphi `$ in $`𝒥_1`$ is trace preserving and $`\lambda 1`$ an eigenvalue and $`A`$ a corresponding eigenvector, then from $`\mathrm{Tr}(A)=\mathrm{Tr}(\varphi (A))=\lambda \mathrm{Tr}(A)`$ we deduce $`\mathrm{Tr}(A)=0`$.
In the finite dimensional case $`dim<\mathrm{}`$ all spaces $`𝒥_p`$ are equal and of course $`\sigma _p(\varphi )`$ is independent of $`p`$. Since all norms $`||||_p`$ are equivalent, by its definition $`r_p(\varphi )`$ is also independent of $`p`$, as it should be.
To formulate the Perron-Frobenius theorem in the present context, we make the following definition, which is the just an adaption of the usual definition. Since any positive map $`\varphi `$ in $`𝒥_p`$ is bounded $`\mathrm{exp}t\varphi `$ is a well defined positive map in $`𝒥_p`$ for all $`t>0`$. Also its inverse $`\mathrm{exp}t\varphi `$ is a well defined bounded map in $`𝒥_p`$, not necessarily positive .
###### Definition 3.1.
A positive map $`\varphi `$ in $`𝒥_p(1p\mathrm{})`$ is positivity improving if $`\varphi (A)>\mathrm{\hspace{0.17em}0}`$ for any $`A0,A0`$. $`\varphi `$ is ergodic if for any $`A0,A0`$ there is $`t_A>0`$ with $`(\mathrm{exp}t_A\varphi )(A)>\mathrm{\hspace{0.17em}0}`$.
If $`\varphi `$ is ergodic then $`(\mathrm{exp}t\varphi )(A)(\mathrm{exp}t_A\varphi )(A)>0`$ for all $`tt_A`$. Obviously $`\varphi `$ is ergodic if it is positivity improving.
A simple necessary criterion for ergodicity is given by
###### Lemma 3.1.
If $`\varphi `$ is ergodic then $`Ker\varphi |_{C_p}=0`$.
###### Proof.
Assume that there is $`0A0`$ with $`\varphi (A)=0`$, such that $`(\mathrm{exp}t\varphi )(A)=A`$ for all $`t`$. If $`A`$ is not positive definite, then we have a contradiction, so let $`A>0`$. We claim there is $`0A^{}0`$ which is not positive definite with $`A^{}A`$. In fact, we may take $`A^{}`$ to be a one-dimensional orthogonal projection onto an eigenvector of $`A`$ times the corresponding eigenvalue, which we choose not to be zero. This gives $`0\varphi (A^{})\varphi (A)=0`$ and we are back to the first situation. ∎
Some remarks concerning the finite dimensional case are in order.
First we will show that when $`dim<\mathrm{}`$ the present definition of ergodicity of $`\varphi `$ is equivalent to the condition $`(1+\varphi )^{dim1}(A)>0`$ for any $`0A0`$, a criterion used in . Obviously this last condition implies ergodicity since $`\mathrm{exp}ta1/(n!)^2\mathrm{min}(1,t)(1+a)^n`$ for all $`n`$ and all $`t,a0`$. As for the converse assume there is $`0A0`$ such that $`(1+\varphi )^{dim1}(A)`$ is not positive definite. So there is an orthogonal projection $`P0,𝕀`$ and $`\lambda >0`$ such that $`\varphi (P)\lambda P`$ (see ). This gives $`(\mathrm{exp}t\varphi )(P)(\mathrm{exp}t\lambda )P`$ for all $`t>0`$, so $`\varphi `$ is not ergodic.
We have the following necessary and simple criterion for $`\varphi `$ to be ergodic.
###### Lemma 3.2.
If $`dim<\mathrm{}`$ then $`\varphi (𝕀)>0`$ for ergodic $`\varphi `$.
The converse is not true as may be seen by taking $`\varphi `$ to be given by (9), where the $`\alpha _i`$’s are taken to be a nontrivial decomposition of unity. Then $`\varphi (𝕀)=𝕀`$ whereas $`\varphi (\alpha _i)=\alpha _i`$ for all $`i`$, such that $`\varphi `$ is not ergodic.
###### Proof.
Assume that $`\varphi (𝕀)0`$ holds such that there is $`\phi `$ with $`\phi ,\varphi (𝕀)\phi =0`$. Then $`\phi ,\varphi (A)\phi =0`$ holds for all selfadjoint $`A`$ due to the bound $`A𝕀AA𝕀`$ and therefore $`\phi ,\varphi ^n(A)\phi =0`$ for all $`n1`$ and all $`A=A^{}`$. This in turn implies $`\phi ,(\mathrm{exp}t\varphi )(\varphi (A))\phi =0`$ for all $`t>0`$ and all $`A=A^{}`$. Since $`\varphi 0`$ by definition, there is $`A0`$ with $`0\varphi (A)0`$ (see the decomposition (2)), so $`\varphi `$ is not ergodic. ∎
Returning to the general case it is clear that $`\varphi _\alpha `$ given by (9) is not ergodic if the $`\alpha _i`$’s satisfy condition (11). Also in the context $`p=1`$ and with a continuous, one parameter family of unitary operators $`U(t)`$ defining a quantum mechanics on $``$ in terms of its infinitesimal generator $`H`$, then $`\varphi _{U(t)}`$ (which describes the time evolution in the Heisenberg picture) for fixed $`t`$ is not ergodic. If the completely positive map $`\varphi `$ in $`𝒥_p`$ is given by (13) in terms of a unitary representation of a compact group, which contains at least one irreducible representation with finite multiplicity, then by Schur’s Lemma $`\varphi `$ is easily seen to be ergodic ( and even positivity improving) if and only if the representation is irreducible. Then also $``$ is of course finite dimensional and the map $`\varphi `$ is a conditional expectation. The following two lemmas are almost obvious.
###### Lemma 3.3.
A completely positive map $`\varphi `$ in $`𝒥_1`$ represented as $`\varphi _{\underset{¯}{\alpha }}`$ is positivity improving if and only for any $`0\phi `$ the closed linear hull of the set of vectors $`\{\alpha _i^{}\phi \}`$ is all of $``$.
###### Proof.
It suffices to consider $`\varphi (P)`$, where $`P`$ is any one-dimensional orthogonal projection (compare the proof of Lemma 3.1). Let $`0\phi _0P`$ be a unit vector, so we have to consider
(17)
$$\phi ,\varphi (P)\phi =\underset{i}{}|\phi ,\alpha _i\phi _0|^2$$
and the claim follows. ∎
This proof shows that in the infinite dimensional case no $`\varphi _{\underset{¯}{\alpha }}`$ with only finitely many non-vanishing $`\alpha _i`$’s can be positivity improving. In fact, for given $`P`$ choose $`0\phi `$ to be orthogonal to all $`\alpha _i\phi _0`$. Then (17) vanishes and therefore $`\varphi (P)`$ is not positive definite.
Let $`𝒜(\underset{¯}{\alpha })`$ be the algebra (not -algebra) generated by the $`\alpha _i^{}`$’s.
###### Lemma 3.4.
A completely positive map $`\varphi `$ in $`𝒥_1`$ represented as $`\varphi _{\underset{¯}{\alpha }}`$ is ergodic if and only if every non-zero vector $`\phi `$ in $``$ is cyclic for $`𝒜(\underset{¯}{\alpha })`$, i.e. the strong closure of the linear space $`𝒜(\underset{¯}{\alpha })\phi `$ is all of $``$.
These two results apply to all Kraus representations $`\varphi _{\underset{¯}{\alpha }}`$ of $`\varphi `$.
###### Proof.
We use the same notation as in the proof of the previous lemma. In addition write $`\alpha _{\underset{¯}{i}}=\alpha _{i_1}\alpha _{i2}..\alpha _{i_n}`$ for $`\underset{¯}{i}=(i_1,i_2,..i_n)`$ and set $`|\underset{¯}{i}|=n`$. Then
$$\phi ,(\mathrm{exp}t\varphi )(P)\phi =\underset{n=0}{\overset{\mathrm{}}{}}\frac{t^n}{n!}\underset{\underset{¯}{i}:|\underset{¯}{i}|=n}{}|\phi ,\alpha _{\underset{¯}{i}}\phi _0|^2$$
and again the claim follows. ∎
The next theorem is part of the Perron-Frobenius theorem in the present context. It is part because it assumes that $`r_p(\varphi )`$ is an eigenvalue, a standard assumption one has to make in the infinite dimensional context if no other information is available.
###### Theorem 3.1.
Let the positive map $`\varphi `$ in $`𝒥_p(1p\mathrm{}`$) be ergodic and assume that $`r_p(\varphi )=\varphi _p`$. If $`r_p(\varphi )`$ is an eigenvalue of $`\varphi `$, then this eigenvalue is simple and the eigenvector $`A`$ may be chosen to be positive definite.
###### Remark 3.1.
According to usual terminology such a vector is called a Perron- Frobenius vector. When $`dim<\mathrm{}`$ then $`r_p(\varphi )`$ is independent of $`p`$ and is always an eigenvalue. This case has been covered by Evans and Hoegh-Krohn with an application to quantum Monte-Carlo processes and by Groh . At present we do not know whether in the infinite dimensional case the condition $`r_p(\varphi )=\varphi _p`$ can be dropped.
For the proof of the theorem we need the following simple
###### Lemma 3.5.
For any $`B>0`$ and $`0B^{}0`$ in $`𝒥_p,\mathrm{\hspace{0.17em}1}p\mathrm{}`$ the following strict inequality
(18)
$$BB^{}_p<B+B^{}_p$$
holds.
The condition $`B>0`$ cannot be weakened to $`B0`$ as the choice $`B=P0`$ and $`B^{}=Q0`$ for two finite dimensional orthogonal projections with $`PQ=0`$ shows.
###### Proof.
The estimate (18) for the case $`p=\mathrm{}`$ is trivial, so let $`p<\mathrm{}`$. Since $`BB^{}𝒥_p`$ is self-adjoint, there is a complete set orthonormal eigenvectors $`\phi _n`$ with real eigenvalues $`\lambda _n`$ of $`BB^{}`$. With $`\lambda _n=\phi _n,(BB^{})\phi _n`$ we obtain
(19)
$$BB^{}_p^p=\underset{n}{}|\lambda _n|^p=\underset{n}{}|\phi _n,B\phi _n\phi _n,B^{}\phi _n|^p.$$
Observe that
(20) $`|\phi _n,B\phi _n\phi _n,B^{}\phi _n|^p`$ $``$ $`\mathrm{max}(\phi _n,B\phi _n^p,\phi _n,B^{}\phi _n^p)`$
$``$ $`(\phi _n,B\phi _n+\phi _n,B^{}\phi _n)^p`$
$`=`$ $`\phi _n,(B+B^{})\phi _n^p`$
$``$ $`\phi _n,(B+B^{})^p\phi _n`$
holds, since $`B>0`$ and $`B^{}0`$. In the last estimate in (20) we used the estimate (7).
Furthermore $`\phi _n,B\phi _n>0`$ holds for all $`n`$. Also $`\phi _n,B^{}\phi _n>0`$ for at least one $`n`$. Otherwise we would also have $`\phi _n,B^{}\phi _n^{}=0`$ for all $`n,n^{}`$ by Schwarz inequality and hence $`B^{}=0`$, contradicting the assumption. Thus the first inequality in (20) is actually strict for at least one $`n`$. Inserting (20) into (19) and using this last observation proves (18). ∎
$`ProofoftheTheorem`$: Let $`A0`$ be any eigenvector for $`\varphi `$ with eigenvalue $`r_p(\varphi )`$. By an observation made above, we may assume that $`A`$ is self-adjoint. Observe the following estimate, valid for any $`t>0`$.
(21) $`\mathrm{exp}tr_p(\varphi )A_p`$ $`=`$ $`\mathrm{exp}tr_p(\varphi )A_p=(\mathrm{exp}t\varphi )(A)_p`$
$`=`$ $`|(\mathrm{exp}t\varphi )(A)|_p(\mathrm{exp}t\varphi )(|A|)_p`$
$``$ $`\mathrm{exp}t||\varphi ||_p|A|_p=\mathrm{exp}t||\varphi ||_p||A||_p.`$
Here we have used the estimate (6) in combination with (3) for the map $`\mathrm{exp}t\varphi `$.
Since by assumption $`r_p(\varphi )=\varphi _p`$, all inequalities in (21) actually have to be equalities. In particular for all $`t>0`$
(22) $`(\mathrm{exp}t\varphi )(A)_p`$ $`=`$ $`(\mathrm{exp}t\varphi )(A_+)(\mathrm{exp}t\varphi )(A_{})_p`$
$`=`$ $`(\mathrm{exp}t\varphi )(|A|)_p=(\mathrm{exp}t\varphi )(A_+)+(\mathrm{exp}t\varphi )(A_{})_p.`$
If $`A_+=0`$ or $`A_{}=0`$ there is nothing to prove so assume that $`A_\pm 0`$. By the ergodicity of $`\varphi `$ we have $`(\mathrm{exp}t\varphi )(A_\pm )>0`$ for all large $`t>0`$. But then (22) contradicts (18). So we must have either $`A_+=0`$ or $`A_{}=0`$ and by replacing $`A`$ by $`A`$ if necessary we may without restriction assume that $`A=|A|=A_+`$. Again by ergodicity we have $`\mathrm{exp}tr_p(\varphi )A=(\mathrm{exp}t\varphi )(A)>0`$ for all large $`t`$ such that $`A>0`$. The proof of non-degeneracy is now easy. Assume there are two linearly independent eigenvectors $`A_1>0`$ and $`A_2>0`$, which we may normalize to $`A_1_p=A_2_p=1`$. Then $`0A_1A_2`$ is also an eigenvector and by our previous discussion we must have either $`A_1>A_2`$ or $`A_2>A_1`$. This would imply that $`A_1_p>A_2_p`$ or $`A_2_p>A_1_p`$ respectively, again an easy consequence of (6) and its proof. This is a contradiction, so we must have $`A_1=A_2`$, thus concluding the proof of the theorem.
## 4. Existence of eigenvalues
In this section we want in particular to establish conditions for positive maps $`\varphi `$ in $`𝒥_p`$, which are sufficient to show that $`r_p(\varphi )`$ is an eigenvalue. We start with a preparation. For any $`A()`$ we write $`A=U_A|A|`$ with unique $`U_A`$ for its polar decomposition. More precisely $`U_A`$ is isometric on $`\overline{\mathrm{Ran}|A|}=(KerA)^{}`$, i.e. $`U_A^{}U_A\phi =\phi `$ for $`\phi \overline{\mathrm{Ran}|A|}=(KerA)^{}`$, and $`U_A\phi =0`$ for $`\phi (\mathrm{Ran}|A|)^{}=KerA`$.
The next two lemmas replace the Kadison-Schwarz inequality (see e.g. the second reference in and for other early references) used in the context of $`C^{}`$-algebras in and in the context of von Neumann algebras in . Since the spaces $`𝒥_p,p<\mathrm{}`$ are not algebras these two lemmas will provide the appropriate substitute.
###### Lemma 4.1.
Let the map $`\varphi `$ in $`𝒥_p,\mathrm{\hspace{0.33em}1}p\mathrm{}`$ be 2-positive. Then for all $`\phi ,\phi ^{}`$ and all $`A𝒥_p`$
(23)
$$|\phi ^{},\varphi (A)\phi |^2\phi ,\varphi (|A|)\phi \phi ^{},\varphi \varphi _{U_A}(|A|)\phi ^{}.$$
Also if not all terms in (23) are vanishing, then there is equality if and only if both relations
(24)
$$\left(\begin{array}{cc}\varphi (|A|)& \varphi (A^{})\\ \varphi (A)& \varphi \varphi _{U_A}(|A|)\end{array}\right)\left(\begin{array}{cc}\phi ^{},\varphi \varphi _{U_A}(|A|)\phi ^{}& \phi \\ \phi ^{},\varphi (A)\phi & \phi ^{}\end{array}\right)=0$$
and
(25)
$$\left(\begin{array}{cc}\varphi (|A|)& \varphi (A^{})\\ \varphi (A)& \varphi \varphi _{U_A}(|A|)\end{array}\right)\left(\begin{array}{cc}& \phi ,\varphi (A)^{}\phi ^{}\phi \\ & \phi ,\varphi (|A|)\phi \phi ^{}\end{array}\right)=0$$
hold.
For $`A=A^{}`$ estimate (23) follows from (3), since then $`U_A^{}`$ commutes with $`A`$ such that $`U_A|A|U_A^{}=AU_A^{}=U_A{}_{}{}^{}A=U_A^{}U_A|A|=|A|`$.
###### Proof.
Here and in the proof of the next lemma we will mimic and extend Exercise 3.4 in , page 39 in the present context. Set $`A_1=|A|^{1/2},A_2=|A|^{1/2}U_A^{}`$ such that $`A_1^2=|A|,A_2^{}A_1=A,A_2^{}A_2=U_A|A|U_A^{}=\varphi _{U_A}(|A|)`$ are all in $`𝒥_p`$. Consider
(26)
$$\left(\begin{array}{cc}A_1& A_2\\ 0& 0\end{array}\right)^{}\left(\begin{array}{cc}A_1& A_2\\ 0& 0\end{array}\right)=\left(\begin{array}{cc}|A|& A^{}\\ A& \varphi _{U_A}(|A|)\end{array}\right)0$$
in $`^2`$, which defines an element in $`𝒞_p()`$. By assumption
(27)
$$\varphi _2\left(\begin{array}{cc}|A|& A^{}\\ A& \varphi _{U_A}(|A|)\end{array}\right)=\left(\begin{array}{cc}\varphi (|A|)& \varphi (A^{})\\ \varphi (A)& \varphi \varphi _{U_A}(|A|)\end{array}\right)0.$$
Consider the following linear transformation $`T`$ on $`^2`$ (equipped with the natural scalar product) given as the $`2\times 2`$ matrix
(28)
$$T=\left(\begin{array}{cc}\phi ,\varphi (|A|)\phi & \phi ,\varphi (A)^{}\phi ^{}\\ \phi ^{},\varphi (A)\phi & \phi ^{},\varphi \varphi _{U_A}(|A|)\phi ^{}\end{array}\right)$$
such that for all $`a_1,a_2`$
(29)
$$\left(\begin{array}{c}a_1\\ a_2\end{array}\right),T\left(\begin{array}{c}a_1\\ a_2\end{array}\right)=\left(\begin{array}{c}a_1\phi \\ a_2\phi ^{}\end{array}\right),\left(\begin{array}{cc}\varphi (|A|)& \varphi (A^{})\\ \varphi (A)& \varphi \varphi _{U_A}(|A|)\end{array}\right)\left(\begin{array}{c}a_1\phi \\ a_2\phi ^{}\end{array}\right)0$$
with the obvious notation for the scalar product in $``$. A linear transformation in $`^2`$ is $`0`$ if and only if its trace and its determinant are both $`0`$. Hence the bound (23) follows. This discussion also easily gives the last claim in Lemma 4.1. Indeed, if $`T0`$, its determinant is equal to zero if and only if there is a non-zero eigenvector of $`T`$ with eigenvalue zero and this is the case if and only if equality in (23) holds and then such an eigenvector of $`T`$ is given either as
$$\left(\begin{array}{c}\phi ^{},\varphi \varphi _{U_A}(|A|)\phi ^{}\\ \phi ^{},\varphi (A)\phi \end{array}\right)$$
or as
$$\left(\begin{array}{c}\phi ,\varphi (A)^{}\phi ^{}\\ \phi ,\varphi (|A|)\phi \end{array}\right).$$
More precisely, if all matrix elements of $`T`$ are non-vanishing then these two eigenvectors have non-zero entries and are proportional. If $`T0`$ with $`detT=0`$ and if one diagonal element is vanishing, such that also the off diagonal elements are vanishing, then exactly one of these vectors is the null vector. In view of (29) this concludes the proof of Lemma 4.1. ∎
###### Lemma 4.2.
Let the map $`\varphi `$ in $`𝒥_p,\mathrm{\hspace{0.33em}1}p\mathrm{}`$ be 2-positive. Then the bound
(30)
$$\varphi (A)_p\varphi (|A|)_p^{1/2}\varphi \varphi _{U_A}(|A|)_p^{1/2}$$
holds for all $`A𝒥_p`$. In particular $`AKer\varphi `$ if $`|A|Ker\varphi `$ or $`|A|Ker\varphi \varphi _{U_A}`$.
###### Proof.
Choosing $`\phi ^{}=U_{\varphi (A)}\phi `$ in (23) we have the estimate
(31)
$$\phi ,|\varphi (A)|\phi \phi ,\varphi (|A|)\phi ^{1/2}U_{\varphi (A)}\phi ,\varphi \varphi _{U_A}(|A|)U_{\varphi (A)}\phi ^{1/2}$$
for all $`\phi `$. Thus (30) follows trivially from (31) in case $`p=\mathrm{}`$. So let $`1p\mathrm{}`$. Take $`\phi _n`$ to be an orthonormal basis in $``$ diagonalizing $`|\varphi (A)|`$ with eigenvalues $`\mu _n0`$. Then we obtain
(32) $`\varphi (A)_p^p`$ $`=`$ $`{\displaystyle \underset{n}{}}\mu _n^p={\displaystyle \underset{n}{}}\phi _n,|\varphi (A)|\phi _n^p`$
$``$ $`{\displaystyle \underset{n}{}}\phi _n,\varphi (|A|)\phi _n^{p/2}U_{\varphi (A)}\phi _n,\varphi \varphi _{U_A}(|A|)U_{\varphi (A)}\phi _n^{p/2}`$
$``$ $`\left({\displaystyle \underset{n}{}}\phi _n,\varphi (|A|)^p\phi _n\right)^{1/2}`$
$`\left({\displaystyle \underset{n}{}}U_{\varphi (A)}\phi _n,(\varphi \varphi _{U_A}(|A|))^pU_{\varphi (A)}\phi _n\right)^{1/2}`$
$``$ $`\varphi (|A|)_p^{p/2}\varphi \varphi _{U_A}(|A|)_p^{p/2}.`$
Here we have used (31) for the first inequality. The second inequality is a consequence of Schwarz inequality and (7). In what follows this will become important, since we will encounter the situation when all inequalities are actually equalities. Also we have used the fact that with $`\phi _n`$ being an orthonormal basis also $`U_{\varphi (A)}\phi _n`$ is a possibly incomplete set of orthonormal vectors and which is responsible for the last inequality. ∎
$`\lambda `$ is called a peripheral eigenvalue of the positive map $`\varphi `$ in $`𝒥_p`$ if $`\lambda `$ and hence also $`\overline{\lambda }`$ is an eigenvalue with $`|\lambda |=r_p(\varphi )`$.
###### Theorem 4.1.
Let the map $`\varphi `$ in $`𝒥_p,\mathrm{\hspace{0.33em}1}p<\mathrm{}`$ be 2-positive with $`r_p(\varphi )=\varphi _p`$. Assume that $`\lambda `$ is a peripheral eigenvalue with eigenvector $`A`$. Then $`r_p(\varphi )`$ is also an eigenvalue with eigenvector $`|A|`$. In particular, if $`\varphi `$ is 2-positive and compact with $`r_p(\varphi )=\varphi _p`$, then $`r_p(\varphi )`$ is an eigenvalue.
Here and in contrast to the previous discussions the case $`p=\mathrm{}`$ is not covered by the present discussion. Since $`𝒥_{p=\mathrm{}}`$ is a $`C^{}`$-algebra, this case, however, is covered by the discussion in . Then and similarly in the context of von Neumann algebras one can actually prove more. In fact, if the spectrum is rescaled via $`\lambda \lambda /r_p(\varphi )`$, then the peripheral eigenvalues form a discrete sub-group of the unit circle group, see and the first reference in .
###### Proof.
We want to prove that $`|A|0`$ is also an eigenvector with eigenvalue $`|\lambda |=r_p(\varphi )`$, i.e. $`\varphi (|A|)=r_p(\varphi )|A|`$. Since $`|\varphi (A)|=|\lambda A|=r_p(\varphi )|A|`$ (such that also $`U_{\varphi (A)}=(\lambda /|\lambda |)U_A`$ and hence $`\varphi _{U_{\varphi (A)}}=\varphi _{U_A}`$ by the uniqueness of $`U_A`$) it suffices to show that
(33)
$$\varphi (|A|)\phi _n=\mu _n\phi _n$$
holds for all $`n`$, where the $`\phi _n`$ are as in the previous lemma. Also by the previous lemma we have
(34) $`r_p(\varphi )A_p`$ $`=`$ $`\lambda A_p=\varphi (A)_p=|\varphi (A)|_p`$
$``$ $`\varphi (|A|)_p^{1/2}\varphi \varphi _{U_A}(|A|)_p^{1/2}`$
$``$ $`\varphi _pA_p.`$
Here we have used the fact that $`\varphi _{U_A}_p1`$. By the assumption $`r_p(\varphi )=\varphi _p`$ we must have equality. In particular this implies that
(35)
$$|\varphi (A)|_p=\varphi (|A|)_p^{1/2}\varphi \varphi _{U_A}(|A|)_p^{1/2}.$$
Inspection of the proof of Lemma 4.2 shows that we must have equality in (31) when $`\phi =\phi _n`$ for all $`n`$. The relations (24) and (25) in Lemma 4.1 then imply that we must have
(36)
$$\phi _n^{},\varphi \varphi _{U_A}(|A|)\phi _n^{}\varphi (|A|)\phi _n=\phi _n^{},\varphi (A)\phi _n\varphi (A)^{}\phi _n^{}$$
and
(37)
$$\phi _n,\varphi (|A|)\phi _n\varphi \varphi _{U_A}(|A|)\phi _n^{}=\phi _n,\varphi (A)^{}\phi _n^{}\varphi (A)\phi _n$$
for all $`n`$ with $`\phi _n^{}=U_{\varphi (A)}\phi _n`$. Assume $`n`$ to be such that $`\mu _n>0`$. Since $`\phi _n\mathrm{Ran}|\varphi (A)|`$ we have $`U_{\varphi (A)}^{}\phi _n^{}=\phi _n`$ and $`\varphi (A)^{}\phi _n^{}=|\varphi (A)|\phi _n=\mu _n\phi _n`$ for all such $`n`$. Similarly $`\varphi (A)\phi _n=U_{\varphi (A)}|\varphi (A)|\phi _n=\mu _n\phi _n^{}`$. Hence we may rewrite (36) and (37) as
(38)
$$\phi _n^{},\varphi \varphi _{U_A}(|A|)\phi _n^{}\varphi (|A|)\phi _n=\mu _n^2\phi _n$$
and
(39)
$$\phi _n,\varphi (|A|)\phi _n\varphi \varphi _{U_A}(|A|)\phi _n^{}=\mu _n^2\phi _n^{}$$
whenever $`\mu _n>0`$.
Furthermore we recall that we made use of Schwarz inequality in the estimate (32). For this to be an equality we claim that we actually must have
(40)
$$\phi _n,\varphi (|A|)\phi _n=\phi _n^{},\varphi \varphi _{U_A}(|A|)\phi _n^{}$$
for all $`n`$. Indeed, the Schwarz inequality in this context is an equality if and only if both sides of (40) are proportional with a common proportionality factor for all $`n`$. By the equality (35) this proportionality factor has to be equal to 1, thus proving (40).
If $`\mu _n=0`$ then $`\phi _n^{}=0`$ and hence by (40) $`\varphi (A)\phi _n=0`$. Thus we have established (33) in the case when $`\mu _n=0`$. If $`\mu _n>0`$ then the left hand sides of (38) and (39) are both non-vanishing since then $`\phi _n^{}0`$. Therefore $`\phi _n`$ is an eigenvector of $`\varphi (|A|)`$ with eigenvalue
$`\phi _n,\varphi (|A|)\phi _n`$ $`=`$ $`{\displaystyle \frac{\mu _n^2}{\phi _n^{},\varphi \varphi _{U_A}(|A|)\phi _n^{}}}`$
$`=`$ $`{\displaystyle \frac{\mu _n^2}{\phi _n,\varphi (|A|)\phi _n}}.`$
This gives (33) in the case when $`\mu _n>0`$, thus concluding the proof of the theorem, since for compact operators the spectrum is pure point. ∎
The use of (7) in the second inequality in (32) does not give any additional mileage when $`p1`$ since the $`\phi _n`$ are eigenvectors of $`\varphi (|A|)`$ and the $`\phi _n^{}`$ eigenvectors of $`\varphi \varphi _{U_A}`$ and hence (7) gives an equality for the case at hand.
The strategy of the proof was motivated by the following observation in the classical context. With the notation introduced after (3) assume that $`S\underset{¯}{z}=\lambda \underset{¯}{z}`$. This gives $`|\lambda ||\underset{¯}{z}|=|\lambda \underset{¯}{z}|=|S\underset{¯}{z}|S|\underset{¯}{z}|`$. Taking its usual Hilbert space norm in $`^n`$ we obtain $`|\lambda ||\underset{¯}{z}|S|\underset{¯}{z}|S|\underset{¯}{z}|`$. So if $`|\lambda |=S`$ we must have $`S|\underset{¯}{z}|=S|\underset{¯}{z}|`$.
The next result extends Theorem 3.1.
###### Theorem 4.2.
Let the map $`\varphi `$ in $`𝒥_p,\mathrm{\hspace{0.33em}1}p<\mathrm{}`$ be 2-positive and ergodic with $`r_p(\varphi )=\varphi _p`$. Assume that $`\lambda `$ is a peripheral eigenvalue. Then $`\lambda `$ is a simple eigenvalue and the corresponding eigenvector $`A`$ is a normal operator.
###### Proof.
Assume there are two eigenvectors $`A_1`$ and $`A_2`$ in $`𝒥_p`$ with eigenvalue $`\lambda `$, normalized to $`A_1_p=|A_1|_p=A_2_p=|A_2|_p=1`$. By Theorem 3.1, its proof and by Theorem 4.1 we have that $`0<|A_1|=|A_2|`$ is an eigenvector with eigenvalue $`r_p(\varphi )`$. In particular both $`U_{A_1}`$ and $`U_{A_2}`$ are unitaries. Since $`A_1\mathrm{exp}i\tau A_2`$ is also an eigenvector for the same eigenvalue for all $`0\tau 2\pi `$ by the same argument we must have
(41)
$$|A_1\mathrm{exp}i\tau A_2|=A_1\mathrm{exp}i\tau A_2_p|A_1|.$$
The function $`\tau A_1\mathrm{exp}i\tau A_2_p`$ is easily seen to be continuous. Hence there is $`\tau _0`$ such that
(42)
$$A_1\mathrm{exp}i\tau _0A_2_p=\underset{0\tau 2\pi }{\mathrm{min}}A_1\mathrm{exp}i\tau A_2_p.$$
Take the square of (41) and use (42) to obtain
$$|A_1\mathrm{exp}i\tau _0A_2|^2|A_1\mathrm{exp}i\tau A_2|^2.$$
Due to the relation $`\overline{\mathrm{Ran}|A_1|}=`$ this gives
$$(U_{A_1}\mathrm{exp}i\tau _0U_{A_2})\phi ^2(U_{A_1}\mathrm{exp}i\tau U_{A_2})\phi ^2$$
which written out in turn gives
(43)
$$\mathrm{}\left(e^{i\tau }U_{A_1}\phi ,U_{A_2}\phi \right)\mathrm{}\left(e^{i\tau _0}U_{A_1}\phi ,U_{A_2}\phi \right)$$
for all $`0\tau 2\pi `$ and all $`\phi `$. We distinguish two possible cases. First if $`A_1\mathrm{exp}i\tau _0A_2=0`$ and hence $`U_{A_1}\mathrm{exp}i\tau _0U_{A_2}=0`$ there is nothing to prove. Otherwise $`U=U_{A_1}^1\mathrm{exp}i\tau _0U_{A_2}`$ is a unitary operator $`𝕀`$ and hence by the spectral theorem for unitary operators there is $`\phi `$ such that $`\mathrm{}(\phi ,U\phi )<|\phi ,U\phi |`$. Choose $`\tau `$ such that $`\mathrm{exp}i(\tau \tau _0)\phi ,U\phi =|\phi ,U\phi |`$. This contradicts (43), thus concluding the proof of the first part of the theorem. As for the last part $`U_A`$ is unitary, as already noted. Relation (39) shows that $`\varphi _{U_A}(|A|)`$ is also an eigenvector of $`\varphi `$ with eigenvalue $`r_p(\varphi )`$. Hence $`\varphi _{U_A}(|A|)=|A|`$ by uniqueness and normalization. So $`U_A`$ commutes with $`|A|`$ and therefore $`A`$ is normal with $`U_A^{}=U_A^{}`$ and $`|A^{}|=|A|`$, where we recall that $`A^{}`$ is an eigenvector with eigenvalue $`\overline{\lambda }`$. ∎
The next result is of relevance in quantum physics. We recall that for completely positive, trace preserving maps $`\varphi `$ in $`𝒥_1`$ one has $`r_1(\varphi )=\varphi _1=1`$. So we immediately obtain the following existence result for Perron-Frobenius vectors.
###### Corollary 4.1.
Let the map $`\varphi `$ in $`𝒥_1`$ be completely positive, trace preserving, compact and ergodic. Then there is a unique density matrix $`\rho >0`$ which is left invariant under $`\varphi `$.
A converse of this result is given by Theorem 5.1 below.
The following result is of the type needed in the standard context of Monte Carlo simulations in lattice models of statistical mechanics. It states that, starting from any non-negative initial configuration, by iterations of an ergodic up-dating procedure (heat bath method) one eventually arrives at the desired equilibrium configuration which typically is given by a Gibbs distribution and which by construction of the up-dating is a Perron-Frobenius vector for the up-dating procedure.
###### Corollary 4.2.
Let the map $`\varphi `$ in $`𝒥_2`$ be 2-positive, selfadjoint, compact and ergodic with Perron-Frobenius vector $`A>0`$ normalized to $`A_2=1`$. Then for any $`B𝒥_2`$
(44)
$$s\underset{t\mathrm{}}{lim}\mathrm{exp}(t(\varphi \varphi _2))(B)=A,B_2A.$$
If $`0B0`$ then the right hand side of (44) is non-vanishing. It would be interesting to find an analogous formulation when $`p2`$.
###### Proof.
We use familiar arguments. Let the $`A_j`$ form an orthonormal basis of eigenvectors for $`\varphi `$ arranged such that $`A_1=A`$ and $`\sigma _{j+1}\sigma _j`$. Here the $`\sigma _j`$’s are the eigenvalues, i.e. $`\varphi (A_j)=\sigma _jA_j`$. Since $`\sigma _1=\varphi _2`$ is a simple eigenvalue, $`\sigma _1>\sigma _2`$. By Plancherel’s theorem we have $`B=_jb_jA_j`$ with $`b_j=A_j,B_2`$ and $`_jb_j^2=B_2^2`$. Obviously
$$\mathrm{exp}(t(\varphi \varphi _2))(B)=\underset{j}{}b_je^{t(\sigma _j\sigma _1)}A_j$$
holds such that
(45)
$$(\mathrm{exp}(t(\varphi \varphi _2)))(B)A,B_2A_2e^{t(\sigma _1\sigma _2)}B_2$$
and the claim follows. ∎
By a slight modification of the proof using the spectral representation of $`\varphi `$ one may replace the compactness condition by the condition that $`\varphi _2`$ is an isolated eigenvalue, i.e. $`\varphi _2`$ is separated by a distance $`m^2`$ (the “mass gap” in physical language) from the remainder of the spectrum of $`\varphi `$, which otherwise may be arbitrary. Then the estimate (45) is still valid with $`\sigma _1\sigma _2`$ being replaced by $`m^2`$.
## 5. Examples
In this section we will provide non-trivial examples when $`dim=\mathrm{}`$. The issue is to show that all 4 conditions in Corollary 4.1 can be fulfilled simultaneously. In the finite dimensional case the compactness criterion is automatically satisfied and then it is easy to construct non-trivial examples. In general the existence of completely positive, trace preserving, compact maps in $`𝒥_1`$ is clarified by
###### Lemma 5.1.
Let $`\underset{¯}{\alpha }=\{\alpha _i\}_i`$ be such that $`_i\alpha _i^{}\alpha _i=𝕀`$ and all $`\alpha _i()`$ are compact. Then the map $`\varphi _{\underset{¯}{\alpha }}`$ in $`𝒥_1`$ is compact.
###### Proof.
We recall that $`\varphi _{\underset{¯}{\alpha }}`$ is the norm limit of $`\varphi _{\underset{¯}{\alpha }_N}`$ as $`N\mathrm{}`$. Since the space of compact operators is closed w.r.t. the norm topology it suffices to prove that all $`\varphi _{\underset{¯}{\alpha }_N}`$ are compact. By assumption to each $`i`$ there is a sequence $`\alpha _{i,k}`$ of finite rank operators in $`)`$ such that
$$\underset{k\mathrm{}}{lim}\alpha _i\alpha _{i,k}=\underset{k\mathrm{}}{lim}\alpha _i^{}\alpha _{i,k}^{}=0.$$
Since $`\alpha _{i,k}^{}`$ is of finite rank $`\varphi _{\underset{¯}{\alpha }_{N,k}}`$ with $`\underset{¯}{\alpha }_{N,k}=(\alpha _{1,k},\mathrm{},\alpha _{N,k})`$ is also of finite rank and in particular compact. Using the a priori estimate $`ABC_1AB_1C`$ and a standard $`ϵ/3`$ argument gives
$$\underset{k\mathrm{}}{lim}\varphi _{\underset{¯}{\alpha }_N}\varphi _{\underset{¯}{\alpha }_{N,k}}__1=0.$$
We are now prepared to provide non-trivial examples, which resulted from a discussion with D. Buchholz. Let $`\psi _i,i`$ be any orthonormal basis in $``$ and $`c_{ik}>0`$ any set of numbers which satisfy the condition $`_ic_{ik}^2=1`$ for all $`k`$. In the Dirac notation we define the following family $`\underset{¯}{\alpha }=\{\alpha _{ik}\}_{1i,k<\mathrm{}}`$ of rank 1 operators
$$\alpha _{ik}=c_{ik}|\psi _i\psi _k|,$$
such that
$$\alpha _{ik}^{}=c_{ik}|\psi _k\psi _i|,$$
In other words $`\alpha _{ik}`$ is the map $`\phi c_{ik}\psi _k,\phi \psi _i`$. We set
$$\varphi _{\underset{¯}{\alpha }}(A)=\underset{i,k}{}\alpha _{ik}A\alpha _{ik}^{}$$
By assumption we have
$$\underset{i,k}{}\alpha _{ik}^{}\alpha _{ik}=\underset{i,k}{}c_{ik}^2|\psi _k\psi _k|=\underset{k}{}|\psi _k\psi _k|=𝕀,$$
so $`\varphi _{\underset{¯}{\alpha }}`$ is completely positive, compact and trace preserving. We claim that it is also positivity improving by Lemma 3.3. Indeed, given $`0\phi `$, assume there is $`0\phi _0`$ such that
$$0=\phi _0,\alpha _{ik}^{}\phi =c_{ik}\phi _0,\psi _k\psi _i,\phi $$
holds for all $`i`$ and $`k`$. Since by assumption $`c_{ik}>0`$ this contradicts the fact that the $`\psi _i`$ form an orthonormal basis. So this $`\varphi _{\underset{¯}{\alpha }}`$ satisfies all conditions of Corollary 4.1. The resulting Perron-Frobenius vector is given as the density matrix
(46)
$$\rho =\underset{i}{}\rho _i|\psi _i\psi _i|,\rho _i>0,\underset{i}{}\rho _i=1,$$
where the $`\rho _i`$ form the components of the Perron-Frobenius vector with eigenvalue 1 for the matrix $`c_{ik}^2`$, i.e. $`_kc_{ik}^2\rho _k=\rho _i`$. As a special case consider the choice $`c_{ik}=c_i>0`$ for all $`k`$ with $`_ic_i^2=1`$. Then $`\rho _i=c_i^2`$.
This discussion also gives the following converse to Corollary 4.1, which in its standard formulation forms one of the starting points for Monte Carlo simulations in lattice models of statistical mechanics and (euclidean) quantum field theory including lattice gauge theories.
###### Theorem 5.1.
Given an arbitrary density matrix $`\rho >0`$, there is a completely positive, trace preserving, compact and ergodic map $`\varphi `$ in $`𝒥_1`$ for which $`\rho `$ is the resulting Perron-Frobenius vector. In addition $`\varphi `$ may be chosen to be idempotent (i.e. $`\varphi `$ satisfies the characteristic equation $`\varphi (1\varphi )=0`$) and $`\sigma _1(\varphi )=\{0,1\}`$.
###### Proof.
Take the $`\psi _i`$’s to form an orthonormal basis in which $`\rho `$ is diagonal, i.e. such that $`\rho `$ takes the form (46) and choose $`\varphi =\varphi _{\underset{¯}{\alpha }}`$ with $`\underset{¯}{\alpha }`$ as above with $`c_{ik}=\rho _i^{1/2}`$. Then we have $`\varphi (A)=\mathrm{Tr}A\rho `$ for all $`A𝒥1`$, such that $`\varphi `$ is idempotent. Also the kernel of $`\varphi `$ consists of all elements which may be written in the form $`A\mathrm{Tr}A\rho `$. Moreover let $`\lambda 0,1`$. Then for given $`A^{}`$ the equation $`(\lambda \varphi )(A)=A^{}`$ has a solution in $`A`$ of the form
$$A=\frac{1}{\lambda (1\lambda )}\mathrm{Tr}A^{}\rho +\frac{1}{\lambda }A^{}.$$
As in the corresponding ( finite dimensional) standard context $`\varphi `$ with prescribed Perron-Frobenius vector is of course not unique. We provide another example which in spirit is much closer to the usual “local” up-grading procedure in Monte Carlo simulations. Set $`p_{\nu ,i}=\rho _{\nu +i}/(\rho _{i+1}+\rho _i+\rho _{i1})`$ for $`\nu =0,\pm 1`$ and with the convention $`\rho _0=0`$. Then $`_{\nu =0,\pm 1}p_{\nu ,i}=1`$ for all $`i`$ and $`p_{\nu ,i}>0`$ for all $`\nu `$ and $`i`$ unless $`\nu =1`$ and $`i=1`$. Set
$$T_{\nu ,i}=p_{\nu ,i}^{1/2}|\psi _{\nu +i}\psi _i|$$
and define $`\varphi _T`$ as
$$\varphi _T(A)=\underset{\nu ,i}{}T_{\nu ,i}AT_{\nu ,i}^{}.$$
An easy calculation shows that $`\varphi _T`$ is completely positive, trace preserving, compact and ergodic.
Acknowledgments The author has profited from discussions with D. Buchholz, B. Kümmerer, M.B. Ruskai and E. Størmer.
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# Nonadiabatic Contribution to the Quasiparticle Self-Energy in Systems with Strong Electron-Phonon Interaction
## I Introduction
The standard description of electron-phonon interaction (EPI) in metals is based on the so-called Born-Oppenheimer ( or adiabatic) theorem. Its main statement says that electrons are not sensitive to the motion of ions and are influenced only by their static electric field. One can to say that ”fast” electrons, due to electroneutrality, follow ”slow” ions. Narrow bands and a strong electron-phonon interaction make their velocities comparable in order of magnitude. An appropriate mathematical description is the standard Feynman-Dyson perturbation theory, where phonons are considered to be uneffected by the EPI. However, there are some processes which in higher orders violate this approximation. Mathematically this can be taken into account by virtue of the so-called ”vertex corrections”. These effects were considered by Migdal who showed that their contributions are small, of the order of a small parameter $`\mathrm{\Omega }_{ph}/\stackrel{~}{W}`$, where $`\stackrel{~}{W}`$ is a characteric energy, defined by the bandwidth $`W`$ or the Fermi energy $`E_F`$.
It was noted by Migdal himself that at $`\stackrel{}{q}=0`$ the vertex corrections are not small (so they are important for optical conductivity and Raman scattering). This case was considered in detail by Engelsberg and Schrieffer . Migdal’s theorem is also violated for a one-dimensional Fermi surface .
Traditionally nonadiabatic corrections were considered to be small because the Migdal parameter $`\mathrm{\Omega }_{ph}/\stackrel{~}{W}`$ was small . But recently some materials (for example, fullerenes and high-$`T_c`$ superconductors) where $`\mathrm{\Omega }_{ph}\stackrel{~}{W}`$ were discovered. So Migdal’s theorem about the smallness of the nonadiabatic corrections can be violated. This gives grounds for not only discussing these corrections, but even examining an antiadiabatic limit $`\mathrm{\Omega }_{ph}\stackrel{~}{W}`$ . But even with small $`\mathrm{\Omega }_{ph}/\stackrel{~}{W}`$ there are contradictory conclusions about the importance of such corrections in the normal and superconducting state (the authors of Refs. consider them to be negligible, while the authors of think the opposite is true). The most intriguing question is whether the violation of adiabaticity leads to an increase of the interelectron interaction (and, as a result, to an increase of $`T_c`$). There are contradictory results about the effect of the vertex corrections on $`T_c`$. For example, Takada reported some increase of $`T_c`$ for small parameters $`\mathrm{\Omega }_{ph}/E_F`$. Even more significant increase was claimed in . At the same time, in Refs. a decrease of $`T_c`$ was obtained due to the vertex corrections.
One can distinguish between the two sorts of the nonadiabatic corrections to the electron self-energy: 1) the vertex ones, resulted from the higher order corrections of perturbation theory, as vertex function $`\mathrm{\Gamma }1`$ even at an infinite electron (or hole) band $`\mathrm{}<\epsilon <+\mathrm{}`$ , and 2) the corrections due to the finite band width, $`W<\epsilon <W`$ .
Nonadiabatic corrections in the superconducting state enter the expression for $`T_c`$ in two ways: indirectly through renormalization of $`Z`$ and directly in the equation for the order parameter $`\mathrm{\Delta }`$. At the same time, analytical calculation of anomalous (crossing) diagrams is rather difficult. Therefore before proceeding to calculations in the superconducting state we want to clarify the role of such corrections in the normal state. The main goal of this paper is to compare different approaches to calculation of nonadiabatic corrections to the electron quasiparticle self-energy in the normal state (or more precisely to the mass renormalization factor $`Z`$). There are two methods to take nonadiabatic corrections into account. One of them, which can be called Migdal’s one, is based on the solution to the Bethe-Salpeter equation for the vertex function in the ladder approximation. In the lowest approximation, the first correction to the unity vertex is determined by the diagram in Fig.1. Then with the obtained vertex function the self-energy $`\mathrm{\Sigma }`$ is calculated . This traditional method does not allow to take into account higher-order diagrams analytically because of complexity of the integration over momenta.
In 1989 the authors of the Ref. considered a non-ladder approximation, using the Ward identity. Later Y. Takada used the same idea in his method , which he called the gauge-invariant self-consistent (GISC) method. In this method, based on the Ward identity
$`i\omega _\nu \mathrm{\Gamma }(i\omega _n,i\omega _ni\omega _\nu ,\stackrel{}{k},\stackrel{}{k}\stackrel{}{q})\stackrel{}{q}\stackrel{}{\mathrm{\Gamma }}(i\omega _n,i\omega _ni\omega _\nu ,\stackrel{}{k},\stackrel{}{k}\stackrel{}{q})`$ (1)
$`=`$ $`G^1(i\omega _n,\stackrel{}{k})G^1(i\omega _ni\omega _\nu ,\stackrel{}{k}\stackrel{}{q}),`$ (2)
the vector term is neglected and the scalar vertex function is chosen to be a functional of the self-energy, which is supposed to be independent from momentum. Recently there appeared another method , aiming at improving the GISC approach as it takes into account the momentum dependence of the vertex function. So, when comparing these methods with the Migdal’s one, we also discuss the validity of these assumptions. In fact, for the self-energy we calculate all the diagrams up to, and including the second order. The parameter $`\mathrm{\Omega }_{ph}/Wm_0`$ is taken to be small.
For simplicity, we consider the model of the Fermi liquid with the usual electron-phonon Hamiltonian and Einstein phonon spectrum with the so-called Eliashberg spectral function $`\alpha ^2F_E(\mathrm{\Omega })=\frac{1}{2}\lambda \mathrm{\Omega }_{ph}\delta (\mathrm{\Omega }\mathrm{\Omega }_{ph})`$, where $`\mathrm{\Omega }_{ph}`$ is the phonon frequency independent from momentum, and $`\lambda `$ is the electron-phonon interaction constant. It is assumed that the density of states is constant $`N(0)`$ for $`W<\epsilon <W`$ where $`2W`$ is the bandwidth. The chemical potential is zero, which corresponds to a half-filled band. It is assumed that $`\lambda 1`$.
We are interested in the quasiparticle self-energy $`\mathrm{\Sigma }`$ on the imaginary axis which is determined by the diagram in Fig.1. The phonon Green’s function is $`D_0(i\omega _\nu )=\mathrm{\Omega }_{ph}^2/(\omega _\nu ^2+\mathrm{\Omega }_{ph}^2)`$, where $`\omega _\nu =2\nu \pi T`$, and the electron Green’s function is $`G(i\omega _n,\stackrel{}{p})=(i\omega _n\epsilon _\stackrel{}{p}\mathrm{\Sigma }(i\omega _n,\stackrel{}{p}))^1`$, where $`\omega _n=(2n+1)\pi T`$, $`\epsilon _\stackrel{}{p}`$ is the bare electron spectrum.
The vertex function itself (two-point one) is only of academic interest. We are, however, interested in the quasiparticle self-energy $`\mathrm{\Sigma }`$ which affects many physical observables. Assuming that $`\mathrm{\Sigma }(i\omega _n,\stackrel{}{k})`$ weakly depends on $`\stackrel{}{k}`$, we take $`|\stackrel{}{k}|=p_F=const`$, and then
$`\mathrm{\Sigma }(i\omega _n)`$ $`=`$ $`{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑\mathrm{\Omega }\alpha ^2F_E(\mathrm{\Omega })T{\displaystyle \underset{\omega _\nu }{}}{\displaystyle \frac{2\mathrm{\Omega }}{\omega _\nu ^2+\mathrm{\Omega }^2}}{\displaystyle \frac{1}{N(0)}}`$ (4)
$`\times {\displaystyle \underset{\stackrel{}{q}}{}}G(i\omega _ni\omega _\nu ,\stackrel{}{k}\stackrel{}{q})\mathrm{\Gamma }(i\omega _n,i\omega _ni\omega _\nu ,\stackrel{}{q}).`$
The methods differ in the choice of the vertex function $`\mathrm{\Gamma }`$, and this will be discussed later. One should note that $`\mathrm{\Gamma }(\stackrel{}{q})`$ enters the integral for $`\mathrm{\Sigma }`$, and all $`\stackrel{}{q}`$’s contribute to it. For simplicity we will obtain results for $`T=0`$ and the variables will be changed as follows: $`\omega _n\omega `$, $`\omega _\nu \nu `$, remaining on the imaginary axis.
## II Migdal’s method
In the method which can be called Migdal’s, the first correction to the unity vertex function $`\mathrm{\Gamma }^{(1)}`$ is considered. This is shown in Fig1. It has been estimated in many papers , but now it is becoming especially significant because the cases were found where as the parameter $`\lambda \mathrm{\Omega }_{ph}/W`$ may be not small.
The diagram in Fig.1 corresponds to the expression
$`\mathrm{\Gamma }^{(2)}(i\omega _n,i\omega _ni\omega _\nu ,\stackrel{}{q})`$ $`=`$ $`T{\displaystyle \underset{\omega _\nu ^{}}{}}{\displaystyle \underset{\stackrel{}{q}^{}}{}}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑\mathrm{\Omega }{\displaystyle \frac{\alpha ^2F_E(\mathrm{\Omega })2\mathrm{\Omega }}{N(0)(\omega _\nu ^{}^2+\mathrm{\Omega }^2)}}`$ (6)
$`\times G^{(0)}(i\omega _ni\omega _\nu ^{},\stackrel{}{p}\stackrel{}{q}^{})G^{(0)}(i\omega _ni\omega _\nu ^{}i\omega _\nu ,\stackrel{}{p}\stackrel{}{q}^{}\stackrel{}{q}).`$
Assuming that $`|\stackrel{}{p}\stackrel{}{q}^{}|p_F`$ and expanding $`\epsilon (\stackrel{}{p}\stackrel{}{q}^{}\stackrel{}{q})\epsilon (\stackrel{}{p}\stackrel{}{q}^{})qV_F\mathrm{cos}\theta `$ we get for $`T=0`$
$`\mathrm{\Gamma }^{(2)}(i\omega ,i\omega i\nu ,\stackrel{}{q})`$ $`=`$ $`\lambda {\displaystyle \frac{1}{4}}{\displaystyle _1^1}𝑑\eta {\displaystyle \frac{1}{im^{}+Q\eta }}`$ (9)
$`\times \{\mathrm{ln}{\displaystyle \frac{1+im+\frac{1}{m_0}}{1im+\frac{1}{m_0}}}+\mathrm{ln}{\displaystyle \frac{1im}{1+im}}+\mathrm{ln}{\displaystyle \frac{1+imim^{}}{1im+im^{}}}`$
$`+\mathrm{ln}{\displaystyle \frac{1im+im^{}Q\eta +\frac{1}{m_0}}{1+imim^{}+Q\eta +\frac{1}{m_0}}}\},`$
where $`m=\omega /\mathrm{\Omega }_{ph}`$, $`m^{}=\nu /\mathrm{\Omega }_{ph}`$, $`Q=qV_F/\mathrm{\Omega }_{ph}`$, $`m_0=\mathrm{\Omega }_{ph}/W`$.
This is a nonanalytical function at $`\omega 0`$, $`\stackrel{}{q}0`$, i.e. the result depends on the order of taking the limits. Let us define the dynamical and static vertex functions
$`\mathrm{\Gamma }_d`$ $`=`$ $`\mathrm{\Gamma }(\stackrel{}{q}=0,\nu 0,\omega );`$ (10)
$`\mathrm{\Gamma }_s`$ $`=`$ $`\mathrm{\Gamma }(\stackrel{}{q}0,\nu =0,\omega ).`$ (11)
From (9) one can get
$`\mathrm{\Gamma }_d^M=\mathrm{\Gamma }_M(\stackrel{}{q}=0,\nu 0,\omega )=\lambda [{\displaystyle \frac{1}{1+(\frac{\omega }{\mathrm{\Omega }_{ph}})^2}}{\displaystyle \frac{1+\frac{W}{\mathrm{\Omega }_{ph}}}{(1+\frac{W}{\mathrm{\Omega }_{ph}})^2+(\frac{\omega }{\mathrm{\Omega }_{ph}})^2}}],`$
$`\mathrm{\Gamma }_s^M=\mathrm{\Gamma }_M(\stackrel{}{q}0,\nu =0,\omega )=\lambda {\displaystyle \frac{1+\frac{W}{\mathrm{\Omega }_{ph}}}{(1+\frac{W}{\mathrm{\Omega }_{ph}})^2+(\frac{\omega }{\mathrm{\Omega }_{ph}})^2}}.`$
It is evident that the nonanalyticity mentioned above gives
$`\mathrm{\Gamma }_d^M\mathrm{\Gamma }_s^M=\lambda {\displaystyle \frac{1}{1+(\frac{\omega }{\mathrm{\Omega }_{ph}})^2}}.`$
At $`V_F|\stackrel{}{q}||\nu |`$
$`\mathrm{\Gamma }^{(2)}(i\omega ,i\omega i\nu ,\stackrel{}{q})`$ (12)
$`=`$ $`\lambda {\displaystyle \frac{\mathrm{\Omega }_{ph}}{\nu }}(\mathrm{arctan}{\displaystyle \frac{\omega }{\mathrm{\Omega }_{ph}}}\mathrm{arctan}{\displaystyle \frac{\omega \nu }{\mathrm{\Omega }_{ph}}}\mathrm{arctan}{\displaystyle \frac{\omega }{W+\mathrm{\Omega }_{ph}}}`$ (13)
$`+\mathrm{arctan}{\displaystyle \frac{\omega \nu }{W+\mathrm{\Omega }_{ph}}})`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Sigma }^{(1)}(i\omega )\mathrm{\Sigma }^{(1)}(i\omega i\nu )}{i\nu }},`$ (14)
which satisfies the Ward identity (2).
At $`V_F|\stackrel{}{q}||\nu |`$ one has $`\mathrm{\Gamma }^{(2)}(i\omega ,i\omega i\nu ,\stackrel{}{q})\lambda \mathrm{\Omega }_{ph}/W,`$ which complies with Refs. .
To calculate the self-energy we expand the last term of (9) in $`Q`$, leaving the first two terms as they are. Then
$`\mathrm{\Gamma }^{(2)}(i\omega ,i\omega i\nu ,\stackrel{}{q})`$ $`=`$ $`\lambda {\displaystyle \frac{\mathrm{\Omega }_{ph}}{qV_F}}\mathrm{arctan}{\displaystyle \frac{qV_F}{\nu }}`$ (17)
$`\times (\mathrm{arctan}{\displaystyle \frac{\omega }{\mathrm{\Omega }_{ph}}}\mathrm{arctan}{\displaystyle \frac{\omega \nu }{\mathrm{\Omega }_{ph}}}\mathrm{arctan}{\displaystyle \frac{\omega }{W+\mathrm{\Omega }_{ph}}}+\mathrm{arctan}{\displaystyle \frac{\omega \nu }{W+\mathrm{\Omega }_{ph}}})`$
$`\lambda {\displaystyle \frac{\frac{\mathrm{\Omega }_{ph}^2}{qV_FW}(1+\frac{\mathrm{\Omega }_{ph}}{W})(\frac{qV_F}{\mathrm{\Omega }_{ph}}\frac{\nu }{\mathrm{\Omega }_{ph}}\mathrm{arctan}\frac{qV_F}{\nu })}{[1+2\frac{\mathrm{\Omega }_{ph}}{W}+(\frac{\mathrm{\Omega }_{ph}}{W})^2(1+(\frac{\omega }{\mathrm{\Omega }_{ph}})^22\frac{\omega \nu }{\mathrm{\Omega }_{ph}^2}+(\frac{\nu }{\mathrm{\Omega }_{ph}})^2)]}}.`$
In Fig. 2a we show $`\mathrm{\Gamma }^{(2)}(\nu )`$ (numerical (9) and approximate analytical (17) results) for the case $`\omega =0,`$ as well as the result of the numerical integration of Eq.(4). One can see that there is a good agreement between these plots, so the expansion in small $`Q`$ is well justified.
The self-energy (4) consists then of three terms:
$`\mathrm{\Sigma }=\mathrm{\Sigma }^{(1)}+\mathrm{\Sigma }_v^{(2)}+\mathrm{\Sigma }_r^{(2)}.`$ (18)
The first order term is
$$\mathrm{\Sigma }^{(1)}(i\omega )=i\lambda \mathrm{\Omega }_{ph}\mathrm{arctan}\frac{\omega }{\mathrm{\Omega }_{ph}}+i\lambda \mathrm{\Omega }_{ph}\mathrm{arctan}\frac{\omega }{W+\mathrm{\Omega }_{ph}}.$$
(19)
The vertex correction $`\mathrm{\Sigma }_v^{(2)}`$ is obtained when one substitutes $`G=G^{(0)}`$ and $`\mathrm{\Gamma }=\mathrm{\Gamma }^{(2)}`$ (17) into (4)
$`\mathrm{\Sigma }_v^{(2)}(i\omega )=\mathrm{\Sigma }_v^{(2)a}(i\omega )+\mathrm{\Sigma }_v^{(2)b}(i\omega ),`$ (20)
$`\mathrm{\Sigma }_v^{(2)a}(i\omega )=i\lambda ^2{\displaystyle \frac{\mathrm{\Omega }_{ph}}{p_FV_F}}d_v^a(m,m_0,m_0^{}),`$
$`m=\omega /\mathrm{\Omega }_{ph}`$, $`m_0=\mathrm{\Omega }_{ph}/W`$, $`m_0^{}=\mathrm{\Omega }_{ph}/(p_FV_F)`$.
$`d_v^a(m,m_0,m_0^{})={\displaystyle \frac{1}{2\pi }}{\displaystyle \underset{\mathrm{}}{\overset{+\mathrm{}}{}}}{\displaystyle \frac{dy}{y^2+1}}\mathrm{arctan}{\displaystyle \frac{1}{m_0(my)}}`$ (21)
$`\times (\mathrm{arctan}m\mathrm{arctan}(my)\mathrm{arctan}{\displaystyle \frac{m}{\frac{1}{m_0}+1}}+\mathrm{arctan}{\displaystyle \frac{my}{\frac{1}{m_0}+1}})`$ (22)
$`\times (2\mathrm{arctan}{\displaystyle \frac{2}{ym_0^{}}}{\displaystyle \frac{1}{2}}ym_0^{}\mathrm{ln}({\displaystyle \frac{4}{y^2m{}_{}{}^{}{}_{0}{}^{}^2}}+1)).`$ (23)
$`\mathrm{\Sigma }_v^{(2)b}(i\omega )=i\lambda ^2{\displaystyle \frac{\mathrm{\Omega }_{ph}}{p_FV_F}}d_v^b(m,m_0,m_0^{}),`$
where
$`d_v^b(m,m_0,m_0^{})={\displaystyle \frac{1}{2\pi }}m_0^2(1+m_0){\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑m^{}{\displaystyle \frac{1}{m^2+1}}\mathrm{arctan}{\displaystyle \frac{1}{m_0(m^{}m)}}`$ (24)
$`\times \{{\displaystyle \frac{2}{m{}_{}{}^{}{}_{0}{}^{}^2}}{\displaystyle \frac{2m^{}}{m_0^{}}}\mathrm{arctan}{\displaystyle \frac{2}{m^{}m_0^{}}}+{\displaystyle \frac{m^{}^2}{2}}\mathrm{ln}[1+({\displaystyle \frac{2}{m^{}m_0^{}}})^2]\}.`$ (25)
For $`\omega /\mathrm{\Omega }_{ph}1`$, $`\mathrm{\Omega }_{ph}/W1`$ it gives
$`\mathrm{\Sigma }_v^{(2)}(i\omega )=\omega \lambda ^2{\displaystyle \frac{\mathrm{\Omega }_{ph}}{p_FV_F}}i[{\displaystyle \frac{\pi ^2}{8}}+({\displaystyle \frac{p_FV_F}{W}})^2].`$ (26)
The plot of $`\mathrm{\Sigma }_v^{(2)}(i\omega )`$ is shown in Fig.2b. The numerical calculations (the exact integration with Eq.(9)) give similar results for the small $`\omega `$’s . The existing discrepancy for larger frequencies is due to our approximation for $`\mathrm{\Gamma }`$.
Another term is the so-called rainbow diagram. As the inner part of the rainbow diagram corresponds to $`\mathrm{\Sigma }^{(1)}(i\omega i\nu )`$, the second order rainbow diagram equals
$`\mathrm{\Sigma }_r^{(2)}(i\omega )={\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑\mathrm{\Omega }\alpha ^2F_E(\mathrm{\Omega })T{\displaystyle \underset{\omega _\nu }{}}{\displaystyle \frac{2\mathrm{\Omega }}{\omega _\nu ^2+\mathrm{\Omega }^2}}{\displaystyle \frac{1}{N(0)}}`$ (27)
$`\times {\displaystyle \underset{\stackrel{}{q}}{}}[G(i\omega _ni\omega _\nu ,\stackrel{}{k}\stackrel{}{q})]^2\mathrm{\Sigma }^{(1)}(i\omega _ni\omega _\nu ),`$ (28)
which results in
$`\mathrm{\Sigma }_r^{(2)}(i\omega )=i\lambda ^2\mathrm{\Omega }_{ph}d_r(m,m_0),\text{ where}`$
$`d_r(m,m_0)`$ $`=`$ $`{\displaystyle \frac{1}{m_0(m^4+2m^2+1+\frac{1}{m_0^4}+2(\frac{m}{m_0})^2\frac{2}{m_0^2})}}`$ (31)
$`\times (m\mathrm{ln}{\displaystyle \frac{m^2+2^2}{m^2+(2+\frac{1}{m_0})^2}}+(\mathrm{arctan}{\displaystyle \frac{m}{2}}\mathrm{arctan}{\displaystyle \frac{m}{2+\frac{1}{m_0}}})(m^21+{\displaystyle \frac{1}{m_0^2}})`$
$`2m\mathrm{ln}{\displaystyle \frac{m_0+1}{m_0+2}}).`$
At $`\omega /\mathrm{\Omega }_{ph}1`$, $`\mathrm{\Omega }_{ph}/W1`$, it equals
$`\mathrm{\Sigma }_r^{(2)}(i\omega )=i{\displaystyle \frac{1}{2}}\lambda ^2\omega {\displaystyle \frac{\mathrm{\Omega }_{ph}}{W}}.`$ (32)
For $`\omega /\mathrm{\Omega }_{ph}1`$, $`\mathrm{\Omega }_{ph}/W1`$, $`\mathrm{\Omega }_{ph}/(p_FV_F)1`$ one can get
$`Z_M1{\displaystyle \frac{\mathrm{\Sigma }(i\omega )}{i\omega }}=1+\lambda {\displaystyle \frac{\pi ^2}{8}}\lambda ^2{\displaystyle \frac{\mathrm{\Omega }_{ph}}{p_FV_F}}{\displaystyle \frac{1}{2}}\lambda ^2{\displaystyle \frac{\mathrm{\Omega }_{ph}}{W}}\lambda ^2{\displaystyle \frac{\mathrm{\Omega }p_FV_F}{W^2}}.`$
One can see that the last expression has three terms with different denominators. When, for example, $`p_FV_F\mathrm{\Omega }_{ph}`$, the main contribution comes from the term $`\frac{\pi ^2}{8}\lambda ^2\mathrm{\Omega }_{ph}/p_FV_F`$ (even for an infinite bandwidth $`W`$). In Fig.2b there are shown for comparison analytical and numerical plots for the self energy correction $`\mathrm{\Sigma }_v^{(2)}`$. It is seen that the approximation used to calculate the correction is well justified.
If one assumes $`W=p_FV_F`$,
$$Z_M1+\lambda 2.7\lambda ^2\frac{\mathrm{\Omega }_{ph}}{W}.$$
(33)
One can see that these corrections lower the renormalization function. Had this tendency persist in all higher orders, this would have resulted in an instability (see below).
## III Cai, Lei, and Xie’s approximation
Following the paper, the authors of Ref. neglected the momentum dependence of the vertex function, considering $`qp_F\nu /W`$, and interpolated it between zero-frequency and infinite-frequency limits:
$`\text{}\mathrm{\Gamma }_{CLX}^{(2)}(i\omega _n,i\omega _n^{})=(1+2{\displaystyle 𝑑\omega \frac{\omega \alpha ^2F(\omega )}{\omega ^2+(\omega _n\omega _n^{})^2}\mathrm{\Lambda }_0\frac{2(\mathrm{\Lambda }_{\mathrm{}}/\mathrm{\Lambda }_0)}{\omega _m^2+\omega _n^2+2(\mathrm{\Lambda }_{\mathrm{}}/\mathrm{\Lambda }_0)}})^11,`$
where $`\mathrm{\Lambda }_0(\omega )=\frac{0.293\omega }{\omega +0.667W}`$, and $`\mathrm{\Lambda }_{\mathrm{}}(\omega )/\mathrm{\Lambda }_0(\omega )=\frac{4\sqrt{2}}{3}\frac{0.667W+\omega }{0.586}.`$
This results in an effective interaction
$`V_{eph}(\omega \omega ^{})={\displaystyle \frac{1}{N(0)}}{\displaystyle \frac{\lambda \mathrm{\Omega }^2}{(\omega \omega ^{})^2+\mathrm{\Omega }^2}}\times {\displaystyle \frac{1}{1+\frac{\lambda \mathrm{\Omega }^2}{(\omega \omega ^{})^2+\mathrm{\Omega }^2}\times \frac{0.293\mathrm{\Omega }}{\mathrm{\Omega }+0.667W}}}.`$
Combined with the contribution from the rainbow diagram, which is the same as in the Migdal method (32), this gives
$`Z_{CLX}1+\lambda 1.38\lambda ^2{\displaystyle \frac{\mathrm{\Omega }}{W}}.`$
The non-ladder approximation which was also discussed in was developed in detail by Takada, and this is considered below.
## IV Kostur and Mitrović’s approximation
There has been a series of papers, e.g. Ref. , which considered the effect of the vertex corrections on $`T_c`$ not only for an isotropic EPI, but also for spin (antiferromagnetic) fluctuations with a pronounced scattering at the wave vector $`\stackrel{}{Q}^{}=(\pi ,\pi )`$. In the paper the authors used the following equations for the self-energy $`\mathrm{\Sigma }(\stackrel{}{k},i\omega _n)=i\omega _n[1Z_n(\stackrel{}{k})]`$:
$`i\omega _nZ_n^{(2)}(\stackrel{}{k})`$ $`=`$ $`{\displaystyle \frac{T^2}{4W}}{\displaystyle \underset{n^{}n^{\prime \prime }}{}}\lambda (nn^{})\lambda (nn^{\prime \prime })`$ (36)
$`\times {\displaystyle _W^W}dϵ_𝐤^{}{\displaystyle _W^W}dϵ_{𝐤^{\prime \prime \prime }}{\displaystyle _W^W}dϵ_{𝐤^{\prime \prime }}M(𝐤,𝐤^{\prime \prime \prime };𝐤^{},𝐤^{\prime \prime }){\displaystyle \frac{i\omega _n^{}i\omega _{n^{\prime \prime }}i\omega _{n^{\prime \prime \prime }}}{Z_n^{}Z_{n^{\prime \prime }}Z_{n^{\prime \prime \prime }}}}`$
$`\times {\displaystyle \frac{1}{[i\omega _n^{}^2(ϵ_𝐤^{}/Z_n^{})^2][i\omega _{n^{\prime \prime }}^2(ϵ_{𝐤^{\prime \prime }}/Z_{n^{\prime \prime }})^2][i\omega _{n^{\prime \prime \prime }}^2(ϵ_{𝐤^{\prime \prime \prime }}/Z_{n^{\prime \prime \prime }})^2]}},`$
where $`\lambda (m)=_0^{\mathrm{}}𝑑\mathrm{\Omega }\alpha ^2F(\mathrm{\Omega })\frac{2\mathrm{\Omega }}{\nu _m^2+\mathrm{\Omega }^2}`$, and $`M(𝐤,𝐤^{\prime \prime \prime };𝐤^{},𝐤^{\prime \prime })`$ is a geometrical factor being a complicated function of momenta which is considered in detail in Ref.. For an isotropic EPI and a three-dimensional spherical Fermi-surface it is approximated by its value on the Fermi surface $`M(𝐤,𝐤^{\prime \prime \prime };𝐤^{},𝐤^{\prime \prime })=1`$. It gives
$`Z_n=1+{\displaystyle \frac{\pi T}{|\omega _n|}}{\displaystyle \underset{n^{}}{}}\lambda (nn^{})\mathrm{\Gamma }_{KM}(n,n^{})s_ns_n^{}a_n^{},\text{ where}`$
$`\lambda (m)=\lambda \frac{\mathrm{\Omega }_{ph}^2}{\nu _m^2+\mathrm{\Omega }_{ph}^2}`$, $`a_n=2/\pi \mathrm{arctan}(W/Z_n|\omega _n|)`$, $`s_n=\text{sign }\omega _n`$, and correction to the vertex function equals
$`\mathrm{\Gamma }_{KM}^{(2)}(n,n^{})={\displaystyle \frac{\pi ^2T}{2W}}{\displaystyle \underset{n^{\prime \prime }}{}}\lambda (nn^{\prime \prime })s_{n^{}+n^{\prime \prime }n}s_{n^{\prime \prime }}a_{n^{}+n^{\prime \prime }n}a_{n^{\prime \prime }}.`$ (37)
One can find that
$`\mathrm{\Gamma }_{KM}^{(2)}(n,n^{})=1\lambda {\displaystyle \frac{\pi }{4}}{\displaystyle \frac{\mathrm{\Omega }_{ph}}{W}}(\pi 2|\mathrm{arctan}{\displaystyle \frac{\omega _n^{}}{\mathrm{\Omega }_{ph}}}\mathrm{arctan}{\displaystyle \frac{\omega _n}{\mathrm{\Omega }_{ph}}}|)`$
and so does not depend from $`q`$. However, for $`qp_F`$ which give the main contribution to the self-energy it presents the correct order of the correction to the vertex function (compare with Eq.(17)) and the self-energy. From Eq. (2) one has for $`\omega /\mathrm{\Omega }_{ph}1`$, $`\mathrm{\Omega }_{ph}/W1`$
$`Z_{KM}=1+\lambda +\lambda ^2{\displaystyle \frac{\mathrm{\Omega }_{ph}}{W}}({\displaystyle \frac{\pi ^2}{8}}{\displaystyle \frac{1}{2}})1+\lambda 1.73\lambda ^2{\displaystyle \frac{\mathrm{\Omega }_{ph}}{W}}.`$
This is similar to Eq. (33), but has a different numerical coefficient, due to neglecting the $`\stackrel{}{q}`$-dependence in the vertex function (37).
## V C. Grimaldi, L. Pietronero, and S. Strässler’s approximation
In the paper the authors started with the same equation as in Migdal method, but made a series of approximate assumptions (for example, expanded in small $`q`$’s to take integrals) and gave the following estimate for the vertex function:
$`\mathrm{\Gamma }_{GPS}^{(2)}(\omega ,Q,\mathrm{\Omega },W)`$ $`=`$ $`{\displaystyle \frac{\lambda }{Q}}[[\mathrm{arctan}m\mathrm{arctan}{\displaystyle \frac{m}{\frac{1}{m_0}+1}}]\mathrm{arctan}{\displaystyle \frac{Q}{m}}`$ (39)
$`[Qm\mathrm{arctan}{\displaystyle \frac{Q}{m}}]{\displaystyle \frac{(\frac{1}{m_0}+1)[(\frac{1}{m_0}+1)^2+2m^2]}{[(\frac{1}{m_0}+1)^2+m^2]^2}}].`$
Here they set one of the external electronic frequencies to zero. In this case this result is similar to that in the Migdal’s approximation (17). However, to calculate $`\mathrm{\Sigma }`$ one needs dependences on both external frequencies.
There are some indications (see, e.g. ) that for a small hole doping strong Coulomb correlations renormalize the EPI, giving rise to the strong forward (small-$`q`$) scattering peak, while the backward scattering is strongly suppressed. With this idea in mind, in Ref. the electron-phonon coupling constant was assumed to have a cut-off in a momentum space $`|g_{\stackrel{}{p},\stackrel{}{k}}|^2=g^2(2k_F/q_c)^2\theta (q_c|\stackrel{}{p}\stackrel{}{k}|)`$. An approximation for $`q0`$ was used and then $`q_c=Q_c2p_F`$ was set to $`q_c=2p_F`$ giving the following result for $`Z`$:
$`Z_{CGPS}1+\lambda {\displaystyle \frac{\pi }{4}}\lambda ^2{\displaystyle \frac{\mathrm{\Omega }_{ph}}{W}}{\displaystyle \frac{1}{Q_c^2}}.`$
For $`Q_c=1`$ it becomes
$`Z_{CGPS}1+\lambda 0.8\lambda ^2{\displaystyle \frac{\mathrm{\Omega }_{ph}}{W}},`$
which gives the correct order of the correction to the self-energy , but underestimates it (compare with Eq.(33)). Probably this is due to the expansion in small $`q`$’s.
## VI Takada’s GISC method
According to Takada’s gauge-invariant self-consistent (GISC) method, the vertex function can be chosen as a functional of the self-energy $`\mathrm{\Gamma }_T=\mathrm{\Gamma }_T[\mathrm{\Sigma }_T]`$. This choice is based on the Ward identity (2) valid for all $`\omega _\nu `$’s and $`\stackrel{}{q}`$’s. In Ref. Takada proposes neglecting the vector term in the Ward identity and choosing
$`\mathrm{\Gamma }_T[\mathrm{\Sigma }]=\mathrm{\Gamma }_T(i\omega _n,i\omega _ni\omega _\nu )=1+({\displaystyle \frac{\mathrm{\Sigma }_T(i\omega _n)}{2i\omega _n}}{\displaystyle \frac{\mathrm{\Sigma }_T(i\omega _ni\omega _\nu )}{2(i\omega _ni\omega _\nu )}}),`$
which corresponds to the estimates of Refs. in this limit. If this method worked it would allow one to calculate also superconducting diagrams without complicated integrations over momenta.
The set of equations (1,2) is suggested to be solved by iterations. At the first step $`G=G^{(0)}`$, $`\mathrm{\Gamma }=\mathrm{\Gamma }_T^{(1)}=1`$ give $`\mathrm{\Sigma }_T=\mathrm{\Sigma }_T^{(1)}=\mathrm{\Sigma }^{(1)}`$ from Eq.(2). At the second step $`G=G[\mathrm{\Sigma }_T^{(1)}]`$ ; and the correction $`\mathrm{\Gamma }_T^{(2)}`$ for $`\mathrm{\Omega }_{ph}/W1`$ gives
$`\mathrm{\Sigma }_T^{(2)}(i\omega )=i\lambda ^2\mathrm{\Omega }_{ph}\{\begin{array}{c}\hfill (\frac{1}{4}+\frac{\pi ^2}{16})\frac{\omega }{\mathrm{\Omega }_{ph}},\text{ }|\omega |\mathrm{\Omega }_{ph},\\ \hfill \frac{\pi ^2}{4}\frac{\mathrm{\Omega }_{ph}}{\omega },\text{ }|\omega |\mathrm{\Omega }_{ph}\end{array}.`$ (42)
For small frequencies it results in
$`Z_T=1+\lambda +\lambda ^2({\displaystyle \frac{\pi ^2}{16}}+{\displaystyle \frac{1}{4}})1+\lambda +0.9\lambda ^2,`$
which has the incorrect sign and order of magnitude of the vertex correction (compare with (33)).
During this calculation, the correction to the vertex $`\mathrm{\Gamma }_T^{(2)}`$ is of the order of $`\lambda `$ for any $`q`$, which contradicts the Migdal’s theorem and can be valid only for $`q0`$ (see e.g. Ref. ). So in Takada’s method the correction $`\mathrm{\Sigma }_T^{(2)}`$ does not have the small parameter $`\mathrm{\Omega }_{ph}/W`$. This can be explained by the fact that the region $`V_F|\stackrel{}{q}||\nu |`$ where $`\mathrm{\Gamma }\lambda `$ gives in fact a small contribution to $`\mathrm{\Sigma }`$. If, on the other hand, $`V_F|\stackrel{}{q}||\nu |`$, the vertex $`\mathrm{\Gamma }`$ is of the order of $`\lambda \mathrm{\Omega }_{ph}/W`$, which gives $`\mathrm{\Sigma }^{(2)}\lambda \mathrm{\Omega }_{ph}^2/W`$ (see (26)).
To explain this result, one can find the vector vertex function $`\stackrel{}{\mathrm{\Gamma }}`$ in the lowest approximation and show that it really cannot be neglected. The diagram in Fig. 1 corresponds to $`\stackrel{}{\mathrm{\Gamma }}^{(2)}`$ if $`\mathrm{\Gamma }^{(1)}`$ is substituted by $`\stackrel{}{\mathrm{\Gamma }}^{(1)}`$ and $`\mathrm{\Gamma }^{(2)}(i\omega _n,i\omega _ni\omega _\nu ,\stackrel{}{q})`$ \- by $`\stackrel{}{\mathrm{\Gamma }}^{(2)}(i\omega ,i\omega i\nu ,\stackrel{}{q})`$. Using electron-hole symmetry, in our model $`\stackrel{}{\mathrm{\Gamma }}(i\omega ,i\omega i\nu ,\stackrel{}{k},\stackrel{}{k}\stackrel{}{q})=\stackrel{}{\mathrm{\Gamma }}^{(1)}+\stackrel{}{\mathrm{\Gamma }}^{(2)}\text{, where }\stackrel{}{\mathrm{\Gamma }}^{(1)}\frac{\stackrel{}{k}}{m},`$
$`\stackrel{}{q}\stackrel{}{\mathrm{\Gamma }}^{(2)}(i\omega ,i\omega i\nu ,\stackrel{}{q})`$ $`=`$ $`\lambda {\displaystyle \frac{1}{4}}{\displaystyle _1^1}𝑑\eta {\displaystyle \frac{\mathrm{\Omega }Q\eta }{im^{}+Q\eta }}`$ (45)
$`\times \{\mathrm{ln}{\displaystyle \frac{1+im+\frac{1}{m_0}}{1im+\frac{1}{m_0}}}+\mathrm{ln}{\displaystyle \frac{1im}{1+im}}+\mathrm{ln}{\displaystyle \frac{1+imim^{}}{1im+im^{}}}`$
$`+\mathrm{ln}{\displaystyle \frac{1im+im^{}Q\eta +\frac{1}{m_0}}{1+imim^{}+Q\eta +\frac{1}{m_0}}}\}.`$
Neglecting the $`Q`$-dependence in the last term of (45) one can get
$`\stackrel{}{q}\stackrel{}{\mathrm{\Gamma }}^{(2)}=i\lambda \mathrm{\Omega }_{ph}(1{\displaystyle \frac{\nu }{qV_F}}\mathrm{arctan}{\displaystyle \frac{qV_F}{\nu }})(\mathrm{arctan}{\displaystyle \frac{\omega }{\mathrm{\Omega }_{ph}}}\mathrm{arctan}{\displaystyle \frac{\omega \nu }{\mathrm{\Omega }_{ph}}}).`$ (46)
We can find regions where scalar or vector terms in the Ward identity dominate.
If one uses at $`\mathrm{\Omega }_{ph}/W1`$ the expression (17) for $`\mathrm{\Gamma }^{(2)}`$, one can find that the terms $`\nu \mathrm{\Gamma }^{(2)}`$ and $`\stackrel{}{q}\stackrel{}{\mathrm{\Gamma }}^{(2)}`$ from (2) are of the same order if $`\nu /qV_F=0.43`$. If, however, $`\nu /qV_F<0.43`$, the vector term dominates and cannot be neglected.
## VII F. Cosenza, L. De Cesare, and M. Fusco Girard’s improvement of the GISC method
Taking into account a criticism of the GISC method , F. Cosenza, L. De Cesare, and M. Fusco Girard tried to improve the GISC method and did not neglect the vector term in the Ward identity. They employed the simplest choice of the solution to the Ward identity
$$\mathrm{\Gamma }_{CDG}^{(2)}(\stackrel{}{k},i\omega _n,\stackrel{}{k^{}},i\omega _n^{})\frac{(i\omega _n^{}i\omega _n)[\mathrm{\Sigma }(\stackrel{}{k},i\omega _n)\mathrm{\Sigma }(\stackrel{}{k^{}},i\omega _n^{})]}{(i\omega _n^{}i\omega _n)^2|\alpha (\stackrel{}{k},\stackrel{}{k^{}})|^2|\stackrel{}{k}\stackrel{}{k^{}}|^2},$$
(47)
where $`\alpha (\stackrel{}{k},\stackrel{}{k^{}})=(\stackrel{}{q}/2+\stackrel{}{k})/m\stackrel{}{V_F}`$. With $`\mathrm{\Sigma }(\stackrel{}{k},i\omega _n)`$ taken from Eq.(19) one gets for $`|\stackrel{}{q}|=|\stackrel{}{k}\stackrel{}{k^{}}|p_F`$
$`\mathrm{\Gamma }_{CDG}^{(2)}(\stackrel{}{k},i\omega _n,\stackrel{}{k^{}},i\omega _n^{})\lambda ({\displaystyle \frac{\mathrm{\Omega }_{ph}}{W}})^2`$
which is different from the Migdal’s estimation $`\mathrm{\Gamma }^{(2)}\lambda \frac{\mathrm{\Omega }_{ph}}{W}`$ (see (17)).
Substituting approximation (47) into the equation for the self-energy (4) one can get
$`\mathrm{\Sigma }_{CDG}^{(2)}(i\omega 0)0.25i\lambda ^2\omega ({\displaystyle \frac{\mathrm{\Omega }_{ph}}{W}})^2(\mathrm{ln}{\displaystyle \frac{W}{\mathrm{\Omega }_{ph}}})^2`$
which is different from (26) and gives
$`Z_{CDG}1+\lambda 0.25\lambda ^2({\displaystyle \frac{\mathrm{\Omega }_{ph}}{W}})^2(\mathrm{ln}{\displaystyle \frac{W}{\mathrm{\Omega }_{ph}}})^2.`$
This can be explained by the fact that the Ward identity is a one equation for two functions (the scalar and vector vertex functions) and thus allows for multiple solutions . So there is no particular preference in using (47) instead of any other solution.
## VIII Discussion
So in the paper we considered consistently the contribution of nonadiabatic effects to the electron self-energy in the normal state at $`T=0`$ in second orders of the EPI constant $`\lambda `$. Several methods of taking into account nonadiabatic corrections were compared which can be summarized in the Table ($`p_FV_F=W`$).
The results following from the approximations of V. N. Kostur and B. Mitrović , E. Cappelluti, C. Grimaldi, L. Pietronero, and S. Strässler give correct estimations of the order of magnitude and sign of the nonadiabatic corrections to the self-energy. The analytical results were compared with numerical calculations.
It is shown that the GISC method and its generalization give for the self-energy overestimated and underestimated results respectively. This is connected with the fact that the leading contribution to $`\mathrm{\Sigma }`$ comes from the region $`qv_F\nu `$ where $`\mathrm{\Gamma }\lambda \mathrm{\Omega }_{ph}/W`$. This means that in the Ward identity one may not neglect the vector term as proposed by Takada. In the framework of the standard (Migdal’s) approach, in the lowest order, there is a calculable parameter $`\nu /qv_F`$ ( the ratio of the frequency of an incoming phonon to its momentum) for which the contributions of the vector and scalar terms become of the same order. They are equal when $`\nu /qv_F=0.43`$. So for the most metals, where the main contribution to the self energy comes from the region $`\nu /qv_F1`$, the GISC method may not be applied. One should also be careful when choosing a solution of the Ward identity, so as not to underestimate the scalar term, like in Ref. . It seems that GISC method can work in systems with long range interaction, for example, in doped semiconductors.
The difference in the renormalization factors $`Z`$ in the Table are due with the different approximations used in the calculations of the momentum and frequency dependence of vertex function. We can illustrate this in Fig. 3 where we represent the results for $`\mathrm{\Gamma }^{(2)}`$ in the Takada’s ($`q`$-independent) approximation, the Migdal approach for different $`q`$ (Eqs. (4,27) of the present paper), and the $`q`$-independent Kostur, Mitrović vertex function. We see that if the former overestimates the vertex corrections, the latter underestimates them.
For a stable system the condition $`Im\mathrm{\Sigma }(i\omega _n)<0`$ for $`\omega _n>0`$ must be satisfied (it is equivalent to $`Z(i\omega _n)>1`$). However, as one can see from the Table, the nonadiabatic corrections reduce the renormalization function $`Z`$, favoring the tendency towards instability. The critical $`\lambda _{crit}`$ at which it takes place (see Eqs. (10-16)) is given by
$`\lambda _{crit}={\displaystyle \frac{\mathrm{arctan}m\mathrm{arctan}\frac{m}{1/m_0+1}}{m_0d_v(m,m_0)+d_r(m,m_0))}},`$
for $`m0`$, where $`d_v(m,m_0)`$ and $`d_r(m,m_0)`$ were defined above. The dependence of $`\lambda `$ on the parameter $`\mathrm{\Omega }_{ph}/W`$ is shown in Fig.4. Thus, for a stable system there are constraints on $`\lambda `$ for a fixed Migdal parameter. This can explain, for example, the existence of PbBi with $`\lambda 3`$. At the same time, in some papers, e.g. , it is stated that $`\lambda 1`$ independently from the Migdal parameter. As one sees from Fig. 4, it may be the case only for $`\mathrm{\Omega }_{ph}W`$.
This possible violation of the analytical properties of the one-particle Green’s function could result in corresponding nonanaliticity of the two-particle Green’s function, and thus in a charge response function. This can lead to a charge instability.
It follows from our analysis that different approximations to the vertex function give in the normal state quantitatively (and sometimes even qualitatively!) different estimates for the self-energy. This makes the conclusions of these theories doubtful, as the same approaches were used there to calculate $`T_c`$ in the superconducting state. This means that only direct Migdal-type calculations can give an answer to the question: Do nonadiabatic effects enhance critical temperature of the superconducting transition?
## IX Acknowledgment
We thank T. Dahm, O. Gunnarsson, M.L. Kulić, V.V. Losyakov, I.I. Mazin, E.G. Maksimov, and N. Schopohl for helpful discussions.
Table
| Method: | $`Z^{(2)}`$ |
| --- | --- |
| Migdal (numerical), with $`\mathrm{\Gamma }^{(2)}`$ from Eq.(4) | $`2.7\lambda ^2\mathrm{\Omega }_{ph}/W`$ |
| Migdal (analytical), with $`\mathrm{\Gamma }^{(2)}`$ from Eq.(7) | $`2.7\lambda ^2\mathrm{\Omega }_{ph}/W`$ |
| CLX, Ref. | $`1.38\lambda ^2\mathrm{\Omega }_{ph}/W`$ |
| KM, Ref. | $`1.73\lambda ^2\mathrm{\Omega }_{ph}/W`$ |
| CGPS, Ref. | $`0.8\lambda ^2\mathrm{\Omega }_{ph}/W`$ |
| Takada’s GISC, Ref. | +$`0.9\lambda ^2`$ |
| CDG, Ref. | $`0.25\lambda ^2(\mathrm{\Omega }_{ph}/W)^2(\mathrm{ln}W/\mathrm{\Omega }_{ph})^2`$ |
Figure Captions
Fig.1. Equation for the vertex function $`\mathrm{\Gamma }^{(2)}`$ in Migdal’s method. The outgoing lines are shown for clarity and are not included in the definition of $`\mathrm{\Gamma }`$.
Fig. 2. a) Plot of $`\mathrm{\Gamma }^{(2)}(\omega =0,i\nu )`$ (Eq.(4), dashed line) and $`\mathrm{\Gamma }_{appr}^{(2)}(\omega =0,i\nu )`$ (Eq.(7), solid line) for $`\lambda =1,\mathrm{\Omega }_{ph}/W=0.1,\omega =0`$, $`q/p_F=0;0,2;0,4;0,6;0,8`$. b) Corresponding frequency dependence of $`\mathrm{\Sigma }_v^{(2)}(i\omega )`$ (dashed line) and $`\mathrm{\Sigma }_{v,appr}{}_{}{}^{(2)}(i\omega )`$ (solid line).
Fig.3. $`\mathrm{\Gamma }^{(2)}(\omega =0,\nu )`$in the Takada’s approximation (long dashed line); the Migdal approach for $`q/p_F=0,0.5,1`$ (solid lines, from top to bottom) ; and the Kostur, Mitrović vertex function ( short dashed line ).
Fig. 4. $`\lambda _{crit}`$ at which $`\mathrm{\Sigma }(i\omega )=0`$ vs the parameter $`\mathrm{\Omega }_{ph}/W.`$
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# Super-Calogero-Moser-Sutherland systems and free super-oscillators : a mapping
## I Introduction
The rational CMS Hamiltonian is described by $`N`$ particles interacting with each other through an inverse square interaction and all particles are subjected to a common confining harmonic force. This model is exactly solvable and the eigen-values, including the degeneracy at each level, are exactly identical to the spectrum of $`N`$ free oscillators, except for a constant shift in the ground state energy. There were enough indications in the literature in different context that a very close connection between the CMS and the free oscillators model might exists. In fact, it has been shown recently that the rational CMS Hamiltonian is equivalent to that of free oscillators through a similarity transformation, confirming all previous speculations. This equivalence has enriched our understanding of the model and also became a very useful tool for studying different aspects of CMS, like eigen functions, integrability, and symmetry algebra, in a new way.
The supersymmetric version of the rational CMS system has also been studied in the literature in different context. The zero fermion sector of the supersymmetric CMS describes the usual CMS model, while the $`N`$ fermion sector describes the CMS model at a shifted value of the coupling constant . Such relation between the zero and the $`N`$ fermion sector of the model is due to ‘shape invariance’ of the Hamiltonian, which is a very popular and useful concept in studying quantum mechanics with one degree of freedom . For other sectors with fermion numbers ranging from one to $`N1`$, however, no such trivial identifications with the usual CMS can be made. Hamiltonian in these sectors are in fact related to CMS with internal degrees of freedom.
The supersymmetric rational CMS model (SRCMSM) of $`A_{N+1}`$-type is exactly solvable in both supersymmetry-preserving and supersymmetry-breaking phases. The spectrum in the supersymmetry-preserving phase is again identical to that of the free super-oscillators. It might be recalled at this point that the supersymmetry is always preserved in the super-oscillator model, once the convention for choosing the ground state in either zero or $`N`$ fermion sector has been made. Thus, the spectrum of SRCMSM in the supersymmetry-breaking phase, has no counter-part in the super-oscillator model. However, it has some similarity with the spectrum of the super-oscillator model modulo a constant shift in the ground state energy. It is intriguing at this point to ask, whether or not the SRCMSM, at least in the supersymmetry-preserving phase, can be shown to be equivalent to free super-oscillators through a similarity transformation, much akin to its non-supersymmetric version.
The purpose of this paper is to show that the SRCMSM of $`A_{N+1}`$-type is indeed equivalent to free super-oscillators through a similarity transformation. This equivalence is valid only in supersymmetry-preserving phase. This explains the identicalness of the spectrum of SRCMSM and free super-oscillators. The eigen functions of these two models are of-course different from each other and we outline a method to construct eigen functions of the SRCMSM from permutationally invariant super-oscillator basis functions. In case, one chooses a basis function which is not symmetric under the combined exchange of bosonic and fermionic coordinates, the corresponding eigenfunction of the SRCMSM is not normalizable. This is due to highly correlated nature of the many-body inverse-square interaction and this has also been observed in the usual CMS model.
We also prescribe on constructing eigen-spectrum of the SRCMSM in the supersymmetry-breaking phase, from the known super-oscillator basis by making use of a duality property of the model. In particular, we construct a new super-Hamiltonian, which differs from the SRCMSM by the fermionic number operator and a constant. This implies that any eigen-function of this dual model is also a valid eigen-function of the SRCMSM. Of course, the corresponding energy eigen-values are different from each other. We show through a similarity transformation that this dual Hamiltonian is again equivalent to a free super-oscillator Hamiltonian. It turns out that the eigen-spectrum of the SRCMSM obtained from this super-oscillator model via the dual Hamiltonian indeed correctly describes the supersymmetry-breaking phase of the model.
The symmetry algebra of the super-oscillator model is well understood in terms of a set of bosonic and fermionic operators. We define a set of such operators for the SRCMSM, which are obtained from the corresponding operators in the super-oscillator model through the inverse similarity transformation. This enables us to study the symmetry algebra of SRCMSM in a simple way, leading to the construction of the complete eigen-spectrum algebraically.
As a generalization of these results, we show that a wide class of models whose bosonic many-body potential is a homogeneous function of degree $`2`$ and all the particles are restricted to move on a line by a common confining harmonic force, are equivalent to the super-oscillator model through a similarity transformation. These Hamiltonians are characterized by a dynamical $`OSp(2|2)`$ supersymmetry. The SRCMSM associated with different root-structures of the Lie-algebra appear as a special small subset of this class. Though the equivalence is valid at the operator level for the general inverse-square potential, one must show that the complete set of eigen-functions as well as eigen-values of such models are indeed obtained from the super-oscillator model. The equivalence relation at the operator level acts as a necessary condition, while the construction of the complete eigen-spectrum and associated wave-functions from the super-oscillator basis is sufficient to claim such relation between these two models. We show that both the necessary and the sufficient conditions are certainly satisfied by the SRCMSM of $`A_{N+1}`$ and $`BC_{N+1}`$ types. However, it appears that all other cases have to be treated individually.
We organize the paper in the following way. We first give an overview of the supersymmetric quantum mechanics with many degrees of freedom in the next section. We mostly review the known results in a way which will become useful for our subsequent discussions. In Sec. III, we consider the $`A_{N+1}`$-type SRCMSM and show its equivalence to free super oscillator model. We first show the equivalence for the supersymmetry-preserving phase in Sec. III.A and outline a method to construct the eigen-spectrum from the known super-oscillator basis. Similar study for the supersymmetry-breaking phase has been discussed in Sec. III.B, using a duality property of the model. We generalize these results to SRCMSM associated with other root-structures of the Lie algebra in Sec. IV. Finally, in Sec. V, we summarize and discuss the implications of these results. We show how the dynamical $`OSp(2|2)`$ supersymmetry is realized by these systems in Appendix A.
## II Supersymmetric Quantum Mechanics with many degrees of freedom: brief review
The supercharge $`Q`$ and its conjugate $`Q^{}`$ are defined as,
$$Q=\underset{i=1}{\overset{N}{}}\psi _i^{}a_i,Q^{}=\underset{i=1}{\overset{N}{}}\psi _ia_i^{},$$
(1)
where the fermionic variables $`\psi _i`$’s satisfy the Clifford algebra,
$$\{\psi _i,\psi _j\}=0=\{\psi _i^{},\psi _j^{}\},\{\psi _i,\psi _j^{}\}=\delta _{ij},i,j=1,2,\mathrm{},N.$$
(2)
The operators $`a_i(a_i^{})`$’s are analogous to bosonic annihilation ( creation ) operators. They are defined in terms of the momentum operators $`p_i=i\frac{}{x_i}`$ and the superpotential $`W(x_1,x_2,\mathrm{},x_N)`$ as,
$$a_i=p_iiW_i,a_i^{}=p_i+iW_i,W_i=\frac{W}{x_i},$$
(3)
and satisfy the following commutation relations among themselves,
$$[a_i,a_j]=0=[a_i^{},a_j^{}],[a_i,a_j^{}]=[a_j,a_i^{}]=2W_{ij},W_{ij}=\frac{^2W}{x_ix_j}.$$
(4)
Note that, by construction, $`W_i`$’s satisfy the so called ‘zero-curvature condition’ $`_iW_j=_jW_i`$. Also, for translationally invariant superpotential, these $`W_i`$’s satisfy the ‘sum to zero’ condition, $`_iW_i=0`$. These two properties are useful ingredients in studying the usual CMS model.
The supersymmetric Hamiltonian is defined in terms of the supercharges as,
$`H`$ $`=`$ $`{\displaystyle \frac{1}{2}}\{Q,Q^{}\}`$ (5)
$`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle \underset{i}{}}\{a_i,a_i^{}\}+{\displaystyle \frac{1}{4}}{\displaystyle \underset{i,j}{}}[a_i,a_j^{}][\psi _i^{},\psi _j].`$ (6)
The Hamiltonian commutes with both $`Q`$ and $`Q^{}`$. The ground state of $`H`$ is annihilated by both $`Q`$ and $`Q^{}`$. Thus, the ground states are given by,
$$\varphi _0=e^W|0>,\varphi _N=e^W|\overline{0}>,$$
(7)
where the fermionic vacuum $`|0>`$ and its conjugate $`|\overline{0}>`$ in the $`2^N`$ dimensional fermionic Fock space are defined as,
$$\psi _i|0>=0,\psi _i^{}|\overline{0}>=0.$$
(8)
The first equation of (8) defines the zero-fermion sector, while the second one defines the $`N`$ fermion sector. In case, either $`\varphi _0`$ or $`\varphi _N`$ is normalizable, the supersymmetry is preserved with zero ground state energy. On the other hand, the supersymmetry is broken if neither $`\varphi _0`$ nor $`\varphi _N`$ is normalizable. The ground state energy in this case is positive-definite.
## III Equivalence : Rational CMS of $`A_{N+1}`$-type and free oscillators
The superpotential for the $`A_{N+1}`$-type SRCMSM is given by,
$$W=\lambda ln\underset{i<j}{}x_{ij}+\frac{1}{2}\underset{i}{}x_i^2,x_{ij}=x_ix_j.$$
(9)
The first term produces the many-body inverse square interaction, while the second term generates the term responsible for harmonic confinement. The Hamiltonian (6), with the above choice of $`W`$, has the following form,
$`H`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i}{}}{\displaystyle \frac{^2}{x_i^2}}+{\displaystyle \frac{1}{2}}\lambda (\lambda 1){\displaystyle \underset{ij}{}}x_{ij}^2+{\displaystyle \frac{1}{2}}{\displaystyle \underset{i}{}}x_i^2{\displaystyle \frac{1}{2}}N\left(1+\lambda (N1)\right)`$ (10)
$`+`$ $`{\displaystyle \underset{i}{}}\psi _i^{}\psi _i+\lambda {\displaystyle \underset{ij}{}}x_{ij}^2\left(\psi _i^{}\psi _i\psi _i^{}\psi _j\right).`$ (11)
The Hamiltonian $`H`$ is permutationally invariant under the combined exchange of bosonic and fermionic coordinates. Observe that the zero-fermion sector of (11) describes the usual CMS, apart from a constant equal to its ground state energy. The ground state of SRCMSM has the well-known form,
$`\mathrm{\Phi }`$ $`=`$ $`e^W|0>`$ (12)
$`=`$ $`{\displaystyle \underset{i<j}{}}x_{ij}^\lambda e^{\frac{1}{2}_ix_i^2}|0>.`$ (13)
Note that $`\mathrm{\Phi }`$ is normalizable for $`\lambda >\frac{1}{2}`$. However, a stronger criteria that each momentum operator $`p_i`$ is self-adjoint for the wave-functions of the form $`\mathrm{\Phi }`$ requires $`\lambda >0`$. The supersymmetry is preserved for $`\lambda >0`$, while it is broken for $`\lambda <0`$.
### A Supersymmetry-preserving phase
Now we would like to show that the Hamiltonian (11) is equivalent to the free super-oscillators model through a similarity transformation. In order to do so, let us first consider the following transformation,
$`H_1`$ $`=`$ $`e^WHe^W`$ (14)
$`=`$ $`{\displaystyle \underset{i}{}}\left(x_i{\displaystyle \frac{}{x_i}}+\psi _i^{}\psi _i\right)S,`$ (15)
$`S`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i}{}}{\displaystyle \frac{^2}{x_i^2}}+\lambda {\displaystyle \underset{ij}{}}x_{ij}^1{\displaystyle \frac{}{x_i}}\lambda {\displaystyle \underset{ij}{}}x_{ij}^2\left(\psi _i^{}\psi _i\psi _i^{}\psi _j\right).`$ (16)
The total fermion number operator $`N_f=_i\psi _i^{}\psi _i`$ commutes with the Hamiltonian $`H`$. The fermionic part of $`H_1`$ is identical to that of $`H`$. Thus, $`N_f`$ commutes with $`H_1`$ and hence, also with $`S`$. Making use of the following identities,
$$[\underset{i}{}x_i\frac{}{x_i},S]=2S,[\underset{i}{}\left(x_i\frac{}{x_i}+\psi _i^{}\psi _i\right),S]=2S,$$
(17)
we find,
$`[H_1,e^{\frac{S}{2}}]=Se^{\frac{S}{2}}`$ (18)
$`H_2=e^{\frac{S}{2}}H_1e^{\frac{S}{2}}`$ (19)
$`={\displaystyle \underset{i}{}}\left(x_i{\displaystyle \frac{}{x_i}}+\psi _i^{}\psi _i\right).`$ (20)
The transformed Hamiltonian $`H_2`$ is nothing but the supersymmetric generalization of the Euler operator. The connection of $`H`$ with the free super-oscillators is apparent from the expression of $`H_2`$. In particular, we get the familiar supersymmetric $`N`$ particle free oscillators model in the following way,
$`H_{sho}`$ $`=`$ $`e^{\frac{1}{2}_ix_i^2}e^{\frac{1}{4}_i\frac{^2}{x_i^2}}H_2e^{\frac{1}{4}_i\frac{^2}{x_i^2}}e^{\frac{1}{2}_ix_i^2}`$ (21)
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i}{}}\left({\displaystyle \frac{^2}{x_i^2}}+x_i^2\right)+{\displaystyle \underset{i}{}}\psi _i^{}\psi _i{\displaystyle \frac{N}{2}}.`$ (22)
This shows the equivalence between SRCMSM and the free super-oscillators.
#### 1 Construction of eigen-functions
The eigen-spectrum of (11) can be constructed either from (20) or (22). We prefer to work with Eq. (20). If $`P_{n,k}`$ is an eigen-function of (20) with the eigen-value $`E_{n,k}`$, then, $`H`$ has the same eigen-value $`E_{n,k}`$ with the eigen-function given by,
$$\chi =e^We^{\frac{S}{2}}P_{n,k}|0>.$$
(23)
We have to choose $`P_{n,k}`$ to be a permutationally symmetric polynomial of $`x_i`$ and $`\psi _i`$, under the combined exchange of the bosonic and fermionic coordinates. Otherwise, the action of $`S`$ on $`P_{n,k}`$ produces non-vanishing singular terms, thereby, making $`\chi `$ non-normalizable. It is worth recalling at this point that similar constraint on $`P_{n,k}`$ has been noticed also for the usual CMS case, reflecting the highly correlated nature of these systems. The highly correlated nature of this model is also present in the supersymmetric version.
There are many choices for the polynomial $`P_{n,k}`$. Let us choose the following form of $`P_{n,k}`$,
$$P_{n,k}=r^{2n}\underset{i}{}x_i^{k1}\psi _i^{},r^2=\underset{i}{}x_i^2,$$
(24)
as the $`N_f=1`$ solution of $`H_2`$ with $`E_{n,k}=2n+k`$. The quantum numbers $`n`$ and $`k1`$ are nonnegative integers. It can be checked easily that the action of $`S^m`$ on $`P_{n,k}`$ does not produce any singularity for positive $`m`$. Let us first consider the action of $`S`$ on $`P_{n,k}`$,
$`SP_{n,k}`$ $`=`$ $`b_1r^{2(n1)}{\displaystyle \underset{i}{}}x_i^{k1}\psi _i^{}+b_2r^{2n}{\displaystyle \underset{i}{}}x_i^{k3}\psi _i^{}+\lambda r^{2n}{\displaystyle \underset{ij}{}}{\displaystyle \underset{l=0}{\overset{k3}{}}}(kl2)x_i^{kl3}x_j^l\psi _i^{},`$ (25)
$`b_1`$ $`=`$ $`n\left[N+2\lambda N(N1)+2(n+k2)\right],b_2={\displaystyle \frac{1}{2}}(k1)(k2).`$ (26)
The first two terms on the right hand side of the first equation in (26) has the same form as that of $`P_{n,k}`$, except for powers of $`r`$ and $`x_i`$. Thus, the contribution of these two terms to $`S^2P_{n,k}`$ can not contain a singular term. The third term has a different form than $`P_{n,k}`$. A term like this, which has a general form,
$$\eta =\underset{i_1i_2\mathrm{},i_N}{}x_{i_1}^{k_1}x_{i_2}^{k_2}\mathrm{}x_{i_N}^{k_N}\psi _{i_1}^{},$$
(27)
keeps on appearing on each successive operation of $`S`$ on the left hand side of the first equation of (26). The integers $`k_i`$’s are determined in terms of $`k`$. However, we keep them as arbitrary nonnegative integers in (27). The first term of $`S`$, a generalized Laplacian operator $`=_i\frac{^2}{x_i^2}`$ , acting on $`\eta `$ can not produce any singularity. Further, we have the following identity,
$`S^{}\eta `$ $`=`$ $`\left(S{\displaystyle \frac{1}{2}}{\displaystyle \underset{i}{}}{\displaystyle \frac{^2}{x_i^2}}\right)\eta `$ (28)
$`=`$ $`\lambda {\displaystyle \underset{i_1i_2\mathrm{}i_N}{}}{\displaystyle \underset{j(i_p)}{}}{\displaystyle \underset{p=1}{\overset{N}{}}}{\displaystyle \underset{l=0}{\overset{k_p2}{}}}\beta _{k_p,l}x_{i_1}^{k_1}x_{i_2}^{k_2}\mathrm{}x_{i_p}^{k_pl2}x_j^l\mathrm{}x_{i_N}^{k_N}\psi _{i_1}^{},`$ (30)
$`\beta _{k_1,l}=k_1l1,\beta _{k_p,l}={\displaystyle \frac{k_p}{2}}forp2.`$
Note that $`S^{}\eta `$ has the same form as that of $`\eta `$, once the summation over the indices $`j,p`$ and $`l`$ has been performed. This proves that $`S^mP_{n,k}`$ can not contain a singular term, instead terminates as a finite degree polynomial. Thus, the well-behaved eigen-functions $`\chi `$ of $`H`$ can be constructed from the super-oscillator basis $`P_{n,k}`$. It may be worth mentioning here that the exact solution for $`N_f=1`$ and certain small values of $`k`$, obtained in , can be reproduced in a systematic way from Eqs. (23) and (24).
Similar results for other values of $`N_f`$ can also be obtained. For example, one may choose $`P_{n,k}`$ for an arbitrary $`N_f`$ as,
$$P_{n,k}=\frac{1}{N_f!}r^{2n}\underset{i_1,i_2,\mathrm{},i_{N_f}}{}f_{i_1i_2\mathrm{}i_{N_f}}(x_1,x_2,\mathrm{},x_N)\psi _{i_1}^{}\psi _{i_2}^{}\mathrm{}\psi _{i_{N_f}}^{},$$
(31)
where $`f_{i_1i_2\mathrm{}i_{N_f}}`$ is anti-symmetric under the exchange of any two indices and is a homogeneous function of degree $`kN_f`$. The anti-symmetric nature of $`f`$ ensures that $`P_{n,k}`$ is permutationally invariant under the combined exchange of bosonic and fermionic coordinates. Though we do not present here results concerning normalizability of eigen-functions $`\chi `$ constructed from (31) for arbitrary $`N_f`$, it is expected that the certain specific choices of $`f`$ would indeed produce well-behaved and physically accepted $`\chi `$. This is because of the result that the eigen-value equation of $`H_1`$ has permutationally symmetric polynomials in $`x_i`$ and $`\psi _i`$ as the solution. This implies that the solution for the eigen equation of $`S`$ are also permutationally symmetric polynomials. Thus, the action of $`S^m`$ on these permutationally symmetric polynomials for any positive $`m`$ are not expected to produce singular terms. We outline a method in the next section to construct the eigenstates in an algebraic way.
#### 2 Algebraic structure
The algebraic structure of the super-oscillators can be exploited to construct the eigenstates of $`H`$ in an algebraic way. Consider the following set of operators,
$`b_i^{}=ip_i={\displaystyle \frac{}{x_i}},b_i^+=2x_i`$ (32)
$`B_n^{}={\displaystyle \underset{i=1}{\overset{N}{}}}T^1b_i^^nT,B_n^+={\displaystyle \underset{i=1}{\overset{N}{}}}T^1b_i^{+^n}T,T=e^{\frac{S}{2}}e^W`$ (33)
$`F_n^{}=T^1\left({\displaystyle \underset{i}{}}\psi _ib_i^{^{n1}}\right)T,F_n^+=T^1\left({\displaystyle \underset{i}{}}\psi _i^{}b_i^{+^{n1}}\right)T,`$ (34)
$`q_n^{}=T^1\left({\displaystyle \underset{i}{}}\psi _i^{}b_i^^n\right)T,q_n^+=T^1\left({\displaystyle \underset{i}{}}\psi _ib_i^{+^n}\right)T.`$ (35)
Note that we are using a particular form of $`b_i^{}`$ and $`b_i^+`$, such that $`[b_i^{},b_j^+]=2\delta _{ij}`$. This choice has been made to make one to one correspondence between the usual annihilation (creation) operator of the harmonic oscillator and the $`b_i^{}(b_i^+)`$. In particular, it can be checked easily,
$$ib_i^{}=t^1a_h^{}t,ib_i^+=t^1a_h^+t,a_h^\pm =p_i\pm ix_i,t=e^{\frac{1}{2}_ix_i^2}e^{\frac{1}{4}_i\frac{^2}{x_i^2}}.$$
(36)
The operators in (35) satisfy the following algebra among themselves.
$`\{F_m^+,F_n^+\}=0,[B_m^+,F_n^+]=0,[B_m^+,B_n^+]=0,`$ (38)
$`\{q_1^{},F_n^+\}=0,\{q_1^+,F_n^+\}=B_n^+,[H,F_n^+]=nF_n^+,`$
$`[q_1^{},B_n^+]=2nF_n^+,[q_1^+,B_n^+]=0,[H,B_n^+]=nB_n^+.`$ (39)
This is also the algebra of the corresponding operators of super-oscillators. Thus, the eigen-functions can be created in a similar way by acting different powers of $`B_n^+`$ and $`F_n^+`$ on the ground state. In particular,
$$\chi _{n_1\mathrm{}n_N\nu _1\mathrm{}\nu _N}=\underset{k=1}{\overset{N}{}}B_k^{+^{n_k}}F_k^{+^{\nu _k}}\mathrm{\Phi },$$
(40)
is the eigenfunction with the eigen-value $`E=_{k=1}^Nk(n_k+\nu _k)`$. The bosonic quantum numbers $`n_k`$’s are nonnegative integers, while the fermionic quantum numbers $`\nu _k`$’s are either $`0`$ or $`1`$. Note that a set of $`N`$ independent super-oscillators with the frequencies $`1,2,\mathrm{},N`$ have the same energy $`E`$. Thus, the spectrum of SRCMSM is identical to that of $`N`$ independent super-oscillators with the frequencies $`1,2,\mathrm{},N`$.
A particular realization of the operators $`B_2^+`$, $`B_3^+`$, $`F_2^+`$ and $`F_3^+`$ was obtained in . One can easily check that the explicit forms of these operators found in , are indeed identical to those obtained from (35). This equivalence is valid modulo an overall normalization factor. Thus, we have given a systematic way to determine $`B_n^+`$ and $`F_n^+`$ for arbitrary $`n`$. It might be noted here that the particular basis we choose for the definitions of these operators is over-complete. However, one may always choose a basis similar to one given in to avoid the over-completeness.
### B Supersymmetry-breaking phase
Consider the following supercharges,
$$\stackrel{~}{Q}=\underset{i}{}\psi _i\left(p_ii\stackrel{~}{W}_i\right),\stackrel{~}{Q}^{}=\underset{i}{}\psi _i^{}\left(p_i+i\stackrel{~}{W}_i\right),\stackrel{~}{W}=\lambda ln\underset{i<j}{}x_{ij}+\frac{1}{2}\underset{i}{}x_i^2.$$
(41)
These supercharges can be obtained from Eqs. (1) and (9) by making $`\lambda \lambda `$ and $`\psi _i\psi _i^{}`$. The dual Hamiltonian $`H_d=\frac{1}{2}\{\stackrel{~}{Q},\stackrel{~}{Q}^{}\}`$ differs from $`H`$ by the fermionic number operator $`N_f`$ and a constant. In particular,
$`H_d`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i}{}}{\displaystyle \frac{^2}{x_i^2}}+{\displaystyle \frac{1}{2}}\lambda (\lambda 1){\displaystyle \underset{ij}{}}x_{ij}^2+{\displaystyle \frac{1}{2}}{\displaystyle \underset{i}{}}x_i^2+{\displaystyle \frac{1}{2}}N\left(1+\lambda (N1)\right)`$ (42)
$``$ $`{\displaystyle \underset{i}{}}\psi _i^{}\psi _i+\lambda {\displaystyle \underset{ij}{}}x_{ij}^2\left(\psi _i^{}\psi _i\psi _i^{}\psi _j\right),`$ (43)
$`H`$ $`=`$ $`H_d+2N_fN\left(1+\lambda (N1)\right).`$ (44)
The ground state of $`H_d`$ is in the $`N`$ fermion sector,
$$\stackrel{~}{\mathrm{\Phi }}=e^{\stackrel{~}{W}}|\overline{0}>=\underset{i<j}{}x_{ij}^\lambda e^{\frac{1}{2}_ix_i^2}|\overline{0}>,$$
(45)
which is normalizable for $`\lambda <\frac{1}{2}`$. A stronger criteria that each momentum operator $`p_i`$ is self-adjoint for wave-functions of the form $`\stackrel{~}{\mathrm{\Phi }}`$ determines $`\lambda <0`$. The supersymmetric phase of $`H_d`$ is described by $`\lambda <0`$. The wave-function $`\stackrel{~}{\mathrm{\Phi }}`$ is also an eigen-state of $`H`$ with positive energy. This is, in fact, the ground state of $`H`$ in the supersymmetry-breaking phase. The complete spectrum of $`H`$ in this phase can be obtained from $`H_d`$ by making use of the second equation of (44).
We get the super-oscillator Hamiltonian under the following transformations,
$`\stackrel{~}{H}_2`$ $`=`$ $`e^{\frac{S}{2}}e^{\stackrel{~}{W}}H_de^{\stackrel{~}{W}}e^{\frac{S}{2}}`$ (46)
$`=`$ $`{\displaystyle \underset{i}{}}\left(x_i{\displaystyle \frac{}{x_i}}\psi _i^{}\psi _i\right)+N`$ (47)
$`\stackrel{~}{H}_{sho}`$ $`=`$ $`e^{\frac{1}{2}_ix_i^2}e^{\frac{1}{4}_i\frac{^2}{x_i^2}}\stackrel{~}{H}_2e^{\frac{1}{4}_i\frac{^2}{x_i^2}}e^{\frac{1}{2}_ix_i^2}`$ (48)
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i}{}}\left({\displaystyle \frac{^2}{x_i^2}}+x_i^2\right){\displaystyle \underset{i}{}}\psi _i^{}\psi _i+{\displaystyle \frac{N}{2}}.`$ (49)
Note the difference between $`H_{sho}`$ and $`\stackrel{~}{H}_{sho}`$. The ground state is in the $`N_f=0`$ sector for the former case, while it is in the $`N_f=N`$ sector for the latter one. This is expected also, since the original many-body Hamiltonians $`H`$ and $`H_d`$ have ground states in the $`N_f=0`$ and $`N_f=N`$, respectively.
We use the first equation of (49) to construct eigen-spectrum of $`H`$. The eigen-function is given by,
$$\widehat{\mathrm{\Phi }}=e^{\stackrel{~}{W}}e^{\frac{S}{2}}\stackrel{~}{P}_{n,k}|\overline{0}>,$$
(50)
where $`\stackrel{~}{P}_{n,k}`$ is a permutationally invariant polynomial under the combined exchange of $`x_i`$ and $`\psi _i`$. We may choose $`\stackrel{~}{P}_{n,k}`$ to have the same form as $`P_{n,k}`$, except for the replacement $`\psi _i^{}\psi _i`$. Following the discussions on the supersymmetry-preserving phase in Sec. III.A.1, it can be checked easily that this choice of $`\stackrel{~}{P}_{n,k}`$ results in well-behaved, normalizable eigenfunction for $`H`$.
The complete eigenstates can also be constructed with the help of bosonic creation operator $`\widehat{B}_n^+`$ and the fermionic creation operator $`\widehat{F}_n^+`$. We define,
$$\widehat{B}_n^+=\underset{i}{}\widehat{T}^1b_i^{+^n}\widehat{T},\widehat{F}_n^+=\widehat{T}^1\left(\underset{i}{}\psi _ib_i^{+^{n1}}\right)\widehat{T},\widehat{q}_n^+=\widehat{T}^1\left(\underset{i}{}\psi _i^{}b_i^{+^n}\right)\widehat{T},$$
(51)
with $`\widehat{T}=e^{\frac{S}{2}}e^{\widehat{W}}`$. The eigenstates are,
$$\widehat{\mathrm{\Phi }}_{n_1,\mathrm{},n_N,\nu _1,\mathrm{},n_N}=\underset{k=1}{\overset{N}{}}\widehat{B}_k^{+^{n_k}}\widehat{F}_k^{+^{\nu _k}}\stackrel{~}{\mathrm{\Phi }},$$
(52)
with the eigen-values, $`E=N(1\lambda (N1))+_{k=1}^N(kn_k+(k2)n_k)`$. The bosonic quantum numbers $`n_k`$’s are non-negative integers, while the fermionic quantum numbers $`\nu _k`$’s are either $`0`$ or $`1`$.
## IV Generalization
We have constructed a similarity transformation which shows the equivalence between the SRCMSM and free super-oscillators. The particular SRCMSM we considered is associated with the $`A_{N+1}`$ type root-structure of the Lie algebra. SRCMSM associated with other root structures also can be shown to be equivalent to free super-oscillators. Instead of considering each model separately, we prove below a general result, which is applicable to all types of SRCMSM and also to a new class of rational models considered in having nearest-neighbor and next-nearest-neighbor interactions. In particular, we consider a super-Hamiltonian $``$ whose bosonic many-body potential is a homogeneous function of degree $`2`$ and all the particles are confined on the line by a common harmonic oscillator potential. It is worth recalling at this point that all types of SRCMSM and models considered in , indeed satisfy this criteria. We construct a similarity transformation which shows the equivalence between $``$ and free super-oscillators $`H_{sho}`$.
Let us decompose the superpotential $`W`$ in terms of superpotentials for the many-body interaction and the harmonic term as,
$$W=w+\frac{1}{2}\underset{i}{}x_i^2,w=lnG(x_1,x_2,\mathrm{},x_N),$$
(53)
where $`G`$ is a homogeneous function of any arbitrary positive degree $`d`$,
$$\underset{i}{}x_i\frac{G}{x_i}=dG.$$
(54)
This property of $`G`$ ensures that each $`w_i`$ is a homogeneous function of degree $`1`$ and hence, the bosonic potential is always homogeneous function of degree $`2`$, apart from the harmonic term. The Hamiltonian is given by,
$$=\frac{1}{2}\underset{i}{}\left[\frac{^2}{x_i^2}+w_i^2+w_{ii}+x_i^2\right](d+\frac{N}{2})+\underset{i}{}\psi _i^{}\psi _i\underset{i,j}{}w_{ij}\psi _i^{}\psi _j.$$
(55)
This Hamiltonian has a dynamical $`OSp(2|2)`$ supersymmetry. The full $`OSp(2|2)`$ algebra and the operators realizing this algebra are given in Appendix-A.
The bosonic sub-algebra $`O(2,1)\times U(1)`$ of $`OSp(2|2)`$ is present for a wide class of Hamiltonians $``$, due to the constraint (54) on the superpotential. This class of Hamiltonians having $`O(2,1)\times U(1)`$ symmetry can even be made larger by adding a term $`T`$ having the following properties,
$$[N_f,T]=0,[\underset{i}{}x_i\frac{}{x_i},T]=2T,$$
(56)
to the Hamiltonian $``$. However, the new Hamiltonian $`^{}=+T`$ will not be supersymmetric anymore for general $`T`$. It is worth mentioning at this point that the presence of the symmetry algebra $`O(2,1)\times U(1)`$ is enough to show the equivalence between $`^{}`$ and free super-oscillators. The supersymmetry of the Hamiltonian does not play any role. In other words, our results are valid even if the $`OSp(2|2)`$ symmetry of $`^{}`$ is lost, but, has only $`O(2,1)\times U(1)`$ symmetry. However, we restrict our discussions in this paper to $`OSp(2|2)`$ supersymmetric Hamiltonian $``$ only.
Observe that $``$ can be transformed to a new Hamiltonian $`_1`$ under the following similarity transformation,
$`_1`$ $`=`$ $`e^We^W`$ (57)
$`=`$ $`{\displaystyle \underset{i}{}}\left(x_i{\displaystyle \frac{}{x_i}}+\psi _i^{}\psi _i\right)\widehat{S}`$ (58)
$`\widehat{S}`$ $`=`$ $`{\displaystyle \underset{i}{}}\left({\displaystyle \frac{1}{2}}{\displaystyle \frac{^2}{x_i^2}}+w_i{\displaystyle \frac{}{x_i}}\right)+{\displaystyle \underset{i}{}}w_{ii}\psi _i^{}\psi _i+{\displaystyle \underset{ij}{}}w_{ij}\psi _i^{}\psi _j.`$ (59)
The total fermion number $`N_f=_i\psi _i^{}\psi _i`$ commutes with $`\widehat{S}`$, $`[N_f,\widehat{S}]=0`$. The commutation relation between the Euler operator $`E=_ix_i\frac{}{x_i}`$ and $`\widehat{S}`$ is given by,
$$[E,\widehat{S}]=2\widehat{S},[H_2,\widehat{S}]=[E+N_f,\widehat{S}]=2\widehat{S}.$$
(60)
The homogeneity property (54) of $`G`$ has been used in deriving the above equations. Now it is easy to show that $`_1`$ is transformed to $`H_2`$ under the following transformation,
$$H_2=e^{\frac{1}{2}\widehat{S}}_1e^{\frac{1}{2}\widehat{S}}.$$
(61)
The super-oscillator Hamiltonian can be obtained from (61) by using the same transformation as used in equation (22).
One might wonder at this point that any super-Hamiltonian with the superpotential $`W`$ described by (53) and (54) is exactly solvable, due to its equivalence to free super oscillators through similarity transformations. We would like to point out that this may not be true always, because, merely showing the equivalence of different models is not sufficient for such conclusions. We have to make sure that the similarity transformation, which is responsible for such equivalence, keeps the original Hamiltonian in its own Hilbert space. Thus, as a check, one should show that the complete spectrum and the corresponding well-behaved, normalizable eigen-functions of $``$ can be constructed from $`H_{sho}`$ or $`H_2`$ through inverse similarity transformation. The equation (61) act as a necessary condition, while the construction of the complete spectrum and associated well-behaved eigen-functions of the original Hamiltonian from the super-oscillator model is sufficient to claim the equivalence between these two Hamiltonians. We discuss these points below with the example of $`BC_{N+1}`$-type SRCMSM.
### A $`BC_{N+1}`$-type SRCMSM and super-half-oscillator
The superpotential for the $`BC_{N+1}`$-type SRCMSM is described by,
$$G(\lambda ,\lambda _1,\lambda _2)=\underset{i<j}{}\left(x_i^2x_j^2\right)^\lambda \underset{k}{}x_k^{\lambda _1}\underset{l}{}(2x_l)^{\lambda _2},$$
(62)
where $`\lambda `$, $`\lambda _1`$ and $`\lambda _2`$ are arbitrary parameters. The $`D_{N+1}`$-type model is described by $`\lambda _1=\lambda _2=0`$, while $`\lambda _1=0(\lambda _2=0)`$ describes $`C_{N+1}(B_{N+1})`$-type Hamiltonian. Without loss of any generality, we restrict our discussions to the $`B_{N+1}`$-type Hamiltonian only. The Hamiltonian is given by,
$`H_{B_{N+1}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i}{}}{\displaystyle \frac{^2}{x_i^2}}+{\displaystyle \frac{1}{2}}\lambda (\lambda 1){\displaystyle \underset{ij}{}}\left[x_{ij}^2+(x_i+x_j)^2\right]+{\displaystyle \frac{1}{2}}\lambda _1(\lambda _11){\displaystyle \underset{i}{}}x_i^2`$ (65)
$`+{\displaystyle \frac{1}{2}}{\displaystyle \underset{i}{}}x_i^2{\displaystyle \frac{1}{2}}N\left[1+2\lambda (N1)+\lambda _1\right]+{\displaystyle \underset{i}{}}\psi _i^{}\psi _i+\lambda _1{\displaystyle \underset{i}{}}\psi _i^{}\psi _ix_i^2`$
$`+\lambda {\displaystyle \underset{ij}{}}\left[x_{ij}^2\left(\psi _i^{}\psi _i\psi _i^{}\psi _j\right)+(x_i+x_j)^2\left(\psi _i^{}\psi _i+\psi _i^{}\psi _j\right)\right].`$
The many-body potential is not translationally invariant like $`A_{N+1}`$-type SRCMSM. Each particle interacts with the images of all other particles and also with itself. This kind of Hamiltonians are suitable for describing systems with boundaries. We choose to work in the $`0<x_1<x_2<\mathrm{}<x_N`$ sector of the phase space. Solutions in other sectors can be obtained by using the fact that the Hamiltonian is permutationally invariant under the combined exchange of $`x_i`$ and $`\psi _i`$. The Hamiltonian also has a very interesting discrete symmetry. It is invariant under any pair $`(x_i,\psi _i)(x_i,\psi _i)`$. This reflection symmetry has a consequence on the spectrum.
The ground-state of (65) in the supersymmetric phase is given by,
$$\mathrm{\Phi }=\underset{i<j}{}\left(x_i^2x_j^2\right)^\lambda \underset{k}{}x_k^{\lambda _1}e^{\frac{1}{2}_ix_i^2},$$
(66)
with $`\lambda ,\lambda _1>0`$. We would like to emphasize here that $`\mathrm{\Phi }`$ is normalizable for $`\lambda ,\lambda _1>\frac{1}{2}`$. However, a stronger criteria that each momentum operator $`p_i`$ is self-adjoint for the wave-function of the form $`\mathrm{\Phi }`$ has been imposed. This requires $`\lambda `$ and $`\lambda _1`$ to be positive definite. The supersymmetry-breaking phase of the $`BC_{N+1}`$-type model has a richer structure than the $`A_{N+1}`$-type model. In the parameter space of $`\lambda `$ and $`\lambda _1`$, there are three regions for which the supersymmetry is broken. They are, (i) $`\lambda <0`$, $`\lambda _1<0`$, (ii) $`\lambda <0`$, $`\lambda _1>0`$ and (iii) $`\lambda >0`$, $`\lambda _1<0`$. We first discuss the spectrum in the supersymmetric phase in the next section. The spectrum in the supersymmetry-breaking phase will be discussed subsequently.
#### 1 Supersymmetric phase
The complete spectrum of $`H_{B_{N+1}}`$ is described by a subset of the spectrum of super-oscillators,
$$E_{B_{N+1}}=2(n+k+N_f),E_{sho}=2n+k+N_f.$$
(67)
At a first thought, this observation might lead to a wrong conclusion regarding the validity of the similarity transformation for $`B_{N+1}`$-type SRCMSM. This apparent contradiction is removed once the discrete reflection symmetry of the $`H_{B_{N+1}}`$ is imposed on the eigen-functions of $`H_2`$. In particular, we have to choose,
$$\widehat{P}_{n,k}=\frac{1}{N_f!}r^{2n}\underset{i_1,i_2,\mathrm{},i_{N_f}}{}f_{i_1i_2\mathrm{}i_{N_f}}(x_1,x_2,\mathrm{},x_N)(x_{i_1}\psi _{i_1}^{})(x_{i_2}\psi _{i_2}^{})\mathrm{}(x_{i_{N_f}}\psi _{i_{N_f}}^{}),$$
(68)
where $`f`$ is anti-symmetric under the exchange of any two indices and a homogeneous function of degree $`2k`$. Note that $`\widehat{P}_{n,k}`$ is invariant under, (a) $`(x_i,\psi _i^{})(x_j,\psi _j^{})`$ and (b) $`(x_i,\psi _i^{})(x_i,\psi _i^{})`$. With this choice of $`\widehat{P}_{n,k}`$, the eigen-value $`E_{sho}^{}`$ of $`H_2`$ is identical with $`E_{B_{N+1}}`$, $`E_{B_{N+1}}=E_{sho}^{}=2(n+k+N_f)`$. Also, the action of $`\widehat{S}^m`$ on $`\widehat{P}_{n,k}`$ does not produce any singularity for positive $`m`$. Thus, $`H_{B_{N+1}}`$ is equivalent to a set of free super-half-oscillators. An explanation on the use of the term ‘super-half-oscillator’ is in order. Note that both $`E_{B_{N+1}}`$ and $`E_{sho}^{}`$ are always even for any integer $`n`$, $`k`$ and $`N_f`$. On the other hand, there is no such restriction on $`E_{sho}`$. It can be both even and odd. Thus, $`E_{sho}^{}`$ or $`E_{BN+1}`$ describes only half of the spectrum described by $`E_{sho}`$. This is because the eigen-value $`E_{sho}`$ is for $`N`$ super-oscillators defined on the full-line. On the contrary, the super-Hamiltonian $`H_{B_{N+1}}`$ is defined only on the positive half-line and, hence, the $`E_{B_{N+1}}`$ or $`E_{sho}^{}`$ corresponds to the eigen-value of a set of free super-oscillators on the half-line. Thus, in analogy with the similar problem for a single particle oscillator Hamiltonian, we use the term ‘super-half-oscillator’.
The eigen-spectrum also can be constructed in an algebraic way. We define the creation and annihilation operators as,
$$_n^+=𝒯^1\underset{i}{}b_i^{+^{2n}}𝒯,_n^+=𝒯^1\underset{i}{}\psi _i^{}b_i^{+^{2n1}}𝒯,𝒯=e^{\frac{\widehat{S}}{2}}e^W.$$
(69)
Note that these operators are invariant under (a) and (b). Thus, the eigen-functions obtained by operating these operators on the ground-state also are invariant under (a) and (b). The eigen-states are obtained as,
$$\chi _{n_1\mathrm{}n_N\nu _1\mathrm{}\nu _N}=\underset{k=1}{\overset{N}{}}_k^{+^{n_k}}_k^{+^{\nu _k}}\mathrm{\Phi },$$
(70)
with the energy $`=_{k=1}^N2k(n_k+\nu _k)`$. The bosonic quantum numbers are non-negative integers, while the fermionic quantum numbers are $`0`$ or $`1`$. Note that the energy $``$ can be interpreted as that of $`N`$ independent super-half-oscillators with the frequencies $`1,2\mathrm{},N`$.
#### 2 Supersymmetry-breaking phase
The eigen-spectrum of the Hamiltonian in the region (i) can be obtained in a similar way as described in Sec. III.B. In particular, we construct a dual Hamiltonian $`H_{B_{N+1}}^d`$ from $`H_{B_{N+1}}`$ by the transformations, $`\psi _i\psi _i^{}`$, $`\lambda \lambda `$ and $`\lambda _1\lambda _1`$. The relation between these two Hamiltonians is given by,
$$H_{B_{N+1}}=H_{B_{N+1}}^d+2N_fN\left[1+2\lambda (N1)+\lambda _1\right].$$
(71)
Using this relation, the complete eigen-spectrum of $`H_{B_{N+1}}`$ in the supersymmetry-breaking phase can be obtained. In particular, the bosonic and fermionic creation operators can be obtained from (69) by replacing $`\lambda \lambda `$, $`\lambda _1\lambda _1`$ and $`\psi _i\psi _i^{}`$. These operators acting on the ground state of $`H_{B_{N+1}}^d`$ produces the eigenstates of $`H_{B_{N+1}}`$ with the eigen-value, $`E=N(12\lambda (N1)\lambda _1)+_{k=1}^N2(kn_k+(k1)\nu _k)`$. This method does not work for the regions (ii) and (iii) in a straightforward way.
The method for constructing eigen-spectrum in the regions described by (ii) and (iii) are similar. We first study the Hamiltonian in the region (ii). The ground state wave-function in this region is given by,
$$\psi (\lambda ,\lambda _1)=e^\theta |\overline{0}>=\underset{i<j}{}\left(x_i^2x_j^2\right)^\lambda \underset{k}{}x_k^{1+\lambda _1}e^{\frac{1}{2}_ix_i^2}|\overline{0}>,$$
(72)
with the ground-state energy $`E=\frac{3N}{2}2\lambda N(N1)`$. It may be noted here that $`\psi (1\lambda ,\lambda _11)`$ is also an exact eigenstate in the $`N_f=0`$ sector. However, the associated energy eigen-value is greater than $`E`$ for $`N3`$. We would also like to point out here that the particular form of $`\psi `$ is due to the shape invariance of the model, relating $`N_f=0`$ sector to $`N_f=N`$ sector . Now we introduce a new Hamiltonian $`H_3`$, which is related to $`H_{B_{N+1}}`$ by the following relation,
$$H_{B_{N+1}}=H_3+2N_f\frac{N}{2}2\lambda N(N1).$$
(73)
The above relation is similar to (71). However, unlike $`H_{B_{N+1}}^d`$ or $`H_d`$ for the $`A_{N+1}`$-type model, $`H_3`$ is not supersymmetric. Thus, we can not use the methods of supersymmetric theory to determine the ground-state energy of $`H_3`$. Instead, we find by inspection that $`\psi (\lambda ,\lambda _1)`$ is the zero-energy eigenstate of $`H_3`$. Now one can check easily,
$$\stackrel{~}{H}_2=e^{\frac{\widehat{S}}{2}}e^\theta H_3e^\theta e^{\frac{\widehat{S}}{2}},$$
(74)
where $`\widehat{S}`$ can be calculated from (59) for the choice of $`w`$ as $`w=lnG(\lambda ,\lambda _1,\lambda _2=0)`$. The bosonic and fermionic creation operators can be obtained from (69) by replacing $`\lambda \lambda `$, $`\lambda _1\lambda _1`$ and $`\psi _i\psi _i^{}`$. These operators acting on $`\psi (\lambda ,\lambda _1)`$ produces the eigenstates of $`H_{B_{N+1}}`$ in region (ii) with the eigen-value, $`E=N(\frac{3}{2}2\lambda (N1))+_{k=1}^N2(kn_k+(k1)\nu _k)`$. Finally, we mention that the ground-state wave-function in the region (iii) is $`\psi ((1+\lambda ),(1\lambda _1))`$ with the ground-state energy $`E=\frac{3N}{2}2\lambda _1N(N1)`$. Rest of the analysis in this region can be done in a straightforward way.
## V Summary and Discussions
We have constructed a similarity transformation which maps the SRCMSM Hamiltonian of $`A_{N+1}`$-type to that of a supersymmetric free harmonic oscillators. This equivalence is valid only in the supersymmetry-preserving phase of SRCMSM. We have outlined methods for the construction of eigen-functions of SRCMSM from the eigen-functions of super-oscillators. Even though there is no equivalence between SRCMSM in supersymmetry-breaking phase and super-oscillators, we are able to construct eigen-spectrum in this phase by using a duality property of the model. We observed that only those eigen functions of the free super-oscillators, which are symmetric under the combined exchange of both bosonic and fermionic coordinates, produce normalizable wave-function for the SRCMSM. This has also been observed in the pure bosonic case. Thus, this brings out the highly correlated nature of these systems.
We have generalized these results to a wide class of super-Hamiltonians whose bosonic many-body interaction is a homogeneous function of degree $`2`$ and all the particles are subjected to a common harmonic confinement. These Hamiltonians are characterized by a dynamical $`OSp(2|2)`$ supersymmetry. Though this equivalence is certainly valid at the operator level, it turns out that the individual super-Hamiltonians should be analyzed carefully to see if the similarity transformation is keeping the Hamiltonian in its original Hilbert space or not. As a check to ascertain this, one should be able to construct the complete set of eigen-values and associated eigen functions of the original Hamiltonian from the super-oscillator model. We discussed the $`BC_{N+1}`$-type SRCMSM as an example and showed its equivalence to half of the spectrum of super-oscillators. To the best of our knowledge, this is the first instance in the literature where the complete spectrum and the eigenstates of the $`BC_{N+1}`$-type SRCMSM has been obtained.
The SRCMSM associated with the root structures other than $`A_{N+1}`$ and $`BC_{N+1}`$ type have not been touched upon in this paper. We believe that permutationally symmetric super-oscillator basis with additional symmetry requirements coming from the specific nature of the root structure, as in the case of $`BC_{N+1}`$-type model, would produce the complete spectrum and associated well-behaved eigen-functions of these models. An universal formulation of the method described here, valid for SRCMSM associated with all the root structures along the line of investigations carried out in , is desirable.
The equivalence relation (61) is valid for a wide class of super-models realizing $`OSp(2|2)`$ supersymmetry. The CMS systems form only a small subset. We have seen that Eq. (61) act as a necessary condition for the equivalence between the original Hamiltonian and super-oscillators. The sufficient condition for the equivalence is to construct the complete spectrum and associated wave-functions of the original Hamiltonian from the super-oscillator basis. Thus, it would be interesting to construct new exactly solvable super-models using the method described here.
###### Acknowledgements.
I would like to thank Tetsuo Deguchi for a careful reading of the manuscript and comments. This work is supported by a fellowship (P99231) of the Japan Society for the Promotion of Science.
## A $`OSp(2|2)`$ super-algebra
We show in this appendix that the Hamiltonian $``$ in (55) has dynamical $`OSp(2|2)`$ supersymmetry. We first define the following four supercharges,
$`q={\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{N}{}}}\psi _i^{}(p_i+iw_iix_i),q^{}={\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{N}{}}}\psi _i(p_iiw_i+ix_i),`$ (A1)
$`\stackrel{~}{q}={\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{N}{}}}\psi _i\left(p_iiw_iix_i\right),\stackrel{~}{q}^{}={\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{N}{}}}\psi _i^{}\left(p_i+iw_i+ix_i\right).`$ (A2)
The Hamiltonian $``$ is given in terms of $`q`$ and $`q^{}`$as, $`=2\{q,q^{}\}`$. The dual Hamiltonian $`^d`$ can be constructed in terms of $`\stackrel{~}{q}`$ and $`\stackrel{~}{q}^{}`$, $`^d=2\{\stackrel{~}{q},\stackrel{~}{q}^{}\}`$. We define the following operators,
$`h={\displaystyle \frac{1}{2}}\left(+^d\right),U={\displaystyle \frac{1}{2}}\left(^d\right),`$ (A3)
$`_2^{}=_0{\displaystyle \frac{1}{4}}{\displaystyle \underset{i}{}}x_i^2{\displaystyle \frac{1}{4}}\left(N+2E\right),_2^+=_0{\displaystyle \frac{1}{4}}{\displaystyle \underset{i}{}}x_i^2+{\displaystyle \frac{1}{4}}\left(N+2E\right),`$ (A4)
$`_0={\displaystyle \frac{1}{4}}{\displaystyle \underset{i}{}}\left(p_i^2+w_i^2+w_{ii}2{\displaystyle \underset{j}{}}w_{ij}\psi _i^{}\psi _j\right).`$ (A5)
The bosonic operators $`_2^\pm `$ and $`h`$ satisfies the following relations,
$$[h,_2^\pm ]=\pm 2_2^\pm ,[_2^{},_2^+]=h.$$
(A6)
The commutator relation $`[E,_0]=2_0`$ has been used in deriving the above equations. The $`U(1)`$ generator $`U`$ commutes with $`_2^\pm `$ and $`h`$.
The non-vanishing anticommutators among $`q`$, $`q^{}`$, $`\stackrel{~}{q}`$ and $`\stackrel{~}{q}^{}`$ are,
$$\{q,q^{}\}=\frac{1}{2}(h+U),\{\stackrel{~}{q},\stackrel{~}{q}^{}\}=\frac{1}{2}(hU),\{q^{},\stackrel{~}{q}^{}\}=_2^+,\{q,\stackrel{~}{q}\}=_2^{}.$$
(A7)
Observe that the relation,
$$=_d+2U,$$
(A8)
which is useful in determining the spectrum in the supersymmetry-breaking phase, follows easily from the first two equations of (A7).
The other non-vanishing commutators are,
$`[_2^+,q]=\stackrel{~}{q}^{},[_2^+,\stackrel{~}{q}]=q^{},[_2^{},\stackrel{~}{q}^{}]=q,[_2^{},q^{}]=\stackrel{~}{q},`$ (A9)
$`[h,q^{}]=q^{},[h,q]=q,[h,\stackrel{~}{q}]=\stackrel{~}{q},[h,\stackrel{~}{q}^{}]=\stackrel{~}{q}^{},`$ (A10)
$`[U,\stackrel{~}{q}]=\stackrel{~}{q},[U,\stackrel{~}{q}^{}]=\stackrel{~}{q}^{},[U,q^{}]=q^{},[U,q]=q.`$ (A11)
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# 1 The statement of the Problem
## 1 The statement of the Problem
The contemporary views on the evolution of globular clusters include the following effects:
* Evaporation of stars.
* The radius (tidal) of a cluster is determined by the condition of compensation of the cluster’s gravity and the field of the Galaxy.
* Core collapse.
* Core collapse may be halted by one or several hard binaries.
* Tidal interaction or shocks of the Galaxy can accelerate the the core collapse.
The time scales of these processes for typical globular clusters are of the order of billion years.
Consider now a binary system both components of which are globular clusters, i.e. when we have a binary globular cluster, particulary a cose binary globular cluster. Our problem here is to reveal heuristically the principal effects which arise when two ordinary globular clusters can be components of a binary system.
The problem of binary star clusters can have direct relation to the galaxies with double nuclei, if assuming that its both components are supermassive star clusters.
## 2 Tidal Radii of Clusters
Consider a binary system with globular clusters A and B of masses $`M_\mathrm{A}`$ and $`M_\mathrm{B}`$, mutual distance $`a`$ and orbital period $`P`$ by relation
$$a=9.4910^8P^{2/3}M^{1/3},$$
$`(1)`$
where the summary mass $`M=M_\mathrm{A}+M_\mathrm{B}`$ is in units of solar mass, $`P`$ in days and $`a`$ in pc. For a binary of globular clusters, for example, with $`M=10^6M_{}`$, and intercomponent distance $`a=300`$ pc, from (1) we have for the orbital period: $`P=510^8`$ years, i.e. comparable with the dynamical time of the galaxy. During that time scale such a binary system can perform two-three tens of revolutions around the mass center of the galaxy.
However, by reasons which we discuss later, we shall deal with masses of the order of $`10^810^9M_{}`$, i.e. higher as compared with the mass of ordinary globular clusters in our Galaxy. Adopting therefore, $`M=10^8M_{}`$ and $`a=300`$ pc we find from (1): $`P=510^7`$ years, i.e. shorter as compared with the dynamical time scale of the galaxy.
The actual sizes of globular clusters in the Galaxy are essentially determined by the tidal field of the Galaxy. For an isolated globular cluster the tidal radius $`r_\mathrm{g}`$, is determined by the gravitation action of the Galaxy of a mass $`M_\mathrm{g}`$ and radius $`R_\mathrm{g}`$; when the cluster is located in the environments of the Galaxy the tidal radius is given by so called Hoerner’s formula (von Hoerner, 1957):
$$r_\mathrm{g}=R_\mathrm{g}\left(\frac{M_\mathrm{A}}{2M_\mathrm{g}}\right)^{1/3}.$$
$`(2)`$
As we see, the tidal radius $`r_\mathrm{g}`$ of a globular cluster depends on the masses both of cluster $`M_\mathrm{A}`$ and the Galaxy $`M_\mathrm{g}`$. At $`M_\mathrm{A}=10^7M_{}`$, for example, for $`M_\mathrm{g}=10^{11}M_{}`$ and $`R_\mathrm{g}=\mathrm{15\hspace{0.17em}000}`$ pc one has $`r_\mathrm{g}=550`$ pc.
The situation is quite different when we deal with a close binary globular cluster. In this case, the tidal radius $`r_\mathrm{t}`$ for one of the components - the cluster A, is determined by the gravitation action of the second cluster B:
$$\frac{GM_\mathrm{B}}{(ar_\mathrm{t})^2}\frac{GM_\mathrm{B}}{a^2}=\frac{GM_\mathrm{A}}{r_\mathrm{t}^2}$$
$`(3)`$
or
$$\frac{1}{(1x)^2}\frac{q}{x^2}=1,$$
$`(4)`$
where $`x`$ and $`q`$ denote
$$x=\frac{r_\mathrm{t}}{a}.q=\frac{M_\mathrm{A}}{M_\mathrm{B}}.$$
$`(5)`$
Eq.(4) is the correct relationship for the determination of $`x`$, i.e. of $`r_t`$. In particular case, when $`x<<1`$, which corresponds “Galaxy - Globular cluster” combination, we shall have from (4): $`x=(q/2)^{1/3}`$ or $`r_t=R_g(q/2)^{1/3}`$, i.e. the Hoerner’s formula (2).
As follows from this relationship the tidal radius $`r_\mathrm{t}`$ for the component A depends on the ratio of masses of both components, but not on the absolute magnitudes of their masses.
In Table I, the numerical values of tidal radius $`r_\mathrm{t}`$, obtained with the help of Eq.(4), are presented for three values $`q=1,2,3`$ and the intercomponent distance $`a=500`$ pc. As we see, in this case the tidal radii are 1.5-2 times smaller than those determined by the gravitational action of the Galaxy. Particularly, for equal masses of both cluster-components, i.e. when $`q=1`$, we have for the tidal radius $`r_t=265`$ pc essentially smaller than we have at “Galaxy - Globular cluster” combination,i.e. $`r_g=550`$ pc.
Thus, we arrive at the first remarkable, and important conclusion: the tidal radius of a globular cluster being a member of a binary cluster system, is determined by the gravitation attraction of the second cluster and not of the Galaxy. Such clusters, therefore, are expected to be more compact than typical isolated globular clusters. In this case the real sizes of globular clusters are determined not by tidal interaction of Galaxy mass but by mutual tidal interactions of globular clusters.
## 3 Roundchrom of a Binary Globular Cluster?
Recently the concept of stellar roundchroms for close binary stellar systems has been proposed to explain the observed anomalously high luminosities of the ultraviolet doublet 2800 MgII of RS CVn type systems (Gurzadyan, 1996, 1997 ab). Physically the roundchrom is a common chromosphere enveloping both components of the binary system having no geometrical contact of their photospheres. The peculiarity of that concept is in the possibility to provide a large emission volume around the secondary component of the binary system. Geometrical and physical parameters of such configurations have been revealed for over 60 RS CVn type close binary stellar systems.
That idea had inspired us to look for an associated phenomenon for a binary globular cluster. As an illustration, in Fig. 1 the configuration of a cluster roundchrom, i.e. an eight-shaped zero-velocity equipotential curve enveloping both components of the binary cluster is shown for the masses of components $`M_\mathrm{A}=410^9M_{}`$, $`M_\mathrm{B}=210^9M_{}`$, i.e. for the mass ratio $`\mu =0.33`$, intercomponent distance $`a=800`$ pc and Jacobi constant $`C=3.94`$, so that a narrow corridor can be formed in the intermediate Lagrangian point $`L_1`$ between components. The dashed circles denote the Roche lobes with radii $`R_\mathrm{L}`$ determined according to the relationship (Eggleton, 1983)
$$R_\mathrm{L}=\frac{0.49q^{2/3}}{0.6q^{2/3}+\mathrm{ln}(1+q^{1/3})},$$
$`(6)`$
yielding $`R_\mathrm{L}(\mathrm{A})=350`$ pc and $`R_\mathrm{L}(\mathrm{B})=175`$ pc.
In the case of stars, owing to the existence of strong outer boundaries – photospheres, one can compare the star’s size with its Roche lobe. In majority of cases the sizes of main components in close binary systems coincide with their Roche lobes. Such a coincidence, again as a rule, is not in the case of secondary components; thus an extensive emission volume is formed between the outer boundaries of stars and their Roche lobes.
Fig. 1
In the case of globular clusters the situation is obviously quite different since one deals with N-body system and hence no clear boundary of the system exists. Nevertheless, the chaotic properties of the stellar orbits (see Gurzadyan, Pfenniger, 1994) ensure that basic effects can arise also for binary globular clusters. Therefore, it appears possible to develop the principal idea of a roundchrom for the case of binary globular cluster supposing for simplicity that the linear sizes of the component-clusters are proportional to $`N_\mathrm{A}`$ and $`N_\mathrm{B}`$, and the outer boundary of the first component is coinciding with its Roche lobe, and is somewhat smaller in the case of the second component.
Then Figure 1 should be interpreted in the following manner. If the outer boundary of the component A is not far from the point $`L_1`$, then one should expect a transfer or flow of stars from A to B. The surrounding volume, i.e. the ”roundchrom” of the secondary will be occupied mainly by these stars, resulting a monotonous decrease of the mass of the primary A with corresponding increase of the mass of the secondary B. As a result, the balance between the kinetic and potential energies cannot remain frozen for each of the clusters, along with the secular variation of their linear sizes.
## 4 The Expansion and Contraction of Clusters at Star Transfer Process
We now aim to determine the rate of the expansion of the second component-cluster (B) of the binary system, and the rate of contraction of the first one (A). More definitely, we hope to deduce the law of the growth of the radius $`R_\mathrm{B}(t)`$ of cluster B and the law of the decreasing of the radius $`R_\mathrm{A}(t)`$ by time, i.e. during the star-transfer process from the first component to the second one.
It is convenient to start from the following basic dynamical principles. When the moment of inertia of the system can be considered as constant within some time interval, i.e. its variation is slow enough, one readily comes to the virial theorem
$$U(t)+2T(t)=0,$$
$`(7)`$
Note, that this relation is obtained not at the infinite time limit, but via averaging within finite time interval.
Considering the slow stationary star transfer process with conservation of the moment of inertia (if the clusters, for example, are far from the last phases of core collapse), we can assume that the energy conservation condition
$$T_0+U_0=T(t)+U(t),$$
$`(8)`$
is fulfilled for each component cluster separately. Similarly, it is well known that for an isolated cluster though the evaporated stars carry away some energy, the total energy of the cluster still can be considered constant within some long enough time interval while estimating the dynamical characteristics of the cluster.
Then, we have the well known relationship between radius $`R`$, velocity dispersion $`v_o`$ and the mass $`M=mN`$ of the cluster
$$v_o^2=\frac{G}{2}\frac{mN}{R},$$
$`(9)`$
where $`N`$ is the number of stars, $`m`$ is the mass of a star.
Now we can obtain the variations of the sizes for each of these components of binary cluster A and B during of star transfer process through the Lagrangian point $`L_1`$ from the component A to B.
First consider the behavior of the component B. For the initial values of kinetic $`T_o`$ and potential energy $`U_o`$ we have
$$T_0=N_\mathrm{B}\frac{mv_0^2}{2}$$
$`(10)`$
$$U_0=\frac{G}{2}\frac{m^2N_\mathrm{B}^2}{R_0},$$
$`(11)`$
where $`R_o`$ is the initial radius of the cluster B.
As a result of star transfer process the cluster B gains an additional kinetic energy
$$\mathrm{\Delta }T=nt\frac{mv_{}^2}{2},$$
$`(12)`$
where $`v_{}`$ is the velocity of stars in the transit point $`L_1`$, and $`n`$ is the flux of stars with a constant rate $`n`$. Thus the total kinetic energy of cluster B at the moment $`t`$ will be
$$T(t)=T_o+\mathrm{\Delta }T=N_\mathrm{B}\frac{mv_o^2}{2}+nt\frac{mv_{}^2}{2}.$$
$`(13)`$
Correspondingly we shall have for the potential energy $`U(t)`$ of the cluster B at the moment $`t`$
$$U(t)=\frac{G}{2}\frac{m^2(N_\mathrm{B}+nt)^2}{R_t},$$
$`(14)`$
where $`R_t`$ is the modified radius of the cluster B.
The condition which we will use is that the virial equilibrium remains valid during this process, i.e. for moment $`t`$. This condition has to be safely fulfilled during the stationary star transfer process due to violent relaxation effects occuring within dynamical (crossing) time scale. Then, substituting (13) and (14) into the virial expression (7) and having in view (9), we obtain for the law of the variation of the radius $`R_t(B)`$ of cluster B by time
$$\frac{R_t(B)}{R_0}=\frac{(1+nt/N_\mathrm{B})^2}{1+(v_{}/v_0)^2nt/N_\mathrm{B}}.$$
$`(15)`$
This same equation can be obtained using the condition of conservation of total energy of each of these clusters during the star transfer process.
Strictly speaking, $`v_{}`$ must be slightly larger than $`v_0`$. However, adopting as a first approximation $`v_{}=v_0`$, we obtain
$$\frac{R_t(B)}{R_0}=1+\frac{n}{N_\mathrm{B}}t$$
$`(16)`$
where the stellar flux rate $`n`$ is considered constant according to the assumption on the stationarity of the star transfer process.
In contrast with cluster B, the cluster A contracts during the stellar transfer process, i.e. its radius should be decreased. Similar relationship can be derived for this case as well, i.e. for the law of decrease of the radius $`R_t(\mathrm{A})`$ of cluster A
$$\frac{R_t(A)}{R_0}=1\frac{n}{N_\mathrm{A}}t.$$
$`(17)`$
However, the star transfer process from the cluster A to the cluster B cannot proceed stationary too long: with the decrease of the diameter of the cluster A, its outer boundary will become more and more farther from the point $`L_1`$. Hence, the initially assumed constant star transfer rate has to be substituted by a decreasing function
$$n_t=n_0e^{\gamma t},$$
$`(18)`$
where $`\gamma `$ in $`yr^{}1`$ characterizes the rate of the variations of the star transfer process from A to B.
In this case the total number of stars $`n(t)`$ transferred from A to B at the moment $`t`$, after starting the transfer process $`(t=0)`$, will be
$$n(t)=\frac{n_0}{\gamma }(1e^{\gamma t}).$$
$`(19)`$
As a result, we shall have for the evolution, i.e. for the growth of the radius of the cluster B at the moment $`t`$ the following relationship, again assuming $`v_{}=v_0`$:
$$\frac{R_B(t)}{R_0}=1+\frac{n_0}{\gamma N_\mathrm{B}}(1e^{\gamma t})$$
$`(20)`$
and for the decrease of the radius of A:
$$\frac{R_A(t)}{R_0}=1\frac{n_0}{\gamma N_\mathrm{A}}(1e^{\gamma t}).$$
$`(21)`$
Above, at the deducation of Eqs. (17) and (21), two extremely case for the rate of star transfer process are considered, namely, with a constant rate and with an exponetially decreasing rate. The real situation, we believe, should be expected between these two limiting cases.
In Fig.2 we give the curves of time variations for both radii, $`R_\mathrm{A}(t)`$ and $`R_\mathrm{B}(t)`$, for the second case, i.e. at exponetially decreasing star transfer rate given by expression (18), for two sets of initial parameters, namely, $`N_\mathrm{A}=510^9`$, $`N_\mathrm{B}=210^9`$, $`\gamma =210^9`$ yr<sup>-1</sup> and $`n_0=1`$, $`n_0=2`$ and $`n_0=3`$ str.yr<sup>-1</sup> in the first case (at left), and $`N_\mathrm{A}=510^8`$, $`N_\mathrm{B}=210^8`$, $`\gamma =10^8`$ yr<sup>-1</sup>, and $`n_0=0.5`$, $`n_0=1`$ and $`n_0=2`$ str.yr<sup>-1</sup> in second (at right).
Fig. 2
The main difference between both Figures is the time scale, $`10^9`$ years in first case and $`10^8`$ years, in the second. While the character itself of the evolution of sizes of binary cluster components, the contraction in one case and expansion in the other, is independent on rate of the star transfer. Note, that the growth of the radius of the cluster B takes place faster than the decrease of the radius of A. At the end of evolution, formally at $`t\mathrm{}`$, the radius of B is doubled while of A is decreased only on 1/5-th of its initial radius.
Also, the volume of roundchrom around component B is decreased monotonocally during the evolution of the system. The volume of the roundchrom around A tests an increasing as well however not so strongly.
It is interesting to note that during the essential evolution of sizes of cluster A, i.e. during of $`2.10^9`$ years, the distance between its outer boundary and the point $`l_1`$ is remained practically unchanged. This important property should be explained also by the slowly motion of the Lagrangian point $`L_1`$ to the direction of the component A synchronously with the decreasing of its mass.
Despite of essential differeces in the time scale, in initial total numbers of stars in clusters as well in the rate of star transfer process, the main character of curves in Fig.2 is nearly the same. This means that we cannot say anything for example about the age of binary system or the phase of evolution or star transfer process of the pair judging only by direct image of that or another interacting binary system in the centre of some galaxies.
As an illustration, an evolution sequence of the sizes of a binary cluster is shown in Fig.3. The initial data at $`t=0`$ are as follows: $`N_\mathrm{A}=510^9`$,
Fig. 3
$`N_\mathrm{B}=210^9`$ with the same mass (solar) of stars in both cases, i.e. $`m=M_{}`$, the ratio of initial radii of both components $`R_\mathrm{A}/R_\mathrm{B}=2.5`$. The configurations of roundchroms correspond to star transfer rate 2 star.yr<sup>-1</sup> and $`\gamma =210^9`$ yr<sup>-1</sup>. The second configuration from above corresponds to $`t=0.510^9`$ yr, the third one – to $`t=210^9`$ yr.
In Fig.3, the roundchrom-cluster concentric configurations allow us to predict the possible existence of the binary globular clusters, particularly in centres of some galaxies, one of components of which in two-envelope form, i.e. with an outer low density ring around the central dense spheric disk.
Returning to the star transfer problem, one should outline the following aspects as well:
a. One can always find wandering stars in the vicinity of the corridor in $`L_1`$, ready for the passage through this corridor.
b. The gravitation attraction of cluster B, should support the accumulation of individual stars-members of the cluster A in the vicinity of the corridor in the point $`L_1`$,
## 5 The Tidal Effect
Independently to the star transfer process from one component to another, each star, say, in the cluster A, individually undergoes gravitation attraction by the cluster B. As a result, the star may acquire an additional velocity. If so, the macrostructure of energy balance of cluster A must be modified due to the variation of its linear size.
To evaluate quantitatively these variations, one has to estimate the amount of the additional kinetic energy gained due to the gravitation attraction of the total mass $`M_\mathrm{B}`$ of B located in its center on the distance $`a`$ from the center of A as is shown in Fig.4.
Fig. 4
Assuming $`N_\mathrm{A}`$ stars of equal mass $`m`$ homogeneously distributed within a spherical cluster A with initial radius $`R_\mathrm{A}`$ and concentration
$$n_0=\frac{3}{4\pi }\frac{N_\mathrm{A}}{R_\mathrm{A}^3},$$
$`(22)`$
we can write down the equation of the motion of a star at the distance $`r`$ from $`M_\mathrm{B}`$
$$m\frac{d\mathrm{\Delta }v}{dt}=G\frac{mM_\mathrm{B}}{r^2},$$
$`(23)`$
which gives for the additional velocity $`\mathrm{\Delta }v`$ of the star
$$\mathrm{\Delta }v=G\frac{M_\mathrm{B}}{r^2}t.$$
$`(24)`$
In view of the gravitational action of the cluster B on individual stars in cluster A, i.e. due to the tidal effect, we can write for the elementary kinetic energy $`dT_{\mathrm{tid}}`$ of a mass included into the element of sphere of radius $`r`$ and thickness $`dr`$
$$dT_{\mathrm{tid}}=dm\frac{\mathrm{\Delta }v^2}{2}=\pi n_0mG^2M_\mathrm{B}^2(1\mathrm{cos}\phi )\frac{t^2}{r^2}dr,$$
$`(25)`$
where $`M_\mathrm{B}=mN_\mathrm{B}`$ and
$$\mathrm{cos}\phi =\frac{a^2R_t^2+r^2}{2ar}.$$
$`(26)`$
After integration of (25) over limits $`aR_t`$ and $`a+R_t`$, we obtain for the additional kinetic energy acquired by the cluster A under the action of cluster B
$$T_{\mathrm{tid}}(t)=\frac{3}{4}G^2m^3\frac{N_\mathrm{A}N_\mathrm{B}^2}{R_t^4}\left(\frac{1}{u^21}\frac{1}{2u}\mathrm{ln}\frac{u+1}{u1}\right)t^2,$$
$`(27)`$
where $`u=a/R_t`$.
Then the energy conservation law should be written in the form
$$T_0+U_0=T(t)+U(t)+T_{\mathrm{tid}}(t),$$
$`(28)`$
where $`T_0`$ and $`U_0`$ are the same as above. For $`T(t)`$ and $`U(t)`$ we have
$$T(t)=\frac{G}{4}\frac{m^2N_\mathrm{A}^2}{R_t},$$
$`(29)`$
$$U(t)=\frac{G}{2}\frac{m^2N_\mathrm{A}^2}{R_t}.$$
$`(30)`$
Substituting (27), (29) and (30) into (28),and in view of (10) and (11), we obtain finally the radius $`R_t`$ of cluster A at the moment of time $`t`$:
$$\frac{N_\mathrm{A}}{R_0}\left(1\frac{R_t}{R_0}\right)=Gm\frac{N_\mathrm{B}^2}{R_t^4}\left(\frac{1}{u^21}\frac{1}{2u}\mathrm{ln}\frac{u+1}{u1}\right)t^2$$
$`(31)`$
with $`u=a/R_t`$, and $`R_0`$ is the initial radius of cluster A.
The relationship (31) defines the law of the variation of cluster’s radius $`R_t`$ with time $`t`$. As we see, the tidal effect of the cluster B leads to the increase of the radius $`R_t`$ of cluster A.
In Fig.5 the curves of the growth of cluster’s relative radius, $`R_t/R_0`$, with $`t`$ calculated with the help of (31), are drawn for three numerical values of $`N_\mathrm{A}`$: $`410^8`$, $`310^8`$ and $`210^8`$ with $`N_\mathrm{B}=10^8`$ and $`m=M_{}`$, cluster radius $`R_0=100`$ pc and distance between the centers of clusters $`a=300`$ pc.
Fig. 5
Note, that the growth of the radius $`R_t/R_0`$ of cluster A due to the tidal effect of the cluster B, occurs rather rapidly – during a time scale of the order of $`10^7`$ yrs. This result seems to be in disagreement with the dynamical nature of the globular clusters. However, as we will show in the next section, the orbital rotation effect can essentially suppress the rapid growth of the radius of cluster A.
The peculiarity of the problem considered above, namely, the role of the tidal process in binary globular clusters, is that one deals not with atoms in the stellar atmospheres of interacting close binaries but with stars. The main difference is in the absence of collisions between stars, i.e. the absence of energy exchange process between stars (particles).
The next step may be, perhaps, the examination of both processes, the tidal and star-transfer, simultaneously. However, as shows the preliminary analysis, the results and basic conclusions will remain unchanged.
## 6 The Role of Orbital Rotation
Each star in both clusters will undergo the action of centrifugal force provoked by the orbital motion the components. This force may compensate in certain degree the gravitational attraction of the other cluster. This factor, as appears, may change radically the above conclusions.
The centrifugal force depends on the tangential velocity $`V_t`$ and the distance to the center of rotation of the binary system. Gravitational attraction provoked by cluster B can be compensated completely by centrifugal force only if the tangential velocity $`V_t`$ of a star in the center of cluster A satisfies the condition
$$V_t=\left(G\frac{M_\mathrm{B}}{a}\right)^{1/2}(1+M_\mathrm{A}/M_\mathrm{B})^{1/2},$$
$`(32)`$
where $`M_\mathrm{A}`$ and $`M_\mathrm{B}`$ are the masses of the clusters, as before.
On the other hand, for the orbital velocity $`V_{\mathrm{or}}`$ of a star we have ( $`P`$ is the orbital period of the binary system)
$$V_{\mathrm{or}}=\frac{\pi a}{P}$$
$`(33)`$
or using the known expression for $`P`$,
$$V_{\mathrm{or}}=\frac{1}{2}\left(G\frac{M_\mathrm{B}}{a}\right)^{1/2}(1+M_\mathrm{A}/M_\mathrm{B})^{1/2}.$$
$`(34)`$
From (34) and (32) we find for the ratio of both velocities
$$\frac{V_{\mathrm{or}}}{V_t}=\frac{1}{2}(1+M_\mathrm{A}/M_\mathrm{B}).$$
$`(35)`$
In particular case when $`M_\mathrm{A}=M_\mathrm{B}`$, we achieve to an important result
$$V_t=V_{\mathrm{or}},$$
i.e. only at equal masses of both clusters we shall have an equilibrium of both types of forces – centrifugal and gravitational; in this case the tidal acceleration of a star towards the cluster A due to the gravitation attraction of cluster B will be compensated by the centrifugal force of an orbital rotation of the system. The further evolution of the sizes of clusters will be determined by star transfer process from one component to another.
Note, that depending on the numerical value of the ratio $`M_\mathrm{A}/M_\mathrm{B}`$ we shall have different relations between $`V_t`$ and $`V_{\mathrm{or}}`$. So, if $`M_\mathrm{A}/M_\mathrm{B}>1`$ we have
$$V_t<V_{\mathrm{or}},$$
and vice versa, when $`M_\mathrm{A}/M_\mathrm{B}<1`$, we shall have
$$V_t>V_{\mathrm{or}}.$$
In the first case, when $`V_t<V_{\mathrm{or}}`$, the centrifugal force will be compensated by the gravitation attraction not completely, hence, in this case the tidal acceleration due to the component B may have some role. Accordingly, some increase of total kinetic energy or decrease of potential energy in the system A and, correspondingly, a decrease of the linear sizes of cluster A may be expected, although on less degree as compared with the limiting case given by relationship (21) or curves in Fig.2.
In the second case when $`V_t>V_{\mathrm{or}}`$, the centrifugal force will dominate over the gravitational attraction, and the final result will be the same as in the first case, i.e. we shall have some increase of kinetic energy or some decrease in linear sizes of the cluster.
The condition $`V_t=V_{\mathrm{or}}`$ can hardly be fulfilled for real cluster binaries, therefore the general conclusion may be formulated in the following manner: the decrease of linear sizes of the main component of a binary cluster at star transfer process is inevitable although at a lower rate as compared with the rate given by (31).
As to the rate of variation of cluster linear size, it depends in which degree the additional velocity $`V_t`$ is larger or smaller than the “thermal” (mean) velocity of stars $`V_0`$ given by relationship
$$V_0^2=\frac{1}{2}\frac{GM_\mathrm{A}}{R_\mathrm{A}}.$$
$`(36)`$
From (32) and (36) we find
$$\frac{V_t}{V_0}=\left(\frac{2R_\mathrm{A}}{a}\frac{M_\mathrm{B}}{M_\mathrm{A}}\right)^{1/2}(1+M_\mathrm{A}/M_\mathrm{B})^{1/2}.$$
$`(37)`$
For example, at $`M_\mathrm{A}/M_\mathrm{B}=2.5`$ and $`a/R_\mathrm{A}=3`$ we have: $`V_t/V_0=0.28`$, i.e. not too small in order to be ignored and not too large in order to be compared with the results presented in Fig.4.
Thus, the star-transfer process from cluster A to cluster B leads to the decrease of the size of cluster A. The tidal effect of the cluster B on cluster A leads to the growth of the sizes of A. These are the consequences of the tidal effect in its“pure” form. However, the rotational effect due to the orbital motion of the clusters A and B, essentially reduces the rate of this process. As a result, the behavior of the size of component A has to be rather complex.
Thus the star-transfer process from cluster A will definitely be accompanied by the decrease of the cluster’s size, however with a rate essentially depended on various physical, dynamical and kinematical parameters of both clusters.
The latter two effects, the tidal and orbital rotation, lead to changes in dynamical state of globular cluster during a time scale of the order of tens of millions years, which is much shorter as compared with the time scale of globular cluster’s evolutionary processes mentioned in Section 1. Thus the clusters in the binary systems must have different evolutionary time scales than the isolated globular clusters.
## 7 On the Binary Nuclei of Galaxies
During the last decade the discoveries of galaxies with double nuclei become more and more common. Due to their form and structure the nuclei are usually interpreted as binary black holes. Nevertheless the the concept that nuclei are compact stellar systems cannot be excluded, i.e. a supermassive globular clusters - of a mass of components of the order of $`10^9M_{}`$ and more. An example of such a double nuclei is shown in Fig.6; this Mrk 273 image is taken by Knapen et all. (1997) via Keck telescope. This image resembles an eight-shaped binary configuration with different sizes of components. At the distance of this galaxy 160 million pc, the intercomponent distance is estimated as 800 pc, with diameters of components 450 pc and 700 pc.
Fig. 6
Note that in Fig. 6 the space between the components reveals certain brightness. According to our concept the point $`L_1`$ located there must provide a corridor for star flow from one component to another, so that the intercomponent space must acquire definite brightness, which can even be roughly estimated. At the flux 1 star per year with stellar velocity 100 km s<sup>-1</sup> we shall have $`10^7`$ stars during $`10^6`$ yrs on 10 pc linear length on this area, thus forming a noticeable background brightness, as maybe indicates Mrk 273.
At present the number of known galaxies with binary nuclei is over one hundred. For some of them, including Mrk 273, the absolute magnitudes of the components $`M_\mathrm{A}`$ and $`M_\mathrm{B}`$ have been obtained (Khachikian, 1998). Then, assuming for these nuclei completely stellar composition of solar mass and solar luminosity, we can evaluate the total number of stars $`N_\mathrm{A}`$ and $`N_\mathrm{B}`$ in the components; the results are presented in Table II.
In the case of Mrk 273, the total number of stars in the components is found $`3.610^9`$ and $`2.010^9`$, i.e. almost of the same order as those used above when drawing the eight-shape curve in Fig.1. Moreover, the derived ratio $`N_\mathrm{A}/N_\mathrm{B}2`$ for this galaxy corresponds, at least qualitatively, to the photographic images of components in Fig. 6. Note, the existence of nuclei with equal number of stars in the components (Mrk 673, ratio 1:1) as well with the ratio up to 1:5 (Mrk 789). It is remarkable that the total number of stars in the binary nuclei at least for these six galaxies is of the same order – $`10^9`$-$`10^{10}`$. Even in view of observational selection, this fact cannot be totally ignored concerning the nature and composition of the binary nuclei of galaxies.
Fig. 7
Another example we have in the case of IC 4553 (Arp 220) with a remarkable binary nuclei structure discovered at radio continuum 4.83 GHz, see Fig. 7 (Baan, Hischick, 1995). The projected separation between the components is 330 pc, the reconstructed separation 466 pc. At orbital period of $`710^6`$ yr for the two nuclei with equal masses, the estimated dynamical mass for both components is $`10^{10}M_{}`$. It is important to note, that in this case no clear hot spots are found in the centers of components within the line emission structure. The spectral properties confirm the starburst dominated nature of the observed emission, and hence, the stellar content of both nuclei.
The results obtained in present article may find an interesting application, particularly, for binary globular clusters in the Large Magellanic Clusters.
## 8 Conclusions
Thus, we formulated and qualitatively considered the binary system with globular clusters as components reveals the importance at least of three aspects on their dynamics:
\- Mutual tidal interaction of both components.
\- The role of the orbital rotation of the binary system.
\- The star transfer process from one component of the system to the other.
The preliminary conclusions can be summarized as follows:
$``$ Globular cluster as a component of a binary system gains new dynamical characteristics and properties;
$``$ Dynamical evolution of a globular cluster as a member of binary system occurs more rapidly as compared with an isolated cluster;
$``$ The tidal radius of a globular cluster in a binary system is smaller as compared with the tidal radius determined by the Galactic field;
$``$ The star transfer process from one component to another in a binary cluster system leads to the reduction of linear size of the first component and the increase of the size of second one;
$``$ In binary globular clusters the tidal interaction leads to the increase of their kinetic energy and, hence, to the expansion of clusters.
$``$ Orbital rotation of a binary cluster system plays a regularing role on the variations of the sizes of the clusters.
Extensive numerical study of this problem can be of remarkable interest.
I thank V.G.Gurzadyan for valuable discussions.
## R e f e r e n c e s
Baan W.A., Haschick A.D. 1995, Ap.J. 454, 745
Eggleton P.P. 1983, Ap.J. 268, 368
Gurzadyan G.A. 1996a, A&SS 241, 211
Gurzadyan G.A. 1996b, Theory of Interplanetary Flights, Gordon & Breach,
Gurzadyan G.A. 1997a, M.N.R.A.S. 290,607
Gurzadyan G.A. 1997b, New Astr. 2, 31
Gurzadyan V.G., Pfenniger D. (Eds.) Ergodic Properties in Stellar Dynamics, Springer-Verlag, 1994.
Hoerner S. von 1957, Ap.J. 125, 451
Khachikian E.Ye.1998, Trans.Armenian Nat.Acad.Sci. 98, 239
Knapen J.H., Laine S., Yates J.A., Robinson A., et al., 1997, Ap.J.Lett. 490, L29
## Figure Captions
Roundchrom structure for a close binary system with star clusters as components with mass ratio $`\mu =0.33`$, masses of components $`M_\mathrm{A}=410^9M_{}`$, $`M_\mathrm{B}=210^9M_{}`$ and intercomponent distance $`a=800`$ pc, $`R_\mathrm{L}`$(A) and $`R_\mathrm{L}`$(B) are the radii of Roche lobes. All linear sizes are in parsecs.
The time-dependent curves of the variations of linear radii $`R_t/R_0`$ of both components, A and B, of binary system with globular clusters as components. For the initial moment, $`t=0`$, it is assumed $`R_t/R_0=1`$. The curves are calculated for three values of star-transfer rate $`n_o`$ from $`0.5`$ up to 3 str.yr<sup>-1</sup> and for two values of gamma, $`2.10^9`$ yr<sup>-1</sup> and $`10^8`$ yr<sup>-1</sup> and the time scale $`10^9`$ yr and $`10^8`$ yr. During the star-transfer process from component A to component B, component A undergoes a decrease of its radius, i.e. $`R_t/R_0<1`$ (lower half) and the component B, the contrary, an increase, $`R_t/R_0>1`$ (upper half).
Three stages of the evolution of globular clusters’ diameter as a function of time $`t`$, i.e. during star-transfer process from component A to component B through the corridor in Lagrangian point $`L_1`$. All three configurations are calculated for one and the same star-flow rate $`n_o=2`$ str.yr<sup>-1</sup> and $`gamma=2.10^9`$ yr<sup>-1</sup>. Initially, at $`t=0`$, the component B has a powerfull roundchrom (pointed area) and component A is practically without a roundchrom.
To the calculations of the tidal effect of the cluster B of a mass $`M_\mathrm{B}`$ on the stars in cluster A.
The estimated time-depended variations of radius, $`R_t/R_0`$, of component A due to the tidal effect for three values of total number of stars $`N_\mathrm{A}`$ in the cluster A, at its radius $`R_0=100`$ pc and intercomponent distance $`a=300`$ pc.
Double nuclei of Mrk 273 (Knapen et al 1997). Components of nuclei are of different masses and different sizes. The separation between components is 730 pc. Note the background in the corridor between components.
Double nuclei of IC 4553 in 4.83 GHz (Baan and Haschick 1995). The separation between nuclei is 330 pc.
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# 1 Introduction
## 1 Introduction
One of the intriguing phenomena, observed at HERA, is the behaviour of the energy dependence of the ratio $`\sigma ^{DD}/\sigma _{tot}`$ in deep inelastic scattering (DIS ). It appears ( see Ref. and Fig. 1 ) that this ratio as a function of energy is almost constant for different masses of the diffractively produced hadrons over a wide range of photon virtualities $`Q^2`$. At present there is no theoretical explanation for this striking experimental observation. The only valid theoretical idea on the market is the quasiclassical gluon field approach (see Ref. ) in which the total as well as the diffractive cross section do not depend on energy.
Recently, Yu. Kovchegov and L. McLerran suggested that the constantcy of the ratio $`\sigma ^{DD}/\sigma _{tot}`$ is closely related to strong shadowing corrections ( SC ) for diffractive production. They derived a formula for this ratio for the case where only a quark \- antiquark pair is produced in diffractive dissociation. On the other hand, K. Golec - Bierat and M. Wüsthoff suggested a phenomenological model which incorporates two main theoretical ideas regarding the transition between “hard” and “soft” processes in QCD : (i) the appearence of a new scale which depends on energy, and is related to the average transverse monentum of a parton in the parton cascade; and (ii) the saturation of the parton density at high energies. This model is successful in describing all the available experimental data on total and diffractive cross sections, including the energy behaviour of $`\sigma ^{DD}/\sigma _{tot}`$ .
The success of the above mentioned papers prompted us to reexamine the energy behaviour of $`\sigma ^{DD}/\sigma _{tot}`$ in perturbative QCD ( pQCD ) in more detail. Our approach is based on two main ideas used to describe the total and diffractive cross section for DIS:
1. The final state of the diffractive processes in HERA kinamatic region , are desribed by the diffraction dissociation of a virtual photon ( $`\gamma ^{}`$ ) into a quark-antiquark pair and quark-antiquark pair plus one extra gluon (see Fig.2 );
2. The Mueller-Glauber approach for calculating SC in the total and diffractive cross sections for DIS, this was used successfully to describe other indications of strong SC in HERA data , such as the $`Q^2`$ \- behaviour of the $`F_2`$ slope and energy behaviour of the diffractive cross section .
The paper is organized as follows. In section 2 we give a simple derivation of the Kovchegov and McLerran formula based on the $`s`$-channel unitarity constraints. We generalize this formula for the case of $`q\overline{q}G`$ final state ( see Fig.2 ) and discuss the relation between this approach and the AGK cutting rules . Section 3 is devoted to a numeric calculation of the ratio $`\sigma ^{DD}/\sigma _{tot}`$ without any restriction on the value of produced masses. In section 4 we discuss how limitations on the mass range change the energy dependence of the ratio. A summary of our results as well as a discussion on future HERA experiments are given in section 5. Appendix fives all formulae and all details of our calculations.
## 2 Shadowing corrections in QCD
### 2.1 Notations and definitions
In this paper we develop further our approach for diffractive production in DIS started in Ref.. We use the same notations and definitions as in Ref. . In this paper we examine the main physical ideas of Ref. , i.e. that the correct degrees of freedom at high enegies ( low $`x`$ ) are colour dipoles, rather than quarks and gluons which appear explicitly in the QCD Lagrangian. The consequence of this hypothesis is that a QCD interaction at high energies does not change the size and energy of a colour dipole. Hence, the majority of our variables and observables are related to the distribution and interaction of the colour dipoles in a hadron.
To facilitate reading the paper we list the notation and definitions which we will use ( see Fig.3 ):
1. $`Q^2`$ denotes the virtuality of the photon in DIS, $`M`$ the produced mass and $`W`$ the energy of the collision in c.m. frame;
2. $`x_P=(Q^2+M^2)/W^2`$ is the fraction of energy carried by the Pomeron ( two gluon ladder in Fig.3 ). Bjorken scaling variable is $`x_B=Q^2/W^2`$;
3. We use the symbol $`x`$ for both $`x_P`$ and $`x_B`$ .
4. $`\beta =Q^2/(Q^2+M^2)=x_B/x_P`$ is the fraction of the Pomeron energy carried by the struck quark;
5. $`k_{}`$ denotes the transverse momentum of the quark, and $`r_{}r`$ the transverse distance between the quark and the antiquark i.e. the size of the colour dipole;
6. $`l_{}`$ is the transverse momentum of the gluon emitted by the quark ( antiquark );
7. $`z`$ is the fraction of the photon momentum in the laboratory frame carried by the quark or antiquark;
8. $`b_t`$ is the impact parameter of the reaction and is the variable conjugated to $`q_{}`$, the momentum transfer from the incoming proton to the recoiled proton. Note that $`t=q_{}^2`$;
9. Our amplitude is normalized so that
$$\frac{d\sigma }{dt}=\pi |f(s=W^2,t)|^2,$$
(2.1)
with the optical theorem given by
$$\sigma _{tot}=\mathrm{\hspace{0.17em}\hspace{0.17em}4}\pi Imf(s,0).$$
(2.2)
10. The scattering amplitude in $`b_t`$ space is defined by
$$a^{el}(s,b_t)=\frac{1}{2\pi }d^2q_{}e^{i\stackrel{}{𝐪_{}}\stackrel{}{𝐛_𝐭}}f(s,t=q_{}^2).$$
(2.3)
11. The $`s`$-channel unitarity constraint then has the form
$$2Ima^{el}(s,b_t)=|a^{el}(s,b_t)|^2+G^{in}(s,b_t),$$
(2.4)
where $`G^{in}`$ denotes the contribution of the all inelastic processes.
Therefore, in the impact parameter representation:
$`\sigma _{tot}=\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}2}{\displaystyle d^2b_tIma^{el}(s,b_t)};`$ (2.5)
$`\sigma _{el}={\displaystyle d^2b_t|a^{el}(s,b_t)|^2};`$ (2.6)
$`\sigma _{in}={\displaystyle d^2b_tG^{in}(s,b_t)}.`$ (2.7)
12. $`x_BG(x_B,Q^2)`$ is the gluon distribution of the nucleon;
13. $`\sigma _{dipole}(x_P,r_{})`$ is the total cross section for dipole - nucleon scattering, and is given by ( see Ref. and references therein )
$$\sigma _{dipole}(x_P,r_{})=\frac{\pi ^2\alpha _S}{3}r_{}^2x_PG(x_P,\frac{4}{r_{}^2}),$$
(2.8)
where $`x_PG(x_P,\frac{4}{r_{}^2})`$ is the number of gluons with $`Q^2=\frac{4}{r_{}^2}`$ and energy $`x_P`$ in the nucleon.
14. To characterize the strength of the colour dipole interaction we introduce
$$\kappa _{dipole}=\frac{\sigma _{dipole}}{\pi R^2}=\frac{\pi ^2\alpha _S}{3\pi R^2}r_{}^2x_PG^{DGLAP}(x_P,\frac{4}{r_{}^2}),$$
(2.9)
where $`R`$ is the nonperturbative parameter which is related to the correlation radius of the gluons in a hadron. We determine its value from high energy phenomenology and HERA experimental data ( see section 3 for discussion ).
To illustrate the physical meaning of Eq. (2.9) we rewrite it in the form:
$$\kappa _{dipole}=\sigma _0\rho (r_{},x),$$
(2.10)
where
$$\rho (r_{},x)=\frac{xG^{DGLAP}(x,\frac{4}{r_{}^2})}{\pi R^2}$$
(2.11)
is the density of colour dipoles of size $`r_{}`$ and energy $`x`$ in the transverse plane. In Eq. (2.10), $`\sigma _0`$ denotes the cross section for the interaction of a dipole of size $`r_{}`$ with a point - like probe
$$\sigma _0=\frac{\pi ^2\alpha _S}{3}r_{}^2.$$
Hence, $`\kappa _{dipole}`$ is a packing factor for dipoles of size $`r_{}`$ in a proton. If $`\kappa _{dipole}`$ is small, we have a diluted gas of colour dipoles in a proton, but at low $`x`$ the gluon density increases and $`\kappa _{dipole}\mathrm{\hspace{0.17em}\hspace{0.17em}1}`$. In a such kinematic region our parton cascade becomes a dense system of colour dipoles which should be treated non-perturbatively;
15. The amplitude for colour dipole scattering on a nucleon is given by
$$a_{dipole}^{el}(s,r_{};b_t)=\frac{\pi ^2\alpha _S}{3}r_{}^2x_PG(x_P,\frac{4}{r_{}^2};b_t),$$
(2.12)
where $`x_PG(x_P,\frac{4}{r_{}^2};b_t)`$ is the number of gluons at fixed impact parameter $`b_t`$.
16. It can be shown ( see Ref. and references therein ) that we can write
$$x_PG^{DGLAP}(x_P,\frac{4}{r_{}^2};b_t)=x_PG^{DGLAP}(x_P,\frac{4}{r_{}})S(b_t),$$
(2.13)
if $`x_PG^{DGLAP}(x_P,\frac{4}{r_{}^2};b_t)`$ satisfies the DGLAP evolution equations.
In Eq. (2.13) $`S(b_t)`$ is the nucleon profile function which is a pure nonperturbative ingredient in our calculations. We assume a Gaussian form for
$$S(b_t)=\frac{1}{\pi R^2}e^{\frac{b_t^2}{R^2}},$$
(2.14)
where $`R`$ has been discussed above;
17. For the exchange of one ladder ( “hard” Pomeron ) as shown in Fig. 3, $`a_{dipole}^{el}`$ can be written as
$$\mathrm{\Omega }_{dipole}^P=a_{dipole}^{el}(x,r_{};b_t)=\kappa _{dipole}^{DGLAP}(x,r_{})e^{\frac{b_t^2}{R^2}}$$
(2.15)
18. $`\mathrm{\Psi }^\gamma ^{}(z,r_{};Q^2)`$ is the wave function of the quark - antiquark pair with the transverse distance $`r_{}`$ between a quark and an antiquark and with a fraction of energy $`z`$ ( colour dipole of the size $`r_{}`$ ). This wave function depends on the polarization of the virtual photon and it has been calculated previously in .
$`\mathrm{\Psi }_L^\gamma ^{}(z,r_{};Q^2)=Qz(1z)K_0(ar_{});`$ (2.16)
$`\mathrm{\Psi }_T^\gamma ^{}(z,r_{};Q^2)=iaK_1(ar_{}){\displaystyle \frac{\stackrel{}{𝐫}_{}}{r_{}}};`$ (2.17)
where $`a^2=z(1z)Q^2+m_q^2`$ and subscripts $`T`$ and $`L`$ denote the transverse and longitudinal polarizations of the photon, respectively.
19. In our calculations we only require the probabilty to find a quark-antiquark pair with the size $`r_{}`$ inside a virtual photon, namely
$`P^\gamma ^{}(z,r_{};Q^2)`$ $`=`$ $`{\displaystyle \frac{\alpha _{em}N_c}{2\pi ^2}}{\displaystyle \underset{f}{}}Z_f^2{\displaystyle \underset{\lambda _1,\lambda _2}{}}\{|\mathrm{\Psi }_T|^2+|\mathrm{\Psi }_L|^2\}`$
$`=`$ $`{\displaystyle \frac{\alpha _{em}N_c}{2\pi ^2}}{\displaystyle \underset{f}{}}Z_f^2\{(z^2+(1z)^2)a^2K_1^2(ar_{})+\mathrm{\hspace{0.17em}4}Q^2z^2(1z)^2K_0^2(ar_{})\}.`$
### 2.2 Shadowing corrections for penetration of $`𝐪\overline{𝐪}`$ pair through the target.
#### 2.2.1 General approach
The physics underlying our approach has been formulated and developed in Refs. . During its passage through the target, the distance $`r_{}`$ between a quark and an antiquark can vary by an amount $`\mathrm{\Delta }r_{}Rk_{}/E`$, where $`E`$ is the pair energy and $`R`$ is the size of the target. Since the quark’s transverse momentum $`k_{}\mathrm{\hspace{0.17em}\hspace{0.17em}1}/r_{}`$, the relation
$$\mathrm{\Delta }r_{}R\frac{k_{}}{E}R\frac{1}{r_{}E}r_{}$$
(2.19)
holds if
$$r^2s\mathrm{\hspace{0.17em}\hspace{0.17em}2}mR,$$
(2.20)
where $`s=W^2=2mE`$ with $`m`$ being the mass of the hadron.
Eq. (2.20) can be rewritten in terms of $`x_P`$, namely,
$$x_P\frac{2}{(1\beta )mR}.$$
(2.21)
From Eq. (2.21) it follows that $`r_{}`$ is a good degree of freedom for high energy scattering.
We can therefore write the total cross section for the interaction of the virtual photon with the target as follows:
$`\sigma _{tot}(\gamma ^{}+p)`$ $`=`$ $`{\displaystyle 𝑑zd^2r_{}P^\gamma ^{}(z,r_{};Q^2)\sigma _{dipole}(x_B,r_{})}`$ (2.22)
$`=`$ $`2{\displaystyle d^2b_t𝑑zd^2r_{}P^\gamma ^{}(z,r_{};Q^2)Ima_{dipole}^{el}(x_B,r_{};b_t)}.`$ (2.23)
The amplitude for diffractive production of a $`q\overline{q}`$ \- pair is equal to
$$a^{DD}(\gamma ^{}+pq+\overline{q})=d^2r_{}\mathrm{\Psi }^\gamma ^{}(z,r_{};Q^2)a_{dipole}^{el}(x_B,r_{};b_t)\mathrm{\Psi }^{q\overline{q}}(k_{},z,r_{}),$$
(2.24)
where $`\mathrm{\Psi }^{q\overline{q}}(k_{},z,r_{})`$ is the wave function of the quark-antiquark pair with fixed momentum $`k_{}`$ and fraction of energy $`z`$. To calculate the total cross section of the diffractive production we should integrate over all $`k_{}`$ and $`z`$. Using the completeness of the $`q\overline{q}`$ wave function one obtains
$$\sigma ^{DD}(\gamma ^{}+pq+\overline{q})=d^2b_t𝑑zd^2r_{}P^\gamma ^{}(z,r_{};Q^2)|a_{dipole}^{el}(x_B,r_{};b_t)|^2.$$
(2.25)
Utilizing the unitarity constraint we obtain a prediction for the ratio
$$\mathrm{}=\frac{\sigma ^{DD}}{\sigma _{tot}}=\frac{d^2b_t𝑑zd^2r_{}P^\gamma ^{}(z,r_{};Q^2)|a_{dipole}^{el}(x_B,r_{};b_t)|^2}{2d^2b_tdzd^2r_{}P^\gamma ^{}((z,r_{};Q^2)Ima_{dipole}^{el}(x_B,r_{};b_t)}.$$
(2.26)
Eq. (2.26) is a general prediction for the ratio $`\mathrm{}`$ . It shows that the diffractive dissociation and total cross sections are related through unitarity. However, Eq. (2.26) is too general to be used for pratical estimates.
Assuming that the dipole - proton amplitude is mainly imaginary at high energy the unitarity constraint of Eq. (2.4) has a general solution
$`a_{dipole}^{el}(x,r_{};b_t)=i\left(\mathrm{\hspace{0.17em}1}e^{\frac{\mathrm{\Omega }(x,r_{};b_t)}{2}}\right);`$ (2.27)
$`G_{dipole}^{in}(x,r_{};b_t)=\mathrm{\hspace{0.17em}\hspace{0.17em}1}e^{\mathrm{\Omega }(x,r_{};b_t)};`$ (2.28)
where $`\mathrm{\Omega }`$ is arbitrary real function.
Using Eq. (2.27) and Eq. (2.28), Eq. (2.26) reduces to
$$\mathrm{}=\frac{\sigma ^{DD}}{\sigma _{tot}}=\frac{d^2b_t𝑑zd^2r_{}P^\gamma ^{}(z,r_{};Q^2)\left(\mathrm{\hspace{0.17em}1}e^{\frac{\mathrm{\Omega }(x,r_{};b_t)}{2}}\right)^2}{2d^2b_t𝑑zd^2r_{}P^\gamma ^{}(z,r_{};Q^2)\left(\mathrm{\hspace{0.17em}1}e^{\frac{\mathrm{\Omega }(x,r_{};b_t)}{2}}\right)}.$$
(2.29)
The main goal of our microscopic approach based on QCD is to find $`\mathrm{\Omega }`$ .
#### 2.2.2 Mueller-Glauber approach
One way to get a more detailed picture of the interaction, is to consider the dipole - proton interaction in the Eikonal model, which is closely related to Mueller - Glauber approach .
The main assumption of this model is to identify the function $`\mathrm{\Omega }`$ in Eq. (2.29) with the exchange of the “hard” Pomeron ( gluon ladder ) given by Eq. (2.15) ( see Fig. 3 ). Since the gluon distribution given by the DGLAP evolution equations originates from inelastic processes of gluon emission, we assume an oversimplified structure of the final states in the Eikonal model, namely, it consists of only a proton and a quark-antiquark pair ( “elastic” scattering ), and an inelastic state with a large number of emitted gluons ($`N_G\mathrm{ln}(1/x)`$ . In particular, we neglect the rich structure of the diffraction dissociation processes and simplify them to the final state of $`p+q\overline{q}`$. For example, we neglect the diffractive production of an excited nucleon in DIS.
We will discuss the accuracy of our approach in the next subsection, where we expand our model to include an excitation of the target. Our accuracy is restricted by the assumption of only a quark - antiquark pair and a nucleon in the final state for diffractive processes. The rough estimate for the contribution of all excitations of the nucleon is
$$\frac{\sigma ^{DD}(\gamma ^{}+pq\overline{q}+N^{})}{\sigma ^{DD}(\gamma ^{}+pq\overline{q}+N)}=\sqrt{\frac{\sigma ^{DD}(p+pp+N^{})}{2\sigma _{el}(p+pp+p)}}\mathrm{\hspace{0.17em}\hspace{0.17em}0.2}÷0.3.$$
(2.30)
consequently, we have to consider the nucleon excitations or, at least, to discuss them in the HERA kinematic region.
Substituting $`\mathrm{\Omega }=\mathrm{\Omega }^P`$ ( see Eq. (2.15) ) in Eq. (2.29) we obtain the Kovchegov - McLerran formula , then for large $`Q^2`$ we have
$$_0^1𝑑zP_T^\gamma ^{}(z,r_{};Q^2)=\frac{4\alpha _{em}N_c}{3\pi ^2Q^2}\times \frac{1}{r_t^4}.$$
(2.31)
We can see from Eq. (2.29) and Eq. (2.31) that
$$\sigma ^{DD}R^2Q^2/<r_{}^2>;$$
$$\sigma _{tot}R^2Q^2/<r_{}^2>;$$
we can evaluate $`<r_{}^2>`$ from the condition
$$\mathrm{\Omega }_{dipole}^P(b_t=0)=\kappa _{dipole}^{DGLAP}(x,<r_{}>)=\mathrm{\hspace{0.17em}\hspace{0.17em}1}.$$
(2.32)
If we assume that $`xG(x,\frac{4}{r_{}^2})`$ depends rather smoothly on $`r_{}^2`$, and substitute $`\frac{4}{r_{}^2}=Q^2`$ as an estimate, we have
$$\frac{1}{<r_{}>}xG(x,Q^2).$$
(2.33)
Eq. (2.33) gives the ratio $`\mathrm{}`$ constant as well as $`\sigma ^{DD}R^2Q^2\times x_BG(x_B,Q^2)`$. This simple estimate, given in Ref. , illustrates that the SC lead to a constant ratio $`\mathrm{}`$, or, vice versa, the constant ratio $`\mathrm{}`$ can be a strong argument for substantial SC. To examine this point we calculate the ratio $`\mathrm{}`$ using the same assumption that $`\frac{4}{r_{}^2}=Q^2`$ in the gluon distribution. Using the result of the explicit calculation in Ref. we obtain:
$$\mathrm{}=\frac{\sigma ^{DD}}{\sigma _{tot}}=\mathrm{\hspace{0.17em}\hspace{0.17em}1}\frac{\mathrm{ln}(\kappa _{dipole})+C+(\mathrm{\hspace{0.17em}1}+\kappa _{dipole})E_1(\kappa _{dipole})+\mathrm{\hspace{0.17em}1}e^{\kappa _{dipole}}}{2[\mathrm{ln}(\frac{\kappa _{dipole}}{2})+C+(\mathrm{\hspace{0.17em}1}+\frac{\kappa _{dipole}}{2})E_1(\frac{\kappa _{dipole}}{2})+\mathrm{\hspace{0.17em}1}e^{\frac{\kappa _{dipole}}{2}}]}.$$
(2.34)
In Fig.4 the ratio $`\mathrm{}`$, given by Eq. (2.34), is plotted as a function of $`\kappa _{dipole}`$. One can see that at large $`\kappa _{dipole}`$ this ratio has a smooth dependence on $`\kappa _{dipole}`$. However, the values of $`\kappa _{dipole}`$ that we are dealing with are not very large ( $`\kappa _{dipole}\mathrm{\hspace{0.17em}\hspace{0.17em}1}`$, see Fig.5, where $`\kappa _{dipole}`$ is calculated using the GRV parameterization for the gluon distribution ). Fig.4-a shows that for $`\kappa _{dipole}=0.22`$ we cannot expect that Eq. (2.34) to yield a more or less constant ratio $`\mathrm{}`$.
However, we would like to draw the reader’s attention to the fact that the ratio differs from the small $`\kappa _{dipole}\mathrm{\hspace{0.17em}\hspace{0.17em}1}`$ limit where $`\mathrm{}\kappa _{dipole}`$.
These simple estimates indicate that the SC are essential, but they are still not sufficiently strong to use the asymptotic formulae. We are in the transition region from low density QCD, described by the DGLAP evolution, to high density QCD where we can use the quasi-classical gluon field approximation . In the transition region the Mueller-Glauber approach is a natural way to obtain reliable estimates of the SC. However, we need to study corrections to Eq. (2.29) more carefully. The most important among them, are the diffractive production of nucleon excitation ( see Eq. (2.30)) and of the $`q\overline{g}G`$ system, which gives the dominant contribution to the diffractive cross section .
#### 2.2.3 Diffractive production of nucleon excitations
A. General approach
We start with a trivial remark, that the nucleon excitations even in DIS are closely related to long distance processes and, therefore, to the “soft” interaction which cannot be determined in QCD. An alternate way of saying this is, to attribute these diffractive processes to nonperturbative QCD, for which at present we only have a phenomenological approach. Theoretically “soft” diffraction can be viewed as a typical quantum mechanical process which occurs, since the hadron states are not diagonal with respect to the strong interaction scattering matrix. In other words diffractive dissociation occurs, because even at high energy hadrons are not the correct degrees of freedom for the strong nonperturbative interaction. Unfortunately, we do not know the correct degrees of freedom and below we will discuss some models for them. We denote by $`n`$ the correct degree of freedom or the set of quantum numbers which characterizes the wave function $`\mathrm{\Psi }_n`$. These function $`\mathrm{\Psi }_n`$ are diagonal with respect to the strong interaction
$$A_{n,n^{}}=<\mathrm{\Psi }_n\left|𝐓\right|\mathrm{\Psi }_n^{}>=A_n\delta _{n,n^{}},$$
(2.35)
where parentheses denote all needed integrations and $`𝐓`$ is the scattering matrix.
Note, that only for amplitudes $`A_n`$ do we have the unitarity constraints in the form of Eq. (2.4), namely,
$$ImA_n^{el}(s,b_t)=|A_n^{el}(s,b_t)|^2+G_n^{in}(s,b_t)$$
(2.36)
which has solutions of Eq. (2.27) and Eq. (2.28) for mainly imaginary $`A_n`$ at high energies:
$`A_n^{el}=i\{\mathrm{\hspace{0.17em}\hspace{0.17em}1}e^{\frac{\mathrm{\Omega }_n(s,b_t)}{2}}\};`$ (2.37)
$`G_n^{in}=\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}1}e^{\mathrm{\Omega }_n(s,b_t)}.`$ (2.38)
The wave function of a hadron is
$$\mathrm{\Psi }_{hadron}=\underset{n=1}{\overset{\mathrm{}}{}}\alpha _n\mathrm{\Psi }_n.$$
(2.39)
For a dipole - hadron interaction the wave function is equal to $`\mathrm{\Psi }_{dipole}(r_{})\times \mathrm{\Psi }_{hadron}`$ before collision. After collision the scattering matrix $`𝐓`$ leads to a new wave function,namely
$$\mathrm{\Psi }_{final}=\mathrm{\Psi }_{dipole}(r_{})\times \underset{n=1}{\overset{\mathrm{}}{}}\alpha _nA_n\mathrm{\Psi }_n.$$
(2.40)
From Eq. (2.40) we obtain the elastic amplitude
$$a_{dipole}^{el}=<\mathrm{\Psi }_{final}|\mathrm{\Psi }_{dipole}(r_{})\times \mathrm{\Psi }_{hadron}>=\underset{n=1}{\overset{\mathrm{}}{}}\alpha _n^2A_n(s,b_t),$$
(2.41)
while for the total cross section of the diffractive nucleon excitations we have
$`\sigma _N^{}^{DD}`$ $`=`$ $`<\mathrm{\Psi }_{final}|\mathrm{\Psi }_{final}>^2<\mathrm{\Psi }_{final}|\mathrm{\Psi }_{dipole}(r_{})\times \mathrm{\Psi }_{hadron}>^2;`$ (2.42)
$`=`$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\alpha _n^2|A_n(s,b_t)|^2\{{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\alpha _n^2A_n(s,b_t)\}^2.`$ (2.43)
Therefore, using Eq. (2.31) instead of Eq. (2.29), we obtain a generalized formula
$$\mathrm{}=\frac{\sigma ^{DD}}{\sigma _{tot}}=\frac{_{n=1}^{\mathrm{}}\alpha _n^2d^2b_t\frac{dr_{}^2}{r_{}^4}\{\mathrm{\hspace{0.17em}1}e^{\frac{\mathrm{\Omega }_n^P(r_{}^2,y;b_t)}{2}}\}^2}{2_{n=1}^{\mathrm{}}\alpha _n^2d^2b_t\frac{dr_{}^2}{r_{}^4}\{\mathrm{\hspace{0.17em}1}e^{\frac{\mathrm{\Omega }_n^P(r_{}^2,y;b_t)}{2}}\}},$$
(2.44)
One can see that this generalized formula has all the attractive features of Eq. (2.29), and at small $`x`$ the ratio tends to $`1/2`$, since the normalization constraint
$$\underset{n=1}{\overset{\mathrm{}}{}}\alpha _n^2=\mathrm{\hspace{0.17em}\hspace{0.17em}1}$$
(2.45)
We can estimate $`\mathrm{\Omega }_n^P`$ using the same Eq. (2.15) here
$$\mathrm{\Omega }_n=\kappa _{dipole}^ne^{\frac{b_t^2}{R_n^2}},$$
(2.46)
where
$$\kappa _{dipole}^n=\frac{\pi ^2\alpha _S}{3\pi R_n^2}r_{}^2x_PG_n(x_P,\frac{4}{r_{}^2}).$$
(2.47)
Therefore, the main difficulty with Eq. (2.46) and Eq. (2.47) is to determine $`R_n`$ and $`xG_n(x,\frac{4}{r_{}^2})`$. We only have direct experimental information for the proton radius and the gluon distribution in the proton. Unfortunately, we cannot evaluate Eq. (2.44) without developing some model for the diffractive excitations.
B. Two channel model for diffractive nucleon excitations.
The main idea of this model is to replace the many final states of diffractively produced hadrons by one state ( effective hadron ). In this case the general Eq. (2.39) reduces to the simple form
$$\mathrm{\Psi }_{hadron}=\alpha _1\mathrm{\Psi }_1+\alpha _2\mathrm{\Psi }_2,$$
(2.48)
with the condition $`\alpha _1^2+\alpha _2^2=\mathrm{\hspace{0.17em}\hspace{0.17em}1}`$ from Eq. (2.45).
The wave function of the produced effective hadron is equal to
$$\mathrm{\Psi }_D=\alpha _2\mathrm{\Psi }_1+\alpha _1\mathrm{\Psi }_2,$$
(2.49)
which is orthogonal to $`\mathrm{\Psi }_{hadron}`$.
Eq. (2.44) can be rewritten in the form
$`\mathrm{}`$ $`=`$ $`{\displaystyle \frac{\sigma ^{DD}}{\sigma _{tot}}}=`$
$`=`$ $`{\displaystyle \frac{d^2b_t\frac{dr_{}^2}{r_{}^4}\left(\alpha _1^2\{\mathrm{\hspace{0.17em}1}e^{\frac{\mathrm{\Omega }_1^P(r_{}^2,y;b_t)}{2}}\}^2+\alpha _2^2\{\mathrm{\hspace{0.17em}1}e^{\frac{\mathrm{\Omega }_2^P(r_{}^2,y;b_t)}{2}}\}^2\right)}{2d^2b_t\frac{dr_{}^2}{r_{}^4}\left(\{\mathrm{\hspace{0.17em}1}e^{\frac{\mathrm{\Omega }_1^P(r_{}^2,y;b_t)}{2}}\}+\alpha _2^2e^{\frac{\mathrm{\Omega }_1^P(r_{}^2,y;b_t)}{2}}\{\mathrm{\hspace{0.17em}1}e^{\frac{\mathrm{\Delta }\mathrm{\Omega }^P(r_{}^2,y;b_t)}{2}}\}\right)}};`$
with $`\mathrm{\Delta }\mathrm{\Omega }=\mathrm{\Omega }_2\mathrm{\Omega }_1`$.
For $`\mathrm{\Omega }_1`$ and $`\mathrm{\Delta }\mathrm{\Omega }`$ we use the parameterization given by Eq. (2.46) and Eq. (2.47), and the following experimental data and phenomenological observations:
1. Data for single diffraction in proton - proton collisions lead to $`\alpha _2^2\mathrm{\hspace{0.17em}\hspace{0.17em}0.2}`$ Ref. ;
2. Data for J/$`\mathrm{\Psi }`$ photoproduction at HERA ( see Fig.6a ) show that $`t`$-dependance is quite different for elastic and inelastic photoproduction. The measured values are $`B_{el}=\mathrm{\hspace{0.17em}\hspace{0.17em}4}GeV^2`$ while $`B_{in}=\mathrm{\hspace{0.17em}\hspace{0.17em}1.66}GeV^2`$. Therefore, we take $`R_1^2=2B_{el}=8GeV^2`$ in $`\mathrm{\Omega }_1`$ and $`R_D^2=2B_{in}=3.32GeV^2`$, where $`R_D^2`$ is the radius in the exponential parameterization for $`\mathrm{\Delta }\mathrm{\Omega }`$ ;
3. The energy behaviour of diffractive J/$`\mathrm{\Psi }`$ photoproduction shows that we can consider this process as a typical “hard” process which occurs at short distances ;
4. We use experimental evidence that the cross section for elastic and inelastic diffractive J/$`\mathrm{\Psi }`$ photoproduction are equal. From this fact we can conclude that
$$\alpha _2(1\alpha _2)(\mathrm{\Delta }\mathrm{\Omega })(b_t=0)^2R_D^2=(1\alpha _2)^2\mathrm{\Omega }_1^2(b_t=0)R_1^2,$$
(2.51)
which gives
$$\mathrm{\Delta }\mathrm{\Omega }(b_t=0)=\sqrt{\frac{(1\alpha _2^2)}{\alpha _2^2}}\frac{R_1}{R_D}\mathrm{\Omega }_1(b_t=0).$$
(2.52)
Substituting in Eq. (2.52) we find $`\mathrm{\Delta }\kappa _{dipole}\mathrm{\hspace{0.17em}\hspace{0.17em}3.4}\kappa _{dipole}^1`$.
In Fig. 7 we display the ratio $`R`$ as a function of $`\kappa _{dipole}`$. Comparing this figure with Fig. 4a we can conclude that in the two channel model, R increases more or less at the same rate as in the Eikonal model, but the value of the ratio depends crucially on the model for the diffractive excitation of the nucleon. Fig. 7b shows the contamination of the total diffractive cross section by the nucleon excitations. One can conclude that the two channel model gives a small fraction of the excitation cross section at sufficiently large $`\kappa _{dipole}`$. ExperimentallyWe thank Henry Kowalski for discussing this data and many problems stimulated by these data with us,$`\sigma _{excitation}^{DD}/\sigma _{elastic}^{DD}`$ = $`35\pm \mathrm{\hspace{0.17em}15}\%`$. Fig.7b supports a low value for this ratio, but this question should be reconsidered after taking into account the diffractive production of the $`q\overline{q}G`$ system.
C. Diffractive production in the Additive Quark Model
As we have mentioned, the main problem of dealing with the “soft” high energy interaction, is to find the correct degrees of freedom, and to incorporate them in the general formalism for high energy scattering. The two channel model gives an estimate of the importance of proton excitation processes, but it is too phenomenological to be instructive and not totally reliable. Here, we consider the nucleon excitation in the Additive Quark Model . In this model the correct degrees of freedom at high energies are the constituent quarks. In spite of a certain naivity this model has not been abandoned and it is included in the standard Donnachie - Landshoff Pomeron approach for “soft” processes at high energies.
Inherent in this model is the assumption, that the $`\gamma ^{}`$ \- constituent quark interactions dominate, while other interactions e.g. the interaction of $`\gamma ^{}`$ with two constituent quarks, are suppressed by factor $`r_Q^2/R_N^2`$ , here $`r_Q`$ is the size of the constituent quark and $`R_N`$ is the radius of the proton. It is obvious that in the AQM we have the same Eq. (2.29) ( or Eq. (2.26) ) where $`\mathrm{\Omega }=\mathrm{\Omega }_Q^P`$ describes the interaction of a colour dipole with the constituent quark. In the AQM the gluon distribution of the constituent quark is equal to $`xG_Q(x,Q^2)=\frac{1}{3}xG_N(x,Q^2)`$ and, therefore
$$\mathrm{\Omega }_Q^P(x,r_{};b_t)=\kappa _{dipole}^Qe^{\frac{b_t^2}{r_Q^2}}=\frac{\pi ^2\alpha _S}{9\pi r_Q^2}r_{}^2xG_N(x,\frac{4}{r_{}^2})e^{\frac{b_t^2}{r_Q^2}}.$$
(2.53)
Consequently, we find $`r_Q`$ using the same AQM to describe the double parton cross section measured by the CDF (see Fig. 7b ). The CDF collaboration has measured the inclusive cross section for the production of two “hard” pairs of jets, with large and almost compensating transverse momenta in each pair, and with similar values of rapidity. Such processes cannot occur in a one parton shower, and only originate from two parton shower interactions as shown in Fig.7b.
The double parton cross section can be written in the form
$$\sigma _{DP}=m\frac{\sigma _{inel}(2jets)\sigma _{inel}(2jets)}{2\sigma _{eff}}$$
(2.54)
where factor $`m`$ is equal to 2 for different pairs of jets, and to 1 for identical pairs. The experimental value of $`\sigma _{eff}=\mathrm{\hspace{0.17em}\hspace{0.17em}14.5}\pm \mathrm{\hspace{0.17em}\hspace{0.17em}1.7}\pm \mathrm{\hspace{0.17em}\hspace{0.17em}2.3}mb`$.
In the AQM (see Fig.7b ) $`\sigma _{eff}`$ can be easily calculated and it is equal to
$$\sigma _{eff}=\mathrm{\hspace{0.17em}\hspace{0.17em}9}\times \mathrm{\hspace{0.17em}2}\pi r_Q^2,$$
(2.55)
where factor 9 reflects the quark counting and $`2\pi r_Q^2`$ comes from the integration over $`b_t`$. Comparing Eq. (2.55) with the experimental value of $`\sigma _{eff}`$ we obtain $`r_Q^2=\mathrm{\hspace{0.17em}\hspace{0.17em}0.66}\pm \mathrm{\hspace{0.17em}0.16}GeV^2`$.
Substituting this result in Eq. (2.53) we find $`\kappa _{dipole}^Q\mathrm{\hspace{0.17em}\hspace{0.17em}5}\kappa _{dipole}^N`$ with the same $`x`$ and $`r_{}`$ dependance. In Fig. 8 one can see the prediction for the ratio $`\sigma ^{DD}/\sigma _{tot}`$. This approach gives the ratio which depends smoothly on $`\kappa _{dipole}`$ for $`\kappa _{dipole}>\mathrm{\hspace{0.17em}0.75}`$.
Comparing Fig.8 and Fig.4a we can conclude that the diffractive production of the nucleon excitations gives about 30% of the total diffractive cross section at $`\kappa _{dipole}\mathrm{\hspace{0.17em}0.2}0.6`$ and becomes very small at $`\kappa _{dipole}\mathrm{\hspace{0.17em}1.5}2`$.
#### 2.2.4 AGK cutting rules and diffractive production
In this section we derive Eq. (2.29) in the Mueller-Glauber approach exploiting the AGK cutting rules . This derivation is more complicated than the previous derivation which was based directly on the $`s`$-channel unitarity constraints. As we intend using the AGK cutting rules for calculating the diffraction production of the $`q\overline{q}G`$ system, we think it instructive to start with a simple example.
The AGK cutting rules provide a prescription of how to calculate the cross sections for the processes with different multiplicities of the produced particles, if one knows the structure of the Pomeron exchange, and the expression for the total cross section in terms of multi Pomeron exchanges . In the Mueller-Glauber approach we have such an expression for the total colour dipole - proton cross section, namely, (see Fig.9)
$$\sigma _{dipole}=\mathrm{\hspace{0.17em}\hspace{0.17em}2}d^2b_t\{\mathrm{\hspace{0.17em}1}e^{\frac{\mathrm{\Omega }^P(x,r_{};b_t)}{2}}\},$$
(2.56)
where $`\mathrm{\Omega }^P`$ is given by Eq. (2.15). We first need to define the Pomeron structure i.e. to specify the kind of inelastic processes that are described by $`\mathrm{\Omega }^P`$. From Eq. (2.15) we see that $`\mathrm{\Omega }^P`$ describes the inelastic processes with large average multiplicity ( $`<n_P>`$ ) of produced partons, since it is related to the DGLAP parton cascade. For example, at large $`Q^2`$ and low $`x`$, $`<n_P>=\mathrm{\hspace{0.17em}\hspace{0.17em}2}\sqrt{\overline{\alpha }_S\mathrm{ln}Q^2\mathrm{ln}(1/x)}\mathrm{\hspace{0.17em}\hspace{0.17em}1}`$. The AGK cutting rules allow us to calculate the cross sections for the processes with average multiplicities 2$`<n_P>`$, 3 $`<n_P>`$ and so on, as well as the process for the diffractive dissociation with multiplicity much smaller that $`<n_P>`$.
Eq. (2.56) can be rewritten in the form
$$\sigma _{dipole}=\mathrm{\hspace{0.17em}\hspace{0.17em}2}d^2b_t\underset{n=1}{\overset{\mathrm{}}{}}C_n(1)^{n+1}\left(\frac{\mathrm{\Omega }^P}{2}\right)^n,$$
(2.57)
where each term corresponds to the exchange of $`n`$ Pomerons. The AGK rules are
$`\sigma _{dipole}^n(k<n_P>)`$ $`=`$ $`{\displaystyle d^2b_tC_n(1)^{nk}\frac{n!}{(nk)!k!}\left(\mathrm{\Omega }^P\right)^n};`$ (2.58)
$`\sigma _{dipole}^n(DD)`$ $`=`$ $`{\displaystyle d^2b_tC_n(1)^n\{\left(\mathrm{\Omega }^P\right)^n\mathrm{\hspace{0.17em}\hspace{0.17em}2}\left(\frac{\mathrm{\Omega }^P}{2}\right)^n\}};`$ (2.59)
where $`\sigma _{dipole}^n(k<n_P>)`$ and $`\sigma _{dipole}^n(DD)`$ are the contributions of $`n`$-Pomeron exchange to the cross section for the process with average multiplicity $`k<n_P>`$, and to the cross section for the diffractive dissociation processes with small multiplicity.
Summing over $`n`$ in Eq. (2.59) we obtain
$`\sigma _{dipole}^{DD}(r_{})`$ $`=`$ $`{\displaystyle d^2b_t\underset{n=1}{\overset{\mathrm{}}{}}C_n(1)^n\{\left(\mathrm{\Omega }^P\right)^n\mathrm{\hspace{0.17em}\hspace{0.17em}2}\left(\frac{\mathrm{\Omega }^P}{2}\right)^n\}};`$ (2.60)
$`=`$ $`{\displaystyle d^2b_t\left(\mathrm{\hspace{0.17em}2}\{\mathrm{\hspace{0.17em}1}e^{\frac{\mathrm{\Omega }^P}{2}}\}\{\mathrm{\hspace{0.17em}1}e^{\mathrm{\Omega }^P}\}\right)};`$ (2.61)
$`=`$ $`{\displaystyle d^2b_t\left(\mathrm{\hspace{0.17em}1}e^{\frac{\mathrm{\Omega }^P}{2}}\right)^2}.`$ (2.62)
One can see that Eq. (2.62) is just the same equation for diffractive production that we have obtained from unitarity ( see Eq. (2.29) ). From Eq. (2.58) we can find the cross section for the production of $`k`$ \- parton showers which is equal to
$$\sigma _{dipole}(k<n_P>)=d^2b_t\frac{\left(\mathrm{\Omega }^P\right)^k}{k!}e^{\mathrm{\Omega }^P}.$$
(2.63)
Eq. (2.63) will be very useful below, when we discuss the diffraction production of the $`q\overline{q}G`$ system.
### 2.3 Cross section for $`𝐪\overline{𝐪}𝐆`$ diffractive production with SC
#### 2.3.1 First correction to the Mueller - Glauber formula for the total cross section
As shown in Fig.9, the Eikonal approach takes into account only rescatterings of the fastest colour dipole. In this subsection we will extend the formalism so as also to include the rescatterings with the target of the two fastest color dipoles: the initial colour dipole and the fastest gluon in Fig.9. Our goal is to include all diagrams of the type shown in Fig.10. Dashed lines in Fig.10 indicate the diffraction dissociation cuts. These diagrams include the diffractive production of $`q\overline{q}G`$ system as well as a quark - antiquark pair.
Since Eq. (2.27) and Eq. (2.28) give the general solution to the unitarity constraint, our problem is to find an expression for $`\mathrm{\Omega }`$ which will be more general that Eq. (2.15) with $`\kappa ^{DGLAP}`$ defined by Eq. (2.9).
The natural generalization of $`\kappa ^{DGLAP}`$ is to substitute in Eq. (2.9) the Mueller-Glauber formula for the gluon structure function , namely<sup>\**</sup><sup>\**</sup>\**It should be stressed that $`xG(x,\frac{4}{r_t^2};b_t)`$ is introduced in a such way that $`d^2b_txG(x,\frac{4}{r_t^2};b_t)=xG(x,\frac{4}{r_t^2})`$.
$$\frac{1}{\pi R^2}xG^{DGLAP}(x,\frac{4}{r_{}^2})e^{\frac{b_t^2}{R^2}}xG^{MG}(x,\frac{4}{r_{}^2};b_t)=$$
(2.64)
$$\frac{4}{\pi ^3}_x^1\frac{dx^{}}{x^{}}_{r_{}^2}^{\mathrm{}}\frac{dr_{}^2}{r_{}^4}\mathrm{\hspace{0.17em}2}\{\mathrm{\hspace{0.17em}1}e^{\frac{\mathrm{\Omega }_G^P(x^{},r_{}^{};b_t)}{2}}\},$$
where
$$\mathrm{\Omega }_G^P=\frac{3\pi ^2\alpha _S}{4\pi R^2}r_{}^2G^{DGLAP}(x,\frac{4}{r_{}^2})e^{\frac{b_t^2}{R^2}}.$$
(2.65)
Eq. (2.64) takes into account the rescattering of the gluon with the target in the Eikonal approach ( see Fig.10), however, the question arises of why we need to only include gluon rescattering for the $`q\overline{q}G`$ system. To understand the physics of Eq. (2.64) it is advantageous to consider the equation that describes the rescattering of all partons. Kovchegov using two principle ideas suggested by A. Mueller proved that the GLR nonlinear equation is able to describe such rescatterings. The principles are:
1. The QCD interaction at high energy does not change the transverse size of interacting colour dipoles, and thus they can be considered as the correct degrees of freedom at high energies;
2. The process of interaction of a dipole with the target has two clear stages:
1. The transition of the dipole into two dipoles, the probability for this is given by
$$|\mathrm{\Psi }(𝐱_{\mathrm{𝟎𝟏}}𝐱_{\mathrm{𝟎𝟐}}+𝐱_{\mathrm{𝟏𝟐}})|^2=\frac{1}{z}\frac{𝐱_{\mathrm{𝟎𝟏}}^\mathrm{𝟐}}{𝐱_{\mathrm{𝟎𝟐}}^\mathrm{𝟐}𝐱_{\mathrm{𝟏𝟐}}^\mathrm{𝟐}},$$
(2.66)
where $`𝐱_{\mathrm{𝐢𝐤}}`$ denotes the size of the dipoles, and $`z`$ the fraction of energy of the initial dipole that the final dipole carries;
2. The interaction of each dipole with the target has an amplitude $`a^{el}(𝐱,b_t,y=\mathrm{ln}(1/x))`$.
The equation is illustrated in Fig.11 and it has the following analytic form:
$$\frac{da_{dipole}^{el}(𝐱_{\mathrm{𝟎𝟏}},b_t,y)}{dy}=\frac{2C_F\alpha _S}{\pi }\mathrm{ln}\left(\frac{𝐱_{\mathrm{𝟎𝟏}}^\mathrm{𝟐}}{\rho ^2}\right)a_{dipole}^{el}(𝐱,b_t,y)+\frac{C_F\alpha _S}{\pi ^2}_\rho d^2𝐱_\mathrm{𝟐}\frac{𝐱_{\mathrm{𝟎𝟏}}^\mathrm{𝟐}}{𝐱_{\mathrm{𝟎𝟐}}^\mathrm{𝟐}𝐱_{\mathrm{𝟏𝟐}}^\mathrm{𝟐}}$$
(2.67)
$$\{2a_{dipole}^{el}(𝐱_{\mathrm{𝟎𝟐}},𝐛_𝐭\frac{\mathrm{𝟏}}{\mathrm{𝟐}}𝐱_{\mathrm{𝟏𝟐}},y)a_{dipole}^{el}(𝐱_{\mathrm{𝟎𝟐}},𝐛_𝐭\frac{\mathrm{𝟏}}{\mathrm{𝟐}}𝐱_{\mathrm{𝟏𝟐}},y)a_{dipole}^{el}(𝐱_{\mathrm{𝟏𝟐}},𝐛_𝐭\frac{\mathrm{𝟏}}{\mathrm{𝟐}}𝐱_{\mathrm{𝟎𝟐}},y)\}.$$
The first term on the r.h.s. of the equation gives the contribution of virtual corrections, which appear in the equation due to the normalization of the partonic wave function of the fast colour dipole ( see Ref. ). The second term describes the decay of the colour dipole of size $`𝐱_{\mathrm{𝟎𝟏}}`$ into two dipoles of sizes $`𝐱_{\mathrm{𝟎𝟐}}`$ and $`𝐱_{\mathrm{𝟏𝟐}}`$, and their interactions with the target in the impulse approximation ( notice factor 2 in Eq. (2.67) ). The third term corresponds to the simultaneous interaction of two produced colour dipoles with the target and describes the Glauber-type corrections for scattering of these dipoles.
The initial condition for this equation is given by Eq. (2.27) with $`\mathrm{\Omega }=\mathrm{\Omega }^P`$ from Eq. (2.15) at $`x=x_0`$.
In DIS the dominant contribution comes from the decay of a small dipole into two large dipoles. Therefore, we can reduce the kernel of Eq. (2.67) to
$$_\rho d^2𝐱_\mathrm{𝟐}\frac{𝐱_{\mathrm{𝟎𝟏}}^\mathrm{𝟐}}{𝐱_{\mathrm{𝟎𝟐}}^\mathrm{𝟐}𝐱_{\mathrm{𝟏𝟐}}^\mathrm{𝟐}}𝐱_{\mathrm{𝟎𝟏}}^\mathrm{𝟐}\pi _{𝐱_{\mathrm{𝟎𝟏}}^\mathrm{𝟐}}^{\frac{1}{\mathrm{\Lambda }_{QCD}^2}}\frac{d𝐱_{\mathrm{𝟎𝟐}}^\mathrm{𝟐}}{(𝐱_{\mathrm{𝟎𝟐}}^\mathrm{𝟐})^2}.$$
(2.68)
We make the first iteration of Eq. (2.67), by substituting the Mueller - Glauber formula for color dipole rescattering ( see Eq. (2.27) with $`\mathrm{\Omega }=\mathrm{\Omega }^P`$ from Eq. (2.15) ), and obtain
$`a_{dipole}^{el}(firstiteration)`$ $`=`$ $`C_F\alpha _S𝐱_{\mathrm{𝟎𝟏}}^\mathrm{𝟐}\{2(1e^{\frac{\mathrm{\Omega }^P}{2}})(1e^{\frac{\mathrm{\Omega }^P}{2}})^2\}`$
$`=`$ $`C_F\alpha _S𝐱_{\mathrm{𝟎𝟏}}^\mathrm{𝟐}\{\mathrm{\hspace{0.17em}1}e^{\frac{2\mathrm{\Omega }^P}{2}}\}.`$
Eq. (2.3.1) gives the Mueller - Glauber formula for the gluon structure function of Eq. (2.64) <sup>††</sup><sup>††</sup>†† Eq. (2.67) is written for large number of colours $`N_c\mathrm{\hspace{0.17em}1}`$. For finite $`N_c`$, $`2\mathrm{\Omega }^P`$ in Eq. (2.3.1) should be replaced by $`2\mathrm{\Omega }^P\frac{9}{4}\mathrm{\Omega }^P`$. . Note, that two assumptions have been made in deriving Eq. (2.3.1): (i) $`b_t𝐱_{\mathrm{𝟏𝟐}}`$ or $`𝐱_{\mathrm{𝟎𝟐}}`$ and (ii) we neglected the first term in Eq. (2.67). Both approximations hold in the so called double log approximation of pQCD .
Eq. (2.64) describes the passage of the $`q\overline{q}G`$ system through the target, as it corresponds to the interaction of two colour dipoles, which is our $`q\overline{q}G`$ system, but not the gluon. However, in the parton cascade in the DIS kinematic region, the gluon always corresponds to two colour dipoles of the same size. By changing $`\mathrm{\Omega }^P\mathrm{\Omega }^{MG}`$ we take into account the fact that initial quark - antiquark pair can fluctuate in the $`q\overline{q}G`$ system many times during the passage through the target. This is illustrated in Fig.10.
Finally, we obtain the following formula for the total cross section:
$$\sigma _{dipole}^{(1)}=d^2b_t2\{\mathrm{\hspace{0.17em}1}e^{\frac{\mathrm{\Omega }^{MG}(x,r_{};b_t)}{2}}\},$$
(2.70)
where
$$\mathrm{\Omega }^{MG}(x,r_{};b_t)=\frac{\pi ^2\alpha _S}{3}r_{}^2xG^{MG}(x,r_{}^2;b_t)$$
(2.71)
with $`xG^{MG}(x,r_{}^2;b_t)`$ given in Eq. (2.64).
#### 2.3.2 Cross section for diffractive production
We would like to obtain the cross section for the diffractive production of both $`q\overline{q}`$ and $`q\overline{q}G`$ final states using the AGK cutting rules. In Fig.10 one can see which cuts in the total cross section are related to the diffracrive processes. They are shown in Fig.10 by dashed lines. First, we have to generalized the AGK cutting rules, since in Eq. (2.58) and Eq. (2.59) we have used the property of one Pomeron exchange i.e. one gluon “ladder” exchange , namely
$$2\mathrm{\Omega }^P=G_{in},$$
(2.72)
where $`G_{in}`$ stands for the inelastic cross section with large multiplicity $`<n_P>`$ ( see Fig.9 ). In Eq. (2.70) $`\mathrm{\Omega }^{MG}`$ itself has a more complicated structure , which can be recovered using the AGK rules of Eq. (2.58) and Eq. (2.59). Eq. (2.59) for $`\mathrm{\Omega }^{MG}`$ gives
$$\sigma _{MG}^{DD}(b_t)=\frac{N_c\alpha _S}{\pi }r_t^2_x^1\frac{dx^{}}{x^{}}_{r_t^2}\frac{dr_{}^2}{r_{}^4}\left(\mathrm{\hspace{0.17em}1}e^{\frac{\mathrm{\Omega }_G^P}{2}}\right)^2,$$
(2.73)
with $`\mathrm{\Omega }_G^P`$ from Eq. (2.65).
We denote by $`\mathrm{\Omega }_{cut}^{MG}`$ the contribution of all processes given by the AGK rules for $`\mathrm{\Omega }^{MG}`$. Using this notation the simple generalization of the AGK cutting rules gives the contributions of different processes to the total cross section of Eq. (2.70);
$`\sigma _{dipole}^{(1)}`$ $`=`$ $`{\displaystyle d^2b_t\underset{n=1}{\overset{\mathrm{}}{}}C_n(1)^{n+1}\left(\frac{\mathrm{\Omega }^{MG}}{2}\right)^n};`$ (2.74)
$`\sigma _{dipole}^{(1)}(n;k)`$ $`=`$ $`{\displaystyle d^2b_tC_n(1)^{nk}\frac{n!}{(nk)!k!}(\mathrm{\Omega }^{MG})^{nk}(\mathrm{\Omega }_{cut}^{MG})^k};`$ (2.75)
$`\sigma _{dipole}^{(1)}(n;0)`$ $`=`$ $`{\displaystyle d^2b_tC_n(1)^n\{\left(\mathrm{\Omega }^{MG}\right)^n\mathrm{\hspace{0.17em}\hspace{0.17em}2}\left(\frac{\mathrm{\Omega }^{MG}}{2}\right)^n\}};`$ (2.76)
where we denote by $`\sigma _{dipole}^{(1)}(n;k)`$ the contribution of $`k`$ \- cut $`\mathrm{\Omega }^{MG}`$ to the $`n`$-th term of Eq. (2.74), $`\sigma _{dipole}^{(1)}(n;0)`$ is the cross section for the diffraction production of a quark - antiquark pair for the $`n`$-th term of Eq. (2.74) ( see Fig.10). The structure of the inelastic processes for each term $`\sigma _{dipole}^{(1)}(n;k)`$ is rather complicated but is well defined by the AGK rules for $`\mathrm{\Omega }^{MG}`$. However, we need only take the cross section of the diffractive process from each $`\mathrm{\Omega }_{cut}^{MG}`$ or, in other words, we should replace $`\mathrm{\Omega }_{cut}^{MG}`$ by $`\sigma _{MG}^{DD}(b_t)`$ in Eq. (2.75). Performing the summation over $`n`$ and $`k`$ we obtain:
$$\sigma _{dipole}^{DD}(q\overline{q}q\overline{q}G)=d^2b_te^{\mathrm{\Omega }^{MG}(x,r_{};b_t)}\left(e^{\sigma _{MG}^{DD}(b_t)}\mathrm{\hspace{0.17em}\hspace{0.17em}1}\right),$$
(2.77)
where $`\mathrm{\Omega }^{MG}`$ defined in Eq. (2.71). Eq. (2.77) gives the contribution to the diffractive dissociation cross section of emission of gluons and $`k`$ is the number of gluons that we summed over. Eq. (2.76) leads to
$$\sigma _{dipole}^{DD}(q\overline{q}q\overline{q})=d^2b_t\left(\mathrm{\hspace{0.17em}1}e^{\frac{\mathrm{\Omega }^{MG}(x,r_{};b_t)}{2}}\right)^2.$$
(2.78)
In this paper we are only interested in production of one extra gluon which corresponds to Eq. (2.75) with $`k`$ = 1. Therefore
$$\sigma _{dipole}^{DD}(q\overline{q}q\overline{q}G)=d^2b_te^{\mathrm{\Omega }^{MG}(x,r_{};b_t)}\sigma _{MG}^{DD}(b_t),$$
(2.79)
which we will use for our numerical calculation.
Finally, collecting Eq. (2.70), Eq. (2.78), and Eq. (2.79) we obtain the generalization of Eq. (2.29), which takes into account the diffractive production of both a quark - antiquark pair, and a quark - antiquark pair plus an extra gluon final states:
$$\mathrm{}=\frac{\sigma ^{DD}}{\sigma _{tot}};$$
(2.80)
$`\sigma ^{DD}={\displaystyle }d^2b_t{\displaystyle }d^2r_{}P^\gamma ^{}(z,r_{};Q^2)\times `$
$`\left(\{\mathrm{\hspace{0.17em}1}e^{\frac{\mathrm{\Omega }^{MG}(x,r_{};b_t)}{2}}\}^2+e^{\mathrm{\Omega }^{MG}(x,r_{};b_t)}r_{}^2{\displaystyle \frac{2\alpha _S(r_{})}{3\pi }}{\displaystyle _x^1}{\displaystyle \frac{dx^{}}{x^{}}}{\displaystyle _{r_{}^2}^{\mathrm{}}}{\displaystyle \frac{dr_{}^2}{r_{}^4}}\{\mathrm{\hspace{0.17em}1}e^{\frac{\mathrm{\Omega }_G^P(x^{},r_{}^{};b_t)}{2}}\}^2\right);`$
$$\sigma _{tot}=\mathrm{\hspace{0.17em}\hspace{0.17em}2}d^2b_td^2r_{}P^\gamma ^{}(z,r_{};Q^2)\left(\mathrm{\hspace{0.17em}\hspace{0.17em}1}e^{\frac{\mathrm{\Omega }^{MG}(x,r_{};b_t)}{2}}\right).$$
We note that Eq. (2.80) was derived in the double log approximation to the DGLAP evolution equations, in which the colour dipoles in the produced $`q\overline{q}G`$ system are much larger than the initial quark - antiquark dipole.
The factor $`e^{\mathrm{\Omega }^{MG}(x,r_{};b_t)}`$ in the second term in Eq. (2.80), leads to the suppression of diffractive production of the the $`q\overline{q}G`$ system. Therefore, in the asymptotic limit at large values of $`\kappa _{dipole}`$, only elastic rescattering of the quark - antiquark pair survives, leading to the value of the ratio $`\sigma ^{DD}/\sigma _{tot}\mathrm{\hspace{0.17em}\hspace{0.17em}1}/2`$. It turns out that this factor is important for numrerical calculations in the HERA kinematic region ( see the next section). It was omitted in Ref.
## 3 Numerical calculation for $`\sigma ^{\mathrm{𝐃𝐃}}/\sigma _{\mathrm{𝐭𝐨𝐭}}`$
In this section we present the numerical result for the ratio of the total diffractive dissociation cross section to the total cross section in DIS. We postpone to the next section the consideration of the influence of the experimental mass cutoff on this ratio.
The main parameters that determine the value of $`\kappa _{dipole}`$, are the value of $`R^2`$ and the value of the gluon distribution. We choose $`R^2=\mathrm{\hspace{0.17em}10}GeV^2`$ since it is the value which is obtained from “soft” high energy phenomenology and is in agreement with HERA data on J/$`\mathrm{\Psi }`$ photoproduction . For $`xG(x,Q^2)`$ we use the GRV’94 parameterization and the leading order solution of the DGLAP evolution equation . We have two reasons for this choice:
1. Our goal in this paper is to understand the influence of the SC on the energy behaviour of the ratio $`\sigma ^{DD}/\sigma _{tot}`$. However, experience shows that we are able to describe almost all HERA data by changing the initial conditions of the DGLAP evolution equations. Unfortunately, we have no theoretical restrictions on this input for any of the parameterizations on the market. On the other hand, we know theoretically that the SC corrections work in a such manner that they alter the initial conditions for the DGLAP evolution, making it impossible to apply them at fixed $`Q^2=Q_0^2`$. With SC we have to solve the DGLAP equations starting with $`Q^2=Q_0^2(x)`$ where $`Q_0^2`$ is the solution of the equation $`\kappa _{dipole}^{DGLAP}(x,Q_0^2(x))=\mathrm{\hspace{0.17em}1}`$. Therefore, we are in controversial situation: we require a GLAP input but, if we take it from the so called global fits, there is a danger that we will incorporate all the effects of the SC in the initial condition of these parameterizations. Only data taken after the 1995 runs are at energies sufficiently high to effect the low $`x`$ behaviour of the initial inputs. GRV’94 is based on the experimental data at rather large $`x`$, so we hope that SC are minimal in this parameterization;
2. The GRV parameterization starts with rather low virtualities ( as low as $`Q^2\mathrm{\hspace{0.17em}0.5}GeV^2`$ ). This is a major weakness of this approach, since one cannot guarantee that only the leading twist contribution is dominant in the DGLAP evolution equations. We agree with this criticism, but the low value of $`Q_0^2`$ leads to the solution of the DGLAP equation which is closer to the leading log approximation, in which we can guarantee the accuracy of our master equation ( see Eq. (2.80) ).
Before we present our numerical results we have to make an important comment. It concerns the substitution in Eq. (2.64), where the iterarted gluon density $`xG^{MG}(x,4/r^2;b)`$ is defined. Due to technical problems, two modifications of the formula are made in the actual numerical calculations. First of all we assume the exponential $`b`$ factorization of $`xG^{MG}(x,4/r^2;b)`$:
$$xG^{MG}(x,4/r^2;b)=\frac{1}{\pi R^2}xG^{MG}(x,4/r^2)e^{\frac{b^2}{R^2}},$$
where
$$xG^{MG}(x,4/r^2)=d^2bxG^{MG}(x,4/r^2;b)$$
(3.1)
We justify the factorization by the following two arguments. The first one is numerical. We actually have checked that factorization holds numerically with satisfactory accuracy. The second argument is that having assumed factorization for $`xG^{DGLAP}`$ on the same basis we can assume it for $`xG^{MG}`$. The second modification, made with Eq. (2.64) is due to the fact that the GRV parameterization we are working with does not satisfy the DGLAP equation in DLA (see Ref. for the discussion of the problem). The authors of suggested a modification of Eq. (3.1), which we implement in all our numerical calculations. The final formula for $`xG^{MG}`$ has the form:
$`xG^{MG}(x,Q^2)={\displaystyle \frac{4}{\pi ^2}}{\displaystyle _x^1}{\displaystyle \frac{dx^{}}{x^{}}}{\displaystyle _{4/Q^2}^{\mathrm{}}}{\displaystyle \frac{dr^2}{r^4}}{\displaystyle 𝑑b^2\mathrm{\hspace{0.17em}2}(1e^{\mathrm{\Omega }^{GRV}(x^{},r^{},b)})}+`$
(3.2)
$`xG^{GRV}(x,Q^2){\displaystyle \frac{\alpha _SN_c}{\pi }}{\displaystyle _x^1}{\displaystyle _{Q_0^2}^{Q^2}}{\displaystyle \frac{dx^{}}{x^{}}}{\displaystyle \frac{dQ^2}{Q^2}}x^{}G^{GRV}(x^{},Q^2).`$
In all the figures representing our numerical results the upper solid line corresponds to the full answer, the widely spaced dashed line to diffractive production of $`q\overline{q}`$ pair, while the narrowly spaced dashed line to the diffractive production of $`q\overline{q}G`$ .
In Fig. 12 we plot the results of our calculations using Eq. (2.80). We have not added any contribution from the nucleon excitations, based on our estimates given in Fig.7b. In Fig.13 we show our predictions for small values of $`Q^2`$. The first observation is the fact that the ratio increases considerably, reaching the value of about 25 %. It should be stressed that even at rather small values of $`Q^2`$, we do not see any sign of saturation of the energy behaviour of this ratio, which is continously increasing in the HERA kinematic region.
It is interesting to consider separately the quark - antiquark diffractive production and the production of the $`q\overline{q}G`$ final state. Both cross sections increase as a function of the energy $`W`$. At small $`Q^2`$ (Fig. 13) the total diffractive cross section is dominated (about 70 %) by the diffractive production of the quark - antiquark pair. The situation changes when higher values of $`Q^2`$ are considered. Then the $`q\overline{q}G`$ final state contibutes even more than the $`q\overline{q}`$ pair. Indeed, at high energy $`W=200GeV`$, $`q\overline{q}G`$ is produced in approximately $`\mathrm{\hspace{0.17em}25}\%`$ of all diffractive events at $`Q^2=\mathrm{\hspace{0.17em}8}GeV^2`$. For smaller values of $`Q^2`$ this fraction decreases, being only 10% at $`Q^2=1GeV^2`$. This is a expected behaviour if the SC play an important role. The results presented show that the contribution of the extra gluon emission is crucial for the predictions, and may lead to a hundred percent enhancement of the ratio.
At small $`Q^2`$, the fraction of $`q\overline{q}G`$ diffractive production to the total cross section increases sufficiently slow due to the suppression factor $`e^{\mathrm{\Omega }^{MG}(x,r_{};b)}`$ in Eq. (2.80). For larger $`Q^2`$, $`\mathrm{\Omega }^{MG}`$ becomes small and the fraction of $`q\overline{q}G`$ diffractive production increases faster with the energy.
In Fig. 14 we illustrate the dependence of our calculation on the value of $`R^2`$. We would like to recall that two values of the radius $`R`$, which we use, have the following physics behind them:
1. $`R^2=\mathrm{\hspace{0.17em}10}GeV^2`$ is related to the Mueller-Glauber approach, which corresponds to the Eikonal model for the nucleon target ;
2. $`R^2=\mathrm{\hspace{0.17em}5}GeV^2`$ is the average radius for the two channel model which has been discussed in subsection 2.2.3(B) ;
One can see from Fig. 14 that the value of the ratio $`\sigma ^{DD}/\sigma _{tot}`$ depends on $`R^2`$. However, the energy dependence is still very pronounced.
Therefore, the general conclusion of this section is that the SC fail to reproduce the constant ratio of $`\sigma ^{DD}/\sigma _{tot}`$ for DIS seen in the HERA kinematic region.
## 4 Ratio $`\sigma ^{\mathrm{𝐃𝐃}}/\sigma _{\mathrm{𝐭𝐨𝐭}}`$ in the mass windows
In this section we examine a possibility that the mass interval will induce an energy independent ratio $`\sigma ^{DD}/\sigma _{tot}`$. As one can see in Fig.1, the experimental measurements were made within some windows in mass.
The full derivation of the mass dependent formulae is presented in the Appendix. The cross sections for the diffractive dissociation production of $`q\overline{q}`$ pair and $`q\overline{q}G`$ parton system with a definite final state mass are given by Eq. (A.4) and Eq. (A.10) respectively. Any mass window can be selected for the mass integrals.
If summation over the whole infinite mass intervals is performed, our formulae for the diffractive dissotiation as well as for the total cross sections should in principle reproduce Eq. (2.80). However, the two sets of formulas are different but consistent with each other in the leading $`\mathrm{log}(1/x)`$ approximation where $`\mathrm{log}(1/x_P)\mathrm{log}(1/x_B)`$. When $`x_P`$ is replaced by $`x_B`$, Eq. (A.4) and Eq. (A.3) reproduce analytically (and numerically) the corresponding expressions in Eq. (2.80). In the numerical computations we use $`x_P`$ instead of $`x_B`$ since this energy variable reflects the real kinematics of the diffraction production process. Since $`x_P=x_B/\beta `$ and the typical values of $`\beta `$ is not very small for $`q\overline{q}`$ system, we do not expect large corrections due to this substitution. The $`q\overline{q}G`$ channel is more sensitive to the kinematic restriction (see Eq. (A.10)). The changes, introduced in Eq. (A.10), concern both the energy variables as well as the integration limits. These cannot be justified in log(1/x) limit which we used in our general formulae, but we have to introduced them to make our calculation reasonable for the diffractive production in the mass window.
Figure 15 presents our results for the mass bins, for which experimental data exists (Fig. 1). The ratio of the diffractive dissociation cross section to the inclusive cross section is plotted as a function of the center of mass energy. The $`q\overline{q}`$ pair (tranverse plus longitudinal) and $`q\overline{q}G`$ contributions are shown separately.
The ratios obtained do not reproduce the experimental data (Fig. 1). The significant energy dependence persists due to the growth of both the $`q\overline{q}`$ and $`q\overline{q}G`$ contributions. In the wide range of the energies our results are smaller than the experimental curves. This observation is quite consistent, since our model excluded the target excitations estimated in the Chapter 2 by about 30%.
As it should be, at small masses the main contributions come from the $`q\overline{q}`$ pair with more than 50% given by the longitudinal part. The $`q\overline{q}G`$ production is suppressed at small masses. Its contibution grows with the mass and dominates at large masses.
Summing all results up to the mass $`M=15`$(GeV) we do not reproduce the inclusive mass results of the previous section. This means that even for $`Q^2=8\mathrm{G}\mathrm{e}\mathrm{V}^2`$ we expect contributions from higher masses. At $`Q^2=8(\mathrm{GeV}^2)`$ these contributions originated from $`q\overline{q}G`$ production. At $`Q^2=60(\mathrm{GeV}^2)`$ both $`q\overline{q}`$ and $`q\overline{q}G`$ will be significant above $`M=15`$(GeV). It is seen that only about 50% of the inclusive DD production is contributed by small masses up to $`M=15`$(GeV).
We wish to point that we have completely disregarded possible final state corrections. As an example of such corrections, the master formula Eq. (2.80) does not take into account the second diagram in Fig. 2. This diagram is believed to give a relatively small contribution compared to the first and the third diagrams.
It is worthwhile comparing our model with the Golec-Biernat Wusthoff model , which successefully reproduces the experimental data (Fig. 1). In the Golec-Biernat Wusthoff model the effective dipole cross section $`\sigma _{GW}(x,r)`$, describing the interaction of the $`q\overline{q}`$ pair with a nucleon has the form:
$`\sigma _{GW}(x,r_{})=\sigma _0[1\mathrm{exp}(r_{}^2/(4R_0^2(x)))];R_0(x)=(x/x_0)^{\lambda /2}(\mathrm{GeV}^1);`$ (4.1)
$`\sigma _0=\mathrm{\hspace{0.17em}23.03}(\mathrm{mb});\lambda =\mathrm{\hspace{0.17em}0.288};x_0=\mathrm{\hspace{0.17em}3.04}\mathrm{\hspace{0.17em}10}^4.`$
In this model, the diffractive dissociation cross section is given by the squaring of $`\sigma _{CW}`$ in Eq. (4.1):
$$\sigma _{GW}^{DD}(x,r_{})=\sigma _{CW}^2/(16\pi B_D);B_D=\mathrm{\hspace{0.17em}6}\mathrm{GeV}^2.$$
(4.2)
A comparison between CW model and the present work model is presented in Fig. 16.
We found a significant difference between the two models . The advantage of G-W model is that this model takes into account in the simplest way a new scale: saturation momentum $`Q_s^2(x)=4/R_0^2(x)`$, but in doing so, this model loses its correspondence with the DGLAP evolution equation. Our approach has a correct matching with the DGLAP evolution and we expected that we would be able to describe experimental data better than the G-W model. It turns out ( see Fig.16-a and Fig. 16-b ) that we do not reproduce the $`\sigma ^{DD}/\sigma _{tot}`$ ratio in contrast to the G-W model, mostly due to our ‘improvement’ in the region of small $`r_{}^2`$.
The second remark is the substantial difference in the way we describe the $`q\overline{q}G`$ state. We failed to find a correspondence between our formula for $`q\overline{q}G`$ production which follows from the AGK cutting rules, and the G-W description of this process. However, our failure to fit the experimental data is mostly due to a large difference in the dipole cross section, rather than in the different treatment of the $`q\overline{q}G`$ diffractive production.
Figs.17-1 and 17-2 show our dipole cross sections at different energies. The main difference with the Golec-Biernat and Wusthoff model is the energy rise of the cross sections which follows from $`b_t`$ dependance included in the Eikonal formulae but neglected in the Golec-Biernat and Wusthoff model. Fig.17-3 - Fig.17-10 show whast distances are essential in our calculations. One can see that the main contributions in diffractive cross sections stem from longer distances than in the total cross sections as have been predicted theoretically . As has been known for long time the typical distances in diffractive cross section is of the order of the saturation scale $`Q_s(x)`$ ($`\kappa _{dipole}(r_t=2/Q_s(x))=1`$ ) in contrast with total cross section were typical distances are much shorter , about $`1/Q`$.. Therefore, one of the reason why we failed to describe the ratio of interest could be that our model cannot describe the dipole cross section in the vicinity of the saturation scale. However, it was demonstrated in AGL papers ( see Ref. ) that our Eikonal model gives a good approximation to the correct non-linear evolution equations at that particular distances which are essential accordingly to Fig.17.
## 5 Summary and discussions
In this paper we developed an approach to the diffractive dissociation in DIS , based on three main ideas:
1. The dominant contribution to diffractive dissociation processes in DIS stems from rather short distances and, therefore, we can use pQCD to describe them;
2. The final state in diffractive dissociation can be simplified in the HERA kinematic region by only considering the production of $`q\overline{q}`$ and $`q\overline{q}G`$ ( );
3. The shadowing corrections which are essential for the description of the diffractive processes, can be taken into account, using Mueller-Glauber approach .
Using pQCD and Mueller-Glauber approach we derived a generalization of Kovchegov-McLerran formula for the ratio $`\sigma ^{DD}/\sigma _{tot}`$ ( see Eq. (2.80) ) which is applicable to the HERA experimental data on diffractive production at HERA. However, we found that Eq. (2.80) cannot describe the approximate energy independence of the ratio $`\sigma ^{DD}/\sigma _{tot}`$ , observed experimentally. Our attempts to introduce the experimental cuts for diffractive production does not change this pessimistic conclusion.
Therefore, we believe, this paper is a strong argument that the nonperturbative QCD contribution is essential for diffractive production and our approach, based on pQCD, should be reconsidered. However, we showed that the main source of the observed energy dependence arises from rather short distances ( see Figs. 16 and 17 ) where we did not see any nonperturbative correction to the total DIS cross section. In principle, it was pointed out in Ref. that the scale anomaly of QCD generates a nonperturbative contribution at high energy at sufficiently short distances $`r_{}^21/M_0^20.25GeV^1`$ . We intend studying the influence of such nonperturbative corrections in further publications.
We studied in detail the contribution of the excited hadronic states in diffractive production. We found that the contribution of nucleon excitations should depend on $`Q^2`$ and $`x`$ leading to large cross section at bigger values of $`Q^2`$ and higher $`x`$.
We hope that our paper will draw the attention of the high energy community to the beautiful experimental data on the energy dependance of the ratio $`\sigma ^{DD}/\sigma _{tot}`$ which have still not recieved an adequate theoretical explanation, and which can provide a new insight to the importance of nonperturbative corrections at sufficiently short distances.
Acknowledgements: E.L. would like to acknowledge the hospitality extended to him at DESY Theory Group where this work was started.
The research was supported in part by BSF # 9800276 and by the Israel Science Foundation, founded by the Israeli Academy of Science and Humanities.
## Appendix A Appendix
In this Appendix we present a derivation of formulae for the cross sections for the diffractive dissociation production of a $`q\overline{q}`$ pair and the $`q\overline{q}G`$ parton system when a final state mass window is selected. In order to preserve the unitarity relation between the DD cross section and the total cross section, we modify the latter. All the formulae are written in the leading $`\mathrm{log}(1/x)`$ aproximation of pQCD. Below we only present results derived for the transversely polarized photon. The results for the longitudinal part can be obtained by similar treatment.
### A.1 $`q\overline{q}`$ contribution
The $`q\overline{q}`$ DD cross section has the form :
$`\sigma _{qq}^{DD}=\mathrm{\hspace{0.17em}16}\pi {\displaystyle }{\displaystyle \frac{dM^2}{M^2+Q^2}}{\displaystyle \frac{\alpha _{em}N_c}{8(2\pi )^2}}{\displaystyle \underset{f}{}}Z_f^2{\displaystyle }{\displaystyle \frac{d^2k_{}}{(2\pi )^2}}{\displaystyle _0^1}dz\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}4}\pi N_\lambda (M^2+Q^2)\times `$
$`\delta (M^2{\displaystyle \frac{k^2}{z(1z)}})|f_{qq}|^2,`$ (A.1)
where within our model for the dipole interaction the square of the amplitude is written:
$`|f_{qq}|^2={\displaystyle }d^2b{\displaystyle }{\displaystyle \frac{d^2r}{2\pi }}{\displaystyle \frac{d^2r^{}}{2\pi }}\mathrm{\Psi }^\gamma ^{}(r))^{}\mathrm{\Psi }^\gamma ^{}(r^{})\{1e^{\frac{\mathrm{\Omega }^{MG}(x^{},r;b)}{2}}\}\{1e^{\frac{\mathrm{\Omega }^{MG}(x^{},r^{};b)}{2}}\}e^{i\stackrel{}{k}(\stackrel{}{r}\stackrel{}{r}^{})}`$ (A.2)
The notation has been introduced in section 2 (see also Fig. 3), while the parameters are defined as follows.
$$x^{}=(k^2/z(1z)+Q^2)/W^2;N_\lambda =z^2+(1z)^2;a^2=Q^2z(1z)$$
It is important to note that in Eq. (A.2) we introduce a correct energy argument $`x^{}`$ in $`\mathrm{\Omega }`$ since at fixed $`M^2`$ the energy of the dipole-proton interaction is $`\beta W^2`$. The kinematic constraint forced by the delta function sets $`x^{}=x_P`$. Performing the angle integrations we obtain
$`|f_{qq}|^2={\displaystyle d^2b\left[_0^{\mathrm{}}𝑑rraK_1(ar)J_1(kr)\{1e^{\frac{\mathrm{\Omega }^{MG}(x^{},r;b)}{2}}\}\right]^2},`$ (A.3)
Finally for the cross-section we get the result
$`\sigma _{qq}^{DD}=\mathrm{\hspace{0.17em}8}\alpha _{em}{\displaystyle \frac{dM}{M^3}_0^{M/2}\frac{dkk^3N_\lambda }{\sqrt{14k^2/M^2}}𝑑b^2\left[_0^{\mathrm{}}𝑑rraK_1(ar)J_1(kr)\{1e^{\frac{\mathrm{\Omega }^{MG}(x_p,r;b)}{2}}\}\right]^2}`$ (A.4)
with
$$N_\lambda =\mathrm{\hspace{0.17em}\hspace{0.17em}1}\mathrm{\hspace{0.17em}\hspace{0.17em}2}k^2/M^2;a=Qk/M.$$
### A.2 $`q\overline{q}G`$ contribution
Consider the diffractively produced $`q\overline{q}G`$ system with $`z`$ and $`z^{}`$ being the fractions of the energy carried by quark and gluon respectively. We assume that the transverse gluon momentum $`l`$ is much smaller than the quark transverse momentum $`k`$. In the leading $`\mathrm{log}(1/x)`$ approximation of pQCD $`z^{}z`$. The kinematic constraint is dictated by the final state mass:
$$M^2=k^2/z(1z)+l^2/z^{}.$$
The cross section for the $`q\overline{q}G`$ production is
$`\sigma _{qqG}^{DD}=(16\pi ){\displaystyle \frac{\alpha _{em}N_c}{8(2\pi )^2}}{\displaystyle \underset{f}{}}Z_f^2{\displaystyle 𝑑M^2\frac{d^2k}{(2\pi )^2}_0^1𝑑z\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}4}\pi N_\lambda |f_{qqG}|^2},`$ (A.5)
with the square of the amplitude
$`|f_{qqG}|^2={\displaystyle }d^2b{\displaystyle }{\displaystyle \frac{d^2r}{2\pi }}{\displaystyle \frac{d^2r^{}}{2\pi }}\mathrm{\Psi }^\gamma ^{}(r))^{}\mathrm{\Psi }^\gamma ^{}(r^{})e^{\frac{\mathrm{\Omega }^{MG}(x^{},r;b)}{2}}e^{\frac{\mathrm{\Omega }^{MG}(x^{},r^{};b)}{2}}e^{i\stackrel{}{k}(\stackrel{}{r}\stackrel{}{r}^{})}\stackrel{~}{\kappa }/2.`$ (A.6)
$`\stackrel{~}{\kappa }`$ is defined as follows.
$`\stackrel{~}{\kappa }={\displaystyle \frac{2}{\pi ^2}}{\displaystyle _0^z}dz^{}{\displaystyle _r^{\mathrm{}}}{\displaystyle \frac{d^2R}{\pi }}{\displaystyle \frac{d^2R^{}}{\pi }}{\displaystyle _0^{k^2}}{\displaystyle \frac{d^2l}{(2\pi )^2}}\delta (M^2k^2/z(1z)l^2/z^{})e^{i\stackrel{}{l}(\stackrel{}{R}\stackrel{}{R}^{})}\times `$
$`\mathrm{\Psi }_g(R)\mathrm{\Psi }_g^{}(R^{})(1e^{\frac{\mathrm{\Omega }_G^P(x_P,R;b)}{2}})(1e^{\frac{\mathrm{\Omega }_G^P(x_P,R^{};b)}{2}})\sqrt{\alpha _S(r^2)}\sqrt{\alpha _S(r^2)}{\displaystyle \frac{\pi ^2\stackrel{}{r}\stackrel{}{r^{}}}{3}}`$ (A.7)
We use the small $`z^{}`$ approximation of the gluon wave function
$$\mathrm{\Psi }_g^{mn}(R,z^{})\frac{\sqrt{2}}{\sqrt{z^{}}R^2}\left(\delta ^{mn}\mathrm{\hspace{0.17em}2}\frac{R^mR^n}{R^2}\right)$$
Performing the angle integration and removing the delta function by doing the $`z^{}`$ integration we obtain
$`|f_{qqG}|^2={\displaystyle \frac{2}{3}}{\displaystyle \frac{1}{M^2k^2/z(1z)}}{\displaystyle }d^2b{\displaystyle _0^{zM^2k^2/(1z)}}{\displaystyle \frac{d^2l}{(2\pi )^2}}\times `$
$`\left[{\displaystyle \frac{d^2r}{2\pi }aK_1(ar)\sqrt{\alpha _S(r^2)}rJ_1(kr)e^{\frac{\mathrm{\Omega }^{MG}(x^{},r;b)}{2}}_{r^2}^{\mathrm{}}\frac{dR^2}{R^2}J_2(lR)(1e^{\frac{\mathrm{\Omega }_G^P(x_P,R;b)}{2}})}\right]^2`$ (A.8)
As a result of the delta function integration we also find that $`k^2M^2z(1z)`$. It should be stressed that $`r`$ is the size of the initial quark -antiquark pair, while $`R`$ is the size of the produced two colour dipoles. In our approximation, both dipoles have the same size with $`Rr`$. Introducing dimensionless variable $`\stackrel{~}{l}`$
$$l^2=(zM^2k^2/(1z))\stackrel{~}{l}^2$$
(A.9)
we finally arrive at the expression for the cross section:
$`\sigma _{qqG}^{DD}={\displaystyle \frac{\alpha _{em}}{6}}{\displaystyle }dM^2{\displaystyle _0^1}dzzN_\lambda {\displaystyle _0^{M^2z(1z)}}dk^2{\displaystyle }db^2{\displaystyle _0^1}d\stackrel{~}{l}^2\times `$ (A.10)
$`\left[{\displaystyle 𝑑rr^2aK_1(ar)\sqrt{\alpha _S(r^2)}J_1(kr)e^{\frac{\mathrm{\Omega }^{MG}(x^{},r;b)}{2}}_{r^2}^{\mathrm{}}\frac{dR^2}{R^2}J_2(lR)(1e^{\frac{\mathrm{\Omega }_G^P(x_P,R;b)}{2}})}\right]^2.`$
Contrary to the $`q\overline{q}`$ case, in the present expression $`x^{}`$, which is $`(k^2/z(1z)+Q^2)/W^2`$ is not equal to $`x_P`$. The energy variables are different for $`\mathrm{\Omega }^{GM}`$ and $`\mathrm{\Omega }_G^P`$, because they describe different physics. Indeed, factor $`e^{\mathrm{\Omega }^{MG}}`$ is a probability that $`q\overline{q}`$ pair does not interact inelastically before emission of the extra gluon, while $`(1e^{\frac{\mathrm{\Omega }_G^P}{2}})^2`$ stands for diffractive production of $`q\overline{q}G`$ system ( two colour dipoles of size $`R`$ ). For the rescattering of the $`q\overline{q}`$ colour dipole the energy is $`s=W^2x_B`$, while the rescattering of the colour dipoles of size $`R`$ occurs at energy $`s^{}=\beta s`$.
### A.3 Total cross section
In order to preserve the unitarity relation between $`q\overline{q}`$ DD cross section and the total cross section we modify the latter. Similarly to the $`q\overline{q}`$ case we write for the total cross section
$`\sigma _{tot}=\mathrm{\hspace{0.17em}\hspace{0.17em}16}\pi {\displaystyle }{\displaystyle \frac{dM^2}{M^2+Q^2}}{\displaystyle \frac{\alpha _{em}N_c}{8(2\pi )^2}}{\displaystyle \underset{f}{}}Z_f^2{\displaystyle }{\displaystyle \frac{d^2k_{}}{(2\pi )^2}}{\displaystyle _0^1}dz\mathrm{\hspace{0.17em}\hspace{0.17em}4}\pi N_\lambda (M^2+Q^2)\times `$
$`\delta (M^2{\displaystyle \frac{k^2}{z(1z)}})|f_{tot}|^2`$ (A.11)
with the square of the amplitude
$`|f_{tot}|^2`$ $`=`$ $`{\displaystyle d^2b\frac{d^2r}{2\pi }\frac{d^2r^{}}{2\pi }(\mathrm{\Psi }^\gamma ^{}(r))^{}\mathrm{\Psi }^\gamma ^{}(r^{})\{2e^{\frac{\mathrm{\Omega }^{MG}(x^{},r;b)}{2}}e^{\frac{\mathrm{\Omega }^{MG}(x^{},r^{};b)}{2}}\}e^{i\stackrel{}{k}(\stackrel{}{r}\stackrel{}{r}^{})}}=`$
$`{\displaystyle d^2b_0^{\mathrm{}}𝑑rr𝑑r^{}r^{}a^2K_1(ar)J_1(kr)K_1(ar^{})J_1(kr^{})\{2e^{\frac{\mathrm{\Omega }^{MG}(x^{},r;b)}{2}}e^{\frac{\mathrm{\Omega }^{MG}(x^{},r^{};b)}{2}}\}}`$
One of the space integrations can be performed analytically noting that
$`{\displaystyle _0^{\mathrm{}}}𝑑rrK_1(ar)J_1(kr)={\displaystyle \frac{k}{a(k^2+a^2)}}`$ (A.13)
Substituting $`a=Qk/M`$ we finally obtain the total cross section
$`\sigma _{tot}=\mathrm{\hspace{0.17em}\hspace{0.17em}16}\alpha _{em}{\displaystyle }{\displaystyle \frac{dM}{M^2}}{\displaystyle \frac{Q}{Q^2+M^2}}{\displaystyle _0^{M/2}}k^3dk{\displaystyle \frac{1\mathrm{\hspace{0.17em}2}k^2/M^2}{\sqrt{14k^2/M^2}}}\times `$
$`{\displaystyle 𝑑b^2𝑑rrK_1(ar)J_1(kr)\{1e^{\frac{\mathrm{\Omega }^{MG}(x_P,r;b)}{2}}\}}`$ (A.14)
In the above expression for the total cross section the mass integration should be carried out over the whole infinite mass interval. The result obtained is consistent with the previous expression for the total cross section. If we change $`x_P`$ to $`x_B`$ and perform the $`M`$ integration we reproduce the old result written in Eq. (2.80). A useful equality we use is
$`{\displaystyle 𝑑M\frac{M^2}{Q^2+M^2}J_1(Mr)}=QK_1(Qr)`$ (A.15)
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# Quantum no-deleting principle and some of its implications
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## Abstract
Unmeasureability of a quantum state has important consequences in practical implementation of quantum computers. Like copying, deleting of an unknown state from among several copies is prohibited. This is called no-deletion prinicple. Here, we present a no deleting principle for qudits. We obtain a bound on $`N`$-to-$`M`$ deleting and show that the quality of deletion drops exponentially with the number of copies to be deleted. In addition, we investigate conditional, state-dependent and approximate quantum deleting of unknown states. We prove that unitarity does not allow us to delete copies from an alphabet of two non-orthogonal states exactly. Further, we show that no-deleting principle is consistent with no-signalling.
PACS Numbers: 03.67.-a, 03.65.Bz, 89.70.+c \]
Linearity of quantum theory unveils that we cannot duplicate an unknown quantum state accurately. Unitarity of quantum evolution shows that non-orthogonal quantum states cannot be perfectly copied . However, they can be can be copied exactly by a unitary and measurement process . The possibility of copying an unknown state approximately by deterministic cloning machines was also proposed. Further, non-orthogonal states can evolve into a linear superposition of multiple copies by a novel cloning machine .
In the quantum information and computation era it is important to know what we can do with the vast amount of information contained in an unknown state and what we cannot. For example, given several (a finite number of) copies of an unknown state we can partly estimate it , we can swap it and we can teleport it . But can we delete a copy of an unknown state from a collection of several copies?
The quantum deletion we study here is not same as the erasure. When we wish to get rid of the last bit (either classical or quantum) of information it is called primitive erasure. This can be achieved by spending certain amount of energy (thermodynamically irreversible), known as Landauers’s erasure principle . We will consistently use the term “deleting” to refer to a uncopying-type of operation as opposed to primitive erasure. We have recently shown that unlike classical information in quantum theory the perfect deletion of an unknown qubit from a collection of two or more qubits is an impossible operation. The basic linear structure of quantum theory puts severe limitations on the complete deleting of the quantum information of an unknown state . In this letter we generalise our no-deleting principle to qudits (qudit is a $`d`$-dimensional quantum system). We obtain a bound on the maximum limit of $`N`$-to-$`M`$ quantum deletion. We study so-called conditional quantum deleting and state-dependent, approximate deleting. We prove a no-deletion theorem for non-orthogonal states using unitarity. Interestingly, we show that the quantum no-deleting principle is consistent with no-signalling.
Quantum deleting of qudits: Consider several copies (say $`N`$) of an unknown quantum state $`|\mathrm{\Psi }`$ each in a $`d`$-dimensional Hilbert space $`=𝒞^d`$. Our $`N`$ copies $`|\mathrm{\Psi }^N`$ live in a smaller dimensional subspace, which is the symmetric subspace of $`^N`$. It contains states that are symmetric under interchange of any pair of qudits. The aim of the quantum deleting machine is to delete one or more number of qudits from a collection of two or more qudits all in the same state. In a sense, we intend to construct a machine which appears to perform the ‘reverse’ of cloning operation (but as we will see later, strictly it is not so). In general the quantum deleting operation is defined for $`N`$ unknown states $`|\mathrm{\Psi }^N`$ such that the linear operator acts on the combined Hilbert space and deletes $`(NM)`$ copies and keeps $`M`$ copies intact. It is defined by
$`|\mathrm{\Psi }^N|A|\mathrm{\Psi }^M|\mathrm{\Sigma }^{(NM)}|A_\mathrm{\Psi },`$ (1)
where $`|\mathrm{\Sigma }`$ is the blank state of a qudit, $`|A`$ is the initial and $`|A_\mathrm{\Psi }`$ is the final state of the ancilla which may in general depend on $`|\mathrm{\Psi }`$.
For simplicity, let us consider a $`2`$-to-$`1`$ quantum deleting machine for qudits. Let $`\{|\psi _i,i=1,2,\mathrm{}d\}`$ be an arbitrary orthonormal basis of a qudit ($`\psi `$ is labelling the choice of basis). For pedagogical reason consider a scenario where only Hilbert spaces involving the initial copies will be investigated. Thus, the deleting machine is a linear operator that acts jointly on the copies of qudits, given by
$`|\psi _i_a|\psi _i_b|\psi _i_a|\mathrm{\Sigma }_b.`$ (2)
Note that if the inputs are different then the deleting machine can yield an arbitrary state such as
$`|\psi _i_a|\psi _j_b|\mathrm{\Phi }_{ij}^\psi _{ab}(ij).`$ (3)
However, if we send two copies of an arbitrary qudit in a state $`|\mathrm{\Psi }=_{i=1}^dc_i|\psi _i`$ (with $`c_i`$’ s being all unknown complex numbers), then by linearity we have
$`|\mathrm{\Psi }_a|\mathrm{\Psi }_b{\displaystyle \underset{i=1}{\overset{d}{}}}c_i^2|\psi _i_a|\mathrm{\Sigma }_b+{\displaystyle \underset{ij=1_{ij}}{\overset{d}{}}}c_ic_j|\mathrm{\Phi }_{ij}^\psi _{ab}.`$ (4)
But ideally we would require $`|\mathrm{\Psi }_a|\mathrm{\Psi }_b_ic_i|\psi _i_a|\mathrm{\Sigma }_b`$. Since ideal output and actual output states are different linearity does not allow us to delete an unknown copy of a quantum state in any finite dimensional Hilbert space.
We can also prove the quantum no-deletion principle by including ancilla. However, when we include ancilla we have to exclude swapping of an unknown copy onto ancilla as a proper deletion. The reason for doing so is first, swapping of an unknown copy onto ancilla is just hiding of quantum information and second, if we allow swapping then the result reduces to primitive erasure and the extra copies that are available at our disposal have played no role in deleting mechanism. Intuitively one would say that if we have a large number of copies of an unknown quantum state we know more about the state (as only in the limit of infinite number of copies we know the state exactly). So extra copies should help in deleting an unknown state. Now including ancilla let us define the action of deleting machine on orthogonal qudits as given by
$`|\psi _i_a|\psi _i_b|A_c|\psi _i_a|\mathrm{\Sigma }_b|A_{\psi _i}_c,`$ (5)
where $`|A_{\psi _i}`$’s need not be orthogonal. When the inputs are in different state then the deleting machine can yield an arbitrary state such as
$`|\psi _i_a|\psi _j_b|A_c|\mathrm{\Phi }_{ij}^\psi _{abc}.`$ (6)
It should be noted that because of (5) and (6) quantum deleting machine is not reverse of cloning machine. If we send two copies of an arbitrary qudit in a state such as $`|\mathrm{\Psi }`$, then by linearity we have
$`|\mathrm{\Psi }_a|\mathrm{\Psi }_b|A_c{\displaystyle \underset{i}{}}c_i^2|\psi _i_a|\mathrm{\Sigma }_b|A_{\psi _i}_c`$ (7)
$`+{\displaystyle \underset{ij_{ij}}{}}c_ic_j|\mathrm{\Phi }_{ij}^\psi _{abc},`$ (8)
which is a qudratic polynomial in $`c_i`$’s. But ideally if we would be able to delete a copy of an unknown qudit then the above equation must reduce to $`_ic_i|\psi _i_a|\mathrm{\Sigma }_b|A_\mathrm{\Psi }_c`$ for all $`c_i`$’s. Note that once we exclude swapping, the ancilla state $`|A_\mathrm{\Psi }`$ is independent of the input state $`|\mathrm{\Psi }`$. Since we know that $`|\mathrm{\Phi }_{ij}^\psi _{abc}`$ can never depend on $`c_i`$’s, therefore the only solution this allows is $`|\mathrm{\Phi }_{ij}^\psi _{abc}=|\psi _i_a|\mathrm{\Sigma }_b|A_{\psi _j}_c`$ and $`|A_\mathrm{\Psi }_c=_ic_j|A_{\psi _j}_c`$. Since the final actual output state has to be normalised for all $`c_i`$’s we can see that all $`|A_{\psi _j}`$’s have to be orthonormal. Therefore, linearity allows only swapping of an unknown qudit onto the $`d`$-dimensional subspace of the ancilla. That is the unknown quantum state is just hidden in the deleting machine state which can of course be retrieved by a unitary transformation. Hence, linearity does not allow us to delete an unknown state against a copy, even in the presence of ancilla. We can only move quantum information around but we cannot delete it completely. This is called quantum no-deletion principle for qudits.
Bounds on $`N`$-to-$`M`$ quantum deleting: Below we restrict our discussions for qubits but they can be generalised to qudits. Consider a new scenario where Victor owns a “qubit company” and prepares copies of a qubit. Let $`N`$ copies of a qubit given to us live in a symmetric subspace of $`(N+1)`$ dimensional Hilbert space of the full $`2^N`$-dimensional Hilbert space. The $`N`$ copies of an unknown qubit $`|\mathrm{\Psi }=\alpha |0+\beta |1`$ can be written as
$`|\mathrm{\Psi }^N=\alpha ^N|0^N+\beta ^N|1^N+{\displaystyle \underset{k=1}{\overset{N1}{}}}f_k(\alpha ,\beta )|k,`$ (9)
where $`|k`$’s are $`(N1)`$ orthogonal bit string states living in the symmetric subspace.
Our deleting machine for orthogonal qubits is defined by $`|0^N|A|0^M|\mathrm{\Sigma }^{(NM)}|A_0,|1^N|A|1^M|\mathrm{\Sigma }^{(NM)}|A_1,`$ and $`|k|A|k^{}`$, where $`|k^{}`$ is final state of the symmetric $`N`$ qubits state and the ancilla. If we send $`N`$ copies of an unknown qubit through our deleting machine this will yield
$`|\mathrm{\Psi }^N|A\alpha ^N|0^M|\mathrm{\Sigma }^{(NM)}|A_0+`$ (10)
$`\beta ^N|1^M|\mathrm{\Sigma }^{(NM)}|A_1+{\displaystyle \underset{k=1}{\overset{N1}{}}}f_k(\alpha ,\beta )|k^{}=|\mathrm{\Psi }_{\mathrm{actual}}.`$ (11)
Now, Victor wants to test how good is our quantum deleting machine by evaluating the worst case success of performing deletion operation. As Victor knows what the state is he can delete $`(NM)`$ copies perfectly and keep $`M`$ copies intact. This ideal operation is given by
$`|\mathrm{\Psi }^N|A|\mathrm{\Psi }^M|\mathrm{\Sigma }^{(NM)}|A_\mathrm{\Psi }=(\alpha ^M|0^M+\beta ^M|1^M`$ (12)
$`+{\displaystyle \underset{j=1}{\overset{M1}{}}}g_j(\alpha ,\beta )|j)|\mathrm{\Sigma }^{(NM)}|A_\mathrm{\Psi }=|\mathrm{\Psi }_{\mathrm{ideal}}.`$ (13)
Therefore, the error introduced by the quantum deleting machine can be calculated from $`=1𝒬=1|\mathrm{\Psi }_{\mathrm{actual}}|\mathrm{\Psi }_{\mathrm{ideal}}|`$. If the ideal and the actual states are identical then obviously there is no error. The quantity $`𝒬`$ called quality function is bounded by
$`𝒬|\alpha |^{(N+M)}+|\beta |^{(N+M)}`$ (14)
$`+[1(|\alpha |^{2N}+|\beta |^{2N})]^{\frac{1}{2}}[1(|\alpha |^{2M}+|\beta |^{2M})]^{\frac{1}{2}}.`$ (15)
The rhs of (11) can be optimised and the optimal value of $`𝒬`$ is given by
$`𝒬_{\mathrm{opt}}={\displaystyle \frac{2}{2^{(N+M)/2}}}+\sqrt{\left(1{\displaystyle \frac{2}{2^N}}\right)\left(1{\displaystyle \frac{2}{2^M}}\right)}.`$ (16)
As expected the function $`𝒬_{\mathrm{opt}}`$ is one for $`N=M`$, so there is no error. For $`N1`$ deleting (meaning keeping a single copy and deleting $`(N1)`$ copies, the quality function $`𝒬_{\mathrm{opt}}=\frac{1}{2^{(N1)/2}}`$. Interestingly, the quality goes down exponentially with number of copies we would like to delete, hence the error increases. Therefore, it is difficult to delete more and more number of copies. From (12) it also follows that if we are given an infinite number of copies and asked to delete a single copy, then we can accomplish it without any error. This is in accordance with our fundamental understanding about quantum information. For $`2`$-to-$`1`$ deleting the optimal quality is $`0.70`$.
Conditional quantum deleting of qubits: To study how the quantum information is distributed among various subsystems and imperfections introduced during a deletion process, we introduce a special class of deleting machines the so-called conditional deleting. If the two input qubits are identical then machine deletes a copy and if they are different then it allows them to pass through without any change. For orthogonal qubits we define it as $`|0|0|A|0|\mathrm{\Sigma }|A_0,|1|1|A|1|\mathrm{\Sigma }|A_1,|0|1|A|0|1|A,`$ and $`|1|0|A|1|0|A`$, where $`|A`$ is the initial state and $`|A_0,|A_1`$ are the final states of ancilla. Notice that if the conditional deletion for orthogonal qubits has to work it is necessary to include the ancilla. Then for an arbitrary qubit the deleting operation will create the following state
$`|\mathrm{\Psi }|\mathrm{\Psi }|A=[\alpha ^2|00+\beta ^2|11+\alpha \beta (|01+|10)]|A`$ (17)
$`\alpha ^2|0|\mathrm{\Sigma }|A_0+\beta ^2|1|\mathrm{\Sigma }|A_1+\alpha \beta (|01+|10)|A`$ (18)
$`=|\mathrm{\Psi }_{\mathrm{out}}.`$ (19)
However, ideally for an arbitrary qubit the deleting machine should have created $`|\mathrm{\Psi }|\mathrm{\Psi }|A|\mathrm{\Psi }|\mathrm{\Sigma }|A_\mathrm{\Psi }`$. Since the final output states in ideal and actual cases are different linearity does not allow to conditionally delete an unknown quantum state.
State-dependent and approximate deleting machine: Since it is impossible to delete an unknown state perfectly we may ask how well one can do the above operation. Here, we discuss the approximate deleting of an unknown state and the fidelity of a state-dependent quantum deleting machine. In (13) if we wish to delete a qubit then the ancilla state should belong to a three dimensional Hilbert space. For a normalised output state in (13) (and hence unitary) we need $`|A,|A_0,`$ and $`|A_1`$ to be orthogonal to each other. The reduced density matrix of the two qubits $`ab`$ after the deleting operation is given by
$`\rho _{ab}=\mathrm{tr}_c(|\mathrm{\Psi }_{\mathrm{out}}\mathrm{\Psi }_{\mathrm{out}}|)=|\alpha |^4|00||\mathrm{\Sigma }\mathrm{\Sigma }|`$ (20)
$`+|\beta |^4|11||\mathrm{\Sigma }\mathrm{\Sigma }|+2|\alpha |^2|\beta |^2|\psi ^+\psi ^+|,`$ (21)
where $`|\psi ^+=\frac{1}{\sqrt{2}}(|01+|10)`$ is one of the four maximally entangled states. The reduced density matrix for the qubit in the mode $`b`$ will be
$`\rho _b=\mathrm{tr}_a(\rho _{ab})=(12|\alpha |^2|\beta |^2)|\mathrm{\Sigma }\mathrm{\Sigma }|+|\alpha |^2|\beta |^2I,`$ (22)
where $`I`$ is the identity matrix in two dimensional Hilbert space. Thus the reduced density matrix of the qubit in the mode $`b`$ is a mixed state which contains the error due to imperfect deleting. The fidelity of deleting can be defined as $`F_b=\mathrm{\Sigma }|\rho _b|\mathrm{\Sigma }=(1|\alpha |^2|\beta |^2)`$. This shows that for either $`\alpha =0`$ and $`\beta =1`$ or $`\alpha =1`$ and $`\beta =0`$ the fidelity of deleting is maximum. For an equal superposition of qubit state the fidelity reaches $`\frac{3}{4}`$ which is the maximum limit for deleting an unknown qubit. The average fidelity of deleting is given by $`\overline{F}_b=𝑑\mathrm{\Omega }F_b=\frac{5}{6}0.83`$, where $`\mathrm{\Omega }=\mathrm{sin}\theta d\theta d\varphi `$.
We can see how good the state of the qubit in mode $`a`$ is after both the qubits have passed through a quantum deleting machine. The reduced density matrix of this mode is given by
$`\rho _a=\mathrm{tr}_b(\rho _{ab})=|\alpha |^4|00|+|\beta |^4|11|+|\alpha |^2|\beta |^2I.`$ (23)
The fidelity of the qubit in mode $`a`$ is $`F_a=\mathrm{\Psi }|\rho _a|\mathrm{\Psi }=(12|\alpha |^2|\beta |^2)`$. For an equal superposition of qubit state the fidelity is $`\frac{1}{2}`$. The average fidelity in this case is $`\overline{F}_a=𝑑\mathrm{\Omega }F_a=\frac{2}{3}0.66`$. This shows that the first mode of the qubit is not faithfully retained during the deleting operation. It is, in fact, less than the actual deleting mode. This shows that linearity of quantum theory neither allows us to delete an unknown state perfectly nor does retain the original state of the other qubit. We can compare the quantum deleting operation to that of the quantum cloning operation defined by Wootters and Zurek . In the cloning operation the reduced density matrix of both the modes are same . Therefore, the average fidelities of both the modes are found to be $`\frac{2}{3}`$. However, as we have shown here the fidelity of the two modes are different for the deleting operation. This again suggests us that the quantum deleting machine is not the reverse of quantum cloning machine.
Quantum deleting of non-orthogonal states: In some physical situations two qubits need not be in orthogonal states nor in completely arbitrary states but they could be chosen secretly from a set containing non-orthogonal states. Though each of a copy from two copies of two orthogonal states can be perfectly deleted, we show here that the same cannot be done for two non-orthogonal states. Suppose we have two copies of the two of the non-orthogonal states $`|\mathrm{\Psi }_i,(i=1,2)`$ with a finite scalar product between them. We ask if there is a unitary operator which can delete one of the copy by keeping the other intact. For simplicity and clarity we work without attaching an ancilla to the qubits. For two copies of distinct non-orthogonal states the deleting machine is a unitary operator which acts on the combined Hilbert space of two qubits and would create the following transformation: $`|\mathrm{\Psi }_1|\mathrm{\Psi }_1|\mathrm{\Psi }_1|\mathrm{\Sigma },|\mathrm{\Psi }_2|\mathrm{\Psi }_2|\mathrm{\Psi }_2|\mathrm{\Sigma },|\mathrm{\Psi }_1|\mathrm{\Psi }_2|\mathrm{\Psi }_1|\mathrm{\Psi }_2,`$ and $`|\mathrm{\Psi }_2|\mathrm{\Psi }_1|\mathrm{\Psi }_2|\mathrm{\Psi }_1`$. Since unitary evolution must preserve the inter inner products we have several conditions to be satisfied simultaneously. These restriction are $`\mathrm{\Psi }_1|\mathrm{\Psi }_2^2=\mathrm{\Psi }_1|\mathrm{\Psi }_2`$, $`\mathrm{\Psi }_1|\mathrm{\Psi }_2=\mathrm{\Sigma }|\mathrm{\Psi }_2`$, $`\mathrm{\Sigma }|\mathrm{\Psi }_2=1`$, $`\mathrm{\Sigma }|\mathrm{\Psi }_1=1`$, and $`\mathrm{\Psi }_2|\mathrm{\Psi }_1=\mathrm{\Sigma }|\mathrm{\Psi }_1`$. These can be satisfied only if $`|\mathrm{\Psi }_1=|\mathrm{\Psi }_2=|\mathrm{\Sigma }`$, which means a contradiction as there are no non-trivial states being processed by the machine. Thus, copies of non-orthogonal states cannot be deleted by a unitary machine.
Perfect deleting and signalling: We know that if we could clone an arbitrary state we can send superluminal signals . The natural question is whether no-deleting principle is also consistent with no-signalling. At prima face, it looks that it may not be so. But on the other hand linear structure of quantum mechanics and no-signalling are consistent with each other and no-deleting is a consequence of linearity. So there should be some nontrivial link between no-deleting and no-signalling. Below we show how this works.
Suppose Alice and Bob share two pairs of EPR singlets and Alice has particles $`1`$ and $`3`$ and Bob has $`2`$ and $`4`$. Since the singlet state is invariant under local unitary operation $`UU`$ it is same in all basis. Let us write the combined state of the system in an arbitrary qubit basis $`\{|\psi =\mathrm{cos}\theta |0+\mathrm{sin}\theta |1,|\overline{\psi }=\mathrm{sin}\theta |0\mathrm{cos}\theta |1\}`$
$`|\mathrm{\Psi }^{}_{12}|\mathrm{\Psi }^{}_{34}={\displaystyle \frac{1}{2}}(|\psi _1|\psi _3|\overline{\psi }_2|\overline{\psi }_4+|\overline{\psi }_1|\overline{\psi }_3|\psi _2|\psi _4`$ (25)
$`|\overline{\psi }_1|\psi _3|\psi _2|\overline{\psi }_4|\psi _1|\overline{\psi }_3|\overline{\psi }_2|\psi _4).`$ (26)
Now if Alice measures her particles $`1`$ and $`3`$ onto the basis $`|\psi _1|\psi _3`$, then the Bob’s particles $`2`$ and $`4`$ are in the state $`|\overline{\psi }_2|\overline{\psi }_4`$. If Alice measures her particles in the basis $`|\overline{\psi }_1|\overline{\psi }_3`$, then Bob’s particles are in the state $`|\psi _2|\psi _4`$. Similarly, one can find the resulting states with other choices of measurements. So whatever measurements Alice does, if she does not convey the measurement result to Bob, then Bob’s particles are in a completely random mixture (i.e. $`\rho _{24}=\frac{I_2}{2}\frac{I_2}{2}`$ ). But suppose Bob has a conditional quantum deleting machine, which can delete an arbitrary state. Then, after Alice does measurement he attaches an ancilla and deletes his copies. The four possible choices for the states of the four particles (with ancilla) are give by
$`|\overline{\psi }_1|\overline{\psi }_3|\psi _2|\psi _4|A|\overline{\psi }_1|\overline{\psi }_3|\psi _2|\mathrm{\Sigma }_4|A_\psi `$ (27)
$`|\psi _1|\psi _3|\overline{\psi }_2|\overline{\psi }_4|A|\psi _1|\psi _3|\overline{\psi }_2|\mathrm{\Sigma }_4|A_{\overline{\psi }}`$ (28)
$`|\overline{\psi }_1|\psi _3|\psi _2|\overline{\psi }_4|A|\overline{\psi }_1|\psi _3|\psi _2|\overline{\psi }_4|A`$ (29)
$`|\psi _1|\overline{\psi }_3|\overline{\psi }_2|\psi _4|A|\psi _1|\overline{\psi }_3|\overline{\psi }_2|\psi _4|A.`$ (30)
If we trace out ancilla and particles $`1`$ and $`3`$ the reduced density matrix for Bob’s particles $`2`$ and $`4`$ is given by
$`\rho _{24}={\displaystyle \frac{1}{4}}(|\psi _2{}_{2}{}^{}\psi ||\mathrm{\Sigma }_{4}^{}{}_{4}{}^{}\mathrm{\Sigma }|+|\overline{\psi }_{2}^{}{}_{2}{}^{}\overline{\psi }||\mathrm{\Sigma }_{4}^{}{}_{4}{}^{}\mathrm{\Sigma }|+`$ (31)
$`|\psi _2_2\psi ||\overline{\psi }_4_4\overline{\psi }|+|\overline{\psi }_2_2\overline{\psi }||\psi _4_4\psi |)`$ (32)
which clearly depends on the choice of basis. This shows that if Alice measures her particles in $`\{|0,|1\}`$ basis then Bob’s particles are in one density matrix. If Alice measures her particles in $`\{|+,|\}`$ basis Bob’s particles are left in a different density matrix. Therefore, if Bob can delete an arbitrary state he can distinguish two statistical mixtures and will allow superluminal signalling. Hence, Bob cannot delete an arbitrary state. Therefore, quantum no-deleting principle is in unison with the principle of no-signalling.
Quantum no-deletion principle being a fundamental limitation on quantum information it ought to have some implications. For example, this may provide special security to copies of files in a quantum computer and possibly in quantum cryptographic protocols– which deserve further investigation in the future. However, constructing a universal quantum deleting machine and obtaining the optimal fidelity of quantum deletion is still an open problem.
Financial support from EPSRC is gratefully acknowledged. We thank useful discussions at various stages with C. Bennett, S. Bose, N. Cerf, A. Chefles and N. Gisin.
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# Initial Data for Numerical Relativity
## 1 Introduction
The goal of numerical relativity is to study spacetimes that cannot be studied by analytic means. The focus is therefore primarily on dynamical systems. Numerical relativity has been applied in many areas: cosmological models, critical phenomena, perturbed black holes and neutron stars, and the coalescence of black holes and neutron stars, for example. In any of these cases, Einstein’s equations can be formulated in several ways that allow us to evolve the dynamics. While Cauchy methods have received a majority of the attention, characteristic and Reggi calculus based methods have also been used. All of these methods begin with a snapshot of the gravitational fields on some hypersurface, *the initial data*, and evolve these data to neighboring hypersurfaces.
The focus of this review is on the initial data needed for Cauchy evolutions of Einstein’s equations. These initial data cannot be freely specified in their entirety. Rather they are subject to certain constraints which must be satisfied. Because of the nonlinearity of Einstein’s equations, there is no unique way of choosing which pieces of the initial data can be freely specified and which are constrained. In §2, I will look at the various formalisms that exist for expressing the initial-value equations. I will use the $`3+1`$ (or ADM) decomposition of Einstein’s equations throughout and I begin the review with a brief introduction of this in §2.1. In §2.2, §2.3, and §2.4, I will explore the most important and widely used decompositions that have been developed.
In the remainder of the review, I will focus on initial data for black holes and neutron stars. Section 3 deals with black-hole initial data, which has received considerable attention over the years. The evolution of black-hole spacetimes is of particular importance because it allows for the study of pure geometrodynamics. The spacetime is either pure vacuum or any matter can be hidden behind an event horizon. Thus, black-hole spacetimes allow us to study the dynamics of gravity alone. In §3.1, I begin with a review of some of the classic analytic solutions of the initial-data equations. Section 3.2 covers current methods for constructing general multi-black-hole initial data sets and also explores some of the limitations of these data sets. Section 3.3 deals with a class of black-hole initial data that has received little attention until recently and which I feel may play an important role in future work on constructing black-hole initial data. Finally, §3.4 examines the issue of constructing black-hole initial data for quasiequilibrium binaries.
This last topic is of extreme importance. Several laser interferometer gravitational wave detectors will become operational in the near future and the coalescence of black-hole binaries is considered to be one of the strongest candidates for detection by the earliest generation of detectors. The chances of detecting these events and then unraveling the information contained in the gravitational wave signals will be greatly increased if we have accurate numerical simulations of black-hole binary coalescence. While post-Newtonian techniques can be used to simulate the inspiral of a compact binary system, the final plunge and coalescence must be simulated numerically. Astrophysically realistic initial data will be needed before these simulations can provide reliable results.
The final plunge and coalescence of neutron-star binaries must also be simulated numerically. Section 4 examines the construction of neutron-star initial data. The major difference between black-hole and neutron-star initial data is the need to deal with the neutron star’s matter. In this section, I will deal with the neutron star matter in general, without considering any particular equation of state. In §4.1, I look at the issue of hydrostatic equilibrium of matter. Section 4.2 takes a brief look at constructing equilibrium initial data for isolated neutron stars. Finally, §4.3 examines the issues involved in constructing quasiequilibrium neutron-star binary initial data. In particular, we look at recent advances in constructing irrotational fluid models.
This article is far from a complete review of all the work that has been done on initial data for numerical relativity. I offer my apologies to all whose work I have not included. In particular, I have not covered any of the issues associated with existence and uniqueness of solutions to the constraint equations, nor dealt extensively with the issue of asymptotic falloff rates. I welcome correspondence on topics that you feel should be covered in this review, references that you think should be included, and the inevitable typographical errors.
### 1.1 Conventions
I will use a 4-metric $`g_{\mu \nu }`$ with signature $`(+++)`$. Following the MTW conventions, I define the Riemann tensor as
$$R^\mu {}_{\nu \alpha \beta }{}^{}_\alpha \mathrm{\Gamma }^\mu {}_{\nu \beta }{}^{}_\beta \mathrm{\Gamma }^\mu {}_{\nu \alpha }{}^{}+\mathrm{\Gamma }^\mu {}_{\sigma \alpha }{}^{}\mathrm{\Gamma }_{}^{\sigma }{}_{\nu \beta }{}^{}\mathrm{\Gamma }^\mu {}_{\sigma \beta }{}^{}\mathrm{\Gamma }_{}^{\sigma }{}_{\nu \alpha }{}^{}.$$
(1)
The Ricci tensor is defined as
$$R_{\mu \nu }R^\sigma {}_{\mu \sigma \nu }{}^{},$$
(2)
and the Einstein tensor as
$$G_{\mu \nu }R_{\mu \nu }\frac{1}{2}g_{\mu \nu }R.$$
(3)
A spatial 3-metric will be written as $`\gamma _{ij}`$ and the Riemann and Ricci tensors associated with it will be defined by (1) and (2). Four-dimensional indices will be denoted by Greek letters, while 3-dimensional indices will be denoted by Latin letters.
To avoid confusion, the covariant derivative and Ricci tensor associated with $`\gamma _{ij}`$ will be written with over-bars—$`\overline{}_j`$ and $`\overline{R}_{ij}`$. I will also frequently deal with an auxiliary 3-dimensional space with a metric that is conformally related to the metric $`\gamma _{ij}`$ of the physical space. The metric for this space will be denoted $`\stackrel{~}{\gamma }_{ij}`$. The covariant derivative and Ricci tensor associated with this metric will be written with tildes—$`\stackrel{~}{}_j`$ and $`\stackrel{~}{R}_{ij}`$. Other quantities that have a conformal relationship to quantities in the physical space will also be written with a tilde over them.
## 2 The Initial-Value Equations
Einstein’s equations, $`G_{\mu \nu }=8\pi GT_{\mu \nu }`$, represent ten independent equations. Since there are ten equations and ten independent components of the four-metric $`g_{\mu \nu }`$, it seems that we have the same number of equations as unknowns. From the definition of the Einstein tensor (3), we see that these ten equations are linear in the second derivatives and quadratic in the first derivatives of the metric. We might expect that these ten second-order equations represent evolution equations for the ten components of the metric. However, a close inspection of the equations reveals that only six of the ten involve second time-derivatives of the metric. The remaining four equations are not evolution equations. Instead, they are constraint equations. The full system of equations is still well posed, however, because of the Bianchi identities
$$_\nu G^{\mu \nu }0.$$
(4)
The four constraint equations appear as a result of the general covariance of Einstein’s theory, which gives us the freedom to apply general coordinate transformations to each of the four coordinates and leave the interval,
$$\mathrm{d}s^2=g_{\mu \nu }\mathrm{d}x^\mu \mathrm{d}x^\nu ,$$
(5)
unchanged.
If we consider Einstein’s equations as a Cauchy problem<sup>1</sup><sup>1</sup>1We could also formulate Einstein’s equations as a characteristic initial-value problem, but we will not pursue that approach in this paper., we find that the ten equations separate into a set of four constraint or initial-value equations, and six evolution or dynamical equations. If the four initial-value equations are satisfied on some spacelike hypersurface, which we can label with $`t=0`$, then the Bianchi identities (4) guarantee that the evolution equations preserve the constraints on neighboring spacelike hypersurfaces.
### 2.1 Initial Data
In the Cauchy formulation of Einstein’s equations, we begin by foliating the 4-dimensional manifold as a set of spacelike, 3-dimensional hypersurfaces (or slices) $`\{\mathrm{\Sigma }\}`$. These slices are labeled by a parameter $`t`$ or, more simply, each slice of the 4-dimensional manifold is a $`t=\text{constant}`$ hypersurface. Following the standard $`3+1`$ decomposition , we let $`n^\mu `$ be the future-pointing timelike unit normal to the slice, with
$$n^\mu \alpha ^\mu t.$$
(6)
Here, $`\alpha `$ is called the *lapse function* (frequently denoted $`N`$ in the literature). The scalar lapse function sets the proper interval measured by observers as they move between slices on a path that is normal to the hypersurface (so-called normal observers):
$$\mathrm{d}s|_{\text{along }n^\mu }=\alpha \mathrm{d}t.$$
(7)
Of course, there is no reason that observers must move along a path normal to the hypersurface. In general, we can define the time vector as
$$t^\mu \alpha n^\mu +\beta ^\mu ,$$
(8)
where
$$\beta ^\mu n_\mu 0.$$
(9)
Here, $`\beta ^\mu `$ is called the *shift vector* (frequently denoted $`N^\mu `$ in the literature).
Because of (9), $`\beta ^\mu `$ has only three independent components and is a *spatial* vector, tangent to the hypersurface on which it resides. At this point, it is convenient to introduce a coordinate system adapted to the foliation $`\{\mathrm{\Sigma }\}`$. Let $`x^i`$ be the spatial coordinates in the slice. The fourth coordinate, $`t`$, is the parameter labeling each slice. With this adapted coordinate system, we find that 3-dimensional coordinate values remain constant as we move between slices along the $`t^\mu `$ direction (8). The four parameters, $`\alpha `$ and $`\beta ^i`$, are a manifestation of the 4-dimensional coordinate invariance, or gauge freedom, in Einstein’s theory. If we let $`\gamma _{ij}`$ represent the metric of the spacelike hypersurfaces, then we can rewrite the interval (5) as
$$\mathrm{d}s^2=\alpha ^2\mathrm{d}t^2+\gamma _{ij}(\mathrm{d}x^i+\beta ^i\mathrm{d}t)(\mathrm{d}x^j+\beta ^j\mathrm{d}t).$$
(10)
In the Cauchy formulation of Einstein’s equations, $`\gamma _{ij}`$ is regarded as the fundamental variable and values for its components must be given as part of a well-posed initial-value problem. Since Einstein’s equations are second order, we must also specify something like a time derivative of the metric. For this, we use the second fundamental form, or extrinsic curvature, of the slice, $`K_{ij}`$, defined<sup>2</sup><sup>2</sup>2A different sign choice for defining the extrinsic curvature is sometimes found in the literature. by
$$K_{ij}\frac{1}{2}_n\gamma _{ij},$$
(11)
where $`_n`$ denotes the Lie derivative along the $`n^\mu `$ direction.
Together, $`\gamma _{ij}`$ and $`K_{ij}`$ are the minimal set of initial data that must be specified for a Cauchy evolution of Einstein’s equations. The metric $`\gamma _{ij}`$ on a hypersurface is induced on that surface by the 4-metric $`g_{\mu \nu }`$. This means that the values $`\gamma _{ij}`$ receives depend on how $`\mathrm{\Sigma }`$ is embedded in the full spacetime. In order for the foliation of slices $`\{\mathrm{\Sigma }\}`$ to fit into the higher-dimensional space, they must satisfy the Gauss–Codazzi–Ricci conditions. Combining these conditions with Einstein’s equations, and using (10), the six evolution equations become
$`_tK_{ij}`$ $`=`$ $`\alpha \left[\overline{R}_{ij}2K_i\mathrm{}K_j^{\mathrm{}}+KK_{ij}8\pi GS_{ij}+4\pi G\gamma _{ij}(S\rho )\right]`$
$`\overline{}_i\overline{}_j\alpha +\beta ^{\mathrm{}}\overline{}_{\mathrm{}}K_{ij}+K_i\mathrm{}\overline{}_j\beta ^{\mathrm{}}+K_j\mathrm{}\overline{}_i\beta ^{\mathrm{}}.`$
Here, $`\overline{}_i`$ is the spatial covariant derivative compatible with $`\gamma _{ij}`$, $`\overline{R}_{ij}`$ is the Ricci tensor associated with $`\gamma _{ij}`$, $`KK_i^i`$, $`\rho `$ is the matter energy density, $`S_{ij}`$ is the matter stress tensor, and $`SS_i^i`$.<sup>3</sup><sup>3</sup>3 The stress-energy tensor is decomposed as $`T_{\mu \nu }=S_{\mu \nu }+n_\mu j_\nu +n_\nu j_\mu +n_\mu n_\nu \rho `$ or, equivalently, $`S_{ij}T_{ij}`$, $`j_in^\mu T_{i\mu }`$, and $`\rho n^\mu n^\nu T_{\mu \nu }`$. Note that the matter terms in Eqs. (2.1), (14), and (15) are defined with respect to a normal observer. This is in contrast to the usual definition of $`\rho `$ with respect to the rest frame of the fluid in hydrodynamics. We have also used the fact that in our adapted coordinate system, $`_t_t`$. The set of second-order evolution equations is completed by rewriting the definition of the extrinsic curvature (11) as
$$_t\gamma _{ij}=2\alpha K_{ij}+\overline{}_i\beta _j+\overline{}_j\beta _i.$$
(13)
Equations (2.1) and (13) are a first-order representation of a complete set of evolution equations for given initial data $`\gamma _{ij}`$ and $`K_{ij}`$. However, the data cannot be freely specified in their entirety. The four constraint equations, following the same procedure outlined above for the evolution equations, become
$$\overline{R}+K^2K_{ij}K^{ij}=16\pi G\rho $$
(14)
and
$$\overline{}_j\left(K^{ij}\gamma ^{ij}K\right)=8\pi Gj^i.$$
(15)
Here, $`\overline{R}\overline{R}_i^i`$ and $`j^i`$ is the matter momentum density. Equation (14) is referred to as the Hamiltonian or scalar constraint, while (15) are referred to as the momentum or vector constraints. Valid initial data for the evolution equations (2.1) and (13) must satisfy this set of constraints. And, as mentioned earlier, the Bianchi identities (4) guarantee that the evolution equations will preserve the constraints on future slices of the evolution.
As we will see below, the Hamiltonian constraint (14) most naturally constrains the 3-metric $`\gamma _{ij}`$, while the momentum constraints (15) naturally constrain the extrinsic curvature $`K_{ij}`$. Taking the constraints into consideration, it seems that the 3-metric has five degrees of freedom remaining, while the extrinsic curvature has three. But we know that the gravitational field in Einstein’s theory has two dynamical degrees of freedom, so we expect that both $`\gamma _{ij}`$ and $`K_{ij}`$ should each have only two free components. The answer to this problem is, once again, the coordinate invariance of Einstein’s theory. This seems strange at first, because we have already used the coordinate invariance of the theory to narrow our scope from the ten components of the 4-metric to the six components of the 3-metric. However, the lapse and shift do not completely specify the coordinate gauge. Rather, they specify how an initial choice of gauge will evolve with the foliation. The metric on a given hypersurface retains full 3-dimensional coordinate invariance, reducing the number of freely specifiable components to two . There also remains one degree of gauge freedom associated with the time coordinate which must be fixed. Each hypersurface represents a $`t=\text{const.}`$ slice of the spacetime, so how the initial hypersurface is embedded in the full 4-dimensional manifold represents our temporal gauge choice. There is no unique way to specify this choice, but it is often convenient to let the trace of the extrinsic curvature $`K`$ represent this temporal gauge choice . Thus, we find that we *are* allowed to choose freely five components of the 3-metric and three components of the extrinsic curvature. However, only *two* of the components for each field represent dynamical degrees of freedom, the remainder are gauge degrees of freedom.
The four constraint equations, (14) and (15), represent conditions which the 3-metric and extrinsic curvature must satisfy. But, they do not specify which components (or combination of components) are constrained and which are freely specifiable. In the weak field limit where Einstein’s equations can be linearized, there are clear ways to determine which components are dynamic, which are constrained, and which are gauge. However, in the full nonlinear theory, there is no unique decomposition. In this case, one must choose a method for decomposing the constraint equations. The goal is to transform the equations into standard elliptic forms which can be solved given appropriate boundary conditions . Each different decomposition yields a unique set of elliptic equations to be solved *and* a unique set of freely specifiable parameters which must be fixed somehow. Seemingly similar sets of assumptions applied to different decompositions can lead to physically different initial conditions.
### 2.2 York–Lichnerowicz Conformal Decompositions
For general initial-data configurations, the most widely used class of constraint decompositions are the York–Lichnerowicz conformal decompositions. At their heart are a conformal decomposition of the metric and certain components of the extrinsic curvature, together with a transverse-traceless decomposition of the extrinsic curvature.
First, the metric is decomposed into a conformal factor $`\psi `$ multiplying an auxiliary 3-metric :
$$\gamma _{ij}\psi ^4\stackrel{~}{\gamma }_{ij}.$$
(16)
The auxiliary 3-metric $`\stackrel{~}{\gamma }_{ij}`$ is often called the conformal or background 3-metric, and it carries five degrees of freedom. Its natural definition is given by
$$\stackrel{~}{\gamma }_{ij}=\gamma ^{\frac{1}{3}}\gamma _{ij},$$
(17)
leaving $`\stackrel{~}{\gamma }=1`$, but we are free to choose any normalization for $`\stackrel{~}{\gamma }`$. Using (16), we can rewrite the Hamiltonian constraint (14) as
$$\stackrel{~}{}^2\psi \frac{1}{8}\psi \stackrel{~}{R}\frac{1}{8}\psi ^5K^2+\frac{1}{8}\psi ^5K_{ij}K^{ij}=2\pi G\psi ^5\rho ,$$
(18)
where $`\stackrel{~}{}^2\stackrel{~}{}^i\stackrel{~}{}_i`$ is the scalar Laplace operator, and $`\stackrel{~}{}_i`$ and $`\stackrel{~}{R}`$ are the covariant derivative and Ricci scalar associated with $`\stackrel{~}{\gamma }_{ij}`$. Equation (18) is a quasilinear elliptic equation for the conformal factor $`\psi `$, and we see that the Hamiltonian constraint naturally constrains the 3-metric.
The conformal decomposition of the Hamiltonian constraint was proposed by Lichnerowicz. But, the key to the full decomposition is the treatment of the extrinsic curvature introduced by York . This begins by splitting the extrinsic curvature into its trace and tracefree parts,
$$K_{ij}A_{ij}+\frac{1}{3}\gamma _{ij}K.$$
(19)
The decomposition proceeds by using the fact that we can covariantly split any symmetric tracefree tensor as follows:
$$𝒮^{ij}(𝕃X)^{ij}+T^{ij}.$$
(20)
Here, $`T^{ij}`$ is a symmetric, transverse-traceless tensor (i.e., $`_jT^{ij}=0`$ and $`T_i^i=0`$) and
$$(𝕃X)^{ij}^iX^j+^jX^i\frac{2}{3}\gamma ^{ij}_{\mathrm{}}X^{\mathrm{}}.$$
(21)
After separating out the transverse-traceless portion of $`𝒮^{ij}`$, what remains, $`(𝕃X)^{ij}`$, is referred to as its “longitudinal” part. We now want to apply this transverse-traceless decomposition to the tracefree part of the extrinsic curvature $`A_{ij}`$. However, the conformal decomposition of the metric leaves us with at least two ways to proceed.
The goal of the decomposition is to produce a coupled set of elliptic equations to be solved with some prescribed boundary conditions. We have already reduced the Hamiltonian constraint to an elliptic equation being solved on a *background* space in terms of differential operators that are compatible with the conformal 3-metric. In the end, we want to reduce the momentum constraints to a set of elliptic equations based on differential operators that are compatible with the same conformal 3-metric. But, the longitudinal operator (21) can be defined with respect to *any* metric space. In particular, it is natural to consider decompositions with respect to both the physical and conformal 3-metrics.
#### 2.2.1 Conformal Transverse-Traceless Decomposition
Let us first consider decomposing $`A^{ij}`$ with respect to the conformal 3-metric . As we will see, when certain assumptions are made, this decomposition has the advantage of producing a simpler set of elliptic equations that must be solved. The first step is to define the conformal tracefree extrinsic curvature $`\stackrel{~}{A}^{ij}`$ by
$$A^{ij}\psi ^{10}\stackrel{~}{A}^{ij}\text{or}A_{ij}\psi ^2\stackrel{~}{A}_{ij}.$$
(22)
Next, the transverse-traceless decomposition is applied to the conformal extrinsic curvature,
$$\stackrel{~}{A}^{ij}(\stackrel{~}{𝕃}X)^{ij}+\stackrel{~}{Q}^{ij}.$$
(23)
Note that the longitudinal operator $`\stackrel{~}{𝕃}`$ and the symmetric, transverse-tracefree tensor $`\stackrel{~}{Q}^{ij}`$ are both defined with respect to covariant derivatives compatible with $`\stackrel{~}{\gamma }_{ij}`$.
Applying equations (16), (19), (21), (22), and (23) to the momentum constraints (15), we find that they simplify to
$$\stackrel{~}{\mathrm{\Delta }}_𝕃X^i=\frac{2}{3}\psi ^6\stackrel{~}{}^iK+8\pi G\psi ^{10}j^i,$$
(24)
where
$$\stackrel{~}{\mathrm{\Delta }}_𝕃X^i\stackrel{~}{}_j(𝕃X)^{ij}=\stackrel{~}{}^2X^i+\frac{1}{3}\stackrel{~}{}^i(\stackrel{~}{}_jX^j)+\stackrel{~}{R}_j^iX^j,$$
(25)
and we have used the fact that
$$\overline{}_j𝒮^{ij}=\psi ^{10}\stackrel{~}{}_j(\psi ^{10}𝒮^{ij})$$
(26)
for any symmetric tracefree tensor $`𝒮^{ij}`$.
In deriving equation (24), we have also used the fact that $`\stackrel{~}{Q}^{ij}`$ is transverse (i.e. $`\stackrel{~}{}_j\stackrel{~}{Q}^{ij}=0`$). However, in general, we will not know if a given symmetric tracefree tensor, say $`\stackrel{~}{M}^{ij}`$, is transverse. By using (20) we can obtain its transverse-traceless part $`\stackrel{~}{Q}^{ij}`$ via
$$\stackrel{~}{Q}^{ij}\stackrel{~}{M}^{ij}(\stackrel{~}{𝕃}Y)^{ij},$$
(27)
and using the fact that if $`\stackrel{~}{Q}^{ij}`$ is transverse, we find
$$\stackrel{~}{}_j\stackrel{~}{Q}^{ij}0=\stackrel{~}{}_j\stackrel{~}{M}^{ij}\stackrel{~}{\mathrm{\Delta }}_𝕃Y^i.$$
(28)
Thus, Eqs. (27) and (28) give us a general way of *constructing* the required symmetric transverse-traceless tensor from a general symmetric traceless tensor.
Using the linearity of $`\stackrel{~}{𝕃}`$, we can rewrite (23) as
$$\stackrel{~}{A}^{ij}=(\stackrel{~}{𝕃}V)^{ij}+\stackrel{~}{M}^{ij},$$
(29)
where
$$V^iX^iY^i.$$
(30)
Similarly, using the linearity of $`\stackrel{~}{\mathrm{\Delta }}_𝕃`$, we can rewrite (24) as
$$\stackrel{~}{\mathrm{\Delta }}_𝕃V^i=\frac{2}{3}\psi ^6\stackrel{~}{}^iK\stackrel{~}{}_j\stackrel{~}{M}^{ij}+8\pi G\psi ^{10}j^i.$$
(31)
By solving directly for $`V^i`$, we can combine the steps of decomposing $`\stackrel{~}{M}^{ij}`$ with that of solving the momentum constraints.
After applying (19) and (22) to the Hamiltonian constraint (18), we obtain the following full decomposition, which I will list together here for convenience:
$`\gamma _{ij}`$ $`=`$ $`\psi ^4\stackrel{~}{\gamma }_{ij}`$
$`K^{ij}`$ $`=`$ $`\psi ^{10}\stackrel{~}{A}^{ij}+\frac{1}{3}\psi ^4\stackrel{~}{\gamma }^{ij}K`$
$`\stackrel{~}{A}^{ij}`$ $`=`$ $`(\stackrel{~}{𝕃}V)^{ij}+\stackrel{~}{M}^{ij}`$ (32)
$`\stackrel{~}{\mathrm{\Delta }}_𝕃V^i\frac{2}{3}\psi ^6\stackrel{~}{}^iK`$ $`=`$ $`\stackrel{~}{}_j\stackrel{~}{M}^{ij}+8\pi G\psi ^{10}j^i`$
$`\stackrel{~}{}^2\psi \frac{1}{8}\psi \stackrel{~}{R}\frac{1}{12}\psi ^5K^2+\frac{1}{8}\psi ^7\stackrel{~}{A}_{ij}\stackrel{~}{A}^{ij}`$ $`=`$ $`2\pi G\psi ^5\rho `$
In the decomposition given by (2.2.1), we are free to specify a symmetric tensor $`\stackrel{~}{\gamma }_{ij}`$ as the conformal 3-metric, a symmetric tracefree tensor $`\stackrel{~}{M}^{ij}`$, and a scalar function $`K`$. Then, with given matter energy and momentum densities, $`\rho `$ and $`j^i`$, and appropriate boundary conditions, the coupled set of constraint equations for $`\psi `$ and $`V^i`$ are solved. Finally, given the solutions, we can construct the physical initial data, $`\gamma _{ij}`$ and $`K^{ij}`$.
The decomposition outlined above has the interesting property that if we choose $`K`$ to be constant and if the momentum density vanishes<sup>4</sup><sup>4</sup>4Or, if we solve the momentum constraint in terms of a background momentum density $`\stackrel{~}{j}^i\psi ^{10}j^i`$, the momentum constraint decouples with nonvanishing $`\stackrel{~}{j}^i`$., then the momentum constraint equations fully decouple from the Hamiltonian constraint. As we will see later, this simplification has proven to be useful.
#### 2.2.2 Physical Transverse-Traceless Decomposition
Alternatively, we can decompose $`A^{ij}`$ with respect to the physical 3-metric . We decompose the extrinsic curvature as
$$A^{ij}(\overline{𝕃}W)^{ij}+Q^{ij}.$$
(33)
In this case, the longitudinal operator $`\overline{𝕃}`$ and the symmetric transverse-tracefree tensor $`Q^{ij}`$ are both defined with respect to covariant derivatives compatible with $`\gamma _{ij}`$.
Applying equations (16), (19), (21), (33), and (26) to the momentum constraint (15), we find that it simplifies to
$$\stackrel{~}{\mathrm{\Delta }}_𝕃W^i+6(\stackrel{~}{𝕃}W)^{ij}\stackrel{~}{}_j\mathrm{ln}\psi =\frac{2}{3}\stackrel{~}{}^iK+8\pi G\psi ^4j^i,$$
(34)
where we have used the fact that
$$(\overline{𝕃}W)^{ij}=\psi ^4(\stackrel{~}{𝕃}W)^{ij}.$$
(35)
As in the previous section, we will obtain the symmetric transverse-traceless tensor $`Q^{ij}`$ from a general symmetric tracefree tensor $`\stackrel{~}{M}^{ij}`$ by using (20). In this case, we take
$$Q^{ij}\psi ^{10}\stackrel{~}{M}^{ij}(\overline{𝕃}Z)^{ij},$$
(36)
and use the fact that $`Q^{ij}`$ is transverse, to obtain
$$\stackrel{~}{\mathrm{\Delta }}_𝕃Z^i+6(\stackrel{~}{𝕃}Z)^{ij}\stackrel{~}{}_j\mathrm{ln}\psi =\psi ^6\stackrel{~}{}_j\stackrel{~}{M}^{ij}.$$
(37)
Again, we can define
$$V^iW^iZ^i,$$
(38)
and use the linearity of $`\stackrel{~}{𝕃}`$ and $`\stackrel{~}{\mathrm{\Delta }}_𝕃`$ to combine the process of obtaining the transverse-traceless part of $`M^{ij}`$ and solving the momentum constraints. We obtain the following full decomposition, which I will list together here for convenience<sup>5</sup><sup>5</sup>5Note that in (2.2.2) we have *not* used the usual conformal scaling of the extrinsic curvature given in equation (22).:
$`\gamma _{ij}`$ $`=`$ $`\psi ^4\stackrel{~}{\gamma }_{ij}`$
$`K^{ij}`$ $`=`$ $`\psi ^4\left(\stackrel{~}{A}^{ij}+\frac{1}{3}\stackrel{~}{\gamma }^{ij}K\right)`$
$`\stackrel{~}{A}^{ij}`$ $`=`$ $`(\stackrel{~}{𝕃}V)^{ij}+\psi ^6\stackrel{~}{M}^{ij}`$ (39)
$`\stackrel{~}{\mathrm{\Delta }}_𝕃V^i+6(\stackrel{~}{𝕃}V)^{ij}\stackrel{~}{}_j\mathrm{ln}\psi `$ $`=`$ $`\frac{2}{3}\stackrel{~}{}^iK\psi ^6\stackrel{~}{}_j\stackrel{~}{M}^{ij}+8\pi G\psi ^4j^i`$
$`\stackrel{~}{}^2\psi \frac{1}{8}\psi \stackrel{~}{R}\frac{1}{12}\psi ^5K^2+\frac{1}{8}\psi ^5\stackrel{~}{A}_{ij}\stackrel{~}{A}^{ij}`$ $`=`$ $`2\pi G\psi ^5\rho `$
In the decomposition given by (2.2.2), we are again free to specify a symmetric tensor $`\stackrel{~}{\gamma }_{ij}`$ as the conformal 3-metric, a symmetric tracefree tensor $`\stackrel{~}{M}^{ij}`$, and a scalar function $`K`$. Then, with given matter energy and momentum densities, $`\rho `$ and $`j^i`$, and appropriate boundary conditions, the coupled set of constraint equations for $`\psi `$ and $`V^i`$ are solved. Finally, given the solutions, we can construct the physical initial data, $`\gamma _{ij}`$ and $`K^{ij}`$.
Notice that, while very similar to the decomposition from §2.2.1, the sets of equations are distinctly different. In general, if we make the same choices for the freely specifiable data in both decompositions (i.e., we choose $`\stackrel{~}{\gamma }_{ij}`$, $`\stackrel{~}{M}^{ij}`$, and $`K`$ the same), we will produce two different sets of initial data. Both will be equally valid solutions of the constraint equations, but they will have distinct physical properties.
There is at least one exception to this. Assume we have a valid set of initial data $`\gamma _{ij}`$ and $`K_{ij}`$, which satisfies the constraint equations (14) and (15). For any everywhere-positive function $`\mathrm{\Psi }`$, we define our freely specifiable data as follows:
$`\stackrel{~}{\gamma }_{ij}`$ $``$ $`\mathrm{\Psi }^4\gamma _{ij}`$
$`\stackrel{~}{M}^{ij}`$ $``$ $`\mathrm{\Psi }^{10}\left(K^{ij}\frac{1}{3}\gamma ^{ij}K\right)`$ (40)
$`K`$ $``$ $`K_i^i.`$
Then the solution to both sets of equations, assuming we use correct boundary conditions, will be $`\psi =\mathrm{\Psi }`$ and $`V^i=0`$, which yields the original data as the solution for each decomposition.
### 2.3 Thin-Sandwich Decomposition
The two general initial-value decompositions outlined in §2.2.1 and §2.2.2 require identical freely specified data ($`\stackrel{~}{\gamma }_{ij}`$, $`\stackrel{~}{M}^{ij}`$, and $`K`$), yet they usually produce different physical initial data. One shortcoming of these approaches is that they provide no direct insight into how to choose the freely specifiable data. All of the data are determined by schemes that involve only a single spacelike hypersurface. The resulting constraint equations are independent of the kinematical variables $`\alpha `$ and $`\beta ^i`$ that govern how the coordinates move through spacetime, and thus there is no connection to dynamics. York’s thin-sandwich decomposition takes a different approach by considering the evolution of the metric between two neighboring hypersurfaces (the thin sandwich). This decomposition is very similar to an approach originated by Wilson , but is somewhat more general. Perhaps the most attractive feature of this decomposition is the insight it yields into the choice of the freely specifiable data.
The decomposition begins with the standard conformal decomposition of the 3-metric (16). However, we next make use of the evolution equation for the metric (13) in order to connect the 3-metrics on the two neighboring hypersurfaces. Label the two slices by $`t`$ and $`t^{}`$, with $`t^{}=t+\delta t`$, then $`\gamma _{ij}^{}=\gamma _{ij}+\left(_t\gamma _{ij}\right)\delta t`$. We would like to specify how the 3-metric evolves, but we do not have full freedom to do this. We know we can freely specify only the conformal 3-metric, and similarly, we are free to specify only the evolution of the conformal 3-metric. We make the following definitions:
$$u_{ij}\gamma ^{\frac{1}{3}}_t(\gamma ^{\frac{1}{3}}\gamma _{ij}),$$
(41)
$$\stackrel{~}{u}_{ij}_t\stackrel{~}{\gamma }_{ij},$$
(42)
and
$$\stackrel{~}{\gamma }^{ij}\stackrel{~}{u}_{ij}0.$$
(43)
The latter definition is made for convenience, so that we can treat $`\psi `$, $`\stackrel{~}{\gamma }_{ij}`$, and $`\stackrel{~}{u}_{ij}`$ as regular scalars and tensors instead of as scalar- and tensor-densities within this thin-sandwich formalism.
The conformal scaling of $`u_{ij}`$ follows directly from (16), (41), (42), (43), and the identity that, for any small perturbation, $`\delta \mathrm{ln}\gamma =\gamma ^{ij}\delta \gamma _{ij}`$. The result is
$$u_{ij}=\psi ^4\stackrel{~}{u}_{ij},$$
(44)
which relies on the useful intermediate result that<sup>6</sup><sup>6</sup>6 Equations (43) and (45) take the place of defining $`\stackrel{~}{\gamma }=1`$ and the subsequent necessity of treating $`\psi `$ and $`\stackrel{~}{\gamma }_{ij}`$ as densities.
$$_t\stackrel{~}{\gamma }=0.$$
(45)
Equation (41) represents the tracefree part of the evolution of the 3-metric, so (13) becomes
$$u^{ij}=2\alpha A^{ij}+(\overline{𝕃}\beta )^{ij}.$$
(46)
Using the conformal scalings (22), (35), and (44), we obtain
$$\stackrel{~}{A}^{ij}=\frac{\psi ^6}{2\alpha }\left((\stackrel{~}{𝕃}\beta )^{ij}\stackrel{~}{u}^{ij}\right).$$
(47)
York has pointed out that it is natural to use the following conformal rescaling of the lapse:
$$\alpha =\psi ^6\stackrel{~}{\alpha }.$$
(48)
This rescaling follows naturally from the “slicing function” that replaces the usual lapse ($`\alpha =\sqrt{\gamma }\stackrel{~}{\alpha }`$) which has been critical in solving several problems . It also results in the natural conformal scaling (22) postulated for the tracefree part of the extrinsic curvature. Substituting (48) into (47) yields what is taken as the definition of the tracefree part of the conformal extrinsic curvature,
$$\stackrel{~}{A}^{ij}\frac{1}{2\stackrel{~}{\alpha }}\left((\stackrel{~}{𝕃}\beta )^{ij}\stackrel{~}{u}^{ij}\right).$$
(49)
Because the tracefree extrinsic curvature satisfies the normal conformal scaling, the Hamiltonian constraint will take on the same form as in (2.2.1). However, the momentum constraint will have a very different form. Combining equations (16), (19), (21), (22), and (49) with the momentum constraint (15), we find that it simplifies to
$$\stackrel{~}{\mathrm{\Delta }}_𝕃\beta ^i(\stackrel{~}{𝕃}\beta )^{ij}\stackrel{~}{}_j\mathrm{ln}\stackrel{~}{\alpha }=\frac{4}{3}\stackrel{~}{\alpha }\psi ^6\stackrel{~}{}^iK+\stackrel{~}{\alpha }\stackrel{~}{}\left(\frac{1}{\stackrel{~}{\alpha }}\stackrel{~}{u}^{ij}\right)+16\pi G\stackrel{~}{\alpha }\psi ^{10}j^i.$$
(50)
Let us, for convenience, group together all the equations that constitute the thin-sandwich decomposition:
$`\gamma _{ij}`$ $`=`$ $`\psi ^4\stackrel{~}{\gamma }_{ij}`$
$`K^{ij}`$ $`=`$ $`\psi ^{10}\stackrel{~}{A}^{ij}+\frac{1}{3}\psi ^4\stackrel{~}{\gamma }^{ij}K`$
$`\stackrel{~}{A}^{ij}`$ $`=`$ $`{\displaystyle \frac{1}{2\stackrel{~}{\alpha }}}\left((\stackrel{~}{𝕃}\beta )^{ij}\stackrel{~}{u}^{ij}\right)`$ (51)
$`\stackrel{~}{\mathrm{\Delta }}_𝕃\beta ^i(\stackrel{~}{𝕃}\beta )^{ij}\stackrel{~}{}_j\mathrm{ln}\stackrel{~}{\alpha }+\frac{4}{3}\stackrel{~}{\alpha }\psi ^6\stackrel{~}{}^iK`$ $`=`$ $`\stackrel{~}{\alpha }\stackrel{~}{}\left(\frac{1}{\stackrel{~}{\alpha }}\stackrel{~}{u}^{ij}\right)+16\pi G\stackrel{~}{\alpha }\psi ^{10}j^i`$
$`\stackrel{~}{}^2\psi \frac{1}{8}\psi \stackrel{~}{R}\frac{1}{12}\psi ^5K^2+\frac{1}{8}\psi ^7\stackrel{~}{A}_{ij}\stackrel{~}{A}^{ij}`$ $`=`$ $`2\pi G\psi ^5\rho `$
In this decomposition (2.3), we are free to specify a symmetric tensor $`\stackrel{~}{\gamma }_{ij}`$ as the conformal 3-metric, a symmetric tracefree tensor $`\stackrel{~}{u}^{ij}`$, a scalar function $`K`$, *and* the scalar function $`\stackrel{~}{\alpha }`$. Solving this set of equations with appropriate boundary conditions yields initial data $`\gamma _{ij}`$ and $`K_{ij}`$ on a *single* hypersurface. However, we also know the following. *If* we chose to use the shift vector obtained from solving (50) and the lapse from (48) via our choice of $`\stackrel{~}{\alpha }`$ and our solution to the Hamiltonian constraint, then the rate of change of the physical 3-metric is given by
$`_t\gamma _{ij}`$ $`=`$ $`u_{ij}+\frac{2}{3}\gamma _{ij}\left(\overline{}_k\beta ^k\alpha K\right)`$
$`=`$ $`\psi ^4\left[\stackrel{~}{u}_{ij}+\frac{2}{3}\stackrel{~}{\gamma }_{ij}\left(\stackrel{~}{}_k\beta ^k+6\beta ^k\stackrel{~}{}_k\mathrm{ln}\psi \psi ^6\stackrel{~}{\alpha }K\right)\right].`$
This direct information about the consequences of our choices for the freely specifiable data is something not present in the previous decompositions. As we will see later, this framework has been used to construct initial data that are in quasiequilibrium.
### 2.4 Stationary Solutions
When there is sufficient symmetry present, it is possible to construct initial data that are in true equilibrium. These solutions possess at least two Killing vectors, one that is timelike at large distances and one that is spatial, representing an azimuthal symmetry. When these symmetries are present, solving for the initial data produces a global solution of Einstein’s equations and the solution is said to be *stationary*. The familiar Kerr-Neumann solution for rotating black holes is an example of a stationary solution in vacuum. Stationary configurations supported by matter are also possible, but the matter sources must also satisfy the Killing symmetries, in which case the matter is said to be in hydrostatic equilibrium .
The basic approach for finding stationary solutions begins by simplifying the metric to take into account the symmetries. Many different forms have been used for the metric (cf. Refs. ). I will use a decomposition that makes comparison with the previous decompositions straightforward. First, define the interval as
$$\mathrm{d}s^2=\psi ^4\mathrm{d}t^2+\psi ^4\left[A^2(\mathrm{d}r^2+r^2\mathrm{d}\theta ^2)+B^2r^2\mathrm{sin}^2\theta (\mathrm{d}\varphi +\beta ^\varphi \mathrm{d}t)^2\right].$$
(53)
This form of the metric can describe *any* stationary spacetime. Notice that the lapse is related to the conformal factor by
$$\alpha =\psi ^2,$$
(54)
and that the shift vector has only one component
$$\beta ^i=(0,0,\beta ^\varphi ).$$
(55)
I have used the usual conformal decomposition of the 3-metric (16) and have written the conformal 3-metric with two parameters as
$$\stackrel{~}{\gamma }_{ij}=\left(\begin{array}{ccc}A^2& 0& 0\\ 0& A^2r^2& 0\\ 0& 0& B^2r^2\mathrm{sin}^2\theta \end{array}\right).$$
(56)
The four functions $`\psi `$, $`\beta ^\varphi `$, $`A`$, and $`B`$ are functions of $`r`$ and $`\theta `$ only.
The equations necessary to solve for these four functions are derived from the constraint equations (14) and (15), *and* the evolution equations (2.1) and (13). For the evolution equations, we use the fact that $`_t\gamma _{ij}=0`$ and $`_tK_{ij}=0`$. The metric evolution equation (13) defines the extrinsic curvature in terms of derivatives of the shift
$$K_{ij}\frac{1}{2\alpha }(\overline{}_i\beta _j+\overline{}_j\beta _i).$$
(57)
With the given metric and shift, we find that $`K=0`$ and the divergence of the shift also vanishes. This means we can write the tracefree part of the extrinsic curvature as
$$A^{ij}=\psi ^{10}\stackrel{~}{A}^{ij}=\frac{1}{2}\psi ^2(\stackrel{~}{𝕃}\beta )^{ij}.$$
(58)
We find that the Hamiltonian and momentum constraints take on the forms given by the thin-sandwich decomposition (2.3) with $`\stackrel{~}{u}^{ij}=K0`$ and $`\stackrel{~}{\alpha }\psi ^8`$. Only one of the momentum constraint equations is non-trivial, and we find that the constraints yield elliptic equations for $`\psi `$ and $`\beta ^\varphi `$. What remains unspecified as yet are $`A`$ and $`B`$ (i.e., the conformal 3-metric).
The conformal 3-metric is determined by the evolution equations for the traceless part of the extrinsic curvature. Of these five equations, one can be written as an elliptic equation for $`B`$, and two yield complementary equations that can each be solved by quadrature for $`A`$. The remaining equations are redundant as a result of the Bianchi identities.
Of course, the clean separation of the equations I have suggested above is an illusion. All four equations must be solved simultaneously, and clever combinations of the four metric quantities can greatly simplify the task of solving the system of equations. This accounts for the numerous different systems used for solving for stationary solutions.
## 3 Black Hole Initial Data
In this section, we will look at Cauchy initial data that represent one or more black holes in an asymptotically flat spacetime. The majority of these will be either vacuum solutions or solutions of the Einstein-Maxwell equations. With no matter to support the gravitational field, we find that we must usually<sup>7</sup><sup>7</sup>7A black hole can be supported by a compact gravitational wave . use a spacetime with a non-trivial topology. It is certainly possible to construct black-hole solutions supported by matter , but it is often desirable to avoid the complications of matter sources.
This raises a point about solutions of Einstein’s equations which we have not yet mentioned. When constructing solutions of Einstein’s initial-value equations, we are free to specify the topology of the initial-data hypersurface. Einstein’s equations of general relativity place no constraints on the topology of the spacetime they describe or of spacelike hypersurfaces that foliate it. For astrophysical black holes (i.e., black holes in an asymptotically flat spacetime), the freedom in the choice of the topology has relatively minor consequences. The primary effects of different topology choices are hidden within the black hole’s event horizon.
In the sections below, we will explore many of the existing black-hole solutions and the schemes for generating them.
### 3.1 Classic Solutions
#### 3.1.1 Schwarzschild
Of course, the simplest black-hole solution is the Schwarzschild solution. It represents a static spacetime containing a single black hole that connects two causally disconnected, asymptotically flat universes. There are actually many different coordinate representations of the Schwarzschild solutions. The simplest representations are time-symmetric ($`K_{ij}=0`$), and so exist on a “maximally embedded” spacelike hypersurface ($`K=0`$). These choices fix the foliation $`\{\mathrm{\Sigma }\}`$. Spherical symmetry fixes two of the three spatial gauge choices. If we choose an “areal-radial coordinate”<sup>8</sup><sup>8</sup>8Spheres at constant areal-radial coordinate $`r`$ have proper area $`4\pi r^2`$., then the interval is written as
$$\mathrm{d}s^2=\left(1\frac{2M}{r}\right)\mathrm{d}t^2+\left(1\frac{2M}{r}\right)^1\mathrm{d}r^2+r^2(\mathrm{d}\theta ^2+\mathrm{sin}^2\theta \mathrm{d}\varphi ^2).$$
(59)
If we choose an isotropic radial coordinate, then the interval is written as
$$\mathrm{d}s^2=\left(\frac{1\frac{M}{2\stackrel{~}{r}}}{1+\frac{M}{2\stackrel{~}{r}}}\right)^2\mathrm{d}t^2+\left(1+\frac{M}{2\stackrel{~}{r}}\right)^4\left(\mathrm{d}\stackrel{~}{r}^2+\stackrel{~}{r}^2\mathrm{d}\theta ^2+\stackrel{~}{r}^2\mathrm{sin}^2\theta \mathrm{d}\varphi ^2\right).$$
(60)
In both (59) and (60), $`M`$ represents the mass of the black hole as measured at spacelike infinity. Both of these solutions exist on the same foliation of $`t=\text{const.}`$ slices. But, notice that the 3-geometry of the slice associated with (60) is conformally flat, while the 3-geometry associated with (59) is not.
The solution given in (60) is easily generated by any of the methods in §2.2 or §2.3. By choosing a time-symmetric initial-data hypersurface, we immediately get $`K_{ij}=0`$, which eliminates the need to solve the momentum constraints. If we choose the conformal 3-geometry to be given by a flat metric (in spherical coordinates in this case), then the vacuum Hamiltonian constraint (18) becomes
$$\stackrel{~}{}^2\psi =0,$$
(61)
where $`\stackrel{~}{}^2`$ is the flat-space Laplace operator. For the solution $`\psi `$ to yield an asymptotically flat physical 3-metric, we have the boundary condition that $`\psi (\stackrel{~}{r}\mathrm{})=1`$. The simplest solution of this equation is
$$\psi =1+\frac{M}{2\stackrel{~}{r}},$$
(62)
where we have chosen the remaining integration constant to give a mass at infinity of $`M`$.
We now have full Cauchy initial data representing a single black hole. If we want to generate a full solution of Einstein’s equations, we must choose a lapse and a shift vector and integrate the evolution equations (2.1) and (13). In this case, a reasonable approach for specifying the lapse is to demand that the time derivative of $`K`$ vanish. For the case of $`K=0`$, this yields the so-called maximal slicing equation which, for the current situation, takes the form
$$\stackrel{~}{}^2(\alpha \psi )=0.$$
(63)
If we choose boundary conditions so that the lapse is frozen on the event horizon ($`\alpha (\stackrel{~}{r}=M/2)=0`$) and goes to one at infinity, we find that the solution is
$$\alpha =\frac{1\frac{M}{2\stackrel{~}{r}}}{1+\frac{M}{2\stackrel{~}{r}}}.$$
(64)
If we now choose $`\beta ^i=0`$, we find that the left-hand sides of the evolution equations (2.1) and (13) vanish identically, and we have found the static solution of Einstein’s equations given in (60). We can, of course, recover the usual Schwarzschild coordinate solution (59) by using the purely spatial coordinate transformation $`r=\stackrel{~}{r}(1+\frac{M}{2\stackrel{~}{r}})^2`$.
It is interesting to examine the differences in these two representations of the Schwarzschild solution. The isotropic radial coordinate representation is well behaved everywhere except, it seems, at $`\stackrel{~}{r}=0`$. However, even here, the solution is well behaved. The 3-geometry is invariant under the coordinate transformation
$$\stackrel{~}{r}\left(\frac{M}{2}\right)^2\frac{1}{r^{}}.$$
(65)
The event horizon at $`\stackrel{~}{r}=\frac{M}{2}`$ is a fixed-point set of the isometry condition (65) which identifies points in two causally disconnected, asymptotically flat universes. We see that $`\stackrel{~}{r}=0`$ is simply an image of infinity in the other universe .
Given our choice for the lapse (64), which is frozen on the event horizon, we find that the solution can cover only the exterior of the black hole. To cover any of the interior with the lapse pinned to zero at the horizon would require we use a slice that is not spacelike everywhere. This is exactly what happens when the usual Schwarzschild areal-radial coordinate is used. At the event horizon, $`r=2M`$, there is a coordinate singularity, and inside this radius the $`t=\text{const.}`$ hypersurface is no longer spacelike. It is impossible to perform a Cauchy evolution interior to the event horizon using the areal-radial coordinate and the given time slicing.
We find that a Cauchy evolution, using the usual Schwarzschild time slicing that is frozen at the horizon, is capable of evolving only the region exterior to the black hole’s event horizon. Portions of the interior of the black hole can be covered by an evolution that begins with data on a standard Schwarzschild time slice, but the result is not a time-independent solution. As we will see later, there are other slicings of the Schwarzschild spacetime that cover the interior of the black hole and yield time-independent solutions.
#### 3.1.2 Time-Symmetric Multi-Hole Solutions
As we saw in §3.1.1, the simplest approach for generating initial data is to assume time symmetry and let the conformal 3-geometry be flat. We could have generated the exterior solution for Schwarzschild with the areal-radial coordinate by a clever choice for the conformal 3-geometry. The Hamiltonian constraint is still linear, even if $`\stackrel{~}{\gamma }_{ij}`$ is not flat. But there is no obvious motivation for the correct $`\stackrel{~}{\gamma }_{ij}`$ that would yield the desired solution.
One approach for generating a time-symmetric multi-hole solution is straightforward. Brill–Lindquist initial data again assume a flat conformal 3-geometry, and the only non-trivial constraint equation is the Hamiltonian constraint, which again takes the form given in (61). But this time, we use the linearity of the Hamiltonian constraint and choose the solution to be a superposition of solutions with the form of (62). More precisely, we choose the solution
$$\psi =1+\underset{\sigma =1}{\overset{N}{}}\frac{\mu _\sigma }{2|𝐱𝐂_\sigma |}.$$
(66)
Here, $`|𝐱𝐂_\sigma |`$ is a coordinate distance from the point $`𝐂_\sigma `$ in the Euclidean conformal space, and the $`\mu _\sigma `$ are constants related to the masses of the black holes. Assuming the points $`𝐂_\sigma `$ are sufficiently far apart, this solution of the initial-value equations represents $`N`$ black holes momentarily at rest in “our” asymptotically flat universe. As was the case for a single black hole, each singular point in the solution, $`𝐱=𝐂_\sigma `$, represents infinity in a different, causally disconnected universe. In fact, each black hole connects “our” universe to a different universe, so that there are $`N+1`$ asymptotically flat hypersurfaces connected together at the throats of $`N`$ black holes. While we started with a 3-dimensional Euclidean manifold, the requirement that we delete the singular points, $`𝐂_\sigma `$, results in a manifold that is not simply connected. This solution is often referred to as having a topology with $`N+1`$ “sheets”.
Brill–Lindquist initial data are very similar to the Schwarzschild initial data in isotropic coordinates except for one major difference: the solution does not represent two identical universes that have been joined together. The coordinate transformation (65) can still be used to show that each pole in the solution corresponds to infinity on an asymptotically flat hypersurface. But, the solution has $`N+1`$ different asymptotically flat universes connected together, not two, and “our” universe containing $`N`$ black holes cannot be isometric to any of the other universes which contain just one. Interestingly, Misner found that it is possible to construct a solution of the vacuum, time-symmetric Hamiltonian constraint (61) that has two isometric asymptotically flat hypersurfaces connected by $`N`$ black holes. The case of two black holes that satisfies this isometry condition (often called inversion symmetry) is usually referred to as “Misner initial data”. It has an analytic representation in terms of an infinite series expansion. The construction is tedious, and there are several representations of the solution <sup>9</sup><sup>9</sup>9In Ref. , Eq. (B7), the line “$`1`$ for $`n=1`$” should read “$`\alpha ^1`$ for $`n=1`$”..
Like the Brill–Lindquist data, the non-simply connected topology of the full manifold is represented on a Euclidean 3-manifold by the presence of singular points that must be removed. Brill–Lindquist data have $`N`$ singular points, each representing an *image* of infinity as seen through the throat of the black hole connecting our universe to that hole’s other universe. Misner’s data contain an infinite number of singular points for each black hole, each representing an image of one of two asymptotic infinities. The result is seen to represent two identical asymptotically flat universes joined by $`N`$ black holes. The two universes join together at the throats of the $`N`$ black holes, with each throat being a coordinate 2-sphere in the conformal space. Data at any point in one universe are related to data at the corresponding point in the alternate universe via the same isometry condition (65) found in the Schwarzschild case. And, as in the Schwarzschild case, the 2-sphere throat of each black hole forms a fixed-point set of the isometry condition.
Both the Brill–Lindquist data and Misner’s data represent $`N`$ black holes at a moment of time-symmetry (i.e., all the black holes are momentarily at rest). Both are conformally flat and the difference in the topology of the two solutions is hidden from an observer outside the black holes. Yet solutions where the holes are chosen to have the same size and separation yield similar, but physically distinct, solutions .
### 3.2 General Multi-Hole Solutions
Time-symmetric black-hole solutions of the constraint equations such as those described in §3.1 are useful as test cases because they have analytic representations. However, they have very little physical relevance. General time-asymmetric solutions are needed to represent black holes that are moving and spinning<sup>10</sup><sup>10</sup>10Time-asymmetric solutions are also needed to represent time-independent solutions that cover the interior of a black hole.. A few approaches for generating general multi-hole solutions have been explored and, below, we will look at the two approaches which are direct generalizations of the Misner and Brill–Lindquist data. Generalizing these two approaches was the most natural first step toward constructing general multi-hole solutions.
The first approach to be developed generalized the Misner approach . It was attractive because an isometry condition relating the two asymptotically flat universes provides two useful things. First, because the two universes are identical, finding a solution in one universe means that you have the full solution. Second, because the throats are fixed-point sets of the isometry, we can construct boundary conditions on any quantity there. This allows us to excise the region interior to the spherical throats from one of the Euclidean background spaces and solve for the initial data in the remaining volume.
The generalization of the Misner approach seemed preferable to trying to solve the constraints on $`N+1`$ Euclidean manifolds stitched together smoothly at the throats of $`N`$ black holes. However, Brandt and Brügmann realized it was possible to factor out analytically the behavior of the singular points in the Euclidean manifold of the $`N+1`$ sheeted approach. Referred to as the “puncture” method, this approach allows us to rewrite the constraint equations for functions on an $`N+1`$ sheeted manifold as constraint equations for new functions on a simple Euclidean manifold.
Another approach tried, which we will not discuss in detail, avoided the issue of the topology of the initial-data slice entirely. Developed by Thornburg , this approach was based on the idea that only the domain exterior to the apparent horizon of a black hole is relevant. The equation describing the location of an apparent horizon can be rewritten in a form that can be used as a boundary condition for the conformal factor in the Hamiltonian constraint equation. Thus, given a compatible solution to the momentum constraints, this boundary condition can be used to construct a solution of the Hamiltonian constraint in the domain exterior to the apparent horizons of any black holes, with no reference at all to the topology of the full manifold.
#### 3.2.1 Bowen–York Data
The generalization of the Misner approach was developed by York and his collaborators . The approach is often called the “conformal-imaging” method, and the data are usually referred to as “Bowen–York” data. This approach begins with a set of simplifying assumptions that is common to all three of the approaches described above. These assumptions are that
$`K`$ $`=`$ $`0,\text{maximal slicing}`$
$`\stackrel{~}{\gamma }_{ij}`$ $`=`$ $`f_{ij},\text{conformal flatness}`$ (67)
$`\psi |_{\mathrm{}}`$ $`=`$ $`1.\text{asymptotic flatness}`$
Here, $`f_{ij}`$ represents a flat metric in any suitable coordinate system. The assumption of conformal flatness means that the differential operators in the constraints are the familiar flat-space operators. More importantly, if we use the conformal transverse-traceless decomposition (2.2.1), we find that in vacuum the momentum constraints completely decouple from the Hamiltonian constraint.
The importance of this last property stems from the fact that York and Bowen were able to find analytic solutions of this version of the momentum constraints, solutions that represent a black hole with both linear momentum and spin . If we choose $`\stackrel{~}{M}^{ij}=0`$, then the momentum constraints (31) become
$$\stackrel{~}{}^2V^i+\frac{1}{3}\stackrel{~}{}^i\left(\stackrel{~}{}_jV^j\right)=0.$$
(68)
A solution of this equation is
$$V^i=\frac{1}{4r}\left[7P^i+n^in_jP^j\right]+\frac{1}{r^2}ϵ^{ijk}n_jS_k.$$
(69)
Here, $`P^i`$ and $`S^i`$ are vector parameters, $`r`$ is a coordinate radius, and $`n^i`$ is the outward-pointing unit normal of a sphere in the flat conformal space ($`n^i\frac{x^i}{r}`$). $`ϵ^{ijk}`$ is the 3-dimensional Levi-Civita tensor.
This solution of the momentum constraints yields the tracefree part of the extrinsic curvature
$`\stackrel{~}{A}_{ij}`$ $`=`$ $`{\displaystyle \frac{3}{2r^2}}\left[P_in_j+P_jn_i(f_{ij}n_in_j)P^kn_k\right]`$
$`+{\displaystyle \frac{3}{r^3}}\left[ϵ_{ki\mathrm{}}S^{\mathrm{}}n^kn_j+ϵ_{kj\mathrm{}}S^{\mathrm{}}n^kn_i\right].`$
Remarkably, using this solution (3.2.1) and the assumptions in (3.2.1), we can determine the physical values for the linear and angular momentum of any initial data we can construct. The momentum contained in an asymptotically flat initial-data hypersurface can be calculated from the integral
$$\mathrm{\Pi }^i\xi _{(k)}^i=\frac{1}{8\pi }_{\mathrm{}}\left(K_i^j\delta _i^jK\right)\xi _{(k)}^i\mathrm{d}^2S_j,$$
(71)
where $`\xi _{(k)}^i`$ is a Killing vector of the 3-metric $`\gamma _{ij}`$<sup>11</sup><sup>11</sup>11If $`\xi _{(k)}^i`$ is a translational Killing vector, then (71) yields the linear momentum in the direction of that Killing vector. If $`\xi _{(k)}^i`$ is a rotational Killing vector, then (71) yields the corresponding angular momentum.. Since we are not likely to have true Killing vectors, we make use of the asymptotic translational and rotational Killing vectors of the flat conformal space. We find from (71), (3.2.1), and (3.2.1) that the *physical* linear momentum of the initial-data hypersurface is $`P^i`$ and the *physical* angular momentum of the slice is $`S^i`$. Furthermore, because the momentum constraints are linear, we can add any number of solutions of the form of Eq. (3.2.1) to represent a collection of linear and angular momentum sources. The total physical linear momentum of the initial-data slice will simply be the vector sum of the individual linear momenta. The total physical angular momentum cannot be obtained by simply summing the individual spins because this neglects the orbital angular momentum of the various sources. However, the total angular momentum can still be computed without having to solve the Hamiltonian constraint .
The Bowen–York solution for the extrinsic curvature is the starting point for all the general multi-hole initial-data sets we have discussed in §3.2. However, this solution is not *inversion symmetric*. That is, it does not satisfy the isometry condition that any field must satisfy to exist on a two-sheeted manifold like that of Misner’s solution. Fortunately, there is a method of images, similar to that used in electrostatics but applicable to tensors, that can be used to make any tensor inversion symmetric .
For the conformal extrinsic curvature of a single black hole, there are two inversion-symmetric solutions
$`\stackrel{~}{A}_{ij}^\pm `$ $`=`$ $`{\displaystyle \frac{3}{2r^2}}\left[P_in_j+P_jn_i(f_{ij}n_in_j)P^kn_k\right]`$
$`{\displaystyle \frac{3a^2}{2r^4}}\left[P_in_j+P_jn_i+(f_{ij}5n_in_j)P^kn_k\right]`$
$`+{\displaystyle \frac{3}{r^3}}\left[ϵ_{ki\mathrm{}}S^{\mathrm{}}n^kn_j+ϵ_{kj\mathrm{}}S^{\mathrm{}}n^kn_i\right].`$
Here, $`a`$ is the radius of the coordinate 2-sphere that is the throat of the black hole. Of course, this coordinate 2-sphere is the fixed-point set of the isometry and is the surface on which we can impose boundary conditions. Notice that this radius enters the solutions only when we make it inversion symmetric.
When the extrinsic curvature represents more than one black hole, the process for making the solution inversion symmetric is rather complex and results in an infinite-series solution. However, in most cases of interest, the solution converges rapidly and it is straightforward to evaluate the solution numerically .
Given an inversion-symmetric conformal extrinsic curvature, it is possible to find an inversion-symmetric solution of the Hamiltonian constraint . Given our assumptions (3.2.1), the Hamiltonian constraint becomes
$$\stackrel{~}{}^2\psi +\frac{1}{8}\psi ^7\stackrel{~}{A}_{ij}\stackrel{~}{A}^{ij}=0.$$
(73)
The isometry condition imposes a condition on the conformal factor at the throat of each hole. This condition takes the form
$$n_\sigma ^i\stackrel{~}{}_i\psi |_{a_\sigma }=\frac{\psi }{2r_\sigma }|_{a_\sigma },$$
(74)
where $`n_\sigma ^i`$ is the outward-pointing unit-normal vector to the $`\sigma ^{\mathrm{th}}`$ throat and $`a_\sigma `$ is the coordinate radius of that throat. This condition can be used as a boundary condition when solving (73) in the region exterior to the throats.
In addition to boundary conditions on the throats, a boundary condition on the outer boundary of the domain is needed before the quasilinear elliptic equation in (73) can be solved as a well-posed boundary-value problem. This final boundary condition comes from the fact that we want an asymptotically flat solution. This implies that the solution behaves as
$$\psi =1+\frac{E}{2r}+𝒪(r^2),$$
(75)
where $`E`$ is the total ADM energy content of the initial-data hypersurface. Equation (75) can be used to construct appropriate boundary conditions either at infinity or at a large, but finite, distance from the black holes .
#### 3.2.2 Puncture Data
The generalization of the Brill–Lindquist data developed by Brandt and Brügmann begins with the same set of assumptions (3.2.1) as the conformal-imaging approach outlined in §3.2.1. We immediately have (3.2.1) from the solution of the momentum constraints, and we must solve the Hamiltonian constraint, which again takes the form of Eq. (73). At this point, however, the method of solution differs from the conformal-imaging approach.
Based on the time-symmetric solution, it is reasonable to assume that the conformal factor will take the form
$$\psi =\frac{1}{\chi }+u,\frac{1}{\chi }\underset{\sigma =1}{\overset{N}{}}\frac{\mu _\sigma }{2|𝐱𝐂_\sigma |}.$$
(76)
If $`u`$ is sufficiently smooth, (76) implies that the manifold will have the topology of $`N+1`$ asymptotically flat regions just as in the Brill–Lindquist solution. In this case, asymptotic flatness requires that $`u=1+𝒪(r^1)`$.
Substituting (76) into the Hamiltonian constraint (73) yields
$$\stackrel{~}{}^2u+\eta (1+\chi u)^7=0,$$
(77)
where
$$\eta =\frac{1}{8}\chi ^7\stackrel{~}{A}_{ij}\stackrel{~}{A}^{ij}.$$
(78)
Near each singular point, or “puncture”, we find that $`\chi |𝐱𝐂_\sigma |`$. From (3.2.1), we see that $`\stackrel{~}{A}_{ij}\stackrel{~}{A}^{ij}`$ behaves no worse than $`|𝐱𝐂_\sigma |^6`$, so $`\eta `$ vanishes at the punctures at least as fast as $`|𝐱𝐂_\sigma |`$.
With this behavior, Brandt and Brügmann have shown the existence and uniqueness of $`C^2`$ solutions of the modified Hamiltonian constraint (77). The resulting scheme for constructing multiple black hole initial data is very simple. The mass and position of each black hole are parameterized by $`\mu _\sigma `$ and $`𝐂_\sigma `$, respectively. Their linear momenta and spin are parameterized by $`𝐏_\sigma `$ and $`𝐒_\sigma `$ in the conformal extrinsic curvature (3.2.1) used for each hole. Finally, the solution for $`u`$ is found on a simple Euclidean manifold, with no need for any inner boundaries to avoid singularities. This is a great simplification over the conformal-imaging approach, where proper handling of the inner boundary is the most difficult aspect of solving the Hamiltonian constraint numerically .
#### 3.2.3 Problems with These Data
Both the conformal-imaging and puncture methods for generating multiple black hole initial data allow for completely general configurations of the relative sizes of the black holes, as well as their linear and angular momenta. This does not mean that these schemes allow for the generation of *all* desired black-hole initial data. The two schemes rely on specific assumptions about the freely specifiable gravitational data. In particular, they assume $`K=0`$, $`\stackrel{~}{M}^{ij}=0`$, and, most importantly, that the 3-geometry is conformally flat.
These choices for the freely specifiable data are not always commensurate with the desired physical solution. For example, if we choose to use either method to construct a single spinning black hole, we will not obtain the Kerr solution. The Kerr–Newman solution can be written in terms of a quasi-isotropic radial coordinate on a $`K=0`$ time slice . Let $`r`$ denote the usual Boyer–Lindquist radial coordinate and make the standard definitions
$$\rho ^2r^2+a^2\mathrm{cos}^2\theta \text{and}\mathrm{\Delta }r^22Mr+a^2+Q^2.$$
(79)
A quasi-isotropic radial coordinate $`\stackrel{~}{r}`$ can be defined via
$$r=\stackrel{~}{r}\left(1+\frac{M+\sqrt{a^2+Q^2}}{2\stackrel{~}{r}}\right)\left(1+\frac{M\sqrt{a^2+Q^2}}{2\stackrel{~}{r}}\right).$$
(80)
The interval then becomes
$$\mathrm{d}s^2=\alpha ^2\mathrm{d}t^2+\psi ^4\left[e^{2\mu /3}(\mathrm{d}\stackrel{~}{r}^2+\stackrel{~}{r}^2\mathrm{d}\theta ^2)+\stackrel{~}{r}^2\mathrm{sin}^2\theta e^{4\mu /3}(\mathrm{d}\varphi +\beta ^\varphi \mathrm{d}t)^2\right],$$
(81)
with
$`\alpha ^2`$ $`=`$ $`{\displaystyle \frac{\rho ^2\mathrm{\Delta }}{r^2+a^2)^2\mathrm{\Delta }a^2\mathrm{sin}^2\theta }}`$
$`\beta ^\varphi `$ $`=`$ $`a{\displaystyle \frac{r^2+a^2\mathrm{\Delta }}{(r^2+a^2)^2\mathrm{\Delta }a^2\mathrm{sin}^2\theta }}`$ (82)
$`\psi ^4`$ $`=`$ $`{\displaystyle \frac{\rho ^{\frac{2}{3}}((r^2+a^2)^2\mathrm{\Delta }a^2\mathrm{sin}^2\theta )^{\frac{1}{3}}}{\stackrel{~}{r}^2}}`$
$`e^{2\mu }`$ $`=`$ $`{\displaystyle \frac{\rho ^4}{(r^2+a^2)^2\mathrm{\Delta }a^2\mathrm{sin}^2\theta }}`$
We see immediately that the 3-geometry associated with a $`t=\text{const.}`$ hypersurface of (81) is not conformally flat. In fact, Garat and Price have shown that in general there is no spatial slicing of the Kerr spacetime that is axisymmetric, conformally flat, and smoothly goes to the Schwarzschild solution as the spin parameter $`a0`$.
Since the Kerr solution is stationary, the inescapable conclusion is that conformally flat initial data for a single rotating black hole must also contain some nonvanishing dynamical component. When we evolve such data, the system will emit gravitational radiation and eventually settle down to the Kerr geometry . But, it cannot be the Kerr geometry initially, and it is unlikely that the spurious gravitational radiation content of the initial data has any desirable physical properties. Conformally flat initial data for spinning holes contain some amount of unphysical “junk” radiation<sup>12</sup><sup>12</sup>12A similar conclusion is reached for conformally flat data for a single black hole with linear momentum ..
The choice of a conformally flat 3-geometry was originally made for convenience. Combined with the choice of maximal slicing, these simplifying assumptions allowed for an analytic solution of the momentum constraints which vastly simplified the process of constructing black-hole initial data. Yet there has been much concern about the possible adverse physical effects that these choices (especially the choice of conformal flatness) will have in trying to study black-hole spacetimes . While these conformally flat data sets may still be useful for tests of black-hole evolution codes, it is becoming widely accepted that the unphysical initial radiation will significantly contaminate any gravitational waveforms extracted from evolutions of these data. In short, these data are not astrophysically realistic.
The various initial-data decompositions outlined in §2.2 and §2.3 are all capable of producing completely general black-hole initial data sets. The only limitation of these schemes is our understanding of what choices to make for the freely specifiable data *and* the boundary conditions to apply when solving the sets of elliptic equations. All these choices will have a critical impact on the astrophysical significance of the data produced. It is also important to remember that similar choices for the freely specifiable data will result in physically different solutions when applied to the different schemes.
The first studies of black-hole initial data that are not conformally flat were carried out by Abrahams et al. . They looked at the superposition of a gravitational wave and a black hole. Using a form of the conformal metric that allows for so-called Brill waves, they constructed time-symmetric initial data that were not conformally flat and yet satisfied the isometry condition (65) used in the conformal-imaging method of §3.2.1. These data were further generalized by Brandt and Seidel to include rotating black holes with a superimposed gravitational wave. In this case, the data are no longer time-symmetric yet they satisfy a generalized form of the isometry condition so that the solution is still represented on two isometric, asymptotically flat hypersurfaces.
Matzner et al. have begun to move beyond conformally flat initial data for binary black holes<sup>13</sup><sup>13</sup>13A related method has been proposed by Bishop et al. , but their approach is much different and outside the current scope of this review.. Their proposal is to use boosted versions of the Kerr metric written in the Kerr–Schild form to represent each black hole. Thus, an isolated black hole will have no spurious radiation content in the initial data. To construct solutions with multiple black holes, they propose, essentially, to use a linear combination of the single-hole solutions. The resulting metric can be used as the *conformal* 3-metric, the trace of the resulting extrinsic curvature can be used for $`K`$, and the tracefree part of the resulting extrinsic curvature can be used for $`\stackrel{~}{M}^{ij}`$. Their scheme uses York’s conformal transverse-traceless decomposition outlined in §2.2.1, with the boundary conditions of $`\psi =1`$ and $`V^i=0`$ on the horizons of the black holes and conditions appropriate for asymptotic flatness at large distances from the holes.
The approach outlined by Matzner should certainly yield “cleaner” data than the conformally flat data currently available. For the task of specifying data for astrophysical black-hole binaries in nearly circular orbits, it is still true that this new data will not contain the correct initial gravitational wave content. Because the black holes are in orbit, they must be producing a continuous wave-train of gravitational radiation. This radiation will not be included in the method proposed by Matzner et al. Also, it is clear that the boundary conditions being used do not correctly account for the tidal distortion of each black hole by its companion. When the black holes are sufficiently far apart, the radiation from the orbital motion can be computed using post-Newtonian techniques. One possibility for producing astrophysically realistic, binary black-hole initial data is to use information from these post-Newtonian calculations to obtain better guesses for $`\stackrel{~}{\gamma }_{ij}`$, $`\stackrel{~}{M}^{ij}`$, and $`K`$.
### 3.3 Horizon-Penetrating Solutions
We noted in §3.1.1 that the time-independent maximal slicing of Schwarzschild with isotropic coordinates covers only the exterior of the black hole. This is because the time independence of this gauge requires that the lapse vanish on the horizon. It is possible to evolve into the black hole’s interior when starting from initial data constructed in this gauge, but it requires a choice for the lapse that yields a time-dependent solution . The time dependence of such a solution is purely gauge, of course, since the spacetime is static.
It is possible to cover all, or part, of the interior of a single black hole with a time-independent slicing. However, doing so seems to require that we give up the maximal-slicing condition. To cover the interior of the black hole, we need a slicing that passes smoothly through the event horizon. A convenient way to generate such solutions is to begin with the metric in standard ingoing-null coordinates. If we want to consider a rotating and charged black hole, then we use the Kerr–Newman geometry in Kerr coordinates:
$`ds^2`$ $`=`$ $`\left(1{\displaystyle \frac{2MrQ^2}{\rho ^2}}\right)\mathrm{d}\stackrel{~}{V}^2+2\mathrm{d}\stackrel{~}{V}\mathrm{d}r2{\displaystyle \frac{2MrQ^2}{\rho ^2}}a\mathrm{sin}^2\theta \mathrm{d}\stackrel{~}{V}\mathrm{d}\stackrel{~}{\varphi }`$
$`+\rho ^2\mathrm{d}\theta ^22a\mathrm{sin}^2\theta \mathrm{d}r\mathrm{d}\stackrel{~}{\varphi }+{\displaystyle \frac{1}{\rho ^2}}\left[\left(r^2+a^2\right)^2\mathrm{\Delta }a^2\mathrm{sin}^2\theta \right]\mathrm{sin}^2\theta \mathrm{d}\stackrel{~}{\varphi }^2,`$
where $`\rho `$ and $`\mathrm{\Delta }`$ are defined by Eq. (79) and $`\stackrel{~}{V}`$ is the ingoing-null coordinate. This metric is regular at $`r=r_\pm M\pm \sqrt{M^2a^2Q^2}`$, where $`r_+`$ and $`r_{}`$ are the locations of the event horizon and the Cauchy horizon, respectively.
This metric can be put into a form suitable for producing time-independent Cauchy initial data by making coordinate transformations of the general form
$$\mathrm{d}t=\mathrm{d}\stackrel{~}{V}+f(r)\mathrm{d}r,\mathrm{d}\varphi =\mathrm{d}\stackrel{~}{\varphi }+g(r)\mathrm{d}r,$$
(84)
where $`f`$ and $`g`$ are suitably chosen functions of the radial coordinate $`r`$. There are a few particularly significant solutions for the general Kerr–Newman geometry, and I will outline these below, listing the nonzero components of the metric in the standard $`3+1`$ format.
#### 3.3.1 Kerr–Schild Coordinates
A spherical coordinate version of the standard Kerr–Schild coordinate system is obtained from (3.3) by using the coordinate choice (cf. Refs. )
$$t=\stackrel{~}{V}r\text{and}\varphi =\stackrel{~}{\varphi }.$$
(85)
The nonzero components of the lapse, shift, and 3-metric are then given by:
$$\alpha ^2=1+\frac{2MrQ^2}{\rho ^2},$$
(86)
$$\beta ^r=\alpha ^2\frac{2MrQ^2}{\rho ^2},$$
(87)
$$\gamma _{rr}=1+\frac{2MrQ^2}{\rho ^2},$$
(88)
$$\gamma _{r\varphi }=\left[1+\frac{2MrQ^2}{\rho ^2}\right]a\mathrm{sin}^2\theta ,$$
(89)
$$\gamma _{\theta \theta }=\rho ^2,$$
(90)
$$\gamma _{\varphi \varphi }=\left[r^2+a^2+\frac{2MrQ^2}{\rho ^2}a^2\mathrm{sin}^2\theta \right]\mathrm{sin}^2\theta $$
(91)
Cartesian coordinate components can be obtained from these via the standard Kerr–Schild coordinate transformations
$$x+iy=(r+ia)e^{i\varphi }\mathrm{sin}\theta \text{and}z=r\mathrm{cos}\theta .$$
(92)
This yields the implicit definition of $`r`$ from
$$r^4r^2(x^2+y^2+z^2a^2)a^2z^2=0,$$
(93)
with $`r>0`$ and $`r=0`$ on the disk described by $`z=0`$ and $`x^2+y^2a^2`$.
#### 3.3.2 Harmonic Coordinates
Harmonic time slicing is integral to some hyperbolic formulations of general relativity, and a time-independent harmonic slicing of the Kerr–Newman geometry does exist . The harmonic time slicing condition is $`\mathrm{}t=0`$, which can be written
$$\frac{1}{\sqrt{g}}_\mu (\sqrt{g}g^{0\mu })=0.$$
(94)
This equation is satisfied by using the coordinate choice
$$t=\stackrel{~}{V}r+2M\mathrm{ln}\left|\frac{2M}{rr_{}}\right|,\varphi =\stackrel{~}{\varphi }.$$
(95)
The nonzero components of the lapse, shift, and 3-metric are then given by:
$$\alpha ^2=1+\frac{2MrQ^2}{\rho ^2}\left(\frac{r+r_+}{rr_{}}\right)+\frac{r_+^2+a^2}{\rho ^2}\left(\frac{2M}{rr_{}}\right)$$
(96)
$$\beta ^r=\alpha ^2\frac{r_+^2+a^2}{\rho ^2}$$
(97)
$$\beta ^\varphi =\alpha ^2\frac{a}{\rho ^2}\left(\frac{2M}{rr_{}}\right)$$
(98)
$$\gamma _{rr}=\left[2\left(1\frac{2MrQ^2}{\rho ^2}\right)\frac{r+r_+}{rr_{}}\right]\frac{r+r_+}{rr_{}}$$
(99)
$$\gamma _{r\varphi }=\left[1+\frac{2MrQ^2}{\rho ^2}\left(\frac{r+r_+}{rr_{}}\right)\right]a\mathrm{sin}^2\theta $$
(100)
$$\gamma _{\theta \theta }=\rho ^2$$
(101)
$$\gamma _{\varphi \varphi }=\left[r^2+a^2+\frac{2MrQ^2}{\rho ^2}a^2\mathrm{sin}^2\theta \right]\mathrm{sin}^2\theta $$
(102)
Cartesian coordinate components can be obtained from these via the standard Kerr–Schild coordinate transformations (92) and (93). However, for the harmonic slicing, the $`t=\text{const.}`$ hypersurface is spacelike only outside the Cauchy horizon at $`r>r_{}`$.
Fully harmonic coordinates ($`\mathrm{}x^\mu =0`$) can be defined when Cartesian spatial coordinates are used by employing a variation of the standard Kerr–Schild coordinate transformations
$$x+iy=(rm+ia)e^{i\varphi }\mathrm{sin}\theta \text{and}z=(rm)\mathrm{cos}\theta .$$
(103)
This yields the implicit definition of $`r`$ from
$$(rm)^4(rm)^2(x^2+y^2+z^2a^2)a^2z^2=0.$$
(104)
Fully harmonic coordinates are useful because applying a boost to a harmonically sliced black hole yields a solution that satisfies (94) only if the black hole is written in fully harmonic coordinates. In this case, the boosted solution also satisfies the fully harmonic coordinate conditions.
#### 3.3.3 Generalized Painlevé–Gullstrand Coordinates
The Painlevé–Gullstrand gauge choice for the Schwarzschild geometry has been rediscovered many times because of its simple form (cf. Refs. ). It is another time-independent solution, but the 3-geometry is completely flat (not simply conformally flat). The lapse is one in this gauge, and all of the information regarding the curvature of spacetime is contained in the shift. The Painlevé–Gullstrand gauge also has an intuitive physical interpretation . An observer starting at rest at infinity and freely falling will trace out a world line that is everywhere orthogonal to the $`t=\text{const.}`$ hypersurfaces in Painlevé–Gullstrand coordinates.
A generalization of the Painlevé–Gullstrand gauge derived by Doran includes the Kerr spacetime, and the extension of this solution to the full Kerr–Newman spacetime is trivial. In the limit that $`a`$ and $`Q`$ vanish, this solution reduces to the Painlevé–Gullstrand gauge. The coordinate transformation is written most easily as
$$\mathrm{d}t=\mathrm{d}\stackrel{~}{V}\frac{\mathrm{d}r}{1+\sqrt{\frac{2MrQ^2}{r^2+a^2}}},\mathrm{d}\varphi =\mathrm{d}\stackrel{~}{\varphi }\frac{a}{r^2+a^2}\frac{\mathrm{d}r}{1+\sqrt{\frac{2MrQ^2}{r^2+a^2}}}$$
(105)
The nonzero components of the lapse, shift, and 3-metric are then given by:
$$\alpha ^2=1$$
(106)
$$\beta ^r=\alpha ^2\sqrt{\frac{r^2+a^2}{\rho ^2}}\sqrt{\frac{2MrQ^2}{\rho ^2}}$$
(107)
$$\gamma _{rr}=\frac{\rho ^2}{r^2+a^2}$$
(108)
$$\gamma _{r\varphi }=\left[\sqrt{\frac{\rho ^2}{r^2+a^2}}\sqrt{\frac{2MrQ^2}{\rho ^2}}\right]a\mathrm{sin}^2\theta $$
(109)
$$\gamma _{\theta \theta }=\rho ^2$$
(110)
$$\gamma _{\varphi \varphi }=\left[r^2+a^2+\frac{2MrQ^2}{\rho ^2}a^2\mathrm{sin}^2\theta \right]\mathrm{sin}^2\theta $$
(111)
Notice that the lapse remains one, but the 3-geometry is no longer flat when the black hole is spinning. Cartesian coordinate components can be obtained from these via the standard Kerr–Schild coordinate transformations (92) and (93). Like the Kerr–Schild time slicing, a $`t=\text{const.}`$ slice of the generalized Painlevé–Gullstrand gauge remains spacelike for all $`r0`$.
### 3.4 Quasicircular Binary Data
One of the primary driving forces behind the development of black-hole initial data has been the two-body problem of general relativity: the inspiral and coalescence of a pair of black holes. This problem is of fundamental importance. Not only is the relativistic two-body problem the most fundamental dynamical problem of general relativity, it is also considered one of the most likely candidates for observation with the upcoming generation of gravitational wave laser interferometers. Because of the circularizing effects of gravitational radiation damping, we expect the orbits of most tight binary systems to have small eccentricities. It is therefore desirable to have a method that can discern which data sets, within the very large parameter space of binary black-hole initial-data sets, correspond to black-hole binaries in a nearly circular (quasicircular) orbit.
Currently, only one approach has been developed for locating quasicircular orbits in a parameter space of binary black-hole initial data . It is based on the fact that minimizing the energy of a binary system while keeping the orbital angular momentum fixed will yield a circular orbit in Newtonian gravity. The idea does not hold strictly for general relativistic binaries since they emit gravitational radiation and cannot be in equilibrium. However, for orbits outside the innermost stable circular orbit, the gravitational radiation reaction time scale is much longer than the orbital period. Thus it is a good approximation to treat such binaries as an equilibrium system. Called an “effective potential method”, this approach was used originally to find the quasicircular orbits and innermost stable circular orbit (ISCO) for equal-sized nonspinning black holes . In this work, the initial data for binary black holes were computed using the conformal-imaging approach outlined in §3.2.1. The approach was also applied to binary black-hole data computed using the puncture method , where similar results were found. Configurations containing a pair of equal-sized black holes with spin also have been examined . In this case, the spins of the black holes are equal in magnitude, but are aligned either parallel to, or anti-parallel to, the direction of the orbital angular momentum.
The approach defines an “effective potential” based on the binding energy of the binary. The binding energy is defined as
$$E_bE_{\mathrm{ADM}}M_1M_2,$$
(112)
where $`E_{\mathrm{ADM}}`$ is the total ADM energy of the system measured at infinity, and $`M_1`$ and $`M_2`$ are the masses of the individual black holes. Quasicircular orbital configurations are obtained by minimizing the effective potential (defined as the nondimensional binding energy $`E_b/\mu `$ (where $`\mu M_1M_2/(M_1+M_2)`$) as a function of separation, while keeping the the ratio of the masses of the black holes $`M_1/M_2`$, the spins of the black holes $`𝐒_1/M_1^2`$ and $`𝐒_2/M_2^2`$, and the total angular momentum $`J/M_1M_2`$ constant.
This approach is limited primarily by the ambiguity in defining the individual masses of the black holes, $`M_1`$ and $`M_2`$. There is no rigorous definition for the mass of an individual hole in a binary configuration and some approximation must be made here. There is also no rigorous definition for the individual spins of holes in a binary configuration. The problem of defining the individual masses becomes particularly pronounced when the holes are very close together (see Ref. ). The limiting choice of conformal flatness for the 3-geometry also has proven to be problematic. The effects of this choice have been clearly seen in the case of quasicircular orbits of spinning black holes , but it is also believed to be a serious problem for any binary configuration because binary configurations are not conformally flat at the second post-Newtonian order .
To date, the results of the effective-potential method have not matched well to the best result from post-Newtonian approximations <sup>14</sup><sup>14</sup>14We note that for the preferred choice of $`\omega _{\mathrm{static}}9`$ the gauge invariant parameters of the ISCO found in Ref. do not match well with the same parameter determined via the effective-potential method in Ref. . However, if $`0<\omega _{\mathrm{static}}<10`$, the agreement is much better.. It will be interesting to see if the results from post-Newtonian approximations and numerical initial-data sets converge, especially when the approximation of conformal flatness is eliminated.
## 4 Neutron-Star Initial Data
The construction of initial data for neutron stars requires that the state of the neutron-star matter be specified before the gravitational data can be determined. Of course, a solution of the gravitational constraint equations can be found in principal for any given energy density $`\rho `$ and momentum density $`j^i`$. But, with neutron-star solutions, we are usually interested in situations where the matter is in (or nearly in) hydrostatic equilibrium and the gravitational fields are also in (or nearly in) equilibrium.
### 4.1 Hydrostatic Equilibrium
For a neutron star to be in true equilibrium, the spacetime must be stationary as discussed in §2.4. This means that the spacetime possesses both “temporal” and “angular” Killing vectors (cf. Ref. ). If the matter is also to be in equilibrium, then the 4-velocity of the matter $`u^\mu `$ must be a linear combination of these two Killing vectors. If we use coordinates as defined in (53) with the angular Killing vector in the $`\varphi `$ direction, then
$$u^\mu =u^t[1,0,0,\mathrm{\Omega }].$$
(113)
Here, $`u^t`$ and $`\mathrm{\Omega }`$ are functions of $`r`$ and $`\theta `$ only. $`\mathrm{\Omega }`$ is the angular velocity of the matter as measured at infinity.
It is common to define $`v`$ as the relative velocity between the matter and a normal observer (often called a zero angular momentum observer) so that
$$\frac{1}{\sqrt{1v^2}}=n_\mu u^\mu =\alpha u^t.$$
(114)
The velocity $`v`$ is then fixed by the normalization condition $`u^\mu u_\mu =1`$.
If we assume that the matter source is a perfect fluid, then the stress-energy tensor is given by
$$T^{\mu \nu }=(\epsilon +P)u^\mu u^\nu +Pg^{\mu \nu },$$
(115)
where $`\epsilon `$ and $`P`$ are the total energy density and pressure, respectively, as measured in the rest frame of the fluid. The vanishing of the divergence of the stress-energy tensor yields the equation of hydrostatic equilibrium (often referred to as the relativistic Bernoulli equation). In differential form, this is
$$\mathrm{d}P(\epsilon +P)(\mathrm{d}\mathrm{ln}u^tu^tu_\varphi \mathrm{d}\mathrm{\Omega })=0.$$
(116)
If the fluid is barytropic<sup>15</sup><sup>15</sup>15For a barytropic fluid, the entropy per baryon and the fractional abundances of the different nuclear species are determined uniquely by the distribution of baryons. In this case, the total energy density $`\epsilon `$ can be expressed as a function of the pressure $`P`$ (cf. Ref. )., then we can define the relativistic *enthalpy* as
$$h(P)\mathrm{exp}\left[_0^P\frac{dP}{\epsilon +P}\right],$$
(117)
and rewrite the relativistic Bernoulli equation as
$$\mathrm{ln}h(P)=\mathrm{ln}h_0+\mathrm{ln}u^t\mathrm{ln}u_0^t_{\mathrm{\Omega }_0}^\mathrm{\Omega }u^tu_\varphi 𝑑\mathrm{\Omega }.$$
(118)
The constants $`h_0`$, $`u_0^t`$, and $`\mathrm{\Omega }_0`$ are the values their respective quantities have at some reference point, often taken to be the surface of the neutron star at the axis of rotation. When uniform rotation is assumed ($`\mathrm{d}\mathrm{\Omega }=0`$), Eq. (118) is rather easy to solve. The case of differential rotation is somewhat more complicated. An integrability condition of (116) requires that $`u^tu_\varphi `$ be expressible as a function of $`\mathrm{\Omega }`$, so
$$u^tu_\varphi F(\mathrm{\Omega }).$$
(119)
$`F(\mathrm{\Omega })`$ is a specifiable function of $`\mathrm{\Omega }`$ which determines the rotation law that the neutron star must obey .
### 4.2 Isolated Neutron Stars
The simplest models of isolated neutron stars are static (i.e., nonrotating), spherically symmetric models that can be constructed, given a suitable equation of state, by solving the Oppenheimer–Volkoff (OV) equations :
$`{\displaystyle \frac{dP}{dr}}`$ $`=`$ $`{\displaystyle \frac{(\epsilon +P)(m+4\pi r^3P)}{r(r2m)}}`$
$`{\displaystyle \frac{dm}{dr}}`$ $`=`$ $`4\pi r^2\epsilon `$ (120)
$`{\displaystyle \frac{d\alpha }{dr}}`$ $`=`$ $`{\displaystyle \frac{\alpha (m+4\pi r^3\epsilon )}{r(r2m)}}`$
for $`0rR`$, where $`R`$ is the radius of the surface of the star. Here, $`r`$ is an areal radius and $`m(r)`$ is the mass inside radius $`r`$. Exterior to the surface of the star, the metric is the standard Schwarzschild metric as in Eq. (59) with $`Mm(R)`$. Interior to the surface of the star, the metric is
$$\mathrm{d}s^2=\alpha ^2\mathrm{d}t^2+\left(1\frac{2m(r)}{r}\right)^1\mathrm{d}r^2+r^2(\mathrm{d}\theta ^2+\mathrm{sin}^2\theta \mathrm{d}\varphi ^2).$$
(121)
The boundary conditions are that $`m(0)=0`$, $`\epsilon (0)`$ is some chosen constant $`\epsilon _0`$, and $`\alpha (R)=\sqrt{1\frac{2M}{R}}`$. The solutions of this equation form a one-parameter family, parameterized by $`\epsilon _0`$ which determines how relativistic the system is. A method for solving these equations in both the areal coordinate $`r`$ and an isotropic radial coordinate $`\stackrel{~}{r}`$ can be found in Ref. .
More generally, isolated neutron stars will be rotating. If the neutron stars are uniformly rotating, then, for any given equation of state, the solutions form a two-parameter family. These models can be parameterized by their central density, which determines how relativistic they are, and by the amount of rotation. If the models are allowed to have differential rotation, then some rotation law must be chosen.
To construct a neutron-star model, the equations for a stationary solution of Einstein’s equations outlined in §2.4 must be solved self-consistently with the equations for hydrostatic equilibrium of the matter outlined above in §4.1. The equations that must be solved depend on the form of the metric chosen, and numerous formalisms and numerical schemes have been used. An incomplete list of references to work on constructing neutron-star models include . Further review information on neutron-star models can be found in Refs. .
### 4.3 Neutron-Star Binaries
Neutron-star binaries, like any relativistic binary system, cannot exist in a true equilibrium configuration since they must emit gravitational radiation. But, as is true for black-hole binary data, for orbits outside the innermost stable circular orbit, the gravitational radiation reaction time scale is much longer than the orbital period and it is a reasonable approximation to consider the stars to be in a quasiequilibrium state.
A binary configuration obviously lacks the azimuthal symmetry that was assumed in the discussions of stationary solutions of Einstein’s equations in §2.4 and hydrostatic equilibrium in §4.1. Fortunately, the condition of hydrostatic equilibrium requires only the presence of a single, timelike Killing vector. With the assumption that gravitational radiation is negligible, we can assume that the matter is in some equilibrium state as viewed from the reference frame that is rotating along with the binary. That is, if the binary has a constant orbital angular velocity of $`\mathrm{\Omega }`$, then the time vector in the rotating frame $`t^\mu `$ is a Killing vector<sup>16</sup><sup>16</sup>16Bonazzola et al refer to this as a *helicoidal* Killing vector. and it is related to the time vector in the rest frame of the binary $`t^\mu `$ by
$$t^\mu =t^\mu +\mathrm{\Omega }\xi ^\mu ,$$
(122)
where $`\xi ^\mu `$ is a generator of rotations about the rotation axis<sup>17</sup><sup>17</sup>17In spherical coordinates, $`\stackrel{}{\xi }=/\varphi `$. In Cartesian coordinates, rotation about the $`z`$ axis would be represented by $`\xi ^\mu =(0,y,x,0)`$..
Two equilibrium states for the matter have been explored in the literature. The simplest case is that of *co*-rotation, where the 4-velocity of the matter is proportional to $`t^\mu `$. In this case, the matter is at rest in the frame of reference rotating with the binary system, the *corotating reference frame*. The second equilibrium state is that of *counter*-rotation, where there is no rotation in the *rest frame of the binary*. We will explore these two cases further below.
Stationarity of the gravitational field, unlike hydrostatic equilibrium, requires the presence of separate timelike and azimuthal Killing vectors. For the case of a binary system, there is no unique definition of quasiequilibrium. The earliest work on constructing quasiequilibrium solutions of Einstein’s equations stems from work by Wilson and Mathews , and others have explored similar schemes . Although written in slightly different forms, the system of equations for the gravitational fields in all of these schemes are fundamentally identical. While they were developed before the thin-sandwich decomposition of §2.3, the thin-sandwich decomposition (see Eq. (2.3)) offers the easiest way to interpret this approach. We consider ourselves to be in the corotating reference frame so that our time vector is $`t^\mu `$. To make the transition back to the rest frame of the binary as easy as possible, we write the shift vector of our $`3+1`$ decomposition as
$$B^i=\beta ^i+\mathrm{\Omega }\xi ^i,$$
(123)
so that $`\beta ^i`$ remains as the shift vector of the $`3+1`$ decomposition made with respect to the rest frame of the binary system.
The primary assumptions are that the conformal 3-metric $`\stackrel{~}{\gamma }_{ij}`$ is flat, the initial-data slice is maximal so that $`K=0`$, and $`\stackrel{~}{u}^{ij}=0`$. We see from (42) that the last assumption implies that the conformal 3-geometry is *instantaneously* stationary as seen in the corotating reference frame. The final choice that must be made is for the conformally rescaled lapse $`\stackrel{~}{\alpha }`$. An elliptic equation for the lapse can be obtained by demanding the trace of the extrinsic curvature $`K`$ also be instantaneously stationary in the corotating reference frame. This is the so-called *maximal slicing* condition on the lapse. For the particular assumptions we have made here, this equation can be written
$$\stackrel{~}{}^2(\stackrel{~}{\alpha }\psi ^7)=(\stackrel{~}{\alpha }\psi ^7)\left[\frac{7}{8}\psi ^7\stackrel{~}{A}_{ij}\stackrel{~}{A}^{ij}+2\pi G\psi ^4(\rho +2S)\right].$$
(124)
It is interesting to note that, for a conformally flat 3-geometry,
$$(\stackrel{~}{𝕃}B)^{ij}=(\stackrel{~}{𝕃}\beta )^{ij},$$
(125)
so $`\mathrm{\Omega }`$ does not appear in the equations for the gravitational fields except in the matter terms and possibly in boundary conditions. Equations (2.3) and (124) can be solved for the gravitational fields on an initial-data hypersurface, given values for the matter terms and appropriate boundary conditions.
*If* we choose the matter so that it is in hydrostatic equilibrium with respect to the pseudo-Killing vector $`t^\mu `$, then these equations for the gravitational fields will yield data that are in quasiequilibrium in the sense that $`\stackrel{~}{\gamma }_{ij}`$ and $`K`$ are both *instantaneously* stationary.
For corotating binaries, the matter is at rest in the corotating reference frame of the binary. It is *rigidly* rotating and hydrostatic equilibrium is specified by solving the relativistic Bernoulli equation (118), with $`\mathrm{d}\mathrm{\Omega }=0`$, self-consistently with the equations for the gravitational fields.
For counterrotating binaries, the matter is not rotating in the rest frame of the binary. Counterrotating equilibrium configurations can be obtained by assuming the matter to have *irrotational flow* . As long as the flow is *isentropic*, we can express the enthalpy (117) as<sup>18</sup><sup>18</sup>18For isentropic flow, the thermodynamic identity reduces to $`\mathrm{d}h=\frac{1}{\rho _0}\mathrm{d}P`$.
$$h=\frac{\epsilon +P}{\rho _0},$$
(126)
where $`\rho _0`$ is the rest-mass density. For irrotational flow, the vorticity of the fluid (cf. Ref. ) is zero. Combining this with the Euler equation, we find that the 4-velocity of the fluid can be expressed as
$$hu_\mu =_\mu \phi ,$$
(127)
where $`\phi `$ is the velocity potential, or flow field. Equation (127), together with the normalization condition $`u^\mu u_\mu =1`$, automatically satisfies the Euler equation and we are left with the continuity equation which must be satisfied,
$$_\mu (\rho _0u^\mu )=0.$$
(128)
The continuity equation (128), the normalization condition, and Eq. (127) yield
$$^2\phi =(^\mu \phi )_\mu \mathrm{ln}(\rho _0/h)\text{with}h=\sqrt{(^\nu \phi )(_\nu \phi )}.$$
(129)
Stationarity (or quasistationarity) in the corotating reference frame requires that
$$hu_\mu t^\mu =C,$$
(130)
where $`C`$ is a positive constant. Now, in terms of the $`3+1`$ decomposition with $`B^i`$ as the shift vector (see Eq. (123)), we find that the Bernoulli equation is written
$$h^2=(\overline{}^i\phi )\overline{}_i\phi +\frac{1}{\alpha ^2}\left(C+B^j\overline{}_j\phi \right)^2.$$
(131)
The flow field $`\phi `$ must satisfy
$`\overline{}^2\phi B^i\overline{}_i\left({\displaystyle \frac{\lambda }{\alpha ^2}}\right){\displaystyle \frac{\lambda }{\alpha }}K`$ $`=`$ $`\left(\overline{}^i\phi {\displaystyle \frac{\lambda }{\alpha ^2}}B^i\right)\overline{}_i\mathrm{ln}\left({\displaystyle \frac{\alpha \rho _0}{h}}\right)`$ (132)
$`\lambda `$ $``$ $`C+B^i\overline{}_i\phi `$
subject to the boundary condition<sup>19</sup><sup>19</sup>19This boundary condition comes from the fact that the fluid motion at the surface of the star must be tangent to the surface, $`u^\mu _\mu \rho _0=0`$. at the surface of the flow
$$\left(\overline{}^i\phi \frac{\lambda }{\alpha ^2}B^i\right)\overline{}_i\rho _0|_{\mathrm{surf}}=0.$$
(133)
Solving Eqs. (131) and (132) self-consistently with the equations for the gravitational fields yields a counterrotating binary in hydrostatic equilibrium.
As mentioned above, the earliest work on neutron-star binaries was carried out by Wilson and Mathews . Wilson et al describe their approach for generating initial data for equilibrium neutron-star binaries. In these early works, the equation of hydrostatic equilibrium was not used. Rather, an initial guess for the density profile was chosen and the full hydrodynamic system was evolved with viscous damping until equilibrium was reached. During each step of the hydrodynamic evolution, the equations for the gravitational fields were resolved. The resulting data represented neither strictly co- nor counter-rotating binary neutron stars. This work led to the controversial result that each neutron star in the binary may become radially unstable and collapse prior to the merger of the pair of stars. While an error was found in this work with the result that the signature of collapse is significantly weaker, the controversy has not yet been completely resolved.
The first use of corotating hydrostatic equilibrium with the Wilson–Mathews approach for specifying the gravitational fields was by Cook et al for the test case of an isolated neutron star. This approach was then used to study corotating neutron-star binaries by Baumgarte et al and by Marronetti et al. Interestingly, turning-point methods for detecting *secular* instabilities can be applied to the case of corotating binaries .
Corotating binary configurations are relatively easy to construct. However, it is believed that the viscosity of neutron-star matter is not large enough to allow for synchronization of the spin with the orbit . But, if the initial spins of the neutron stars are not too large, close binaries should be well approximated by irrotational models. Bonazzola et al (as corrected by Asada ) developed the first approach for constructing counterrotating binary configurations. However, simpler formulations of irrotational flow were developed independently by Teukolsky and Shibata , and Gourgoulhon showed that all three approaches were equivalent. Numerical solutions of the equations for irrotational flow coupled to the equations for the gravitational fields are more difficult to construct than those for corotation because of the boundary condition (133) on the flow field that must be applied on the surface of each neutron star. This boundary condition is particularly difficult to implement because the location of the surface of the star is not known a priori, and will move as the equations are being solved. The first models of irrotational binary neutron stars were constructed by Bonazzola et al, Marronetti et al, and Uryū and Eriguchi . A full description of the numerical methods used by Bonazzola can be found in Ref. .
## 5 Acknowledgments
This work was supported by NSF grant PHY-9988581 to Wake Forest University. I am especially indebted to James York, Saul Teukolsky, and Stuart Shapiro for their support over the years. I would also like to thank Andrew Abrahams, Thomas Baumgarte, and Mark Scheel for innumerable helpful discussions.
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# 1 Introduction
## 1 Introduction
The realm of integrable systems in 2d statistical physics or (1+1)d QFT has attracted much attention for a long time. Important contributions in this development were the solution of the anisotropic $`sl(2)`$ Heisenberg model ($`XYZ`$-model) by means of the Bethe ansatz and its subsequent reformulation in an algebraic version, the algebraic Bethe ansatz or so called Quantum Inverse Scattering Mathod (QISM) . Generalizations to higher rank groups, which allow to treat models with internal degrees of freedom, were in this context achieved in . The efforts resulted in a plethora of models soluble by means of the QISM, insofar that eigenvalues and eigenvectors of the Hamiltonian were found explicitely. This method also stimulated the investigation of various areas in mathematical physics, such as quantum groups, theory of knots etc. Despite the indisputable achievements of the QISM and the rather simple action of the inverse problem operators, which can be interpreted as creation and annihilation operators for quasiparticles and as generating functions for the conserved quantities respectively, the study of correlation functions and formfactors has proven to be rather intricate. This is partly due to the fact that the solution of the inverse problem (expressing the original microscopic operators by means of the operators figuring in the algebraic Bethe ansatz) has only been achieved recently , and by the fact that the action of the quasiparticle creation and annihilation operators which figure in the construction of the eigenfunctions (Bethe wavevectors) is obscured on the level of the microscopic variables (spin raising and lowering operators) by nonlocal effects arising from polarization clouds (compensating pairs of local raising and lowering spin operators). In a seminal paper Maillet and Sanchez de Santos revealed an application of so called Drinfel’d twists to obtain a basis for the $`sl(2)`$ $`XXX`$ and $`XXZ`$ model which allows to express the creation and annihilation operators in a completely symmetric way with the further advantage of being polarization free, that is being built from the respective quasiclassical Gaudin operators dressed diagonally, thus supressing non-local effects for these operators. A generalized transformation has subsequently been applied to the $`sl(2)`$ $`XYZ`$ model and been used to resolve the nested hierarchy in the Bethevectors of the $`sl(n)`$ $`XXX`$ model .
In this paper we report the construction of suitable Drinfel’d twists for the $`_n\times _n`$ symmetric Belavin model , which can be thought of as a n-state generalization of the $`sl(2)`$ $`XYZ`$ model. Our results generalize the findings for the $`sl(n)`$ $`XXX`$ model and provide a completely symmetric representation of the creation operators as well as a resolution of the hierarchical structure of the Bethe wavevectors.
The paper is organized as follows: Section 2 provides a short survey of the Belavin model and its reformulation as an face-type model. Section 3 deals with the factorizing twists and the computation of the operator valued entries of the monodromy matrix. In section 4 we discuss the Bethe wavevectors and section 5 contains the conclusions.
## 2 Belavin model and the corresponding IRF model
A possible n-state generalization of the eight-vertex model is given by the $`_n\times _n`$ symmetric model of Belavin , whose Boltzmann weights fulfill the Yang–Baxter equation (YBE) ,
$`S(z_1z_2)_{k_1,k_2}^{i_1,i_2}S(z_1z_3)_{j_1,k_3}^{k_1,i_3}S(z_2z_3)_{j_2,j_3}^{k_2,k_3}=S(z_2z_3)_{k_2,k_3}^{i_2,i_3}S(z_1z_3)_{k_1,j_3}^{i_1,k_3}S(z_1z_2)_{j_1,j_2}^{k_1,k_2}`$ (1)
(here and subsequently double indices $`k`$ mean summation over $`0,1,\mathrm{},n1`$ unless stated otherwise). It can be parametrized in the following way <sup>1</sup><sup>1</sup>1The slight difference compared to originates from the normalization of the S-matrix.
$`S(z)_{i,j}^{k,l}=\{\begin{array}{cc}0\text{if}i+jk+l\hfill & \\ \frac{h(w)h(z)\theta ^{(ij)}(z+w)}{h(z+w)\theta ^{(ik)}(w)\theta ^{(kj)}(z)}\text{if}i+j=k+l;\text{mod}n\hfill & \end{array}`$ (4)
with
$`h(z)={\displaystyle \underset{i=0}{\overset{n1}{}}}\theta ^{(i)}(z)\left({\displaystyle \underset{i=1}{\overset{n1}{}}}\theta ^{(i)}(0)\right)^1`$ (5)
where $`\theta ^{(i)}(z)=\theta \left(\genfrac{}{}{0pt}{}{\frac{1}{2}\frac{i}{n}}{\frac{1}{2}}\right)(z)`$ represents the theta function of rational characteristics $`1/2i/n,1/2`$. The theta function of characteristics $`a,b`$ is given as a Fourier series ($`\tau `$ is a fixed complex number in the upper half plane, and $`\mathrm{\Lambda }_\tau +\tau `$ is the lattice generated by $`1`$ and $`\tau `$)
$`\theta \left({\displaystyle \genfrac{}{}{0pt}{}{a}{b}}\right)(z)={\displaystyle \underset{m}{}}exp\left(\pi i\tau (m+a)z+2\pi i(m+a)(z+b)\right)`$ (6)
and has zeros at $`z=(\frac{1}{2}b)+(\frac{1}{2}a)\tau \text{mod}\mathrm{\Lambda }_\tau `$.
The matrix (4) is unitary, $`S(z)S(z)=1\mathrm{I}`$ and obeys the initial condition $`S(0)=P`$ , where $`P`$ is the permutation operator.
There exists a vertex-face map which transforms the Belavin model into a face type model , which in turn can be thought of as a multicomponent generalization of the six-vertex model.
The correspondence is given by
$`S(z_1z_2)M_a^\mu (z_1)M_{a+\widehat{\mu }}^\nu (z_2)=M_{a+\widehat{\nu }^{}}^\mu ^{}(z_1)M_a^\nu ^{}(z_2)R(z_1z_2|a)_{\mu ^{},\nu ^{}}^{\mu ,\nu }`$ (7)
where $`R(a|z)`$ has the form
$`R(z|a)_{\mu ^{},\nu ^{}}^{\mu ,\nu }`$ $`=`$ $`b^{\mu ,\nu }(z|a)\delta _{\mu ^{}\mu }\delta _{\nu ^{}\nu }+c^{\mu ,\nu }(z|a)\delta _{\mu ^{}\nu }\delta _{\nu ^{}\mu }`$ (8)
$`b^{\mu ,\nu }(z|a)`$ $`=`$ $`{\displaystyle \frac{h(z)}{h(z+w)}}{\displaystyle \frac{h(s^\nu s^\mu +w(a\widehat{\nu })^{\nu ,\mu })}{h(s^\nu s^\mu +w(a\widehat{\nu })^{\nu ,\mu }+w)}}{\displaystyle \underset{i=0}{\overset{n1}{}}}{\displaystyle \frac{g_{i\nu }(a)}{g_{i\nu }(a+\widehat{\mu })}}{\displaystyle \underset{i=0}{\overset{n1}{}}}{\displaystyle \frac{g_{i\mu }(a+\widehat{\nu })}{g_{i\mu }(a)}}`$
$`c^{\mu ,\nu }(z|a)`$ $`=`$ $`{\displaystyle \frac{h(z+w)}{h(z)}}{\displaystyle \frac{h(s^\nu s^\mu +w(a\widehat{\nu })^{\nu ,\mu }+w+z)}{h(s^\nu s^\mu +w(a\widehat{\nu })^{\nu ,\mu }+w)}}`$ (9)
$`g_{i\mu }(a)`$ $`=`$ $`\{\begin{array}{cc}1\text{if}i\mu \hfill & \\ h(s^0s^\mu +wa^{0,\mu })\text{if}i=0\hfill & \\ h(s^is^\mu +wa^{i,\mu }w)\text{if}\mathrm{\hspace{0.33em}0}<i<\mu \hfill & \end{array}`$ (13)
where $`s^\mu `$ are arbitrary complex numbers, $`a_{\mu =0}^{n1}\mathrm{\Lambda }_\mu `$ with $`\mathrm{\Lambda }_\mu `$ weights of the affine Lie algebra $`A_{n1}^{(1)}`$, $`\widehat{\mu }=\mathrm{\Lambda }_{\mu +1}\mathrm{\Lambda }_\mu ,\mu =0,1,\mathrm{},n1`$ and $`\mathrm{\Lambda }_0=\mathrm{\Lambda }_n`$. The $`a^{\mu ,\nu }`$ is given in and obeys $`(a\widehat{\nu })^{\mu ,\mu }=0`$.
The intertwining vector $`M_a^\mu (z)`$ is given by
$`M_a^\mu (z)`$ $`=`$ $`{}_{}{}^{t}(\varphi _a^\mu (z)_0,\mathrm{},\varphi _a^\mu (z)_{n1}){\displaystyle \underset{i=0}{\overset{n1}{}}}g_{i\mu }^1(a)`$ (14)
$`\varphi _a^\mu (z)_j`$ $`=`$ $`\theta ^{(j)}(znwa^{\mu ,0}).`$ (15)
Applying a sum of threefold tensorproducts of functions (14) to (1) using the vertex face map (7) and exploiting the explicit index structure of (8) we obtain a modified Yang-Baxter equation
$`R(|z_1z_2|a+\widehat{\alpha }_3)_{\mu _1,\mu _2}^{\alpha _1,\alpha _2}R(z_1z_3|a)_{\beta _1,\mu _3}^{\mu _1,\alpha _3}R(z_2z_3|a+\widehat{\beta }_1)_{\beta _2,\beta _3}^{\mu _2,\mu _3}`$ (16)
$`=`$ $`R(z_2z_3|a)_{\mu _2,\mu _3}^{\alpha _2,\alpha _3}R(z_1z_3|a+\widehat{\mu }_2)_{\mu _1,\beta _3}^{\alpha _1,\mu _3}R(z_1z_2|a)_{\beta _1,\beta _2}^{\mu _1,\mu _2}.`$
This can be written symbolically as (we set $`R(z_iz_j|a)=R_{ij}(a)`$)
$`R_{12}(a+\widehat{\sigma }_3)R_{13}(a)R_{23}(a+\widehat{\sigma }_1)=R_{23}(a)R_{13}(a+\widehat{\sigma }_2)R_{12}(a).`$
The monodromy matrix for a chain with $`N`$ sites associated with the R-matrix (8) is
$`T_{0,1\mathrm{}N}(\lambda |a)=R_{0N}(a+\widehat{\sigma }_1+\mathrm{}+\widehat{\sigma }_{N1})\mathrm{}R_{02}(a+\widehat{\sigma }_1)R_{01}(a)`$ (17)
where $`0`$ denotes the horizontal auxiliary space with spectral parameter $`\lambda `$; $`1,\mathrm{},N`$ label the vertical quantum spaces (with local inhomogeneities $`\left\{z_i\right\}`$) whose tensorproduct constitutes the physical Hilbertspace $`_N`$. An algebraic construction of eigenvalues of the monodromy matrix associated to the matrix (4) was performed in and parallels the procedure for the eight-vertex model in (One has to take into account that our construction of the monodromy matrix (17) differs from the monodromy matrix in by an additional transformation in the physical space).
Subsequently we will focus on the construction of a factorizing $`F`$ matrix for the R-matrix (8) and pursue its consequences for the structure of the monodromy matrix (17).
## 3 The F basis
The factorizing $`F`$-matrix for $`N`$ sites ($`N`$ quantum spaces), being defined by the relation $`R_{1\mathrm{}N}^\sigma (a)=F_{\sigma (1\mathrm{}N)}^1(a)F_{1\mathrm{}N}(a)`$ , turns out to be given by formally the same expression as found in for the $`sl(n)`$ $`XXX`$-model
$`F_{1\mathrm{}N}(a)`$ $`=`$ $`{\displaystyle \underset{\sigma S_N}{}}{\displaystyle \underset{\alpha }{\overset{}{}}}{\displaystyle \underset{i=1}{\overset{N}{}}}P_{\sigma (i)}^{\alpha _{\sigma (i)}}R_{1\mathrm{}N}^{\sigma _\alpha }(z_1,\mathrm{},z_N|a)`$ (18)
where $`P_i^{\alpha _i}`$ projects on the $`\alpha _i`$-th component in the $`i`$-th space and the labels $`\alpha _{\sigma (i)}`$ satisfy the conditions
$`\alpha _{\sigma (i+1)}`$ $``$ $`\alpha _{\sigma (i)}\text{if}\sigma (i+1)>\sigma (i)`$
$`\alpha _{\sigma (i+1)}`$ $`>`$ $`\alpha _{\sigma (i)}\text{if}\sigma (i+1)<\sigma (i).`$ (19)
The modification of the Yang-Baxter equation (16) requires a particular rule for the handling of the parameter $`a`$ in the formation of the intertwining matrix $`R^\sigma (a)`$ (related to the permutation $`\sigma `$), which can be read off from the modified composition law
$`R^{\sigma \sigma _i}(a)`$ $`=`$ $`R_{\sigma (i),\sigma (i+1)}(\stackrel{~}{a}_i)R^\sigma (a)`$
$`\stackrel{~}{a}_i`$ $`=`$ $`a+\widehat{\sigma }_{\sigma (1)}+\mathrm{}+\widehat{\sigma }_{\sigma (i1)}`$ (20)
where $`\sigma _i`$ is the transposition of $`i,i+1`$, and $`\sigma `$ an arbitrary permutation.
$`R^\sigma (a)`$ has the intertwining property
$`R^\sigma (a)T_{0,1\mathrm{}N}(a)=T_{0,\sigma (1)\mathrm{}\sigma (N)}(a)R^\sigma (a+\widehat{\sigma }_h)`$ (21)
where $`\sigma _h`$ is associated with the matrix indices in the space 0.
The matrix $`F_{1\mathrm{}N}(a)`$ satisfies the factorizing equation
$`R_{1\mathrm{}N}^\sigma (a)=F_{\sigma (1\mathrm{}N)}^1(a)F_{1\mathrm{}N}(a).`$ (22)
A proof of the latter equation can be found in and relies on the fact that
$`P_i^\alpha P_j^\alpha R_{ij}=P_i^\alpha P_j^\alpha 1\mathrm{I}_{ij}.`$ (23)
The modification of the composition law induced by the parameter $`a`$ being insignificant.
We will proceed by computing elements of the monodromy matrix in the new basis provided by the F-matrix. The transformation law $`\stackrel{~}{T}_{0,1\mathrm{}N}(a)=F_{1\mathrm{}N}(a)T_{0,1\mathrm{}N}(a)F_{1\mathrm{}N}^1(a+\widehat{\sigma }_h)`$ is enforced by the requirement that the resulting operator is symmetric, i.e. $`\stackrel{~}{T}_{0,1\mathrm{}N}(a)=\stackrel{~}{T}_{0,\sigma (1)\mathrm{}\sigma (N)}(a)`$, which then follows from (21).
The computation of the diagonal element $`T_{n1n1}(\lambda |a)`$ proceeds along the same lines as that in the generalized $`sl(n)`$ $`XXX`$-model in . Let us consider the action of the matrix $`F(a)`$ on $`T_{n1n1}(a)`$:
$`F_{1\mathrm{}N}(a)T_{n1n1}(a)`$ $`=`$ $`{\displaystyle \underset{\sigma S_N}{}}{\displaystyle \underset{\alpha _{\sigma (1)}\mathrm{}\alpha _{\sigma (N)}}{\overset{}{}}}{\displaystyle \underset{i=1}{\overset{N}{}}}P_{\sigma (i)}^{\alpha _{\sigma (i)}}R_{1\mathrm{}N}^\sigma (a)P_0^{n1}T_{0,1\mathrm{}N}(a)P_0^{n1}`$ (24)
$`=`$ $`{\displaystyle \underset{\sigma S_N}{}}{\displaystyle \underset{\alpha _{\sigma (1)}\mathrm{}\alpha _{\sigma (N)}}{\overset{}{}}}{\displaystyle \underset{i=1}{\overset{N}{}}}P_{\sigma (i)}^{\alpha _{\sigma (i)}}P_0^{n1}T_{0,\sigma (1)\mathrm{}\sigma (N)}(a)P_0^{n1}R_{1\mathrm{}N}^\sigma (a+\sigma _h).`$
The specialization to the entry $`(n1,n1)`$ of the auxiliary space here is achieved by the projectors $`P_0^{n1}`$. For the second equality in (24) we have used relation (21) and the obvious fact that $`P_0^{n1}`$ commutes with $`R_{1\mathrm{}N}^\sigma (a)`$. To simplify the following argument we distinguish in the sum $`^{}`$ cases of various multiplicities of the occurrence of the group index $`n1`$:
$`F_{1\mathrm{}N}(a)T_{n1n1}(a)`$ $`=`$ $`{\displaystyle \underset{\sigma S_N}{}}{\displaystyle \underset{k=0}{\overset{N}{}}}{\displaystyle \underset{\alpha _{\sigma (1)}\mathrm{}\alpha _{\sigma (N)}}{\overset{^{}}{}}}{\displaystyle \underset{j=Nk+1}{\overset{N}{}}}P_{\sigma (j)}^{n1}\delta _{\alpha _{\sigma (j)},n1}{\displaystyle \underset{j=1}{\overset{Nk}{}}}P_{\sigma (j)}^{\alpha _{\sigma (j)}}`$ (25)
$`\times P_0^{n1}T_{0,\sigma (1)\mathrm{}\sigma (N)}(a)P_0^{n1}R_{1\mathrm{}N}^\sigma (a+\sigma _h).`$
Let us consider the prefactor of $`R_{1\mathrm{}N}^\sigma (a+\sigma _h)`$ on the r.h.s. of Eq. (25) more closely. Using specific features of the $`R`$-matrices we can rewrite it as follows:
$`{\displaystyle \underset{j=1}{\overset{Nk}{}}}P_{\sigma (j)}^{\alpha _{\sigma (j)}}{\displaystyle \underset{j=Nk+1}{\overset{N}{}}}P_{\sigma (j)}^{n1}P_0^{n1}T_{0,\sigma (1)\mathrm{}\sigma (N)}(a)P_0^{n1}`$ (26)
$`=`$ $`{\displaystyle \underset{j=1}{\overset{Nk}{}}}P_{\sigma (j)}^{\alpha _{\sigma (j)}}\left(R_{0,\sigma (N)}(a+{\displaystyle \underset{m=1}{\overset{N1}{}}}\widehat{\sigma }_{\sigma (m)})\right)_{n1n1}^{n1n1}\mathrm{}\left(R_{0,\sigma (Nk+1)}(a)\right)_{n1n1}^{n1n1}`$
$`\times P_0^{n1}T_{0,\sigma (1)\mathrm{}\sigma (Nk)}(a)P_0^{n1}{\displaystyle \underset{j=Nk+1}{\overset{N}{}}}P_{\sigma (j)}^{n1}`$
$`=`$ $`{\displaystyle \underset{j=1}{\overset{Nk}{}}}P_{\sigma (j)}^{\alpha _{\sigma (j)}}P_0^{n1}T_{0,\sigma (1)\mathrm{}\sigma (Nk)}(a)P_0^{n1}{\displaystyle \underset{j=Nk+1}{\overset{N}{}}}P_{\sigma (j)}^{n1}`$
$`=`$ $`{\displaystyle \underset{i=1}{\overset{Nk}{}}}\left(R_{0\sigma (i)}(a+{\displaystyle \underset{m=1}{\overset{i1}{}}}\widehat{\sigma }_{\sigma (m)})\right)_{n1,\alpha _{\sigma (i)}}^{n1,\alpha _{\sigma (i)}}{\displaystyle \underset{j=1}{\overset{Nk}{}}}P_{\sigma (j)}^{\alpha _{\sigma (j)}}{\displaystyle \underset{j=Nk+1}{\overset{N}{}}}P_{\sigma (j)}^{n1}P_0^{n1}`$
Inserting the r.h.s. of (26) into Eq. (25) one sees that the product $`_i\left(R_{0\sigma (i)}(a+_{m=1}^{i1}\widehat{\sigma }_{\sigma (m)})\right)_{n1,\alpha _{\sigma (i)}}^{n1,\alpha _{\sigma (i)}}`$ creates a diagonal dressing factor for $`\stackrel{~}{T}_{n1n1}(a)`$ and the product of projectors applied to $`R^\sigma (a+\sigma _h)`$ gives $`F_{1\mathrm{}N}(a+\sigma _h)`$. One obtains
$`\stackrel{~}{T}_{n1n1}(\lambda |a)=Y_{n1}(a)_{i=1}^N\text{diag}\{b(\lambda z_i),\mathrm{},b(\lambda z_i),1\}.`$ (27)
with the abbreviations
$`Y_j(a)`$ $`=`$ $`{\displaystyle \underset{m=0}{\overset{j1}{}}}\{{\displaystyle \underset{i=0}{\overset{n1}{}}}{\displaystyle \frac{g_{ij}(\stackrel{~}{a_m}+k_m\widehat{m})}{g_{ij}(\stackrel{~}{a_m})}}{\displaystyle \underset{j_m=1}{\overset{k_m}{}}}[{\displaystyle \underset{i=0}{\overset{n1}{}}}{\displaystyle \frac{g_{im}(\stackrel{~}{a_m}+(j_m1)\widehat{m})}{g_{im}(\stackrel{~}{a_m}+(j_m1)\widehat{m}+\widehat{j})}}\times `$
$`\times {\displaystyle \frac{h(s^ms^j+w[\stackrel{~}{a_m}+(j_m2)\widehat{m}]^{m,j})}{h(s^ms^j+w[\stackrel{~}{a_m}+(j_m2)\widehat{m}]^{m,j}+w)}}]\}`$
$`\stackrel{~}{a_m}`$ $`=`$ $`a+{\displaystyle \underset{i=0}{\overset{m1}{}}}k_i\widehat{i}`$
$`b(\lambda )`$ $`=`$ $`{\displaystyle \frac{h(\lambda )}{h(\lambda +w)}}`$ (28)
and $`k_m`$ gives the multiplicity of the upper matrix labels $`\alpha _i=m`$.
To compute $`T_{n1n2}(\lambda |a)`$ one has to distinguish in the sum $`^{}`$ cases of various multiplicities $`k_{n1}`$ and $`k_{n2}`$ of the occurrence of group indices $`n1`$ and $`n2`$:
$`F_{1\mathrm{}N}(a)T_{n1n2}(a)`$ $`=`$ $`{\displaystyle \underset{\sigma S_N}{}}{\displaystyle \underset{k_{n1}=0}{\overset{N}{}}}{\displaystyle \underset{k_{n2}=0}{\overset{Nk_{n1}}{}}}{\displaystyle \underset{\alpha _{\sigma (1)}\mathrm{}\alpha _{\sigma (N)}}{\overset{^{\prime \prime }}{}}}{\displaystyle \underset{j_{n1}=Nk_{n1}+1}{\overset{N}{}}}P_{\sigma (j_{n1})}^{n1}{\displaystyle \underset{j_{n2}=Nk_{n1}k_{n2}+1}{\overset{Nk_{n1}}{}}}P_{\sigma (j_{n2})}^{n2}`$ (29)
$`\times {\displaystyle \underset{j=1}{\overset{Nk_{n1}k_{n2}}{}}}P_{\sigma (j)}^{\alpha _{\sigma (j)}}P_0^{n1}T_{0,\sigma (1)\mathrm{}\sigma (N)}(a)P_0^{n2}R_{1\mathrm{}N}^\sigma (a+\sigma _h).`$
Evaluating the matrix product in $`T_{0,\sigma (1)\mathrm{}\sigma (N)}(a)`$ leads to
$`{\displaystyle \underset{j_{n1}=Nk_{n1}+1}{\overset{N}{}}}P_{\sigma (j_{n1})}^{n1}{\displaystyle \underset{j_{n2}=Nk_{n1}k_{n2}+1}{\overset{Nk_{n1}}{}}}P_{\sigma (j_{n2})}^{n2}{\displaystyle \underset{j=1}{\overset{Nk_{n1}k_{n2}}{}}}P_{\sigma (j)}^{\alpha _{\sigma (j)}}P_0^{n1}T_{0,\sigma (1)\mathrm{}\sigma (N)}(a)P_0^{n2}`$ (30)
$`=`$ $`{\displaystyle \underset{i=Nk_{n1}k_{n2}+1}{\overset{Nk_{n1}}{}}}\left(R_{0\stackrel{~}{Nk_{n1}}}(\widehat{a}_{Nk_{n1}})\right)_{n1n2}^{n1n2}\mathrm{}\left(R_{0\stackrel{~}{i+1}}(\widehat{a}_{i+1})\right)_{n1n2}^{n1n2}`$
$`\times \left(R_{0\stackrel{~}{i}}(\widehat{a}_i)\right)_{n2n1}^{n1n2}\left(R_{0\stackrel{~}{i1}}(\widehat{a}_{i1})\right)_{n2n2}^{n2n2}\mathrm{}\left(R_{0\stackrel{~}{Nk_{n1}k_{n2}+1}}(\widehat{a}_{Nk_{n1}k_{n2}+1})\right)_{n2n2}^{n2n2}`$
$`\times {\displaystyle \underset{k=1}{\overset{Nk_{n1}k_{n2}}{}}}\left(R_{0\stackrel{~}{k}}(\widehat{a}_k)\right)_{n1\alpha _{\stackrel{~}{k}}}^{n2\alpha _{\stackrel{~}{k}}}E_{n2,n1}^{\stackrel{~}{i}}{\displaystyle \underset{ji}{}}P_{\stackrel{~}{j}}^{\alpha _{\stackrel{~}{j}}}P_{\stackrel{~}{i}}^{n1}E_{n1,n2}^0`$
$`=`$ $`{\displaystyle \underset{i=Nk_{n1}k_{n2}+1}{\overset{Nk_{n1}}{}}}{\displaystyle \underset{l=i+1}{\overset{Nk_{n1}}{}}}b^{n1,n2}(\lambda z_{\stackrel{~}{l}}|\widehat{a}_l)c^{n2,n1}(\lambda z_{\stackrel{~}{i}}|\widehat{a}_i)`$
$`\times {\displaystyle \underset{k=1}{\overset{Nk_{n1}k_{n2}}{}}}b^{n2,\alpha _{\stackrel{~}{k}}}(\lambda z_{\stackrel{~}{k}}|\widehat{a}_k)E_{n2,n1}^{\stackrel{~}{i}}{\displaystyle \underset{ji}{}}P_{\stackrel{~}{j}}^{\alpha _{\stackrel{~}{j}}}P_{\stackrel{~}{i}}^{n1}E_{n1,n2}^0`$
with the abbreviations
$`\widehat{a}_i`$ $`=`$ $`a+{\displaystyle \underset{m=1}{\overset{i1}{}}}\widehat{\sigma }_{\stackrel{~}{m}}`$
$`\stackrel{~}{i}`$ $`=`$ $`\sigma (i)`$
$`\left(E_{a,b}^i\right)_{\beta _i}^{\alpha _i}`$ $`=`$ $`\delta _{\alpha _i,a}\delta _{\beta _i,b}.`$
One notes that in the calculation the index $`\alpha _{\stackrel{~}{i}}`$ has changed from $`n2`$ to $`n1`$. As the distribution of $`\alpha `$’s is therefore no longer consistent with the conditions (19) in the sum $`^{}`$ one has to correct this fact by commuting the site $`\stackrel{~}{i}`$ through all higher sites $`\stackrel{~}{j}`$ with $`\alpha _{\stackrel{~}{j}}=n2`$. So, taking into account (20), one has to insert an additional factor
$`R_{\stackrel{~}{i}\stackrel{~}{Nk_{n1}}}(\widehat{a}_{\stackrel{~}{Nk_{n1}}}+\sigma _h)\mathrm{}R_{\stackrel{~}{i}\stackrel{~}{i+1}}(\widehat{a}_{\stackrel{~}{i+1}}+\sigma _h)`$ (31)
between the projectors and $`R_{1\mathrm{}N}^\sigma (a+\sigma _h)`$ in (29). Because of Eq. (23) no further corrections are neccessary. For the following calculation two equations are needed
$`P_i^{n1}P_j^{n2}\mathrm{\hspace{0.17em}1}\mathrm{I}_{ij}=P_i^{n1}P_j^{n2}\left\{b^{n1,n2}(z_iz_j|a)^1R_{ij}(a){\displaystyle \frac{c^{n2,n1}(z_iz_j|a)}{b^{n1,n2}(z_iz_j|a)}}P_{ij}\right\}`$ (32)
and
$`E_{n1,n}^iP_j^{n1}P_{ij}=E_{n1,n}^jP_i^{n1}`$ (33)
with $`\left(P_{ij}\right)_{\beta _i\beta _j}^{\alpha _i\alpha _j}=\delta _{\alpha _i,\beta _j}\delta _{\alpha _j,\beta _i}`$ the permutation operator in space i and j.
Let us now concentrate on the term with $`i=Nk_{n1}k_{n2}+1`$ in (30) and use (32) to create the needed factor (31). Because of (33) the second term in (32) gives rise to an $`E_{n2,n1}^{\stackrel{~}{j}}`$ with $`\stackrel{~}{j}\stackrel{~}{i}`$. So the only possibility to get an $`E_{n2,n1}^{\stackrel{~}{i}}`$ is to use the first term in (32). Corrections in the other terms with $`j>Nk_{n1}k_{n2}+1`$ cannot lead to an expression with $`E_{n2,n1}^{\stackrel{~}{i}}`$ as $`\stackrel{~}{j}`$ has not to be commuted with the site $`\stackrel{~}{i}`$. So the only term that contains $`E_{n2,n1}^{\stackrel{~}{i}}`$ after the corrections for that special $`R_{1\mathrm{}N}^\sigma `$ in (30) is
$`{\displaystyle \underset{l=i+1}{\overset{Nk_{n1}}{}}}{\displaystyle \frac{b^{n1,n2}(\lambda z_{\stackrel{~}{l}}|\widehat{a}_l)}{b^{n1,n2}(z_{\stackrel{~}{i}}z_{\stackrel{~}{l}}|\widehat{a}_l)}}c^{n2,n1}(\lambda z_{\stackrel{~}{i}}|\widehat{a}_i){\displaystyle \underset{k=1}{\overset{Nk_{n1}k_{n2}}{}}}b^{n2,\alpha _{\stackrel{~}{k}}}(\lambda z_{\stackrel{~}{k}}|\widehat{a}_k)E_{n2,n1}^{\stackrel{~}{i}}`$ (34)
$`=`$ $`{\displaystyle \underset{l=i+1}{\overset{Nk_{n1}}{}}}{\displaystyle \frac{b(\lambda z_{\stackrel{~}{l}})}{b(z_{\stackrel{~}{i}}z_{\stackrel{~}{l}})}}c^{n2,n1}(\lambda z_{\stackrel{~}{i}}|\stackrel{~}{a}_{n2}){\displaystyle \underset{k=1}{\overset{Nk_{n1}k_{n2}}{}}}b^{n2,\alpha _{\stackrel{~}{k}}}(\lambda z_{\stackrel{~}{k}}|\widehat{a}_k)E_{n2,n1}^{\stackrel{~}{i}}.`$
Because of the symmetry of $`\stackrel{~}{T}_{0,1\mathrm{}N}(\lambda |a)`$ all other terms have to be of the same form as (34). After taking into account the projectors the resulting expression is
$`\stackrel{~}{T}_{n1n2}(\lambda |a)`$ $`=`$ $`Y_{n2}(a){\displaystyle \underset{i=1}{\overset{N}{}}}c^{n2,n1}(\lambda z_i|\stackrel{~}{a}_{n2})E_{n2,n1}^i`$ (35)
$`_{ji}^N\text{diag}\{b(\lambda z_j),\mathrm{},b(\lambda z_j),b(\lambda z_j)b^1(z_iz_j),1\}.`$
For the calculation of $`\stackrel{~}{T}_{n1n3}(\lambda |a)`$ one has to distinguish the cases $`n1`$, $`n2`$ and $`n3`$ in the sum $`^{}`$. The only difference compared to $`\stackrel{~}{T}_{n1n2}(\lambda |a)`$ is a term containing a product $`E_{n3,n2}^iE_{n2,n1}^j`$ now showing up in the matrix product in $`T_{0,\sigma (1)\mathrm{}\sigma N}(a)`$. One again has to correct the distribution of $`\alpha ^{}s`$ with the analog of the equations (32) and (33)
$`P_i^{\alpha _i}P_j^{\alpha _j}\mathrm{\hspace{0.17em}1}\mathrm{I}_{ij}=P_i^{\alpha _i}P_j^{\alpha _j}\left\{b^{\alpha _i,\alpha _j}(z_iz_j|a)^1R_{ij}(a){\displaystyle \frac{c^{\alpha _j,\alpha _i}(z_iz_j|a)}{b^{\alpha _i,\alpha _j}(z_iz_j|a)}}P_{ij}\right\}`$ (36)
$`E_{a,b}^iP_j^aP_{ij}=E_{a,b}^jP_i^a`$ (37)
and also with a new relation which has to be taken into account when dealing with the term containing the product $`E_{n3,n2}^iE_{n2,n1}^j`$ :
$`E_{n3,n1}^iP_j^{n2}P_{ij}=E_{n3,n2}^iE_{n2,n1}^j.`$ (38)
This reasoning leads to ($`b_{ik}=b(z_iz_k)`$):
$`\stackrel{~}{T}_{n1n3}(\lambda |a)`$ $`=`$ $`Y_{n3}(a){\displaystyle \underset{i=1}{\overset{N}{}}}c^{n3,n1}(\lambda z_i|\stackrel{~}{a}_{n3})E_{n3,n1}^i`$ (39)
$`_{ki}^N\text{diag}\{b(\lambda z_k),\mathrm{},b(\lambda z_k),b(\lambda z_k)b_{ik}^1,b(\lambda z_k)b_{ik}^1,1\}`$
$`+Y_{n3}(a){\displaystyle \underset{ij}{\overset{N}{}}}f_2(n3,n2;n1)E_{n3,n2}^iE_{n2,n1}^j`$
$`_{ki,j}^N\text{diag}\{b(\lambda z_k),\mathrm{},b(\lambda z_k),b(\lambda z_k)b_{ik}^1,b(\lambda z_k)b_{jk}^1,1\}`$
with
$`f_2(n3,n2;n1)`$ $`=`$ $`c^{n2,n1}(\lambda z_j|\stackrel{~}{a}_{n2})c^{n3,n2}(\lambda z_i|\stackrel{~}{a}_{n3})`$
$``$ $`{\displaystyle \frac{c^{n2,n1}(z_iz_j|\stackrel{~}{a}_{n2})}{b^{n1,n2}(z_iz_j|\stackrel{~}{a}_{n2})}}b^{n1,n2}(\lambda z_j|\stackrel{~}{a}_{n2})c^{n3,n1}(\lambda z_i|\stackrel{~}{a}_{n3})`$
where the second term in (LABEL:f2) has its origin in the second term in (36) and in (38).
Proceeding in an analogous manner one obtains in the general case
$`\stackrel{~}{T}_{n1\alpha }(\lambda |a)`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{\alpha }{}}}{\displaystyle \underset{i_1\mathrm{}i_k}{}}{\displaystyle \underset{n1\alpha <n_1<\mathrm{}n_k<n1}{}}`$
$`E_{n_1,n_2}^{(i_1)}E_{n_2,n_3}^{(i_2)}\mathrm{}E_{n_k,n1}^{(i_k)}f_k(n_1,n_2,\mathrm{},n_k;n1)Y_{n1\alpha }(a)`$
$`_{j\{i_k\}}\text{diag}\{b(\lambda z_j),\mathrm{},\underset{n_2n_1}{\underset{}{b(\lambda z_j)b_{i_1j}^1,\mathrm{},b(\lambda z_j)b_{i_1j}^1}},\mathrm{},\underset{n1n_k}{\underset{}{b(\lambda z_j)b_{i_kj}^1,\mathrm{},b(\lambda z_j)b_{i_kj}^1}},1\}`$
with $`f_k(n_1,n_2,\mathrm{},n_k;n_0)`$ being defined recursively by <sup>2</sup><sup>2</sup>2In the rational $`sl(n)`$ case this recursion relation leads directly to the result in Eq. (42) in .
$`f_k(n_1,n_2,\mathrm{},n_k;n_0)={\displaystyle \frac{c_{i_{k1},i_k}^{n_k,n_0}}{b_{i_{k1},i_k}^{n_0,n_k}}}b_{0,i_k}^{n_0,n_k}f_{k1}(n_1,n_2,\mathrm{},n_{k1};n_0)+c_{0,i_k}^{n_k,n_0}f_{k1}(n_1,n_2,\mathrm{},n_{k1};n_k)`$ (42)
where
$`f_1(n_1;n_0)=c_{0,i_{n_1}}^{n_1,n_0}`$
Remark: As the above procedure to obtain the operators of the monodromy matrix in the F-basis did not rely on the invariance of the monodromy matrix under combined $`sl(n)`$ rotations in the auxilliary and quantum space, the $`sl(n)`$ XXZ-model can be treated in the same way, and we obtain expressions for the $`\stackrel{~}{T}_{n1\alpha }`$ similar to that in (LABEL:Ta), with the diffenence that the $`a`$ dependence vanishes, and the parametrization of the elements of the R-matrix is trigonometric instead of elliptic ($`c(\lambda )=\frac{sinh(w)}{sinh(\lambda +w)},b(\lambda )=\frac{sinh(\lambda )}{sinh(\lambda +w)}`$).
## 4 Bethe wavevectors
Our presentation of the hierarchical Bethe ansatz will be rather sketchy. For details we refer the reader to .
The operators $`T_{n1\alpha }`$ $`(\alpha <n1)`$ act as quasiparticle creation operators and satisfy the Faddeev–Zamolodchikov algebra
$`T_{n1\alpha }(\lambda |a)T_{n1\beta }(\mu |a+\widehat{\alpha })=T_{n1\gamma }(\mu |a)T_{n1\delta }(\lambda |a+\widehat{\gamma })R(\lambda \mu |a)_{\alpha \beta }^{\delta \gamma }`$ (43)
or in components of the R-matrix
$`[T_{n1\alpha }(\lambda |a),T_{n1\alpha }(\mu |a+\widehat{\alpha })]`$ $`=`$ $`0`$ (44)
$`T_{n1\alpha }(\lambda |a)T_{n1\beta }(\mu |a+\widehat{\alpha })`$ $`=`$ $`{\displaystyle \frac{1}{b^{\beta \alpha }(\mu \lambda |a)}}T_{n1\beta }(\mu |a)T_{n1\alpha }(\lambda |a+\widehat{\beta })`$ (45)
$``$ $`{\displaystyle \frac{c^{\beta \alpha }(\mu \lambda |a)}{b^{\beta \alpha }(\mu \lambda |a)}}T_{n1\beta }(\lambda |a)T_{n1\alpha }(\mu |a+\widehat{\beta }).`$
Inspired by , an ansatz for a Bethe vector $`\mathrm{\Psi }_n`$ is given in terms of a linear superposition of products of operators $`T_{n1\alpha }`$ acting on a reference state $`\mathrm{\Omega }_N^{(n)}`$:
$`\mathrm{\Psi }_n(N;\lambda _1,\mathrm{},\lambda _p|a)={\displaystyle \underset{\alpha _1,\mathrm{},\alpha _p}{}}\mathrm{\Phi }_{\alpha _1,\mathrm{},\alpha _p}T_{n1\alpha _1}(\lambda _1|a)T_{n1\alpha _2}(\lambda _2|a+\widehat{\alpha _1})\mathrm{}T_{n1\alpha _p}(\lambda _p|a+{\displaystyle \underset{i=1}{\overset{n1}{}}}\widehat{\alpha _i})\mathrm{\Omega }_N^{(n)}`$
(46)
where the reference state $`\mathrm{\Omega }_N^{(n)}`$ is constituted as a $`N`$-fold tensor product of lowest weight states $`v_n^{(i)}={}_{}{}^{t}(0,\mathrm{},0,1)`$ in $`_n^{(i)}`$
$`\mathrm{\Omega }_N=_{i=1}^Nv_n^{(i)}`$
and the $`\mathrm{\Phi }_{\alpha _1,\mathrm{},\alpha _p}`$ denote some c-number coefficients, which themselves have to be chosen s.t. they are components of a $`sl(n1)`$ wavevector, leading to a nested structure finally giving a $`sl(2)`$ eigenvalue problem.
It is important to note that the reference state is invariant under the $`F`$-transformation
$`F\mathrm{\Omega }_N^{(n)}=\mathrm{\Omega }_N^{(n)}`$
due to the special form of the R-matrix.
We will not impose the Bethe ansatz equations for the spectral parameters $`\{\lambda _i\}`$ which turns the vector (46) into an eigenvector of the transfer matrix, that is we consider the Bethe wavevector being “off-shell” .
In what follows we want to determine the functional form of such vectors, using the explicit form of the operators relevant for the Bethe wavevectors.
In order to clarify the arguments employed in the course of the computation we will present the case of $`sl(2)`$ and $`sl(3)`$ in quite a detail. The generalization to the general case of $`sl(n)`$, $`n>3`$ will then be rather straightforward.
For the $`sl(2)`$ case we have from (35)
$`\stackrel{~}{T}_{10}(\lambda |a)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}c^{0,1}(\lambda z_i|a)\sigma _+^{(i)}_{ji}\left(\begin{array}{cc}b(\lambda z_j)b_{ij}^1& 0\\ 0& 1\end{array}\right)_{[j]}`$ (47)
$`\stackrel{~}{\mathrm{\Psi }}_2(N;\lambda _1,\mathrm{},\lambda _p|a)`$ $`=`$ $`\stackrel{~}{T}_{10}(\lambda _1|a)\stackrel{~}{T}_{10}(\lambda _2|a+\widehat{0})\mathrm{}\stackrel{~}{T}_{10}(\lambda _p|a+(p1)\widehat{0})\mathrm{\Omega }_N^{(2)}`$ (48)
$`=`$ $`{\displaystyle \underset{i_1\mathrm{}i_p}{}}B_p^{(2)}(\lambda _1,\mathrm{},\lambda _p|z_{i_1},\mathrm{},z_{i_p}|a)\sigma _+^{(i_1)}\mathrm{}\sigma _+^{(i_p)}\mathrm{\Omega }_N^{(2)}.`$
The c-number coefficient $`B^{(2)}(\{\lambda _i\}|\{z_i\}|a)`$ is, due to the “diagonally dressed” spin raising operators $`\sigma _+^i`$ in (47), of the form
$`B_p^{(2)}(\lambda _1,\mathrm{},\lambda _p|z_1,\mathrm{},z_p)`$ $`=`$ $`{\displaystyle \underset{\sigma S_p}{}}{\displaystyle \underset{m=1}{\overset{p}{}}}c^{0,1}(\lambda _mz_{\sigma (m)}|a+(m1)\widehat{0}){\displaystyle \underset{l=m+1}{\overset{p}{}}}{\displaystyle \frac{b(\lambda _mz_{\sigma (l)})}{b(z_{\sigma (m)}z_{\sigma (l)})}}.`$
We now turn to the $`sl(3)`$ case. The strategy will rely on the symmetry of the wavevector under the exchange of arbitrary spectral parameters (the verification of this fact follows the same lines as the one in using (16) and (43)) which enables us to concentrate on a particularily simple term in the sum (46), and the repeated use of the Faddeev–Zamolodchikov algebra.
These ideas lead us to propose the following form for the Bethe wavevector
$`\stackrel{~}{\mathrm{\Psi }}_3(N,\lambda _1,\mathrm{},\lambda _{p_0};\lambda _{p_0+1},\mathrm{},\lambda _{p_0+p_1}|a)`$ (50)
$`=`$ $`{\displaystyle \underset{\sigma S_{p_0}}{}}B_{p_1}^{(2)}(\lambda _{p_0+1},\mathrm{},\lambda _{p_0+p_1}|\lambda _{\sigma (1)},\mathrm{},\lambda _{\sigma (p_1)}|a)`$
$`\times `$ $`{\displaystyle \underset{k=1}{\overset{p_1}{}}}{\displaystyle \underset{l=p_1+1}{\overset{p_0}{}}}b^{1,0}(\lambda _{\sigma (k)}\lambda _{\sigma (l)}|a+(k1)\widehat{0}+(lp_11)\widehat{1})^1`$
$`\times `$ $`\stackrel{~}{T}_{21}(\lambda _{\sigma (p_1+1)}|a)\mathrm{}\stackrel{~}{T}_{21}(\lambda _{\sigma (p_0)}|a+(p_0p_11)\widehat{1})`$
$`\times `$ $`\stackrel{~}{T}_{20}(\lambda _{\sigma (1)}|a+(p_0p_1)\widehat{1})\mathrm{}\stackrel{~}{T}_{20}(\lambda _{\sigma (p_1)}|a+(p_0p_1)\widehat{1}+(p_11)\widehat{0})\mathrm{\Omega }_N^{(3)}.`$
Consider a special term in the sum (46) of the form (which is motivated by the fact that the associated coefficient $`\mathrm{\Phi }`$ is especially simple to compute, see below)
$`\stackrel{~}{T}_{20}(\lambda _1|a)\stackrel{~}{T}_{20}(\lambda _1|a+\widehat{0})\mathrm{}\stackrel{~}{T}_{20}(\lambda _{p_1}|a+(p_11)\widehat{0})`$
$`\times `$ $`\stackrel{~}{T}_{21}(\lambda _{p_1+1}|a+p_1\widehat{0})\stackrel{~}{T}_{21}(\lambda _{p_1+2}|a+p_1\widehat{0}+\widehat{1})\mathrm{}\stackrel{~}{T}_{21}(\lambda _{p_0}|a+p_1\widehat{0}+(p_0p_11)\widehat{1})\mathrm{\Omega }_N^{(3)}.`$
Commuting all $`\stackrel{~}{T}_{20}(\lambda _1|a)`$ to the right using the first term in (45) yields an additional factor
$`{\displaystyle \underset{x=1}{\overset{p_1}{}}}{\displaystyle \underset{y=p_1+1}{\overset{p_0}{}}}\left\{b^{1,0}(\lambda _y\lambda _x|a+(x1)\widehat{0}+(yp_11)\widehat{1})\right\}^1.`$
It has to be noted that the associated c-number coefficient $`\mathrm{\Phi }_{0\mathrm{}01\mathrm{}1}^{(2)}`$ in (LABEL:01) is not evaluated in the $`sl(3)`$ F basis. It can however be expressed in the form (LABEL:B2) as it is invariant under the action of the $`sl(2)`$ F-matrix. This is due to the fact that it constitutes a component of the $`sl(2)`$ vector whose labels (a non-decreasing series of $`\alpha _i`$ with respect to the original ordering of sites $`i`$) correspond via (19) to the identity permutation in the definition of the F-matrix (18).
Invoking the exchange symmetry we arrive thus at the formula (50).
The creation operators with respect to the lowest weight state are of the form (cf. (LABEL:Ta))
$`\stackrel{~}{T}_{21}(\lambda |a)`$ $`=`$ $`Y_1(a){\displaystyle \underset{i=1}{\overset{N}{}}}c^{1,2}(\lambda z_i|\stackrel{~}{a_1})E_{1,2}^{(i)}_{ji}\text{diag}\{b(\lambda z_j),b(\lambda z_j)b_{ij}^1,1\}_{[j]}`$ (52)
$`\stackrel{~}{T}_{20}(\lambda |a)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}c^{0,2}(\lambda z_i|\stackrel{~}{a_0})E_{0,2}^{(i)}_{ji}\text{diag}\{b(\lambda z_j)b_{ij}^1,b(\lambda z_j)b_{ij}^1,1\}_{[j]}`$ (53)
$`+`$ $`\text{terms involving}E_{1,2}^{(i)}E_{0,1}^{(j)}.`$
The second term in (53) annihilates the vacuum $`\mathrm{\Omega }_N^{(3)}`$, which is why we did not cite the explicit form of the prefactors accompanying these roots. The form of the creation operators permits us to further simplify the wavevector (50). Taking into account the action of the roots on the respective dressing as well as the fact that both groups of creation operators generate a factor similar to the $`sl(2)`$ problem, we obtain
$`\stackrel{~}{\mathrm{\Psi }}_3(N,\lambda _1,\mathrm{},\lambda _{p_0};\lambda _{p_0+1},\mathrm{},\lambda _{p_0+p_1}|a)`$
$`=`$ $`{\displaystyle \underset{\sigma S_{p_0}}{}}{\displaystyle \underset{i_1\mathrm{}i_{p_0}}{}}{\displaystyle \underset{k=1}{\overset{p_1}{}}}{\displaystyle \underset{l=p_1+1}{\overset{p_0}{}}}b^{1,0}(\lambda _{\sigma (l)}\lambda _{\sigma (k)}|a+(k1)\widehat{0}+(lp_11)\widehat{1})^1b(\lambda _{\sigma (l)}z_{i_k})`$
$`\times `$ $`B_{p_0p_1}^{1,2}(\lambda _{\sigma (p_1+1)},\mathrm{},\lambda _{\sigma (p_0)}|z_{i_{p_1+1}},\mathrm{},z_{i_{p_0}}|a)B_{p_1}^{0,1}(\lambda _{p_0+1},\mathrm{},\lambda _{p_0+p_1}|\lambda _{\sigma (1)},\mathrm{},\lambda _{\sigma (p_1)}|a)`$
$`\times `$ $`B_{p_1}^{0,2}(\lambda _{\sigma (1)},\mathrm{},\lambda _{\sigma (p_1)}|z_{i_1},\mathrm{},z_{i_{p_1}}|a+(p_0p_1)\widehat{1})E_{1,2}^{(i_{p_1+1})}\mathrm{}E_{1,2}^{(i_{p_0})}E_{0,2}^{(i_1)}\mathrm{}E_{0,2}^{(i_{p_1})}\mathrm{\Omega }_N^{(3)}`$
where we defined as a generalization of (LABEL:B2)
$`B_p^{\alpha ,\beta }(\lambda _1,\mathrm{},\lambda _p|z_1,\mathrm{},z_p|a)`$
$`=`$ $`{\displaystyle \underset{\sigma S_p}{}}{\displaystyle \underset{m=1}{\overset{p}{}}}\left\{c^{\alpha ,\beta }(\lambda _mz_{\sigma (m)}|\stackrel{~}{a_\alpha }+(m1)\widehat{\alpha })Y_\alpha (a+(m1)\widehat{\alpha })\right\}{\displaystyle \underset{l=m+1}{\overset{p}{}}}{\displaystyle \frac{b(\lambda _mz_{\sigma (l)})}{b(z_{\sigma (m)}z_{\sigma (l)})}}.`$
We are now in the position to proceed to the general $`sl(n)`$ case.
We start with a special term of (46), where all operators $`\stackrel{~}{T}_{n1\alpha }`$ are to the left of operators $`\stackrel{~}{T}_{n1\beta }`$ if $`\alpha <\beta `$. The associated coefficient $`\mathrm{\Phi }`$ again remains invariant under the action of F, which enables us to express it in a manner similar to (LABEL:B2). We then use the Faddeev-Zamolodchikov algebra to reverse the order of all operators. Once again only the term containing one $`E_{a,b}^i`$ in the expression (LABEL:Ta) of the respective operators $`\stackrel{~}{T}_{n1\alpha }`$ contributes in this special ordering.
The wavevector can be expressed in a form similar to (LABEL:Psi3b) ($`\overline{\alpha }n2\alpha `$)
$`\stackrel{~}{\mathrm{\Psi }}_n(N,p_0,p_1,\mathrm{},p_{n2})={\displaystyle \underset{i_1\mathrm{}i_{p_0}}{}}B_{p_0,p_1,\mathrm{},p_{n2}}^{(n)}(\lambda _1,\mathrm{},\lambda _{p_0+\mathrm{}p_{n2}}|z_{i_1},\mathrm{},z_{i_{p_0}}|a){\displaystyle \underset{\alpha =0}{\overset{n2}{}}}{\displaystyle \underset{j=p_{\overline{\alpha }+1}+1}{\overset{p_{\overline{\alpha }}}{}}}E_{\alpha n1}^{(i_j)}\mathrm{\Omega }_N^{(n)}`$
(56)
with the recursion relation for $`B^{(n)}`$
$`B_{p_0p_1\mathrm{}p_{n2}}^{(n)}(\lambda _1,\mathrm{},\lambda _{p_0+p_1+\mathrm{}p_{n2}}|z_1,\mathrm{},z_{p_0}|a)`$ (57)
$`=`$ $`{\displaystyle \underset{\sigma S_{p_0}}{}}{\displaystyle \underset{\alpha =0}{\overset{n3}{}}}{\displaystyle \underset{\beta =\alpha +1}{\overset{n2}{}}}{\displaystyle \underset{k_\alpha =p_{\overline{\alpha }+1}+1}{\overset{p_{\overline{\alpha }}}{}}}{\displaystyle \underset{l_\beta =p_{\overline{\beta }+1}+1}{\overset{p_{\overline{\beta }}}{}}}b(\lambda _{\sigma (l_\beta )}z_{k_\alpha })`$
$`\times `$ $`b^{\beta ,\alpha }\left(\lambda _{\sigma (l_\beta )}\lambda _{\sigma (k_\alpha )}|a+{\displaystyle \underset{m=0}{\overset{\alpha 1}{}}}(p_{\overline{m}}p_{\overline{m}+1})\widehat{m}+(k_\alpha p_{\overline{\alpha }+1}1)\widehat{\alpha }+(l_\beta p_{\overline{\beta }+1}1)\widehat{\beta }\right)^1`$
$`\times `$ $`{\displaystyle \underset{\gamma =0}{\overset{n2}{}}}B_{p_{\overline{\gamma }}p_{\overline{\gamma }+1}}^{\gamma ,n1}(\lambda _{\sigma (p_{\overline{\gamma }+1}+1)}\mathrm{}\lambda _{\sigma (p_{\overline{\gamma }})}|z_{p_{\overline{\gamma }+1}+1}\mathrm{}z_{p_{\overline{\gamma }}}|a+{\displaystyle \underset{m=\gamma +1}{\overset{n2}{}}}(p_{\overline{m}}p_{\overline{m}+1})\widehat{m})`$
$`\times `$ $`B_{p_1\mathrm{}p_{n2}}^{(n1)}(\lambda _{p_0+1}\mathrm{}\lambda _{p_0+p_1+\mathrm{}+p_{n2}}|\lambda _{\sigma (1)}\mathrm{}\lambda _{\sigma (p_1)}|a)`$
which can be solved explicitely to yield
$`B_{p_0p_1\mathrm{}p_{n2}}^{(n)}(\lambda _1,\mathrm{},\lambda _{p_0+p_1+\mathrm{}+p_{n2}}|z_1,\mathrm{},z_{p_0}|a)`$
$`={\displaystyle \underset{\sigma _0S_{p_0}}{}}{\displaystyle \underset{\sigma _1S_{p_1}}{}}\mathrm{}{\displaystyle \underset{\sigma _{n3}S_{p_{n3}}}{}}{\displaystyle \underset{i=0}{\overset{n3}{}}}{\displaystyle \underset{\alpha _i=0}{\overset{n3i}{}}}{\displaystyle \underset{\beta _i=\alpha _i+1}{\overset{n2i}{}}}{\displaystyle \underset{k_{\alpha _i}=p_{\overline{\alpha _i}+1}+1}{\overset{p_{\overline{\alpha _i}}}{}}}{\displaystyle \underset{l_{\beta _i}=p_{\overline{\beta _i}+1}+1}{\overset{p_{\overline{\beta _i}}}{}}}b(\lambda _{q_{i1}+\sigma _i(l_{\beta _i})}\lambda _{q_{i2}+\sigma _{i1}(k_{\alpha _i})})`$
$`\times `$ $`b^{\beta _i,\alpha _i}\left(\lambda _{q_{i1}+\sigma _i(l_{\beta _i})}\lambda _{q_{i1}+\sigma _i(k_{\alpha _i})}|a+{\displaystyle \underset{m=0}{\overset{\alpha _i1}{}}}(p_{\overline{m}}p_{\overline{m}+1})\widehat{m}+(k_{\alpha _i}p_{\overline{\alpha _i}+1}1)\widehat{\alpha _i}+(l_{\beta _i}p_{\overline{\beta _i}+1}1)\widehat{\beta _i}\right)^1`$
$`\times `$ $`{\displaystyle \underset{\gamma _i=0}{\overset{n2i}{}}}B_{p_{\overline{\gamma _i}}p_{\overline{\gamma _i}+1}}^{\gamma _i,n1}(\lambda _{q_{i1}+\sigma _i(p_{\overline{\gamma _i}+1}+1)}\mathrm{}\lambda _{q_{i1}+\sigma _i(p_{\overline{\gamma _i}})}|\lambda _{q_{i2}+\sigma _{i1}(p_{\overline{\gamma _i}+1}+1)}\mathrm{}\lambda _{q_{i2}+\sigma _{i1}(p_{\overline{\gamma _i}})}|a+{\displaystyle \underset{m=\gamma _i+1}{\overset{n2i}{}}}(p_{\overline{m}}p_{\overline{m}+1})\widehat{m})`$
$`\times `$ $`B_{p_{n2}}^{0,1}\left(\lambda _{q_{n3}+1}\mathrm{}\lambda _{q_{n3}+p_{n2}}|\lambda _{q_{n4}+\sigma _{n3}(1)}\mathrm{}\lambda _{q_{n4}+\sigma _{n3}(p_{n2})}|a\right)`$
where we defined
$$q_i=\underset{j=0}{\overset{i}{}}p_j;q_1=0$$
and
$$\lambda _{\sigma _1(k)}=z_k.$$
By expressing the $`sl(n)`$ wavevector (56) with the help of $`sl(2)`$ building blocks (LABEL:Bn) we have achieved a resolution of the hierarchy.
Remark: Once again the arguments in this section apply to the $`sl(n)`$ $`XXZ`$-model as well, leading to a wavevector for this model which displays the same features as (56) in connection with (LABEL:Bn), without the $`a`$ dependence and with a trigonometric parametrization.
## 5 Conclusion
We accomplished the construction of a factorizing F-matrix for the Belavin model enabling one to construct completely symmetric creation operators which moreover are devoid of non-local effects from polarization clouds. These operators were used to resolve the intricacies of the nested structure of Bethe wavevectors.
In contrast to the $`sl(n)`$ $`XXX`$-model the above method does not rely on an $`sl(n)`$ invariance of the monodromy matrix, which renders the extraction of the elements in the grid of the monodromy matrix much more cumbersome. The only ingredients needed in our computation are the form of the R-matrix (8), i.e. its structure, unitarity and the fact that it constitutes a representation of the permutation group, and the property (23). Thus our findings directly yield the corresponding expression for the generalized $`sl(n)`$ $`XXZ`$-model.
In view of the similarities between the generalized $`sl(n)`$ rational, trigonometric and elliptic model it is tempting to ask whether such a distinguished basis exists for all integrable models in two dimensions. In it was shown that in every integrable two-dimensional quantum field theory there exist semi-local polarization-free generators which are localized in wedge-shaped regions of Minkowski space. It is conceivable that there is a relation between these operators and polarization-free operators in lattice spin models.
We hope that these results might prove useful for the construction of formfactors starting from the microscopical level.
Acknowledgement: We thank H. Boos, R. Flume and R.H. Poghossian for helpful discussions.
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# 1 Introduction
## 1 Introduction
The GHZ theorem provides a powerful test of quantum non-locality, which can be confirmed or refuted by the outcome of just one single experiment . Formulated for three spin 1/2 particles , the argument is based on the anti-commutative nature of the 2x2 spin operators $`\sigma _x`$, $`\sigma _y`$. The values of the three mutually commuting observables
$$\sigma _x^a\sigma _y^b\sigma _y^c\sigma _x^a\sigma _y^b\sigma _y^c,\sigma _y^a\sigma _x^b\sigma _y^c,\sigma _y^a\sigma _y^b\sigma _x^c,$$
(1)
and their product, $`\sigma _x^a\sigma _x^b\sigma _x^c`$, cannot be obtained, consistently, by making local assignments to each of the individual spin operators, $`m_x^I,m_y^I=\pm 1`$, $`I=a,b,c`$. This is not a contradiction of Quantum Mechanics: the state $`|\psi =\frac{1}{\sqrt{2}}(||)`$, for instance, is one of the common eigenstates of the four operators, with eigenvalues $`\lambda _1=\lambda _2=\lambda _3=1`$, $`\lambda _4=1`$, respectively. $`|\psi `$ is a highly correlated (entangled) state of the three parties which has no defined value for $`\sigma _x^I,\sigma _y^I`$.
In this note we address the question of how to generalize the argument to particles of higher spin and find that there are no non-trivial extensions other than direct sums of operators that can be brought into the form $`\sigma _x`$,$`\sigma _y`$ by means of local unitarity transformations. (For odd dimensional Hilbert spaces the direct sum is completed by a one-dimensional submatrix, i.e., a c-number in the diagonal). We give a proof for the cases of spin 1 and 3/2. Similar problems have been addressed in .
Let us look for observables $`A`$, $`B`$ such that $`AB=\omega BA`$ (their hermiticity implies that $`\omega `$ is at most a phase): this is a necessary condition for the commutator relations $`[A_1^aA_2^bA_3^c,B_1^aB_2^bB_3^c]=\mathrm{etc}\mathrm{}=0`$ to hold. As we shall see, all interesting cases correspond to $`\omega =1`$. Without loss of generality, $`A`$ can always be taken diagonal, $`A=\mathrm{diag}(\lambda _1,\lambda _2)`$, for the simplest case s=1/2. The above condition reads
$$AB\omega BA=\left(\begin{array}{cc}(1\omega )\lambda _1b_{11}& (\lambda _1\omega \lambda _2)b_{12}\\ (\lambda _2\omega \lambda _1)b_{12}^{}& (1\omega )\lambda _2b_{22}\end{array}\right)=0.$$
(2)
If $`\omega 1`$, a solution with non-vanishing off-diagonal elements is allowed if $`\omega ^2=1`$, i.e., $`\omega =1`$ This leads to
$$A=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),B=\left(\begin{array}{cc}0& b\\ b^{}& 0\end{array}\right),$$
(3)
which can always be transformed to $`\sigma _x`$ and $`\sigma _y`$, by rotations and adequate normalization. These are the operators of the example (1). For spin 1/2 the set of GHZ operators are in this sense unique.
## 2 Spin one
For higher spins the proof proceeds along the same lines. We find one case of interest, with $`\omega =1`$,
$$A=\left(\begin{array}{ccc}1& & \\ & 1& \\ & & 1\end{array}\right),B=\left(\begin{array}{ccc}0& b& c\\ b^{}& 0& 0\\ c^{}& 0& 0\end{array}\right).$$
(4)
In the basis where $`B`$ is diagonal $`A`$ and $`B`$ read
$$A=\left(\begin{array}{ccc}1& 0& 0\\ 0& 0& 1\\ 0& 1& 0\end{array}\right),B=\sqrt{|b|^2+|c|^2}\left(\begin{array}{ccc}0& & \\ & 1& \\ & & 1\end{array}\right)$$
(5)
which proves the assertion in the case of spin one, as a rotation around $`x`$ brings $`B`$ into the form $`0\sigma _y`$, while $`A`$ is left as $`1\sigma _x`$, up to normalizations.
## 3 Spin 3/2
For spin 3/2, in addition to cases that reduce straightforwardly to those of lower spins, we find:
$$A=\left(\begin{array}{cccc}1& & & \\ & 1& & \\ & & 1& \\ & & & 1\end{array}\right),B=\left(\begin{array}{cccc}0& a& b& c\\ a^{}& 0& 0& 0\\ b^{}& 0& 0& 0\\ c^{}& 0& 0& 0\end{array}\right).$$
(6)
In the basis where $`B`$ is diagonal $`A`$ and $`B`$ read
$$A=\left(\begin{array}{cccc}0& 1& 0& 0\\ 1& 0& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& 1\end{array}\right),B=\sqrt{|a|^2+|b|^2+|c|^2}\left(\begin{array}{cccc}1& & & \\ & 1& & \\ & & 0& \\ & & & 0\end{array}\right)$$
(7)
which is again diagonal in two, 2x2, blocks.
The last case corresponds to
$$A=\left(\begin{array}{cccc}1& & & \\ & 1& & \\ & & 1& \\ & & & 1\end{array}\right),B=\left(\begin{array}{cccc}0& a& 0& b\\ a^{}& 0& c^{}& 0\\ 0& c& 0& d\\ b^{}& 0& d^{}& 0\end{array}\right).$$
(8)
The following list of unitary transformations bring these matrices to the desired form:
a) With
$$F=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 0& 1& 0\\ 0& 1& 0& 0\\ 0& 0& 0& 1\end{array}\right),$$
(9)
$`F^{}=F=F^1`$, we find
$$A^{}=FAF=\left(\begin{array}{cc}I& \\ & I\end{array}\right),B^{}=FBF=\left(\begin{array}{cc}& \\ ^{}& \end{array}\right),$$
(10)
where
$$=\left(\begin{array}{cc}a& b\\ c& d\end{array}\right).$$
b) A unitary transformation of the form $`U=\left(\begin{array}{cc}U_1& \\ & U_2\end{array}\right)`$ leaves $`A^{}`$ invariant and allows to diagonalize $``$
$$A^{\prime \prime }=A^{},B^{\prime \prime }=UB^{}U^{}=\left(\begin{array}{cc}& U_1U_2^{}\\ (U_1U_2^{})^{}& \end{array}\right)=\left(\begin{array}{cccc}& & m& 0\\ & & 0& n\\ m^{}& 0& & \\ 0& n^{}& & \end{array}\right).$$
(11)
We have used the result that the generic matrix $``$ can be brought to a diagonal form with two unitary matrices $`U_1`$, $`U_2`$.
c) Finally, acting with $`F`$ again,
$$A^{\prime \prime \prime }=A,B^{\prime \prime \prime }=\left(\begin{array}{cccc}0& m& & \\ m^{}& 0& & \\ & & 0& n\\ & & n^{}& 0\end{array}\right),$$
(12)
which completes the proof.
## 4 Conclusions
We conclude that the equation $`AB=\omega BA`$ is very restrictive on $`\omega `$ and on the possible forms of A and B; as the Hilbert space dimension increases, with increasing spin, all its solutions for $`\omega 1`$ have $`\omega =1`$ and are essentially direct sums of the two-dimensional $`\sigma _x`$ and $`\sigma _y`$. In this sense there are no solutions that could, in principle, enrich the possibilities opened by the GHZ theorem.
## 5 Acknowledgments
J.S., J.T., R.T. acknowledge the Centro de Ciencias de Benasque for hospitality while this work was beeing done. J.S. also acknowledges the Department ECM for hospitality and financial support. J.T. and R.T. acknowledge financial support by CICYT project AEN 98-0431, CIRIT project 1998 SGR-00026 and CEC project IST-1999-11053.
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# Vector Boson Production at Hadron Colliders: Results from HERWIG and Resummed Calculations1footnote 11footnote 1Talk given at the 22nd annual MRST (Montreal-Rochester-Syracuse-Toronto) Meeting on High Energy Physics (MRST 2000), Rochester, NY, U. S. A., 8-9 May 2000.
## Introduction
The production of vector bosons $`W`$, $`Z`$ and $`\gamma `$ altarelli is one of the most interesting processes in the phenomenology of hadron collisions and provides an environment to test both Quantum Chromodynamics and the Standard Model of electroweak interactions (see lhc ; lh for a review). The lowest-order processes $`q\overline{q}^{}V`$ are not sufficient to make reliable predictions, but the initial-state radiation should be taken into account. Monte Carlo event generators and resummed analytical calculations are available tools to describe the multiple radiation accompanying the incoming hadrons.
Standard Monte Carlo algorithms herwig ; pythia describe the initial-state parton showers in the soft/collinear approximation, but can have ‘dead zones’, where no radiation is permitted. The radiation in these regions can be generated by the use of the exact first-order matrix element. Referring to the HERWIG event generator, matrix-element corrections to Drell–Yan processes have been implemented in corsey , following the general method of sey , and included in the new version HERWIG 6.1 herwig61 .
Another possible approach consists of performing an analytical resummation of the large logarithmic coefficients which multiply the strong coupling constant. Considering the transverse momentum $`q_T`$ distribution, logarithms of the ratio $`m_V/q_T`$, $`m_V`$ being the vector boson mass, arise in calculating higher-order corrections to the Born process. The resummation of these logarithms, which are large at small $`q_T`$, was initially proposed by Dokshitzer, Dyakonov and Troyan (DDT) ddt , then accomplished by Collins, Sterman and Soper (CSS) css . CSS performed the resummation in the space of the impact parameter $`b`$, which is the Fourier conjugate of $`q_T`$. Their results have been implemented numerically in ly ; balazs , while more recent analyses can be found in ellis ; fnr ; kulesza , where the resummation is performed in both $`q_T`$\- and $`b`$-space.
In this paper, we review some results for the $`W/Z`$ transverse momentum distribution according to the HERWIG event generator and resummed calculations.<sup>2</sup><sup>2</sup>2See also huston for a similar comparison for Higgs production at hadron colliders.
## The HERWIG parton shower algorithm
HERWIG simulates the initial-state radiation in hadron collisions according to a ‘backward evolution’ marweb , in which the scale is reduced away from the hard vertex and traces the hard-scattering partons back into the incoming hadrons. The branching algorithm relies on the universal structure of the elementary probability in the leading infrared approximation. The probability of the emission of an additional soft/collinear parton from a parton $`i`$ is given by the general result:
$$dP=\frac{dq_i^2}{q_i^2}\frac{\alpha _S\left(\frac{1z_i}{z_i}q_i\right)}{2\pi }P_{ab}(z_i)dz_i\frac{\mathrm{\Delta }_{S,a}(q_{i\mathrm{max}}^2,q_c^2)}{\mathrm{\Delta }_{S,a}(q_i^2,q_c^2)}\frac{x_i/z_i}{x_i}\frac{f_b(x_i/z_i,q_i^2)}{f_a(x_i,q_i^2)}.$$
(1)
The ordering variable is $`q_i^2=E^2\xi _i`$, where $`E`$ is the energy of the parton that splits and $`\xi _i=\frac{pp_i}{EE_i}`$, with $`p`$ and $`p_i`$ being the four-momenta of the splitting and of the emitted parton respectively; $`z_i`$ is the energy fraction of the outgoing space-like parton with respect to the incoming one; $`P_{ab}(z)`$ is the Altarelli–Parisi splitting function for a parton $`a`$ evolving in $`b`$. In the approximation of massless partons, we have $`\xi _i=1\mathrm{cos}\theta `$, where $`\theta `$ is the emission angle to the incoming hadron direction. For soft emission ($`E_iE`$), ordering according to $`q_i^2`$ corresponds to angular ordering. When the emission is hard, the energy of the radiated parton is similar to that of the splitting parton, so $`q_i^2`$-ordering corresponds to transverse momentum ordering. In (1) $`f_a(x_i,q_i^2)`$ is the parton distribution function for the partons of type $`a`$ in the initial-state hadron, $`x_i`$ being the parton energy fraction. The function
$$\mathrm{\Delta }_S(q_2^2,q_1^2)=\mathrm{exp}\left[\frac{\alpha _S}{2\pi }_{q_1^2}^{q_2^2}\frac{dk^2}{k^2}_{Q_1/Q_2}^{1Q_1/Q_2}𝑑zP(z)\right]$$
(2)
is the Sudakov form factor, expressing the probability of no-resolvable branching in the range $`q_1^2<q^2<q_2^2`$. The ratio of form factors in (1) is therefore the probability of no branching at higher values of $`q_i^2`$. Unitarity dictates that the Sudakov form factor sums up all-order virtual and unresolved contributions. In (1), $`q_{i\mathrm{max}}`$ is the maximum value of $`q`$, fixed by the hard process, and $`q_c`$ is the value at which the backward evolution is terminated, corresponding, in the case of HERWIG, to a cutoff on the transverse momentum of the showering partons. However, since $`q_c`$ is smaller than the minimum scale at which the parton distribution functions are evaluated, an additional cutoff on the evolution variable $`q_i^2`$ has to be set.
If the backward evolution has not resulted in a valence quark, an additional non-perturbative parton emission is generated to evolve back to a valence quark. Such a valence quark has a Gaussian distribution with respect to the non-perturbative intrinsic transverse momentum in the hadron, with a width $`q_{T\mathrm{int}}`$ that is an adjustable parameter and whose default value is zero.
We need finally to specify the showering frame, the variables $`q_i^2`$ and $`z_i`$ being frame-dependent. One can show that, as a result of the $`q^2`$-ordering, the maximum $`q`$-values of two colour connected partons $`i`$ and $`j`$ are related via $`q_{i\mathrm{max}}q_{j\mathrm{max}}=p_ip_j`$, which is Lorentz-invariant. For vector boson production, symmetric limits are set in HERWIG: $`q_{i\mathrm{max}}=q_{j\mathrm{max}}=\sqrt{p_ip_j}`$. Furthermore, the energy of the parton which initiates the cascade is given by $`E=q_{\mathrm{max}}=\sqrt{p_ip_j}`$. It follows that ordering according to $`q^2`$ implies $`\xi <z^2`$.
The region $`\xi >z^2`$ is therefore a ‘dead zone’ for the shower evolution. In such a zone the physical radiation is not logarithmically enhanced, but not completely absent as happens in the standard algorithm. We therefore need to improve the HERWIG parton showers by the use of matrix-element corrections.
## Matrix-element corrections
According to sey , we populate the ‘dead zone’ of the phase space using the exact $`𝒪(\alpha _S)`$ matrix element (hard correction). We also correct the emission in the already-populated region using the first-order result any time an emission is capable of being the ‘hardest so far’ (soft correction), where the hardness of an emission is measured in terms of the transverse momentum of the emitted parton relative to the splitting one.
We consider the process $`q(p_1)\overline{q}^{}(p_2)V(q)g(p_3)`$, define the Mandelstam variables $`\widehat{s}=(p_1+p_2)^2`$, $`\widehat{t}=(p_1p_3)^2`$ and $`\widehat{u}=(p_2p_3)^2`$ and obtain the total phase-space limits
$`m_V^2<`$ $`\widehat{s}`$ $`<s,`$ (3)
$`m_V^2\widehat{s}<`$ $`\widehat{t}`$ $`<\mathrm{\hspace{0.33em}\hspace{0.33em}0},`$ (4)
$`s`$ being the total centre-of-mass energy. We observe that the soft singularity corresponds to $`s=m_V^2`$ and the lines $`\widehat{t}=0`$ and $`\widehat{t}=m_V^2\widehat{s}`$ to collinear gluon emission.
After relating the parton shower variables $`z`$ and $`\xi `$ to $`\widehat{s}`$ and $`\widehat{t}`$, as done in corsey , and setting $`\xi <z^2`$, one can get the HERWIG phase-space limits in terms of $`\widehat{s}`$ and $`\widehat{t}`$. In Fig. 1 we plot the total and the HERWIG phase space for $`\sqrt{s}=200`$ GeV and $`m_V=80`$ GeV; the soft and collinear singularity are inside the HERWIG region.
We calculate the differential cross section with respect to $`\widehat{s}`$ and $`\widehat{t}`$ following the prescription of sey2 , where it is shown that, assuming that the rapidity and the virtuality of the vector boson are fixed by the Born process, the following factorization formula holds:
$$d^2\sigma =\sigma _0\frac{f_{q/1}(\chi _1)f_{\overline{q}^{}/2}(\chi _2)}{f_{q/1}(\eta _1)f_{\overline{q}^{}/2}(\eta _2)}\frac{C_F\alpha _S}{2\pi }\frac{d\widehat{s}d\widehat{t}}{\widehat{s}^2\widehat{t}\widehat{u}}\left[(m_V^2\widehat{u})^2+(m_V^2\widehat{t})^2\right],$$
(5)
where $`f_{q/1}(\chi _1)`$ and $`f_{\overline{q}^{}/2}(\chi _2)`$ are the parton distribution functions of the scattering partons inside the incoming hadrons 1 and 2 for energy fractions $`\chi _1`$ and $`\chi _2`$ in the process $`q\overline{q}^{}Vg`$, while $`f_{q/1}(\eta _1)`$ and $`f_{\overline{q}^{}/2}(\eta _2)`$ refer to the Born process.
A similar treatment holds for the Compton process $`q(p_1)g(p_3)q^{}(p_2)V(q)`$, as discussed in corsey .
The distribution (5) or the equivalent one for the Compton process is implemented to generate events in the missing phase space and in the populated region every time an emission is the hardest so far.
## Transverse momentum distributions: HERWIG results
An interesting observable to study is the vector boson transverse momentum, which is constrained to be $`q_T<m_V`$ in the soft/collinear approximation. After matrix-element corrections, a fraction of events at higher $`q_T`$ is to be expected. In Figs. 2 and 3 we plot the $`W`$ $`q_T`$ distribution at the Tevatron and at the LHC, according to HERWIG 5.9 and HERWIG 6.1, the new version including matrix-element corrections, for $`q_{T\mathrm{int}}=0`$. We see a big effect at large $`q_T`$: after some $`q_T`$ the 5.9 version does not generate events anymore, while we still have a non-zero cross section after matrix-element corrections.
Moreover, a slight suppression can be seen at small $`q_T`$. It is related to the fact that, although we are providing the shower with the tree-level $`𝒪(\alpha _S)`$ corrections, virtual contributions are missing and we still get the leading-order cross section. The enhancement at large $`q_T`$ is therefore compensated by a suppression in the low-$`q_T`$ range.
It is now interesting to compare the HERWIG results with some Tevatron data. In Fig. 4 we compare the HERWIG 6.1 distribution with some DØ data d0 and find reasonable agreement over the whole $`q_T`$ range. As shown in corsey , smearing the HERWIG curve to account for detector effects must be included to achieve this agreement. Also, we do not see any relevant impact of setting $`q_{T\mathrm{int}}=1`$ GeV after detector corrections.
In Fig. 5 we compare the HERWIG 5.9 and 6.1 results, for different values of $`q_{T\mathrm{int}}`$, with some CDF data cdf , already corrected for detector effects. We find good agreement after matrix-element corrections, while the 5.9 version is not able to fit in with the data for $`q_T>50`$ GeV. At low $`q_T`$, the best agreement to the data is obtained for $`q_{T\mathrm{int}}=2`$ GeV, as shown in Fig. 6. While the $`Z`$ distribution strongly depends on $`q_{T\mathrm{int}}`$ at small $`q_T`$, in corsey1 and Fig. 7 it is shown that the ratio $`R`$ of the $`W`$ and $`Z`$ spectra is approximately independent of it.<sup>3</sup><sup>3</sup>3The negative slopes of the plots in Fig. 7 are due to the $`W/Z`$ mass difference.. Such a ratio is one of the main inputs for the $`W`$ mass measurement in hadron collisions and it is good news that it does not depend on unknown non-perturbative parameters.
## Resummed calculations
Another possible approach to study the vector boson transverse momentum distribution consists of resumming the logarithmic terms $`l=\mathrm{log}(m_V^2/q_T^2)`$ in the low-$`q_T`$ range. It is interesting to compare the HERWIG phenomenological results with those of some resummed calculations, in particular ellis and fnr . According to ddt , the resummed differential cross section for $`W`$ production can be written as:
$`{\displaystyle \frac{d^2\sigma }{dm_V^2dq_T^2}}`$ $`=`$ $`\sigma _0{\displaystyle \frac{d}{dq_T^2}}{\displaystyle \underset{q,q^{}}{}}|V_{q\overline{q}^{}}|^2{\displaystyle _0^1}𝑑x_1𝑑x_2\delta (x_1x_2\tau )`$ (6)
$`\times `$ $`[f_{q/1}(x_1,q_T)f_{\overline{q}^{}/2}(x_2,q_T)\mathrm{exp}[S(m_V,q_T)]+(q\overline{q}^{})],`$
where $`V_{q\overline{q}^{}}`$ is the relevant Cabibbo–Kobayashi–Maskawa matrix element and $`\tau =m_V^2/s`$. In (6), $`\mathrm{exp}[S(m_V,q_T)]`$ is a Sudakov-like form factor which resums the large logarithms associated to the initial-state radiation. It reads:
$$S(m_V,q_T)=_{q_T^2}^{m_V^2}\frac{d\mu ^2}{\mu ^2}\left[A(\alpha _S(\mu ^2))\mathrm{log}\frac{m_V^2}{\mu ^2}+B(\alpha _S(\mu ^2))\right],$$
(7)
where $`A(\alpha _S)`$ and $`B(\alpha _S)`$ can be expanded as:
$$A(\alpha _S)=A_1\alpha _S+A_2\alpha _S^2+\mathrm{};B(\alpha _S)=B_1\alpha _S+B_2\alpha _S^2+\mathrm{}$$
(8)
As far as the logarithms which contribute to the resummation are concerned, two conflicting nomenclatures exist. One consists of looking at Sudakov exponent, where the leading logarithms (LL) are $`\alpha _S^nl^{n+1}`$ and the next-to-leading ones (NLL) $`\alpha _S^nl^n`$. It is straightforward to show that the LL contributions are obtained by keeping only the $`A_1`$ term in the expansions (8) while NLL accuracy is achieved by considering $`A_2`$ and $`B_1`$ as well. In this sense, the approach fnr is NLL.
Another classification relies on the expansion of the exponent
$$S(m_V,q_T)=\underset{n}{}c_{n,n+1}\alpha _S^nl^{n+1}+\underset{n}{}c_{n,n}\alpha _S^nl^n,$$
(9)
where the leading term is $`\alpha _Sl^2`$, so that the leading contributions to $`\mathrm{exp}[S(m_V,q_T)]`$ are $`\alpha _S^nl^{2n}`$, terms $`\alpha _S^nl^{2n1}`$ being next-to-leading. This is equivalent to saying that in the differential cross section the LL and NLL contributions are $`(1/q_T^2)\alpha _S^nl^{2n1}`$ and $`(1/q_T^2)\alpha _S^nl^{2n2}`$ respectively. According to this nomenclature, the calculations ellis and kulesza are NNLL and NNNLL respectively.
In the $`b`$-space formalism, following fnr , the differential cross section reads:
$`{\displaystyle \frac{d^2\sigma }{dm_V^2dq_T^2}}`$ $`=`$ $`{\displaystyle \frac{\sigma _0}{4\pi }}{\displaystyle \underset{q,q^{}}{}}|V_{q\overline{q}^{}}|^2{\displaystyle _0^1}𝑑x_1𝑑x_2\delta (x_1x_2\tau ){\displaystyle d^2be^{i\stackrel{}{q_T}\stackrel{}{b}}}`$ (10)
$`\times `$ $`[f_{q/1}(x_1,c_1/b)f_{\overline{q}^{}/2}(x_2,c_1/b)\mathrm{exp}[S(m_V,b)]+(q\overline{q}^{})],`$
where $`c_1`$ and $`c_2`$ are integration constants of order 1 and $`S(m_V,b)`$ is the Sudakov exponent in $`b`$-space.
For high $`b`$ values, i.e. small $`q_T`$, non-perturbative effects are taken into account via a Gaussian function $`F_{NP}=\mathrm{exp}(gb^2)`$, as suggested in ly . In both ellis and fnr the value $`g=3`$ $`\mathrm{GeV}^2`$ is chosen.
Also, in order to allow resummed calculations to be reliable even at large $`q_T`$, we wish to match the calculations of ellis and fnr to the exact $`𝒪(\alpha _S)`$ result. We add the first-order cross section to the resummed result and, in order to avoid double counting, we subtract off the term which they have in common, which is the $`q_T0`$ limit of the exact $`𝒪(\alpha _S)`$ result. According to our prescription, the matching works fine if at the point $`q_T=m_V`$ we have a continuous distribution.
## Comparison of HERWIG and resummed calculations
A detailed and general discussion on the comparison of angular-ordered parton shower algorithms with resummed calculations for Drell–Yan processes was already performed in cmw , where the authors showed that, for $`\tau 1`$, HERWIG always accounts for the term $`A_1`$, corresponding to the leading logarithms in the exponent, and $`B_1`$ as well. Furthermore, one is able to account for the NLL term $`A_2`$ by simply modifying the Altarelli–Parisi splitting function introducing a second-order contribution
$$P_{qq}^{}(\alpha _S,z)=\frac{\alpha _S}{2\pi }C_F\frac{1+z^2}{1z}+\frac{C_F}{2}\left(\frac{\alpha _S}{\pi }\right)^2\frac{K}{1z},$$
(11)
where the $`K`$ factor is given by:
$$K=C_A\left(\frac{67}{18}\frac{\pi ^2}{6}\right)\frac{5}{9}N_f,$$
(12)
$`N_f`$ being the number of flavours, $`C_F=4/3`$ and $`C_A=1/2`$. This is equivalent to redefining the QCD parameter $`\mathrm{\Lambda }`$ to the ‘Monte Carlo’ $`\mathrm{\Lambda }_{MC}`$:
$$\mathrm{\Lambda }_{MC}=\mathrm{\Lambda }\mathrm{exp}(K/4\pi \beta _0),$$
(13)
with $`\beta _0=(11C_A2N_f)/(12\pi )`$. Even after these replacements, the HERWIG algorithm cannot be considered completely accurate at the next-to-leading level, since it is still missing higher-order contributions in the strong coupling constant or the parton distribution functions (see, for instance, the discussion in lh ).
In Fig. 8 we show the $`W`$ transverse momentum distribution at the Tevatron in the low-$`q_T`$ range according to HERWIG 6.1 and the calculations of ellis in $`q_T`$-space and of fnr in $`q_T`$\- and $`b`$-space. The HERWIG curve lies within the range of the resummed calculations, which is a reasonable result, considering that we are actually comparing different approaches. In Fig. 9 we consider the whole $`q_T`$ range, with the resummations matched to the exact first-order amplitude. We find that the matching works fine only for the approach fnr in $`q_T`$-space, the others showing a step at $`q_T=m_W`$, due to uncompensated contributions of order $`\alpha _S^2`$ or higher. The well-matched distribution agrees with the HERWIG 6.1 prediction at large $`q_T`$.
## Conclusions
We studied the initial-state radiation in vector boson production according to the HERWIG event generator and some resummed calculations. In particular, we investigated the effect of the recently-implemented matrix-element corrections to the HERWIG algorithm. We found a big effect of such corrections on $`W/Z`$ transverse momentum distributions at the Tevatron and at the LHC, and good agreement with the DØ and CDF data, with a crucial role played by such corrections in order to be able to fit in with the data at large $`q_T`$. We also found that, even though the spectra at small $`q_T`$ do depend on the intrinsic non-perturbative transverse momentum, the ratio of the $`W`$ to the $`Z`$ spectrum is roughly independent of it. We then considered some resummed calculations, which we matched to the exact $`𝒪(\alpha _S)`$ matrix element, which makes them reliable at large $`q_T`$ as well. We found reasonable agreement of such approaches with HERWIG and fine matching only for the calculation which keeps all the next-to-leading logarithms in the Sudakov exponent in $`q_T`$-space.
Finally, we have to say that the discussed method of improving the initial-state radiation in parton-shower Monte Carlo simulations can be extended to a wide range of interesting processes for the phenomenology of hadron colliders. The implementation of hard and soft corrections to top and Higgs production is in progress.
## Acknowledgements
The presented results have been obtained in collaboration with Mike Seymour. We also acknowledge Lynne Orr for a careful reading of this manuscript. This work was supported by grant number DE-FG02-91ER40685 from the U.S. Dept. of Energy.
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# 1 Introduction
## 1 Introduction
### 1.Star product.
Let us introduce, following a star product that is an associative $`R[[\mathrm{}]]`$ – bilinear product on algebra $`A[[\mathrm{}]]`$ of smooth functions on a finite-dimensional $`C^{\mathrm{}}`$-manifold:
$$\begin{array}{c}fg=fg+\mathrm{}B_1(f,g)+\mathrm{}^2B_2(f,g)+\mathrm{}A[[\mathrm{}]]\end{array}$$
(1)
where $`\mathrm{}`$ is a formal variable and $`B_i(f,g)`$ – bidifferential operators.
Associativity for the $`n^{th}`$ order means: $`(fg)h=f(gh)+O(\mathrm{}^{n+1}).`$
### 2.Gauge group.
There is a natural gauge group which acts on star products:
$$\begin{array}{c}^{},f^{}(\mathrm{})=Df(\mathrm{}),f^{}{}_{}{}^{}g_{}^{}=D(D^1f^{}D^1g^{})\end{array}$$
(2)
where $`D=1+_{i1}\mathrm{}^iD_i`$ and $`D_i`$’s are arbitrary differential operators
$$ff+\mathrm{}D_1f+\mathrm{}^2D_2f+O(\mathrm{}^3)$$
Associativity of the new star product is obvious since $`f^{}{}_{}{}^{}g_{}^{}{}_{}{}^{}h_{}^{}=D(fgh)`$.
It follows from above
$$\begin{array}{c}B_{1}^{}{}_{}{}^{}(f,g)=B_1(f,g)+D_1(fg)fD_1(g)gD_1(f)\end{array}$$
(3)
$`B_1(f,g)`$ can be chosen to be a skew-symmetric bi-vector field (see ). Then, we put $`B_1(f,g)=\alpha ^{ab}_af_bg`$, $`\alpha ^{ab}=\alpha ^{ba}`$. The Poisson bi-vector may depend on $`\mathrm{}`$:
$$\begin{array}{c}\alpha ^{ab}(\mathrm{})=\underset{i0}{}\mathrm{}^i\alpha _i^{ab}\end{array}$$
(4)
The second order term $`O(\mathrm{}^2)`$ in the associativity equation $`f(gh)=(fg)h`$ implies that $`\alpha `$ gives a Poisson structure on $`X`$,
$$f,g,h\{f\{g,h\}\}+\{g,\{h,f\}\}+\{h,\{f,g\}\}=0,$$
where $`\{f,g\}:=\frac{fggf}{\mathrm{}}|_{\mathrm{}=0}`$
### 3.Moyal product.
An example of the star product is the Moyal product:
$$\begin{array}{c}fg=fg+\mathrm{}\vartheta ^{ij}_if_jg+\\ \\ +\frac{\mathrm{}^2}{2!}\vartheta ^{ij}\vartheta ^{kl}_i_kf_j_lg+\frac{\mathrm{}^3}{3!}\vartheta ^{ij}\vartheta ^{kl}\vartheta ^{mn}_i_k_mf_j_l_ng+\mathrm{}=\\ =\underset{n=0}{\overset{\mathrm{}}{}}\frac{\mathrm{}^n}{n!}\underset{i_1,\mathrm{},i_n;j_1,\mathrm{},j_n}{}\underset{k=1}{\overset{n}{}}\vartheta ^{i_kj_k}(\underset{k=1}{\overset{n}{}}_{i_k})(f)\times (\underset{k=1}{\overset{n}{}}_{j_k})(g)=\\ =e^{\mathrm{}\vartheta ^{ij}_{i}^{}{}_{}{}^{(1)}_{j}^{}{}_{}{}^{(2)}}f(x_{(1)})g(x_{(2)})|_{x_{(1)}=x_{(2)}=x}\end{array}$$
(5)
where $`\vartheta ^{ij}`$ is constant and skew-symmetric.
### 4. Kontsevich formula.
In paper the following theorem was stated
### Theorem.
(1) Let $`\alpha `$ be a Poisson bi-vector field in a domain of $`R^d`$. The formula:
$$\begin{array}{c}fg:=\underset{n=0}{\overset{\mathrm{}}{}}\mathrm{}^n\underset{\mathrm{\Gamma }G_n}{}\omega _\mathrm{\Gamma }B_{\mathrm{\Gamma },\alpha }(f,g)\end{array}$$
(6)
defines an associative product. (2) Changing coordinates, one obtains a gauge equivalent star product.
In formula (6) expressions $`B_{\mathrm{\Gamma },\alpha }(f,g)`$ are constructed with the help of Kontsevich’s diagrams $`\mathrm{\Gamma }G_n`$, $`G_n`$ being a set of the $`n^{th}`$-order diagrams and $`\omega _\mathrm{\Gamma }`$ are constant coefficients corresponding to diagrams $`\mathrm{\Gamma }G_n`$. The values of these coefficients can be computed by the prescription given in .
Up to the second order, this formula can be written as follows
$$\begin{array}{c}fg=fg+\mathrm{}\alpha ^{ab}_af_bg++\frac{1}{2}\mathrm{}^2\alpha ^{ab}\alpha ^{cd}_a_cf_b_dg+\\ +\frac{1}{3}\mathrm{}^2\alpha ^{as}_s\alpha ^{bc}(_a_bf_cg+_a_bg_cf)+O(\mathrm{}^3)\end{array}$$
(7)
### 5. Purpose of the paper.
We are going to make a direct check of Kontsevich’s statement that changing coordinates one obtains a gauge equivalent star product, up to the $`\mathrm{}^3`$-order starting from the constant bi-vector (in this case, the star product is Moyal’s one). For the symplectic case, this fact was stated in .
We make a change of variables in the Moyal product, which is nothing but Kontsevich’s star product with the constant Poisson bi-vector (let us denote it $`\vartheta ^{ij}`$), and try to represent it with a help of gauge transformations in the form of (6) in new coordinates $`z^a(x^i)`$ using only the Poisson bi-vector field $`\alpha ^{ab}(\mathrm{})`$ $`\alpha _{0}^{}{}_{}{}^{ab}=\vartheta ^{ij}\frac{z^a}{x^i}\frac{z^b}{x^j}`$. The result should be compared with Kontsevich’s formula for non-constant (non-Moyal) $`\vartheta ^{ij}`$. In this way, we obtain formula (8). In the third order, we see that bi-vector terms can not be rewritten in new coordinates and thus are interpreted as $`\alpha _2^{ab}`$ in formula (4): $`\alpha ^{ab}(\mathrm{})=\underset{i0}{}\mathrm{}^i\alpha _i^{ab}`$. The definition of $`\alpha _2`$ is not unique as we may add to the formula (8) a bi-vector term (for example, $`_s\alpha ^{pt}_p\alpha ^{so}_o_t\alpha ^{ab}_af_bg`$) and thus subtract it from $`\alpha _2`$. However, this kind of terms always correspond to loop diagrams. If one considers the case when there are no loops, this consequently fixes the value of $`\alpha `$. In this case, the answer we are going to get is given by the following formula
$$\begin{array}{c}fg=fg+\mathrm{}\alpha ^{ab}_af_bg+\\ +\mathrm{}^2\left[\frac{1}{2}\alpha ^{ab}\alpha ^{cd}_a_cf_b_dg+\frac{1}{3}\alpha ^{as}_s\alpha ^{bc}(_a_bf_cg+_a_bg_cf)\right]+\\ +\mathrm{}^3[\frac{1}{6}\alpha ^{ab}\alpha ^{cd}\alpha ^{ho}_a_c_hf_b_d_og+\frac{1}{3}\alpha ^{tp}_p\alpha ^{as}_s_t\alpha ^{bc}(_a_cf_bg_a_cg_bf)+\\ +[\frac{2}{3}\alpha ^{dp}_p\alpha ^{as}_s\alpha ^{bc}+\frac{1}{3}\alpha ^{ap}_p\alpha ^{ds}_s\alpha ^{cb}]_a_cf_b_dg+\\ +\frac{1}{6}\alpha ^{as}\alpha ^{ct}_s_t\alpha ^{bd}(_a_b_cf_dg_a_b_cg_df)+\\ +\frac{1}{3}\alpha ^{as}_s\alpha ^{bc}\alpha ^{hd}(_a_b_hf_c_dg_a_b_hg_c_df)]+O(\mathrm{}^4)\end{array}$$
(8)
where
$$\begin{array}{c}\alpha ^{ab}=\vartheta ^{ij}\frac{z^a}{x^i}\frac{z^b}{x^j}+\mathrm{}^2[\frac{1}{3!}\vartheta ^{ij}\vartheta ^{kl}\vartheta ^{mn}\frac{^3z^a}{x^ix^kx^m}\frac{^3z^b}{x^jx^lx^n}\\ \frac{1}{18}S^{spt}_p_s_t\alpha ^{ab}\frac{1}{4}\vartheta ^{ij}\vartheta ^{kl}\frac{^2z^s}{x^ix^k}\frac{^2z^t}{x^jx^l}_s_t\alpha ^{ab}]+O(\mathrm{}^3)\end{array}$$
(9)
and $`S^{abs}`$ is determined in (15). The differential operator in the gauge transformation (2) necessary for obtaining (8),(9) is the following
$$D=1+\mathrm{}^2[\frac{1}{4}\vartheta ^{ij}\vartheta ^{kl}\frac{^2z^a}{x^ix^k}\frac{^2z^b}{x^jx^l}_a_b+$$
$$+\frac{1}{18}\vartheta ^{ij}\vartheta ^{kl}(\frac{^2z^a}{x^ix^k}\frac{z^b}{x^j}\frac{z^c}{x^l}+\frac{^2z^c}{x^ix^k}\frac{z^a}{x^j}\frac{z^b}{x^l}+\frac{^2z^b}{x^ix^k}\frac{z^c}{x^j}\frac{z^a}{x^l})_a_b_c]+O(\mathrm{}^4)$$
We show that $`\alpha (\mathrm{})`$ defined by (9) do not satisfy the Jacobi identity: $`\alpha ^{as}_s\alpha ^{bc}+\alpha ^{cs}_s\alpha ^{ab}+\alpha ^{bs}_s\alpha ca=0`$. Still $`\alpha _0`$ does satisfy the Jacobi identity and we will use this fact to compute coefficients in (8) from the requirement of associativity. These coefficients are in agreement with . It is discussed in section 5.
## 2 Diagram representation
There is a natural way to represent separate terms in Kontsevich’s formula by diagrams.
For the $`n^{th}`$-order term one needs $`n`$ vertices, each vertex containing $`\alpha `$, and two more vertices containing functions $`f`$ and $`g`$. Vertices can be connected by arrows. Each vertex is an origin of two ordered arrows. If an arrow ends in some vertex, it describes a partial derivative acting on this vertex.
The following diagrams correspond to formula (8) (here we do not draw the arrows as all of them point to the right):
Note that there are no loop diagrams in this picture. In particular, there is no structure of the type $`_c\alpha ^{da}_d\alpha ^{bc}`$.
## 3 Calculations up to the second order terms
Let us consider the terms up to the second order ($`\alpha =\alpha _0`$) in the Moyal formula (5) and make a change of variables in it. We are going to obtain formula (7) with the help of gauge transformations. The change of variables gives the following expression for the Moyal product
$$\begin{array}{c}fg=fg+\mathrm{}\alpha ^{ab}_af_bg+\\ +\frac{1}{2}\mathrm{}^2\vartheta ^{ij}\vartheta ^{kl}\left(\frac{^2z^a}{x^ix^k}\frac{f}{z^a}+\frac{z^a}{x^i}\frac{^2f}{z^az^c}\frac{z^c}{x^i}\right)\left(\frac{^2z^b}{x^jx^l}\frac{g}{z^b}+\frac{z^b}{x^l}\frac{^2g}{z^bz^d}\frac{z^d}{x^j}\right)=\\ =fg+\mathrm{}\alpha ^{ab}_af_bg+\frac{1}{2}\mathrm{}^2\alpha ^{ab}\alpha ^{cd}_a_cf_b_dg+\\ +\frac{1}{2}\mathrm{}^2\vartheta ^{ij}\vartheta ^{kl}\frac{^2z^a}{x^ix^k}\frac{^2z^b}{x^jx^l}_af_bg+\\ +\mathrm{}^2\vartheta ^{ij}\vartheta ^{kl}\frac{^2z^c}{x^ix^k}\frac{z^a}{x^l}\frac{z^b}{x^j}(_a_bf_cg+_a_bg_cf)+O(\mathrm{}^3)\end{array}$$
(10)
The last term is a symmetric bi-vector and can be canceled by the following gauge transformation
$$\begin{array}{c}D_{}^{}{}_{2}{}^{}=\frac{1}{4}\vartheta ^{ij}\vartheta ^{kl}\frac{^2z^a}{x^ix^k}\frac{^2z^b}{x^jx^l}_a_b\end{array}$$
(11)
Now we pay attention only to the terms $`^2fg`$, $`^2gf`$ which are
$$\begin{array}{c}\frac{1}{2}P_1=\frac{1}{2}\frac{^2z^c}{x^ix^k}\frac{z^a}{x^l}\frac{z^b}{x^j}(_a_bf_cg+_a_bg_cf)\end{array}$$
(12)
If one puts $`\alpha ^{ab}=\vartheta ^{ij}\frac{z^a}{x^i}\frac{z^b}{x^j}`$ in $`K_0\alpha ^{as}_s\alpha ^{bc}(_a_bf_cg+_a_bg_cf)`$, where $`K_0`$ is constant, this gives
$$\begin{array}{c}K_0\left(\frac{^2z^a}{x^ix^k}\frac{z^c}{x^l}\frac{z^b}{x^j}+\frac{^2z^c}{x^ix^k}\frac{z^a}{x^l}\frac{z^b}{x^j}\right)(_a_bf_cg+_a_bg_cf)=K_0P_1K_0P_2\end{array}$$
(13)
where $`P_2=\frac{^2z^a}{x^ix^k}\frac{z^c}{x^l}\frac{z^b}{x^j}(_a_bf_cg+_a_bg_cf)`$.
We are going now to find a gauge transformation that makes (12) equal to expression (13). Consider
$$\begin{array}{c}D_2=K_1S^{abc}_a_b_c\end{array}$$
(14)
$$\begin{array}{c}S^{abc}=\vartheta ^{ij}\vartheta ^{kl}\left(\frac{^2z^a}{x^ix^k}\frac{z^b}{x^j}\frac{z^c}{x^l}+\frac{^2z^c}{x^ix^k}\frac{z^a}{x^j}\frac{z^b}{x^l}+\frac{^2z^b}{x^ix^k}\frac{z^c}{x^j}\frac{z^a}{x^l}\right)\end{array}$$
(15)
Note that $`S^{abc}`$ is symmetric with respect to $`a,b,c`$. Transformation (14-15) adds to (12) the following terms
$$\begin{array}{c}D_2(fg)fD_2(g)gD_2(f)=3S^{abc}(_a_bf_cg+_a_bg_cf)=3K_1(P_1+2P_2)\end{array}$$
(16)
The equality $`(12)+(16)=(13)`$ looks like
$$\frac{1}{2}P_1+3K_1(P_1+2P_2)=K_0(P_1P_2)$$
$$K_0=\frac{1}{3},K_1=\frac{1}{18}$$
This means we should substitute in formula (3) the gauge transformations $`D_{}^{}{}_{2}{}^{}`$ (11) and $`D_2`$ (14-15), where $`K_1=\frac{1}{18}`$. Note that one can add to the star product in formula (3) the symmetric bi-vector
$$\begin{array}{c}\lambda \mathrm{}^2_d\alpha ^{ac}_c\alpha ^{bd}\end{array}$$
(17)
(where $`\lambda `$ is a constant) which corresponds to loop diagram.
It is easy to get corrections to the expression for $`B_3(f,g)`$ from formula (1) under gauge transformations $`D_2`$ made in the second order:
$$\begin{array}{c}B_3^{}(f,g)=B_3(f,g)+D_2B_1(f,g)B_1(D_2f,g)B_1(f,D_2g)\end{array}$$
(18)
Thus, it is possible to change higher order terms by changing $`\lambda `$. We will put $`\lambda =0`$.
## 4 Calculations up to the third order terms
Let us write down the terms of the third order from the Moyal product in the new coordinates:
$$\begin{array}{c}\frac{1}{6}\vartheta ^{ij}\vartheta ^{kl}\vartheta ^{mn}(\frac{^3f}{z^az^cz^e}\frac{z^e}{x^i}\frac{z^c}{x^k}\frac{z^a}{x^m}+\frac{^2f}{z^az^c}(\frac{^2z^a}{x^kx^m}\frac{z^c}{x^i}+\frac{^2z^a}{x^ix^m}\frac{z^c}{x^k}+\\ +\frac{^2z^a}{x^ix^k}\frac{z^c}{x^m})+\frac{^3z^a}{x^ix^kx^m}\frac{f}{z^a}\left)\right(\frac{^3g}{z^bz^dz^e}\frac{z^e}{x^j}\frac{z^d}{x^l}\frac{z^b}{x^n}+\\ +\frac{^2g}{z^bz^d}(\frac{^2z^b}{x^lx^n}\frac{z^d}{x^j}+\frac{^2z^b}{x^lx^j}\frac{z^d}{x^n}+\frac{^2z^b}{x^nx^j}\frac{z^d}{x^l})+\frac{^3z^b}{x^jx^lx^n}\frac{g}{z^b})\end{array}$$
(19)
Due to (18) we are to add to the above expression
$$\begin{array}{c}D_2B_1(f,g)B_1(D_2f,g)B_1(f,D_2g)+D_2^{}B_1(f,g)B_1(D_2^{}f,g)B_1(f,D_2^{}g)\end{array}$$
(20)
where $`D_2`$ and $`D_2^{}`$ are given by (14-15) and (11).
Again we are going to represent the sum of the Moyal product and gauge terms in diagrams (in the new coordinates).
An important remark is that there is no necessity to make gauge transformations in the third order, since the corrections to the third order terms $`B_3(f,g)`$ coming from the gauge transformation of the second order (20) are enough to represent this expression in the form of Kontsevich’s diagrams.
Now consider an example of computation for the terms of the type $`(^2f^2g)`$. Such terms in the Moyal product (19) are
$$\begin{array}{c}\frac{1}{6}\vartheta ^{ij}\vartheta ^{kl}\vartheta ^{mn}(\frac{^2z^a}{x^kx^m}\frac{z^c}{x^i}+\frac{^2z^a}{x^ix^m}\frac{z^c}{x^k}+\frac{^2z^a}{x^ix^k}\frac{z^c}{x^m})(\frac{^2z^b}{x^lx^n}\frac{z^d}{x^j}+\\ +\frac{^2z^b}{x^lx^j}\frac{z^d}{x^n}+\frac{^2z^b}{x^nx^j}\frac{z^d}{x^l})_a_cf_b_dg=\\ =\vartheta ^{ij}\vartheta ^{kl}\vartheta ^{mn}\left(\frac{1}{2}\frac{^2z^a}{x^ix^k}\frac{^2z^b}{x^jx^l}\frac{z^c}{x^m}\frac{z^d}{x^n}+\frac{^2z^a}{x^ix^k}\frac{^2z^b}{x^lx^n}\frac{z^c}{x^m}\frac{z^d}{x^j}\right)_a_cf_b_dg\end{array}$$
(21)
The gauge terms are
$$\begin{array}{c}D_2^{}:\frac{1}{2}\vartheta ^{ij}\vartheta ^{kl}\vartheta ^{mn}\frac{^2z^a}{x^ix^k}\frac{^2z^b}{x^jx^l}\frac{z^c}{x^m}\frac{z^d}{x^n}_a_cf_b_dg\end{array}$$
(22)
(Note that this gauge transformation cancels the first term in the right side of (21).)
$$\begin{array}{c}D_2:\frac{1}{3}S^{scd}_s\alpha ^{ab}_a_cf_b_dg\end{array}$$
(23)
where $`S^{scd}`$ is defined by (15). One may check that (23) is equal to
$$\begin{array}{c}\frac{1}{3}S^{scd}_s\alpha ^{ab}_a_cf_b_dg=\frac{1}{3}Y_1\frac{1}{3}Y_2\frac{1}{3}Y_3\end{array}$$
(24)
where $`Y_1`$,$`Y_2`$ and $`Y_3`$ are
$$\begin{array}{c}Y_1=\vartheta ^{ij}\vartheta ^{kl}\vartheta ^{mn}\frac{z^a}{x^j}\frac{^2z^d}{x^kx^i}\left(\frac{^2z^c}{x^mx^l}\frac{z^b}{x^n}\frac{^2z^b}{x^mx^l}\frac{z^c}{x^n}\right)\end{array}$$
(25)
$$\begin{array}{c}Y_2=\vartheta ^{ij}\vartheta ^{kl}\vartheta ^{mn}\frac{z^a}{x^j}\frac{z^d}{x^l}\frac{^2z^s}{x^ix^k}\frac{x^o}{z^s}\left(\frac{^2z^c}{x^mx^o}\frac{z^b}{x^n}\frac{^2z^b}{x^mx^o}\frac{z^c}{x^n}\right)\end{array}$$
(26)
$$\begin{array}{c}Y_3=\vartheta ^{ij}\vartheta ^{kl}\vartheta ^{mn}\frac{z^d}{x^j}\frac{^2z^a}{x^kx^i}\left(\frac{^2z^c}{x^mx^l}\frac{z^b}{x^n}\frac{^2z^b}{x^mx^l}\frac{z^c}{x^n}\right)\end{array}$$
(27)
We will try to find an answer in the form
$$\begin{array}{c}[A\alpha ^{dp}_p\alpha ^{as}_s\alpha ^{bc}+B\alpha ^{ap}_p\alpha ^{ds}_s\alpha ^{cb}]_a_cf_b_dg\end{array}$$
(28)
In old coordinates it looks like
$$\begin{array}{c}[A\alpha ^{dp}_p\alpha ^{as}_s\alpha ^{bc}+B\alpha ^{ap}_p\alpha ^{ds}_s\alpha ^{cb}]_a_cf_b_dg=A(Y_2Y_1)+B(Y_3Y_2)\end{array}$$
(29)
where $`A`$ and $`B`$ are constants. Note also that the last term in the right part of (21) is equal to $`Y_3`$ since
$$\vartheta ^{ij}\vartheta ^{kl}\vartheta ^{mn}\frac{^2z^a}{x^ix^k}\frac{^2z^c}{x^lx^n}\frac{z^b}{x^m}\frac{z^d}{x^j}_a_cf_b_dg=0$$
due to the skew symmetry with respect to the interchange $`(ac`$ and $`db)`$. Now we can find $`A`$ and $`B`$
$$Y_3\frac{1}{3}Y_1\frac{1}{3}Y_2\frac{1}{3}Y_3=A(Y_3Y_2)+B(Y_2Y_1)$$
$$A=\frac{2}{3},B=\frac{1}{3}$$
Thus, we have found the $`(^2f^2g)`$-terms in the deformation quantization formula (8):
$$\mathrm{}^3[\frac{2}{3}\alpha ^{dp}_p\alpha ^{as}_s\alpha ^{bc}+\frac{1}{3}\alpha ^{ap}_p\alpha ^{ds}_s\alpha ^{cb}]_a_cf_b_dg$$
All other terms can be found in the same manner. The only exception happens with the $`fg`$-terms. For these terms, one has
$$\begin{array}{c}\text{Moyal term}:\frac{1}{6}\vartheta ^{ij}\vartheta ^{kl}\vartheta ^{mn}\frac{^3z^a}{x^ix^kx^m}\frac{^3z^b}{x^jx^lx^n}\end{array}$$
(30)
$$\begin{array}{c}\text{Gauge}D_2:\frac{1}{18}S^{spt}_p_s_t\alpha ^{ab}_af_bg\end{array}$$
(31)
$$\begin{array}{c}\text{Gauge}D_2^{}:\frac{1}{4}\vartheta ^{ij}\vartheta ^{kl}\frac{^2z^s}{x^ix^k}\frac{^2z^t}{x^jx^l}_s_t\alpha ^{ab}_af_bg\end{array}$$
(32)
These terms can not be canceled by gauge transformations as they are skew symmetric, and none of them can be represented in the form of Kontsevich’s diagrams. Therefore, we are going to treat them as $`\alpha _{2}^{}{}_{}{}^{ab}_af_bg`$. The definition of $`\alpha _2`$ is not unique for the following reason. Formula (32) is contained in $`_p\alpha ^{so}_o\alpha ^{tp}_s_t\alpha ^{ab}_af_bg`$, however, this expression contains other terms. Formula (31) can not be realized in the form of Kontsevich’s diagrams at all because of the third derivative. The Moyal term (30) contains $`^3z^a^3z^b`$ and thus is contained in the only diagram: $`\alpha ^{pt}_t_s\alpha ^{ao}_o_p\alpha ^{bs}`$ but it again contains many other terms.
Thus we have following expression for the Poisson bi-vector
$$\begin{array}{c}\alpha ^{ab}=\vartheta ^{ij}\frac{z^a}{x^i}\frac{z^b}{x^j}+\mathrm{}^2[\frac{1}{3!}\vartheta ^{ij}\vartheta ^{kl}\vartheta ^{mn}\frac{^3z^a}{x^ix^kx^m}\frac{^3z^b}{x^jx^lx^n}\\ \frac{1}{18}S^{spt}_p_s_t\alpha ^{ab}\frac{1}{4}\vartheta ^{ij}\vartheta ^{kl}\frac{^2z^s}{x^ix^k}\frac{^2z^t}{x^jx^l}_s_t\alpha ^{ab}]+O(\mathrm{}^3)\end{array}$$
(33)
So we found $`\alpha =\alpha _0+\mathrm{}^2\alpha _2`$. It is obvious that $`\alpha ^{ab}=\vartheta ^{ij}\frac{z^a}{x^i}\frac{z^b}{x^j}`$ satisfy the Jacobi identity
$$\alpha _{0}^{}{}_{}{}^{as}_s\alpha _{0}^{}{}_{}{}^{bc}+\alpha _{0}^{}{}_{}{}^{cs}_s\alpha _{0}^{}{}_{}{}^{ab}+\alpha _{0}^{}{}_{}{}^{bs}_s\alpha _{0}^{}{}_{}{}^{ca}=0$$
The Poisson bi-vector $`\alpha (\mathrm{})`$ does not satisfy the Jacobi identity in the $`\mathrm{}^2`$-order. This statement may be proved in the following way. One may consider a certain type of terms in the right hand side of the Jacobi identity and show that there is no way to drop them off.
In the $`\mathrm{}^2`$-order, the Jacobi identity looks like
$$\alpha _{0}^{}{}_{}{}^{as}_s\alpha _{2}^{}{}_{}{}^{bc}+\alpha _{2}^{}{}_{}{}^{as}_s\alpha _{0}^{}{}_{}{}^{bc}+\alpha _{0}^{}{}_{}{}^{cs}_s\alpha _{2}^{}{}_{}{}^{ab}+\alpha _{2}^{}{}_{}{}^{cs}_s\alpha _{0}^{}{}_{}{}^{ab}+\alpha _{0}^{}{}_{}{}^{bs}_s\alpha _{2}^{}{}_{}{}^{ca}+\alpha _{2}^{}{}_{}{}^{bs}_s\alpha _{0}^{}{}_{}{}^{ca}=0$$
To see it is not true for $`\alpha `$ from (33) one should have a look on the terms $`z^4z^3z`$. The coordinate $`z`$ can have different types of indices: external $`a,b,c`$ and internal $`o,t,p`$. Thus, there is a number of the above expressions with different types of indices and each of them should satisfy the Jacobi identity in order to have it for the whole $`\alpha `$. This kind of terms appears from the two first terms in (33). The indices have the following structures: $`z^a^4z^b^3z^c`$ from the first term, $`z^a^4z^s^3z^b`$ and $`z^a^4z^b^3z^s`$ from the second one. Each of the above expressions makes the Jacobi identity non-valid. Still, we should remember that $`\alpha _2`$ can be defined in a different way if one adds $`fg`$ terms to the star product (8). The only such a term which gives the same type of terms in the Jacobi identity is $`_p\alpha ^{so}_o\alpha ^{tp}_s_t\alpha ^{ab}`$. It can cancel $`z^a^4z^b^3z^c`$ and add $`z^a^4z^s^3z^b`$, $`z^a^4z^b^3z^s`$ and $`z^a^4z^o^3z^s`$. In this case, the expression $`z^a^4z^o^3z^s`$ violates the Jacobi identity. Note that we did not consider the ”tadpole” diagrams (which contain $`_s\alpha ^{sa}`$), since there are no such terms in Kontsevich’s formula but they save the Jacobi identity neither.
## 5 Values of coefficients from associativity
In this section, we are going to obtain the coefficients in formula (8) from the condition of associativity and the Jacobi identity for $`\alpha _0`$. Associativity for the $`n^{th}`$-order means $`(fg)h=f(gh)+O(\mathrm{}^{n+1})`$. For the second and the third orders, we have (star product is defined by (1)):
$$\begin{array}{c}B_2(fg,h)+B_1(B_1(f,g),h)+B_2(f,g)h=B_2(f,gh)+B_1(f,B_1(g,h))+fB_2(g,h)\end{array}$$
(34)
$$\begin{array}{c}B_3(fg,h)+B_2(B_1(f,g),h)+B_1(B_2(f,g),h)+B_3(f,g)h=\\ =B_3(f,gh)+B_2(f,B_1(g,h))+B_1(f,B_2(g,h))+fB_3(g,h)\end{array}$$
(35)
In the second order condition (34), one may put $`\alpha =\alpha _0`$. In the third order (35), $`B_3(f,g)`$ contains the term depending on $`\alpha _2`$, but this is a bi-vector $`\alpha _{2}^{}{}_{}{}^{ab}_af_bg`$ and, thus, cancels from (35) due to the Leibnitz rule. It means there is no possibility to obtain some knowledge about $`\alpha _2`$ from associativity in the $`\mathrm{}^3`$-order. So we may use the Jacobi identity for $`\alpha `$ in equations (34),(35).
Let us put arbitrary coefficients in formula (8)
$$\begin{array}{c}fg=fg+\mathrm{}\alpha ^{ab}_af_bg+\\ +\mathrm{}^2[A_1\alpha ^{ab}\alpha ^{cd}_a_cf_b_dg+A_2\alpha ^{as}_s\alpha ^{bc}_a_bf_cg+A_3\alpha ^{as}_s\alpha ^{bc}_a_bg_cf]+\\ \mathrm{}^3[C_1\alpha ^{ab}\alpha ^{cd}\alpha ^{ho}_a_c_hf_b_d_og+C_2\alpha ^{tp}_p\alpha ^{as}_s_t\alpha ^{bc}_a_cf_bg+\\ +C_3\alpha ^{tp}_p\alpha ^{as}_s_t\alpha ^{bc}_a_cg_bf+\\ +[C_4\alpha ^{dp}_p\alpha ^{as}_s\alpha ^{bc}+C_5\alpha ^{ap}_p\alpha ^{ds}_s\alpha ^{cb}]_a_cf_b_dg+\\ +C_6\alpha ^{as}\alpha ^{ct}_s_t\alpha ^{bd}_a_b_cf_dg+C_7\alpha ^{as}\alpha ^{ct}_s_t\alpha ^{bd}_a_b_cg_df+\\ +C_8\alpha ^{as}_s\alpha ^{bc}\alpha ^{hd}_a_b_hf_c_dg+C_9\alpha ^{as}_s\alpha ^{bc}\alpha ^{hd}_a_b_hg_c_df]+O(\mathrm{}^4)\end{array}$$
(36)
Equation (34) simply gives $`A_1=\frac{1}{2}`$ and
$$((1A_2)\alpha ^{as}_s\alpha ^{bc}+(1A_3)\alpha ^{cs}_s\alpha ^{ab}+(A_2+A_3)\alpha ^{bs}_s\alpha ^{ca})_af_bg_ch=0$$
As the Jacobi identity is known to be satisfied by $`\alpha _0`$, one obtains
$$1A_2=1A_3=A_2+A_3$$
and thus $`A_2=A_3=\frac{1}{3}`$.
Similarly from equation (35) one may obtain the following results
$$C_1=\frac{1}{6},C_2=C_3=\frac{1}{3},C_4C_5=\frac{1}{3},C_6=C_7=\frac{1}{6}C_8=C_9=\frac{1}{3}$$
These results are in agreement with formula (8). At the same time, they are in agreement with . The only difference is that the authors of write $`C_4=C_5=\frac{1}{6}`$ and, in formula (8), $`C_4=\frac{2}{3},C_5=\frac{1}{3}`$. The both cases satisfy the obtained condition $`C_4C_5=\frac{1}{3}`$. However, in our case the two coefficients $`C_4`$ and $`C_5`$ are not fixed since the two diagrams corresponding to the terms $`^2f^2g`$ are dependent through the Jacobi identity for $`\alpha _0`$. Since the entire $`\alpha `$ does not satisfy the Jacobi identity, these coefficients in formula (8) can not be changed. We expect that our coefficients will better suit next orders calculations.
## 6 Conclusion
In the present paper, we checked, up to the third order in $`\mathrm{}`$ the statement made in that changing coordinates in Kontsevich’s star product leads to gauge equivalent star products. The manifest calculations were presented starting from the star product with the constant Poisson bi-vector (giving the Moyal product). In this way, we obtained formula (8) for the deformation quantization which is agreement with . We also obtained the same result from the requirement of associativity of the star product and the Jacobi identity for $`\alpha _0`$ (which is not valid for $`\alpha (\mathrm{})`$).
## 7 Acknowledgments
Author is grateful to A.Gerasimov, A.Losev, A.Mironov, K.Saraikin and especially to A.Morozov for introducing the problem and discussions. This work was partially supported by RFBR grant N00-02-16530 and the program for support of the scientific schools 00-15-96557.
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# Untitled Document
IFT-00/15
hep-ph/0007255
Two–mixing texture for three active neutrinos<sup>*</sup><sup>*</sup>*Supported in part by the Polish KBN–Grant 2 P03B 052 16 (1999–2000).
Wojciech Królikowski
Institute of Theoretical Physics, Warsaw University
Hoża 69, PL–00–681 Warszawa, Poland
Abstract
The conjecture that among three massive neutrinos $`\nu _1,\nu _2,\nu _3`$ there is no direct mixing between $`\nu _1`$ and $`\nu _3`$ leads to a two–mixing texture for three active neutrinos $`\nu _e,\nu _\mu ,\nu _\tau `$. This texture, much discussed previously, is neatly consistent with the observed deficits of solar $`\nu _e`$’s and atmospheric $`\nu _\mu `$’s, but (without extra mixing with at least one sterile neutrino $`\nu _s`$) predicts no LSND effect for accelerator $`\nu _\mu `$’s. In this option, the masses $`m_1^2\stackrel{<}{}m_2^2m_3^2`$ are readily estimated. The characteristic feature of the two–mixing neutrino texture that only the close neighbours in the hierarchy of massive neutrinos $`\nu _1,\nu _2,\nu _3`$ mix significantly may be physically meaningful. Going out from the notion of mixing matrix we construct an intrinsic occupation–number operator whose eigenvalues 0, 1, 2 numerate the three generations of massive neutrinos. Analogical constructions work also for charged leptons as well as for up and down quarks.
PACS numbers: 12.15.Ff , 14.60.Pq , 12.15.Hh .
July 2000
If one conjectures that in the generic Cabibbo–Kobayashi–Maskawa–type matrix for leptons ,
$$U=\left(\begin{array}{ccc}c_{13}c_{12}& c_{13}s_{12}& s_{13}e^{i\delta }\hfill \\ c_{23}s_{12}s_{13}s_{23}c_{12}e^{i\delta }& c_{23}c_{12}s_{13}s_{23}s_{12}e^{i\delta }& c_{13}s_{23}\hfill \\ s_{23}s_{12}s_{13}c_{23}c_{12}e^{i\delta }& s_{23}c_{12}s_{13}c_{23}s_{12}e^{i\delta }& c_{13}c_{23}\hfill \end{array}\right)$$
(1)
with $`s_{ij}=\mathrm{sin}\theta _{ij}>0`$ and $`c_{ij}=\mathrm{cos}\theta _{ij}0`$, $`(i,j=1,2,3)`$, there is practically no direct mixing of massive neutrinos $`\nu _1`$ and $`\nu _3`$ (i.e., $`\theta _{13}=0`$), then $`U`$ is reduced to the following two–mixing form much discussed previously :
$$U=\left(\begin{array}{ccc}c_{12}& s_{12}& 0\\ c_{23}s_{12}& c_{23}c_{12}& s_{23}\\ s_{23}s_{12}& s_{23}c_{12}& c_{23}\end{array}\right)=\left(\begin{array}{ccc}1& 0& 0\\ 0& c_{23}& s_{23}\\ 0& s_{23}& c_{23}\end{array}\right)\left(\begin{array}{ccc}c_{12}& s_{12}& 0\\ s_{12}& c_{12}& 0\\ 0& 0& 1\end{array}\right).$$
(2)
For the two–mixing option (2) the neutrino mixing formula $`\nu _\alpha =_iU_{\alpha i}\nu _i`$ takes the form
$`\nu _e`$ $`=`$ $`c_{12}\nu _1+s_{12}\nu _2,`$
$`\nu _\mu `$ $`=`$ $`c_{23}(s_{12}\nu _1+c_{12}\nu _2)+s_{23}\nu _3,`$
$`\nu _\tau `$ $`=`$ $`s_{23}(s_{12}\nu _1+c_{12}\nu _2)+c_{23}\nu _3,`$ (3)
while the inverse neutrino mixing formula $`\nu _i=_\alpha U_{\alpha i}^{}\nu _\alpha `$ gives
$`\nu _1`$ $`=`$ $`c_{12}\nu _es_{12}(c_{23}\nu _\mu s_{23}\nu _\tau ),`$
$`\nu _2`$ $`=`$ $`s_{12}\nu _e+c_{12}(c_{23}\nu _\mu s_{23}\nu _\tau ),`$
$`\nu _3`$ $`=`$ $`s_{23}\nu _\mu +c_{23}\nu _\tau .`$ (4)
In the representation, where the charged–lepton mass matrix is diagonal (and thus the corresponding diagonalizing matrix — unit), the lepton mixing matrix $`U=\left(U_{\alpha i}\right)`$ $`(\alpha =e,\mu ,\tau ,i=1,2,3)`$ is, at the same time, the diagonalizing matrix for neutrino mass matrix $`M=\left(M_{\alpha \beta }\right)`$ $`(\alpha ,\beta =e,\mu ,\tau )`$ , $`U^{}MU=\mathrm{diag}(m_1,m_2,m_3)`$ with $`m_1^2m_2^2m_3^2`$, so that $`M=\left(_iU_{\alpha i}U_{\beta i}^{}m_i\right)`$. In this case, the orthogonal two–mixing form (2) of $`U`$ leads to the real and symmetric
$$M=(\begin{array}{ccc}c_{12}^2m_1+s_{12}^2m_2& (m_2m_1)c_{12}s_{12}c_{23}& (m_2m_1)c_{12}s_{12}s_{23}\\ (m_2m_1)c_{12}s_{12}c_{23}& s_{23}^2m_3+c_{23}^2(s_{12}^2m_1+c_{12}^2m_2)& (m_3s_{12}^2m_1c_{12}^2m_2)c_{23}s_{23}\\ (m_2m_1)c_{12}s_{12}s_{23}& (m_3s_{12}^2m_1c_{12}^2m_2)c_{23}s_{23}& c_{23}^2m_3+s_{23}^2(s_{12}^2m_1+c_{12}^2m_2)\end{array}).$$
(5)
Here, as is seen from Eq. (4), the values $`c_{23}=1/\sqrt{2}=s_{23}`$ give maximal mixing of $`\nu _\mu `$ and $`\nu _\tau `$: $`(\nu _\mu \pm \nu _\tau )/\sqrt{2}`$, and then $`c_{12}1/\sqrt{2}s_{12}`$ — a nearly maximal mixing of $`\nu _e`$ and $`(\nu _\mu \nu _\tau )/\sqrt{2}`$: approximately $`[\nu _e\pm (\nu _\mu \nu _\tau )/\sqrt{2}]/\sqrt{2}`$.
From the familiar neutrino oscillation formulae
$$P(\nu _\alpha \nu _\beta )=|\nu _\beta |e^{iPL}|\nu _\alpha |^2=\delta _{\alpha \beta }4\underset{j>i}{}U_{\beta j}^{}U_{\alpha j}U_{\beta i}U_{\alpha i}^{}\mathrm{sin}^2x_{ji},$$
(6)
with
$$x_{ji}=1.27\frac{\mathrm{\Delta }m_{ji}^2L}{E},\mathrm{\Delta }m_{ji}^2=m_j^2m_i^2$$
(7)
($`\mathrm{\Delta }m_{ji}^2`$, $`L`$ and $`E`$ measured in eV<sup>2</sup>, km and GeV, respectively) which is valid for $`U_{\beta j}^{}U_{\alpha j}U_{\beta i}U_{\alpha i}^{}`$ real (CP violation neglected), one infers in the case of two–mixing option (2) that
$`P(\nu _e\nu _e)`$ $`=`$ $`1(2c_{12}s_{12})^2\mathrm{sin}^2x_{21},`$
$`P(\nu _\mu \nu _\mu )`$ $`=`$ $`1(2c_{12}s_{12}c_{23})^2\mathrm{sin}^2x_{21}(2c_{23}s_{23})^2(s_{12}^2\mathrm{sin}^2x_{31}+c_{12}^2\mathrm{sin}^2x_{32})`$
$``$ $`1(2c_{23}s_{23})^2\mathrm{sin}^2x_{32},`$
$`P(\nu _\mu \nu _e)`$ $`=`$ $`(2c_{12}s_{12}c_{23})^2\mathrm{sin}^2x_{21},`$ (8)
where the final step in the second formula is valid when $`\mathrm{\Delta }m_{21}^2\mathrm{\Delta }m_{31}^2\mathrm{\Delta }m_{32}^2`$ or equivalently $`m_1^2m_2^2m_3^2`$.
The first formula (8) is consistent with the observed deficit of solar $`\nu _e`$’s if one applies the vacuum global solution or large–angle MSW global solution or finally LOW global solution with $`(2c_{12}s_{12})^2\mathrm{sin}^22\theta _{\mathrm{sol}}(0.90`$ or 0.79 or 0.91) and $`\mathrm{\Delta }m_{21}^2\mathrm{\Delta }m_{\mathrm{sol}}^2(4.4\times 10^{10}`$ or $`2.7\times 10^5`$ or $`1.0\times 10^7)\mathrm{eV}^2`$, respectively. This gives $`c_{12}^20.5+(0.16`$ or 0.23 or 0.15) and $`s_{12}^20.5(0.16`$ or 0.23 or 0.15), when taking $`c_{12}^2s_{12}^2`$.
The second formula (8) describes correctly the observed deficit of atmospheric $`\nu _\mu `$’s if $`(2c_{23}s_{23})^2\mathrm{sin}^22\theta _{\mathrm{atm}}1`$ and $`\mathrm{\Delta }m_{32}^2\mathrm{\Delta }m_{\mathrm{atm}}^23.5\times 10^3\mathrm{eV}^2`$, since then $`\mathrm{\Delta }m_{21}^2\mathrm{\Delta }m_{31}^2\mathrm{\Delta }m_{32}^2`$ for $`\mathrm{\Delta }m_{21}^2`$ determined as in the case of solar $`\nu _e`$’s. This implies that $`c_{23}^20.5s_{23}^2`$ and $`m_3^23.5\times 10^3\mathrm{eV}^2`$, because $`m_1^2m_2^2m_3^2`$.
Then, the third formula (8) shows that no LSND effect for accelerator $`\nu _\mu `$’s should be observed, $`P(\nu _\mu \nu _e)0`$, since with $`\mathrm{\Delta }m_{21}^2\mathrm{\Delta }m_{\mathrm{sol}}^2(10^{10}`$ or $`10^5`$ or $`10^7)\mathrm{eV}^2\mathrm{\Delta }m^2_{\mathrm{LSND}}1\mathrm{eV}^2`$, one gets $`\mathrm{sin}^2(x_{12})_{\mathrm{LSND}}10^{19}`$ or $`10^9`$ or $`10^{14}\mathrm{sin}^2x_{\mathrm{LSND}}1`$, while $`(2c_{12}s_{12}c_{23})^2(0.90`$ or 0.79 or $`0.91)\times 0.5>\mathrm{sin}^22\theta _{\mathrm{LSND}}10^2`$. As is well known, the confirmation of LSND effect would require (beside three active neutrinos $`\nu _e,\nu _\mu ,\nu _\tau `$) the existence of at least one sterile neutrino $`\nu _s`$ (blind to all Standard Model gauge interactions) , mixing with $`\nu _e`$ through a mass generation mechanism.
In the case of Chooz experiment looking for oscillations of reactor $`\overline{\nu }_e`$’s , where it happens that $`(x_{32})_{\mathrm{Chooz}}=1.27\mathrm{\Delta }m_{32}^2L_{\mathrm{Chooz}}/E_{\mathrm{Chooz}}1`$ for $`\mathrm{\Delta }m_{32}^2\mathrm{\Delta }m_{\mathrm{atm}}^2`$, the first formula (8) leads to $`P(\overline{\nu }_e\overline{\nu }_e)1`$, since $`(x_{21})_{\mathrm{Chooz}}(x_{32})_{\mathrm{Chooz}}1`$ for $`\mathrm{\Delta }m_{21}^2\mathrm{\Delta }m_{\mathrm{sol}}^2`$ ( $`U_{e3}=0`$ in our case). This is consistent with the negative result of Chooz experiment. We can see, however, that for the actual lepton counterpart of Cabibbo–Kobayashi–Maskawamatrix the entry $`U_{e3}`$ may be a potential correction to the two–mixing option (2) ($`|U_{e3}|<0.2`$ according to the estimation in Chooz experiment).
Further on, we will put $`c_{23}=1/\sqrt{2}=s_{23}`$. Then, from Eq. (5) we infer that
$$M=(\begin{array}{ccc}c_{12}^2m_1+s_{12}^2m_2& (m_2m_1)c_{12}s_{12}/\sqrt{2}& (m_2m_1)c_{12}s_{12}/\sqrt{2}\\ (m_2m_1)c_{12}s_{12}/\sqrt{2}& (m_3+s_{12}^2m_1+c_{12}^2m_2)/2& (m_3s_{12}^2m_1c_{12}^2m_2)/2\\ (m_2m_1)c_{12}s_{12}/\sqrt{2}& (m_3s_{12}^2m_1c_{12}^2m_2)/2& (m_3+s_{12}^2m_1+c_{12}^2m_2)/2\end{array}).$$
(9)
Here, $`M_{e\mu }=M_{e\tau }`$, $`M_{\mu \mu }=M_{\tau \tau }`$ and
$`M_{ee}=c_{12}^2m_1+s_{12}^2m_2,M_{ee}`$ $`+`$ $`M_{\mu \mu }M_{\mu \tau }=m_1+m_2,M_{\mu \mu }+M_{\mu \tau }=m_3,`$
$`M_{e\mu }`$ $`=`$ $`(m_2m_1)c_{12}s_{12}/\sqrt{2}.`$ (10)
Assuming that $`M_{ee}=0`$, we get from Eq. (10) the relations $`M_{\mu \mu }=(m_3+m_2+m_1)/2`$, $`M_{\mu \tau }=(m_3m_2m_1)/2`$, $`M_{e\mu }=(s_{12}/c_{12})m_2/\sqrt{2}`$, and
$$\frac{m_1}{m_2}=\frac{s_{12}^2}{c_{12}^2},\mathrm{\Delta }m_{21}^2m_2^2m_1^2=m_2^2\frac{c_{12}^2s_{12}^2}{c_{12}^4}$$
(11)
or
$$m_1=\sqrt{\mathrm{\Delta }m_{21}^2}\frac{s_{12}^2}{\sqrt{c_{12}^2s_{12}^2}},m_2=\sqrt{\mathrm{\Delta }m_{21}^2}\frac{c_{12}^2}{\sqrt{c_{12}^2s_{12}^2}},$$
(12)
when taking $`m_1m_2`$. For instance, applying to Eq. (12) the LOW solar solution i.e., $`s_{12}^20.50.15`$, $`c_{12}^20.5+0.15`$ and $`\mathrm{\Delta }m_{21}^21.0\times 10^7\mathrm{eV}^2`$, we estimate
$$m_12.0\times 10^4\mathrm{eV},m_23.8\times 10^4\mathrm{eV},$$
(13)
while the Super–Kamiokande result $`\mathrm{\Delta }m_{32}^23.5\times 10^3\mathrm{eV}^2`$ leads to the estimation
$$m_35.9\times 10^2\mathrm{eV},$$
(14)
what shows explicitly that $`|m_1|\stackrel{<}{}m_2m_3`$. Thus, in this case
$$M_{ee}=0,M_{\mu \mu }=M_{\tau \tau }3.0\times 10^2\mathrm{eV},M_{e\mu }=M_{e\tau }1.9\times 10^4\mathrm{eV},M_{\mu \tau }3.0\times 10^2\mathrm{eV},$$
(15)
where $`M_{\mu \mu }\stackrel{>}{}M_{\mu \tau }M_{e\mu }`$.
In the general case, carrying out the diagonalization of mass matrix $`M=\left(M_{\alpha \beta }\right)`$ given in Eq. (9), we obtain
$`m_{1,2}`$ $`=`$ $`{\displaystyle \frac{M_{ee}+M_{\mu \mu }M_{\mu \tau }}{2}}\sqrt{\left({\displaystyle \frac{M_{ee}M_{\mu \mu }+M_{\mu \tau }}{2}}\right)^2+2M_{e\mu }^2}`$ (18)
$`=`$ $`\{\begin{array}{c}M_{ee}XM_{e\mu }\sqrt{2}\hfill \\ M_{\mu \mu }M_{\mu \tau }+XM_{e\mu }\sqrt{2}\hfill \end{array},`$
$`m_3`$ $`=`$ $`M_{\mu \mu }+M_{\tau \tau }`$ (19)
and
$$c_{12}=\frac{1}{\sqrt{1+X^2}},s_{12}=\frac{X}{\sqrt{1+X^2}},$$
(20)
where
$`X`$ $``$ $`{\displaystyle \frac{m_1M_{ee}}{M_{e\mu }\sqrt{2}}}={\displaystyle \frac{M_{e\mu }\sqrt{2}}{m_2M_{ee}}}`$ (21)
$`=`$ $`{\displaystyle \frac{M_{ee}M_{\mu \mu }+M_{\mu \tau }}{2M_{e\mu }\sqrt{2}}}+\sqrt{\left({\displaystyle \frac{M_{ee}M_{\mu \mu }+M_{\mu \tau }}{2M_{e\mu }\sqrt{2}}}\right)^2+1}>0.`$
Here, $`0<X<1`$ if $`M_{ee}M_{\mu \mu }+M_{\mu \tau }<0`$. For instance, for LOW solar solution
$$X=\frac{s_{12}}{c_{12}}\sqrt{0.54}=0.73,$$
(22)
showing that then $`M_{ee}M_{\mu \mu }+M_{\mu \tau }<0`$.
In conclusion, the two–mixing texture of three (Dirac or Majorana) active neutrinos $`\nu _\alpha (\alpha =e,\mu ,\tau )`$, described by the formulae (2) and (5), is neatly consistent with the observed solar and atmospheric neutrino deficits, but it predicts no LSND effect whose confirmation should imply, therefore, the existence of at least one sterile neutrino $`\nu _s`$, mixing with $`\nu _e`$. This might be either one extra, light (Dirac or Majorana) sterile neutrino $`\nu _s`$ or one of three conventional, light Majorana sterile neutrinos $`\nu _\alpha ^{(s)}=\nu _{\alpha R}+(\nu _{\alpha R})^c(\alpha =e,\mu ,\tau )`$ existing in this case beside three light Majorana active neutrinos $`\nu _\alpha ^{(a)}=\nu _{\alpha L}+(\nu _{\alpha L})^c(\alpha =e,\mu ,\tau )`$ \[of course, $`\nu _\alpha ^{(a)}=\nu _{\alpha L}`$ and $`\nu _{\alpha L}^{(s)}=(\nu _{\alpha R})^c]`$.
The essential agreement of the observed neutrino oscillations with the two–mixing option (2) for $`U`$ (provided there is really no LSND effect) suggests that the conjecture of absence of direct mixing of massive neutrinos $`\nu _1`$ and $`\nu _3`$, leading to $`U`$ of the form (2), is somehow physically important. This absence tells us that only the close neighbours in the hierarchy of massive neutrinos $`\nu _1,\nu _2,\nu _3`$ mix significantly.
Making use of Gell–Mann matrices (in the space of three generations)
$$\lambda _2=\left(\begin{array}{ccc}0& i& 0\\ i& 0& 0\\ 0& 0& 0\end{array}\right),\lambda _7=\left(\begin{array}{ccc}0& 0& 0\\ 0& 0& i\\ 0& i& 0\end{array}\right)$$
(23)
we can rewrite the two–mixing matrix (2) in the compact form
$$U=U^{(23)}U^{(12)}=e^{i\lambda _7\theta _{23}}e^{i\lambda _2\theta _{12}},$$
(24)
while the generic matrix (1) includes also the phased 13–rotation
$$U^{(13)}=(\begin{array}{ccc}c_{13}& 0& s_{13}e^{i\delta }\\ 0& 1& 0\\ s_{13}e^{i\delta }& 0& c_{13}\end{array})=(\begin{array}{ccc}1& 0& 0\\ 0& 1& 0\\ 0& 0& e^{i\delta }\end{array})e^{i\lambda _5\theta _{13}}(\begin{array}{ccc}1& 0& 0\\ 0& 1& 0\\ 0& 0& e^{i\delta }\end{array}),\lambda _5=(\begin{array}{ccc}0& 0& i\\ 0& 0& 0\\ i& 0& 0\end{array}),$$
(25)
inserted between the previous 23– and 12–rotations, $`U^{(23)}`$ and $`U^{(12)}`$, of closely neighbouring massive neutrinos, $`U=U^{(23)}U^{(13)}U^{(12)}`$ .
Then, in terms of the (truncated) annihilation and creation operators (in generation space)
$$a=\left(\begin{array}{ccc}0& 1& 0\\ 0& 0& \sqrt{2}\\ 0& 0& 0\end{array}\right),a^{}=\left(\begin{array}{ccc}0& 0& 0\\ 1& 0& 0\\ 0& \sqrt{2}& 0\end{array}\right)$$
(26)
satisfying together with the operator
$$n=a^{}a=\left(\begin{array}{ccc}0& 0& 0\\ 0& 1& 0\\ 0& 0& 2\end{array}\right)$$
(27)
the canonical annihilation and creation relations
$$[a,n]=a,[a^{},n]=a^{},$$
(28)
but obeying also the truncation condition
$$a^3=0=a^{\mathrm{\hspace{0.17em}3}},$$
(29)
we can put
$$\lambda _2=\frac{1}{2i}a(aa^{})a^{},\lambda _7=\frac{1}{i\sqrt{2}}a^{}(aa^{})a,\lambda _5=\frac{1}{i\sqrt{2}}(a^2a^{\mathrm{\hspace{0.17em}2}})$$
(30)
in exponents of the factor matrices $`U^{(12)},U^{(23)},U^{(13)}`$, respectively. Other Gell–Mann matrices, absent from $`U`$, can be put in the form
$$\lambda _1=\frac{1}{2}a(a+a^{})a^{},\lambda _6=\frac{1}{\sqrt{2}}a^{}(a+a^{})a,\lambda _4=\frac{1}{\sqrt{2}}(a^2+a^{\mathrm{\hspace{0.17em}2}}),$$
(31)
and
$$\lambda _3=\frac{1}{2}(a^2a^{\mathrm{\hspace{0.17em}2}}aa^{\mathrm{\hspace{0.17em}2}}a),\lambda _8=\frac{1}{\sqrt{3}}(aa^{}a^{}a).$$
(32)
Note that
$$[a,a^{}]=\left(\begin{array}{ccc}1& 0& 0\\ 0& 1& 0\\ 0& 0& 2\end{array}\right)$$
(33)
is not a canonical commutation relation for bosons, though the canonical annihilation and creation relations (25) hold. Note also the formulae
$$a=\frac{1}{2}(\lambda _1+i\lambda _2)+\frac{1}{\sqrt{2}}(\lambda _6+i\lambda _7),a^{}=\frac{1}{2}(\lambda _1i\lambda _2)+\frac{1}{\sqrt{2}}(\lambda _6i\lambda _7).$$
(34)
The (truncated) occupation–number operator (24), appearing naturally in our description of mixing matrix $`U`$, tells us that massive neutrinos of three generations, $`\nu _1,\nu _2,\nu _3`$, can be characterized by its three eigenvalues 0,1,2. In fact,
$$n|n_i=n_i|n_i,a|n_i=\sqrt{n_i}|n_i1,a^{}|n_i=\sqrt{n_i+1}|n_i+1,n_i|n_i=1$$
(35)
with $`|n_i=|\nu _i(n_i=0,1,2,i=1,2,3)`$. Here $`,a|0=0`$ and $`a^{}|2=0`$ i.e., $`|1=0`$ and $`|3=0`$, due to the truncation condition (26). We can see from Eqs. (27) that the matrix $`\lambda _5`$, absent from $`U`$ in the two–mixing form (2) or (21), involves linearly (in contrast to matrices $`\lambda _2`$ and $`\lambda _7`$) two–step transition operators $`a^2`$ and $`a^{\mathrm{\hspace{0.17em}2}}`$ mixing directly $`\nu _1`$ and $`\nu _3`$.
Analogical algebraic constructions work also for three generations of other fundamental fermions: charged leptons as well as up and down quarks, but the corresponding parameters $`c_{ij}`$ and $`s_{ij}`$ take different values. In the representation, where the charged–lepton and up–quark mass matrices are diagonal (and thus the corresponding diagonalizing matrices — unit), the lepton and quark mixing matrices are, at the same time, the neutrino and down–quark diagonalizing matrices (strictly speaking, in the case of quarks $`V=U^{}`$ is the conventional mixing matrix). As is well known, in contrast to neutrinos, in the down–quark case no large mixing appears experimentally: the corresponding $`s_{ij}`$ are always considerably smaller than $`1/\sqrt{2}`$ (the largest of them is $`s_{12}0.22)`$.
References
1. Z. Maki, M. Nakagawa and S. Sakata, Progr. Theor. Phys. 28, 870 (1962).
2. Cf. e.g. F. Feruglio, Acta Phys. Pol. B 31, 1221 (2000); and references therein.
3. Cf. e.g. J.N. Bahcall, P.I. Krastev and A.Y. Smirnov, Phys. Lett. B 477, 401 (2000); hep–ph/0002293.
4. Y. Fukuda et al. (Super–Kamiokande Collaboration), Phys. Rev. Lett. 81, 1562 (1998) \[E. 81, 4279 (1998)\]; 82, 1810 (1999); 82, 2430 (1999).
5. C. Athanassopoulos et al. (LSND Collaboration), Phys. Rev. Lett. 75, 2650 (1995); Phys. Rev. C 54, 2685 (1996); Phys. Rev. Lett. 77, 3082 (1996); 81, 1774 (1998).
6. Cf. e.g. W. Królikowski, Nuovo Cim. A 111, 1257 (1999); A 112, 893 (1999), also hep–ph/9904489; hep–ph/0001023; hep–ph/0004222; and references therein.
7. W. Królikowski, Acta Phys. Pol. B 31, 663 (2000); and references therein.
8. M. Appolonio et al. (Chooz Collaboration), Phys. Lett. B 420, 397 (1998); B 466, 415 (1999).
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# Oscillatory integral operators with low–order degeneracies
## Introduction
Let $`\mathrm{\Omega }_L`$, $`\mathrm{\Omega }_R`$ be open sets in $`^d`$. This paper is concerned with $`L^2`$ bounds for oscillatory integral operators $`T_\lambda `$ of the form
$$T_\lambda f(x)=e^{ı\lambda \mathrm{\Phi }(x,z)}\sigma (x,z)f(z)𝑑z$$
$`1.1`$
where $`\mathrm{\Phi }C^{\mathrm{}}(\mathrm{\Omega }_L\times \mathrm{\Omega }_R)`$ is real-valued, $`\sigma C_0^{\mathrm{}}(\mathrm{\Omega }_L\times \mathrm{\Omega }_R)`$ and $`\lambda `$ is large. We shall also write
$$T_\lambda T_\lambda [\sigma ]$$
to indicate the dependence on the symbol $`\sigma `$.
The decay in $`\lambda `$ of the $`L^2`$ operator norm of $`T_\lambda `$ is determined by the geometry of the canonical relation
$$𝒞=\{(x,\mathrm{\Phi }_x,z,\mathrm{\Phi }_z):(x,z)\mathrm{\Omega }_L\times \mathrm{\Omega }_R\}T^{}\mathrm{\Omega }_L\times T^{}\mathrm{\Omega }_R,$$
$`1.2`$
specifically by the behavior of the projections $`\pi _L:𝒞T^{}\mathrm{\Omega }_L`$ and $`\pi _R:𝒞T^{}\mathrm{\Omega }_R`$ ,
$`\pi _L:(x,z)(x,\mathrm{\Phi }_x(x,z))`$ $`1.3`$
$`\pi _R:(x,z)(z,\mathrm{\Phi }_z(x,z));`$
here $`\mathrm{\Phi }_x`$ and $`\mathrm{\Phi }_z`$ denote the partial gradients with respect to $`x`$ and $`z`$. Note that $`\text{rank }D\pi _L=\text{rank }D\pi _R`$ is equal to $`d+\text{rank }\mathrm{\Phi }_{xz}`$ and that the determinants of $`D\pi _L`$ and $`D\pi _R`$ are equal to
$$h(x,z):=det\mathrm{\Phi }_{xz}(x,z).$$
$`1.4`$
If $`𝒞`$ is locally the graph of a canonical transformation, i.e., if $`h0`$, then $`T_\lambda =O(\lambda ^{d/2})`$ (see Hörmander , ). If the projections have singularities then there is less decay in $`\lambda `$ and in various specific cases the decay has been determined. In dimension $`d=1`$ Phong and Stein obtained a complete description of the $`L^2`$ mapping properties, for the case of real-analytic phase functions. Similar results for $`C^{\mathrm{}}`$ phases (which however missed the endpoints) and related $`L^p`$ estimates for averaging operators in the plane are in . The bounds for oscillatory integral operators in one dimension, with $`C^{\mathrm{}}`$ phases, have recently been substantially improved by Rychkov , so that many endpoint estimates are now available in the $`C^{\mathrm{}}`$ category.
Such general results are not known in higher dimensions even under the assumption of $`\text{rank }\mathrm{\Phi }_{xz}d1`$. We list some known cases. If both projections $`\pi _L`$ and $`\pi _R`$ have fold $`(S_{1,0})`$ singularities then $`T_\lambda =O(\lambda ^{(d1)/21/3})`$ (, , ). If only one of the projections has fold singularities then by we have $`T_\lambda =O(\lambda ^{(d1)/21/4})`$; this is sharp if the other projection is maximally degenerate () but can be improved when that projection satisfies some finite type finite type condition (for sharp results of this sort see Comech ). This one-sided behavior comes up naturally when studying restricted X-ray transforms , , . In the authors began a study of the case of higher one-sided Morin $`(S_{1_r,0})`$ singularities, which are the stable singularities of corank one, and it was shown under suitable additional (”strongness”) assumptions that such estimates can be deduced from sharp estimates for two-sided $`S_{1_{r1},0}`$ singularities. Thus the authors were able to prove that if one projection is a Whitney cusp, i.e., of type $`S_{1,1,0}`$, then $`T_\lambda =O(\lambda ^{(d1)/21/6})`$; again this is only sharp if the other projection is maximally degenerate.
It is conjectured that if one of $`\pi _L`$ or $`\pi _R`$ has $`S_{1_r,0}`$ singularities then $`T_\lambda =O(\lambda ^{(d1)/21/(2r+2)})`$ (for the discussion of some model cases where this is satisfied and sharp see ). Here we take up the case $`r=3`$; such mappings are commonly referred to as swallowtail singularities. In order to prove this result it is crucial to get a sharp result for operators with two-sided cusp singularities.
###### Theorem
(i) Suppose that the only singularities of one of the projections ($`\pi _L`$ or $`\pi _R`$) are Whitney folds, Whitney cusps or swallowtails. Then $`T_\lambda _{L^2L^2}=O(\lambda ^{(d1)/21/8})`$ for $`\lambda 1`$.
(ii) Suppose that the only singularities of both projections $`\pi _L`$ and $`\pi _R`$ are Whitney folds or Whitney cusps. Then $`T_\lambda _{L^2L^2}=O(\lambda ^{(d1)/21/4})`$ for $`\lambda 1`$.
A slightly weaker result than (ii) was recently obtained by Comech and Cuccagna , who proved for two-sided cusp singularities the bound $`T_\lambda _{L^2L^2}C_\epsilon \lambda ^{(d1)/21/4+\epsilon }`$ with $`C_\epsilon \mathrm{}`$ as $`\epsilon 0`$.
We shall prove somewhat more general results about operators of the same “type” but with the stability assumptions weakened. To formulate the hypotheses we review the definition of kernel vector fields for a map. Fix $`n`$-dimensional manifolds $`M,N`$ and points $`P_0M`$ and $`Q_0N`$. Let $`f:MN`$ be a $`C^{\mathrm{}}`$ map with $`f(P_0)=Q_0`$. Let $`𝒰`$ be a neighborhood of $`P`$. A vector field $`V`$ is a kernel field for the map $`f`$ on $`𝒰`$ if $`V`$ is smooth on $`𝒰`$ and if $`Df_PV=det(Df_P)W_{f(P)}`$ for $`P𝒰`$; here $`W`$ is a smooth vector field on $`N`$ defined near $`Q_0=f(P_0)`$ and $`det(Df_P)`$ is calculated with respect to any local systems of coordinates.
Suppose now that $`\text{rank }Df(P_0)n1`$. Then there is a neighborhood of $`P`$ and a nonvanishing kernel vector field $`V`$ for $`f`$ on $`𝒰`$. If $`\stackrel{~}{V}`$ is another kernel field on $`U`$ then $`\stackrel{~}{V}=\alpha Vdet(Df)W`$ in some neighborhood of $`P_0`$, for some vector field $`W`$ and smooth function $`\alpha `$.
This is easy to see by an elementary calculation. Indeed we may choose coordinates $`x=(x^{},x_n)`$ on $`M`$, $`y=(y^{},y_n)`$ on $`N`$ vanishing at $`P_0`$ and $`Q_0`$, respectively, so that $`D_x^{}f=(A,b)`$ and $`D_{x_n}f=(c^t,d)`$ where $`A`$ is an invertible $`(n1)\times (n1)`$ matrix, $`b`$ and $`c`$ are vectors in $`d`$, and $`A,b,c,d`$ depend smoothly on $`x`$. Define the vector field $`V`$ by $`V=_{x_n}A^1b,_x^{}`$. Then clearly $`Df(V)=(dc^tA^1b)_{y_n}`$ and $`detDf=(dc^tA^1b)detA`$; thus $`V`$ is a kernel field. Now assume that $`\stackrel{~}{V}=\beta ^{},_x^{}+\beta _n_{x_n}`$ so that $`Df(\stackrel{~}{V})=det(Df)Z`$ with $`Z=Z^{}+\gamma _n_{y_n}`$, $`Z^{}=\gamma ^{},_y^{}`$; here $`\beta =(\beta ^{},\beta _n)`$ are smooth functions of $`x`$ and $`\gamma =(\gamma ^{},\gamma _n)`$ are smooth functions of $`y`$. Then, at any $`x`$, $`A\beta ^{}+b\beta _n=det(Df)\beta ^{}`$; therefore $`\beta ^{}=A^1b\beta _n+detDfA^1\gamma ^{}`$ and thus $`\stackrel{~}{V}=\beta _nV+(detDf)Z^{}`$ as claimed.
###### Definition Definition
Suppose that $`M`$ and $`N`$ are smooth $`n`$-dimensional manifolds and that $`f:MN`$ is a smooth map with $`dim\text{ker}(Df)1`$ on $`M`$. We say that $`f`$ is of type $`k`$ at $`P`$ if there is a nonvanishing kernel field $`V`$ near $`P`$ so that $`V^j(detDf)_P=0`$ for $`j<k`$ but $`V^k(detDf)_P0`$.
¿From the previous discussion it is clear that this definition does not depend on the choice of the nonvanishing kernel field. If one assumes that $`Df`$ drops rank simply on the singular variety $`\{detDf=0\}`$ (i.e., if $`detDf0`$) then the definition agrees with the one proposed by Comech .
###### Theorem 1.1
Suppose that both $`\pi _L`$ and $`\pi _R`$ are of type $`2`$ on $`𝒞`$. Then for $`\lambda 1`$
$$T_\lambda _{L^2L^2}=O(\lambda ^{(d1)/21/4}).$$
###### Theorem 1.2
Suppose that $`D\pi _L`$ drops rank simply on the singular variety $`\{detD\pi _L=0\}`$ and suppose that $`\pi _L`$ is of type $`3`$ on $`𝒞`$. Then for $`\lambda 1`$
$$T_\lambda _{L^2L^2}=O(\lambda ^{(d1)/21/8}).$$
Of course the analogous statement holds with $`\pi _L`$ replaced by $`\pi _R`$ in Theorem 1.2. As a corollary of both theorems we obtain the sharp endpoint estimate for two-sided cusp and one-sided swallowtail singularities stated above.
###### Remark Remark
The estimates in Theorem 1.1 and Theorem 1.2 are stable under small perturbations of $`\mathrm{\Phi }`$ and $`\sigma `$ in the $`C^{\mathrm{}}`$-topology.
The above theorems imply sharp $`L^2`$-Sobolev estimates for Fourier integral operators (see ). Let $`CT^{}\mathrm{\Omega }_L\{0\}\times T^{}\mathrm{\Omega }_R\{0\}`$ and let $`FI^\mu (\mathrm{\Omega }_L,\mathrm{\Omega }_R;C)`$ (see for the definition and for the reduction of smoothing estimates for Fourier integral operators to decay estimates for oscillatory integral operators). As a corollary of Theorems 1.1 and 1.2 one obtains
###### Theorem 1.3
(i) If both $`\pi _L`$ and $`\pi _R`$ are of type $`2`$, then $`F`$ maps $`L_{\alpha ,\text{comp}}^2`$ to $`L_{\alpha \mu 1/4,\text{loc}}^2`$.
(ii) If one projection ($`\pi _L`$ or $`\pi _R`$) is of type $`3`$ and the rank of its differential drops only simply, then $`F`$ maps $`L_{\alpha ,\text{comp}}^2`$ to $`L_{\alpha \mu 3/8,\text{loc}}^2`$.
###### Remark Remarks
1. As an example of part (i) of Theorem 1.3, consider as in \[9, §6\] a family of curves in $`^4`$ of the form
$$\gamma _x(t)=\text{exp}(tX+t^2Y+t^3Z+t^4W)(x)$$
for smooth vector fields $`X,Y,Z,W`$ on $`^4`$ such that both of the sets of vectors
$$\{X,Y,Z\frac{1}{6}[X,Y],W\frac{1}{4}[X,Z]+\frac{1}{24}[X,[X,Y]]\}$$
are linearly independent at each point $`x`$. Then the generalized Radon transform
$$Rf(x)=_{}f(\gamma _x(t))\chi (t)𝑑t,\chi C_0^{\mathrm{}}(),$$
belongs to $`I^{\frac{1}{2}}(^4,^4;C)`$ with the canonical relation $`C`$ a two–sided cusp, i.e., both $`\pi _L`$ and $`\pi _R`$ are Whitney cusps, and hence it follows from Theorem 1.3(i) that $`R:L_{\alpha ,\text{comp}}^2L_{\alpha +1/4,\text{loc}}^2`$, for all $`\alpha `$, generalizing the well-known fact for the translation–invariant family $`\gamma _x(t)=x+(t,t^2,t^3,t^4)`$.
2. Consider the translation–invariant families of curves in $`^3`$, $`\gamma _x^1(t)=x+(t,t^2,t^4)`$ and $`\gamma _x^2(t)=x+(t,t^3,t^4)`$. Then $`\{\gamma _x^1\}`$ is associated with a canonical relation, $`C^1`$, which is a two–sided cusp, while $`\{\gamma _x^2\}`$ is associated with a canonical relation, $`C^2`$, for which both projections are type 2, but not Whitney cusps. In fact, the singular variety of $`C^2`$ is not smooth: it is a union of two intersecting hypersurfaces, and $`det(D\pi _L)`$ and $`det(D\pi _R)`$ vanish of order two at the intersection and simply away from that intersection. Averaging operators associated with any (not necessarily translation–invariant) sufficiently small $`C^{\mathrm{}}`$ perturbation of either $`\{\gamma _x^1\}`$ or $`\{\gamma _x^2\}`$ will still have both projections of type 2 and hence map $`L_{\alpha ,\text{comp}}^2L_{\alpha +\frac{1}{4},\text{loc}}^2`$.
3. As an instance of part (ii) of Theorem 1.3, let $``$ be the restricted X-ray transform associated to a well–curved line complex $``$ in $`^5`$ (see \[9, §5\] for the definition). Then $`\pi _R`$ has (at most) swallowtail singularities and $``$ maps $`L_{\text{comp}}^2(^5)`$ into $`L_{1/8,\text{loc}}^2()`$. As an example consider a curve $`\alpha \gamma (\alpha )`$ in $`^4`$ with $`\gamma ^{}`$, $`\gamma ^{\prime \prime }`$, $`\gamma ^{\prime \prime \prime }`$ and $`\gamma ^{(4)}`$ being linearly independent at each $`\alpha `$ and consider the X-ray transform associated to the rigid $`5`$-dimensional line complex consisting of lines $`\{\mathrm{}_{x^{},\alpha }:x^{}^4,\alpha \}`$ in $`^5`$ where $`\mathrm{}_{x^{},\alpha }=\{(x^{}+t\gamma (\alpha ),t),t\}`$, and perturbations of this example. For the rigid case the projection $`\pi _L`$ is a blowdown in the sense of or , i.e., it exhibits a maximal degeneracy; this behavior however is not invariant under small perturbations and is not required for Theorem 1.3 to apply.
4. As an example of a restricted $`X`$-ray transform in $`^4`$ which is not well–curved in the sense of , consider the situation as in the previous example, but with $`\gamma `$ replaced by one of the curves $`\gamma ^{(1)}(\alpha )=(\alpha ,\alpha ^2,\alpha ^4)`$ or $`\gamma ^{(2)}(\alpha )=(\alpha ,\alpha ^3,\alpha ^4)`$ in $`^3`$. For both examples $`\pi _R`$ satisfies a type three condition with $`detd\pi _R`$ vanishing simply; however the singularity of $`d\pi _R`$ for the canonical relation associated to the second line complex (defined by $`\gamma ^{(2)}`$) is not of swallowtail type. Again $``$ and perturbations thereof map $`L_{\text{comp}}^2(^4)`$ into $`L_{1/8,\text{loc}}^2()`$.
5. For conormal operators in two dimension the condition of type $`k`$ for $`\pi _L`$ corresponds to a left finite type condition of order $`k+2`$ in the terminology of , and the condition of (exact) type $`k`$ corresponds to the type $`(1,k+1)`$ condition in the terminology of .
## 2. Bounds for operators with two-sided type two conditions
We decompose the operator according to the size of $`det\mathrm{\Phi }_{xz}`$, following Phong and Stein who used this decomposition to estimate operators with fold singularities. Various extensions and refinements are in , , , , , ; in fact we will use the key estimate in as the first step in our proof of Theorem 1.1. As in that work (see also , ) we shall need to localize $`V_Lh`$ and $`V_Rh`$ where $`V_L`$ and $`V_R`$ are nonvanishing kernel vector fields for $`\pi _L`$ and $`\pi _R`$, respectively. We may suppose that the support of $`\sigma `$ is small and choose coordinates $`x=(x^{},x_d)`$, $`z=(z^{},z_d)`$ in $`^{d1}\times `$ vanishing at a reference point $`P^0=(x^0,z^0)`$ so that
$$\mathrm{\Phi }_{x^{}z^{}}(P^0)=I_{d1},\mathrm{\Phi }_{x_dz^{}}(P^0)=0,\mathrm{\Phi }_{x^{}z_d}(P^0)=0.$$
Write $`\mathrm{\Phi }^{z^{}x^{}}:=\mathrm{\Phi }_{x^{}z^{}}^1`$ and $`\mathrm{\Phi }^{x^{}z^{}}:=\mathrm{\Phi }_{z^{}x^{}}^1=(\mathrm{\Phi }_{x^{}z^{}}^t)^1`$. Representatives for the kernel vector fields are then given by
$`V_R`$ $`=_{x_d}\mathrm{\Phi }_{x_dz^{}}\mathrm{\Phi }^{z^{}x^{}}_x^{}`$ $`2.1`$
$`V_L`$ $`=_{z_d}\mathrm{\Phi }_{z_dx^{}}\mathrm{\Phi }^{x^{}z^{}}_z^{}`$
(see and the discussion in the introduction).
Let $`K`$ be a fixed compact set in $`\mathrm{\Omega }_L\times \mathrm{\Omega }_R`$ which contains the support of $`\sigma `$ in its interior. Let $`A_010^{2d}`$ so that
$$\mathrm{\Phi }_{C^5(K)}10^{2d}A_0.$$
$`2.2`$
We also assume that
$$|V_L^2h|A_1^1,|V_R^2h|A_1^1.$$
$`2.3`$
for some $`A_11`$. After additional localization we may assume that $`\sigma `$ is supported on a set of small diameter $`\epsilon `$, for later use we choose
$$\epsilon =10^1\mathrm{min}\{A_0^2,A_1^2\}.$$
$`2.4`$
Let $`\beta _0C^{\mathrm{}}()`$ be an even function supported in $`(1,1)`$, and equal to one in $`(1/2,1/2)`$. Let $`\beta (s)\beta _1(s)=\beta _0(s/2)\beta _0(s)`$ and for $`j1`$ let $`\beta _j(s)=\beta _1(2^{1j}s)=\beta _0(2^js)\beta _0(2^{j+1}s)`$.
We may assume that $`\lambda `$ is large. Let $`\mathrm{}_0=[\mathrm{log}_2(\sqrt{\lambda })]`$, that is the largest integer $`\mathrm{}`$ so that $`2^{\mathrm{}}\lambda ^{1/2}`$. Let then
$`\sigma _{j,k,l}(x,z)`$ $`=\sigma (x,z)\beta (2^lh(x,z))\beta _j(2^{l/2}V_Rh(x,z))\beta _k(2^{l/2}V_Lh(x,z))`$ $`2.5`$
$`\sigma _{j,k,\mathrm{}_0}^0(x,z)`$ $`=\sigma (x,z)\beta _0(2_0^{\mathrm{}}h(x,z))\beta _j(2^{\mathrm{}_0/2}V_Rh(x,z))\beta _k(2^{\mathrm{}_0/2}V_Lh(x,z));`$
thus if $`j,k>0`$ then on the support of $`\sigma _{j,k,l}`$ we have that $`|h|2^l`$, $`|V_Lh|2^{kl/2}`$, $`|V_Rh|2^{jl/2}`$.
Our main technical result sharpens estimates given in ; we use here, as throughout, the notation $`AB`$ to denote inequalities $`ACB`$ with constants $`C`$ independent of $`\lambda ,j,k,l`$.
###### Theorem 2.1
We have the following estimates:
(i) For $`0<l<\mathrm{}_0=[\mathrm{log}_2(\sqrt{\lambda })]`$
$$T_\lambda [\sigma _{j,k,l}]_{L^2L^2}\lambda ^{(d1)/2}\mathrm{min}\{2^{l/2}\lambda ^{1/2};2^{(l+j+k)/2}\}.$$
$`2.6`$
(ii)
$$T_\lambda [\sigma _{j,k,\mathrm{}_0}^0]_{L^2L^2}\lambda ^{(d1)/21/4}2^{(j+k)/2}.$$
$`2.7`$
Given Theorem 2.1 we can deduce Theorem 1.2 by simply summing the estimates (2.6) and (2.7): The bound $`_{j,k}T_\lambda [\sigma _{j,k,\mathrm{}_0}^0]\lambda ^{(d1)/21/4}`$ is immediate. Moreover
$$\underset{0l\mathrm{log}_2(\sqrt{\lambda })}{}\underset{0j,kl/2}{}T_\lambda [\sigma _{j,k,l}]_2I+II$$
where
$`I`$ $`{\displaystyle \frac{{\displaystyle \underset{0l\mathrm{log}_2(\sqrt{\lambda })}{}}{\displaystyle }}{j,k}}`$
$`j+k\mathrm{log}_2(\lambda 2^{2l})2^{l/2}\lambda ^{d/2}`$
$`\lambda ^{(d1)/21/4}{\displaystyle \underset{0l\mathrm{log}_2(\sqrt{\lambda })}{}}(\lambda 2^{2l})^{1/4}[1+\mathrm{log}(\lambda 2^{2l})]^2\lambda ^{(d1)/21/4}`$
and
$`II`$ $`{\displaystyle \frac{{\displaystyle \underset{0l\mathrm{log}_2(\sqrt{\lambda })}{}}2^{l/2}\lambda ^{(d1)/2}{\displaystyle }}{j,k}}`$
$`j+k\mathrm{log}_2(\lambda 2^{2l})2^{(j+k)/2}`$
$`\lambda ^{(d1)/2}{\displaystyle \underset{0l\mathrm{log}_2(\sqrt{\lambda })}{}}2^{l/2}(2^l\lambda ^{1/2})(1+\mathrm{log}_2(\lambda 2^{2l}))\lambda ^{(d1)/21/4}.\mathit{}`$
We make some preliminary observations needed in the proof of Theorem 2.1. In what follows we always make the
Assumption: $`kj`$.
For $`kj`$ apply the corresponding estimates for the adjoint of $`T_\lambda [\sigma _{j,k,l}]`$. For the proof of Theorem 2.1 we may assume, by the known result for one-sided folds , that
$$2^{kl/2}2^{jl/2}\epsilon $$
$`2.8`$
where $`\epsilon `$ is as in (2.4).
### Affine changes of variables
Before starting with estimates we wish to mention the effect of changes of variables on (2.5). Set $`x=x(u)`$ and $`z=z(v)`$ and let $`\mathrm{\Psi }(u,v)=\mathrm{\Phi }(x(u),z(v))`$. Let $`h(x,z)=det\mathrm{\Phi }_{xz}`$ and $`\stackrel{~}{h}(u,v)=det\mathrm{\Psi }_{uv}(u,v)`$ then
$$\stackrel{~}{h}(u,v)=h(x(u),z(v))det\frac{Dx}{Du}det\frac{Dz}{Dv}.$$
If $`V_R=s_i(x,z)_{x_i}`$, $`V_L=t_i(x,z)_{z_i}`$ and $`\stackrel{~}{V}_R=\sigma _i(u,v)_{u_i}`$, $`\stackrel{~}{V}_L=\tau _i(u,v)_{v_i}`$, then $`V_Rg(x(u),z(v))=\stackrel{~}{V}_R(g(x(u),z(v))`$ and $`V_Lg(x(u),z(v))=\stackrel{~}{V}_L(g(x(u),z(v))`$ if and only if $`\stackrel{}{\sigma }=(\frac{Dx}{Du})^1\stackrel{}{s}`$ and $`\stackrel{}{\tau }=(\frac{Dz}{Dv})^1\stackrel{}{t}`$.
In particular if our changes of variables are affine and of the form
$$x(u)=x^0+(u^{}+a^{}u_d,u_d),z(v)=z^0+(v^{}+b^{}v_d,v_d)$$
$`2.9`$
with constant vectors $`a^{},b^{}^{d1}`$ and $`P=(x^0,z^0)`$ and if $`V_L`$ and $`V_R`$ are of the form (2.1) then we have $`\stackrel{~}{V}_R=_{u_d}(a^t+\mathrm{\Phi }_{x_dz^{}}\mathrm{\Phi }^{z^{}x^{}})_u^{}`$ and $`\stackrel{~}{V}_L=_{v_d}(b^t+\mathrm{\Phi }_{z_dx^{}}\mathrm{\Phi }^{x^{}z^{}})_v^{}`$ where all coefficient functions are evaluated at $`(x^0,z^0)+(u^{}+a^{}u_d,u_d,v^{}+b^{}v_d,v_d)`$. Thus by choosing $`a^{}=\mathrm{\Phi }^{x^{}z^{}}(P)\mathrm{\Phi }_{z^{}x_d}(P)`$, $`b^{}=\mathrm{\Phi }^{z^{}x^{}}(P)\mathrm{\Phi }_{x^{}z_d}(P)`$ we achieve that $`(\stackrel{~}{V}_R)_{0,0}=_{u_d}`$, $`(\stackrel{~}{V}_L)_{0,0}=_{v_d}`$.
### Localization
We shall perform various localizations to small boxes in $`(x,z)`$-space.
Let $`P=(x^0,z^0)\mathrm{\Omega }_L\times \mathrm{\Omega }_R_L^d\times _R^d`$ and let $`a_L^d`$ and $`b_R^d`$ be vectors with $`1a_{\mathrm{}},b_{\mathrm{}}2`$ and let $`\pi _a^{}`$, $`\pi _b^{}`$ be the orthogonal projections to the orthogonal complement of $`a`$ in $`_L^d`$ and $`b`$ in $`_R^d`$, respectively. Suppose $`0<\gamma _1\gamma _2<1`$ and $`0<\delta _1\delta _2\epsilon `$ and let
$$\begin{array}{c}B_P^{a,b}(\gamma _1,\gamma _2,\delta _1,\delta _2)=\hfill \\ \hfill \{(x,z):|\pi _a^{}(xx^0)|\gamma _1,|xx^0,a|\gamma _2,|\pi _b^{}(zz^0)|\delta _1,|zz^0,b|\delta _2.\}\end{array}$$
$`2.10`$
###### Definition Definition
We say that $`\chi C_0^{\mathrm{}}`$ is a normalized cutoff function associated to $`B_P^{a,b}(\gamma _1,\gamma _2,\delta _1,\delta _2)`$ if it is supported in $`B_P^{a,b}(\gamma _1,\gamma _2,\delta _1,\delta _2)`$ and satisfies the (natural) estimates
$$|(\pi _a^{}_x)^{m_L}a,_x^{n_L}(\pi _b^{}_z)^{m_R}b,_z^{n_R}\chi (x,z)|\gamma _1^{m_L}\gamma _2^{n_L}\delta _1^{m_R}\delta _2^{n_R}$$
whenever $`m_L+n_L10d`$, $`m_R+n_L10d`$. Here $`(\pi _a^{}_x)^{n_L}`$ stands for any differential operator $`\stackrel{}{u}_1,_x\mathrm{}\stackrel{}{u}_{n_L},_x`$ where the vectors $`u_1,\mathrm{},u_{n_L}`$ are unit vectors perpendicular to $`a`$.
We denote by $`𝒵_P^{a,b}(\gamma _1,\gamma _2,\delta _1,\delta _2)`$ the class of all normalized cutoff functions associated to $`B_P^{a,b}(\gamma _1,\gamma _2,\delta _1,\delta _2)`$.
We shall often localize to boxes of the form (2.10) and consider $`T_\lambda [\zeta \sigma _{j,k,l}]`$ where $`\zeta `$ is a cutoff function which is controlled by an absolute constant times a normalized cutoff function in the above sense.
Suppose now that $`P=(x^0,z^0)`$ and our change of variable is as in (2.9) and that $`a=(a^{},1)`$, $`b=(b^{},1)`$. Suppose that $`\zeta `$ is a normalized cutoff function associated to $`B_P^{a,b}(\gamma _1,\gamma _2,\delta _1,\delta _2)`$. Let $`\stackrel{~}{\zeta }(u,v)=\zeta (x(u),z(v))`$. Then $`\stackrel{~}{\zeta }`$ is supported in $`\stackrel{~}{B}_{(0,0)}^{e_d,e_d}(\stackrel{~}{\gamma }_1,\stackrel{~}{\gamma }_2,\stackrel{~}{\delta }_1,\stackrel{~}{\delta }_2)`$ with $`\stackrel{~}{\gamma }=(1+|a^{}|)\gamma `$, $`\stackrel{~}{\delta }=(1+|b^{}|)\delta `$ and there is a positive constant $`C`$ (independent of $`\gamma `$, $`\delta `$) so that $`C^1\stackrel{~}{\zeta }`$ is a normalized cutoff function associated to $`\stackrel{~}{B}_{(0,0)}^{e_d,e_d}(\stackrel{~}{\gamma }_1,\stackrel{~}{\gamma }_2,\stackrel{~}{\delta }_1,\stackrel{~}{\delta }_2)`$.
Changing variables as in (2.9) in the expression for the operator $`T_\lambda [\zeta \sigma _{j,k,l}]`$ yields that
$$T_\lambda [\zeta \sigma _{j,k,l}]f(z(v))=\stackrel{~}{\zeta }(u,v)\stackrel{~}{\sigma }_{j,k,l}e^{ı\lambda \mathrm{\Psi }(u,v)}𝑑v$$
with $`\stackrel{~}{\sigma }_{j,k,l}(u,v)=\sigma (x(u),z(v))\beta (2^l\stackrel{~}{h}(u,v))\beta _j(2^{l/2}\stackrel{~}{V}_R\stackrel{~}{h}(u,v))\beta _k(2^{l/2}\stackrel{~}{V}_L\stackrel{~}{h}(u,v))`$ and $`\stackrel{~}{h}(u,v)=det\mathrm{\Psi }_{uv}(u,v)`$.
### Basic estimates
We now give estimates for various pieces localized to (thin) boxes which will usually be longer in the directions of the kernel fields $`V_R`$ and $`V_L`$.
In order to formulate our results we start with a definition.
###### Definition Definition
Let $`P=(x^0,z^0)\mathrm{\Omega }_L\times \mathrm{\Omega }_R`$ and let
$$a_P=(\mathrm{\Phi }^{x^{}z^{}}(P)\mathrm{\Phi }_{z^{}x_d}(P),1),b_P=(\mathrm{\Phi }^{z^{}x^{}}(P)\mathrm{\Phi }_{x^{}z_d}(P),1)$$
$`2.11`$
Define, for fixed $`j,k,l`$,
$`𝒜_P(\gamma _1,\gamma _2,\delta _1,\delta _2):=sup\{T_\lambda [\zeta \sigma _{j,k,l}]:\zeta 𝒵_P^{a_P,b_P}(\gamma _1,\gamma _2,\delta _1,\delta _2)\}`$ $`2.12`$$`2.13`$
$`𝒜_P^0(\gamma _1,\gamma _2,\delta _1,\delta _2):=sup\{T_\lambda [\zeta \sigma _{j,k,\mathrm{}_0}^0]:\zeta 𝒵_P^{a_P,b_P}(\gamma _1,\gamma _2,\delta _1,\delta _2)\}.`$
Here $`\mathrm{}_0=[\mathrm{log}_2(\sqrt{\lambda })]`$.
The main estimate in Comech-Cuccagna applies to operators whose kernels are localized to boxes $`B_P^{a_P,b_P}(2^l,2^{jl/2},2^l,2^{kl/2})`$. This result is formulated in (2.14) of the following proposition. The constants implicit in the inequalities below, do not depend on $`j,k,l`$.
###### Proposition 2.2
(i) For $`2^l\lambda ^{1/2}`$, $`kjl/2`$,
$$\underset{P}{sup}𝒜_P(2^l,2^{jl/2},2^l,2^{kl/2})2^{l/2}\lambda ^{d/2}.$$
$`2.14`$
and
$$\underset{P}{sup}𝒜_P(2^l,2^{jl/2},2^l,2^{kl/2})2^{(l+j+k)/2}\lambda ^{(d1)/2}.$$
$`2.15`$
(ii) Let $`l=[\mathrm{log}_2(\sqrt{\lambda })]=\mathrm{}_0`$. Then for $`kj\mathrm{}_0/2`$,
$$\underset{P}{sup}𝒜_P^0(2^\mathrm{}_0,2^{j\mathrm{}_0/2},2^\mathrm{}_0,2^{k\mathrm{}_0/2})\lambda ^{(d1)/21/4}2^{(j+k)/2}.$$
$`2.16`$
Proposition 2.2 is the starting point in our proof and is extended via orthogonality arguments. The basic steps are contained in the following Propositions 2.3-2.5.
In what follows $`N`$, denotes an integer $`10d1`$ and $`l=[\mathrm{log}_2(\sqrt{\lambda })]=\mathrm{}_0`$. Then the following estimates hold uniformly in $`j,k,l`$.
###### Proposition 2.3
(i) For $`2^l\lambda ^{1/2}`$, $`kjl/2`$,
$`\underset{P}{sup}`$ $`𝒜_P(2^{j+kl},2^{kl/2},2^{j+kl},2^{kl/2})`$ $`2.17`$
$`\underset{P}{sup}𝒜_P(2^l,2^{jl/2},2^l,2^{kl/2})+2^{l(2d1)/2}2^{(j+k)/2}(2^{2l}\lambda )^{N/2}.`$
(ii) For $`kj\mathrm{}_0/2`$,
$`\underset{P}{sup}`$ $`𝒜_P^0(2^{j+k\mathrm{}_0},2^{k\mathrm{}_0/2},2^{j+k\mathrm{}_0},2^{k\mathrm{}_0/2})`$ $`2.18`$
$`\underset{P}{sup}𝒜_P^0(2^\mathrm{}_0,2^{j\mathrm{}_0/2},2^\mathrm{}_0,2^{k\mathrm{}_0/2})+\lambda ^{d/2+1/4}2^{(j+k)/2}.`$
###### Proposition 2.4
(i) For $`2^l\lambda ^{1/2}`$, $`kjl/2`$,
$`\underset{P}{sup}`$ $`𝒜_P(2^{jl/2},2^{jl/2},2^{kl/2},2^{kl/2})`$ $`2.19`$
$`\underset{P}{sup}𝒜_P(2^{j+kl},2^{kl/2},2^{j+kl},2^{kl/2})+2^{(j+k)(d1)}2^k2^{l(2d1)/2}(2^{j+k2l}\lambda )^{N/2}.`$
(ii) For $`kj\mathrm{}_0/2`$,
$`\underset{P}{sup}`$ $`𝒜_P^0(2^{j\mathrm{}_0/2},2^{j\mathrm{}_0/2},2^{k\mathrm{}_0/2},2^{k\mathrm{}_0/2})`$ $`2.20`$
$`\underset{P}{sup}𝒜_P^0(2^{j+k\mathrm{}_0},2^{k\mathrm{}_0/2},2^{j+k\mathrm{}_0},2^{k\mathrm{}_0/2})+2^{(j+k)\frac{2(d1)N}{2}+k}\lambda ^{\frac{2d1}{4}\frac{N}{2}}\mathrm{}`$
###### Proposition 2.5
(i) For $`2^l\lambda ^{1/2}`$, $`kjl/2`$,
$$T[\sigma _{j,k,l}]\underset{P}{sup}𝒜_P(2^{jl/2},2^{jl/2},2^{kl/2},2^{kl/2})+2^{(j+kl)d/2}(\lambda 2^{k3l/2})^{N/2}.$$
$`2.21`$
(ii) For $`kj\mathrm{}_0/2`$,
$$T[\sigma _{j,k,\mathrm{}_0}^0]\underset{P}{sup}𝒜_P^0(2^{j\mathrm{}_0/2},2^{j\mathrm{}_0/2},2^{k\mathrm{}_0/2},2^{k\mathrm{}_0/2})+2^{(j+k)d/2kN/2}\lambda ^{d/2N/8}.$$
$`2.22`$
Taking these estimates for granted we can give the
###### Demonstration Proof of Theorem 2.1
Observe that since $`kjl/2`$ and $`2^l\lambda ^{1/2}`$ the quantities $`2^{l(2d1)/2}2^{(j+k)/2}(2^{2l}\lambda )^{N/2}`$, $`2^{(j+k)(d1)/2+kl(2d1)/2}(2^{j+k2l}\lambda )^{N/2})`$ and $`2^{(j+kl)d/2}(\lambda 2^{k3l/2})^{N/2}`$ are all dominated by a constant times $`\lambda ^{(d1)/2}\mathrm{min}\{2^{(l+j+k)/2},2^{l/2}\lambda ^{1/2}\}`$, and a combination of the first parts of the Propositions 2.3-2.5 gives
$`T_\lambda [\sigma _{j,k,l}]`$ $`\underset{P}{sup}𝒜_P(2^l,2^{jl/2},2^l,2^{kl/2})+\lambda ^{(d1)/2}\mathrm{min}\{2^{(l+j+k)/2},2^{l/2}\lambda ^{1/2}\}.`$
We estimate the quantities $`𝒜_P(2^l,2^{jl/2},2^l,2^{kl/2})`$ by Proposition 2.2 and (2.5) follows. (2.6) is proved in the same way, using instead (2.18), (2.20) and (2.22).∎
## 3. Proofs of the Propositions
### Preliminaries
We begin by stating two elementary Lemmas which will be used several times in the proof of Propositions 2.3-5.
###### Lemma 3.1
Suppose that $`\zeta 𝒵_P^{a_P,b_P}(\gamma _1,\gamma _2,\delta _1,\delta _2)`$. Then $`\zeta =_{i=1}^Mc_i\zeta _i`$ where $`\zeta _i𝒵_{Q_i}(\epsilon \gamma _1,\epsilon \gamma _2,\epsilon \delta _1,\epsilon \delta _2)`$ with $`Q_iB_P^{a_P,b_P}(\gamma _1,\gamma _2,\delta _1,\delta _2)`$ and so that $`|c_i|,MC_\epsilon `$ (independent of the specific choice of $`\zeta `$ and $`\gamma `$, $`\delta `$, $`P`$).
###### Demonstration Proof
Immediate.∎
###### Lemma 3.2
Let $`Q=(x^Q,z^Q)B_P^{a_P,b_P}(\gamma _1,\gamma _2,\delta _1,\delta _2)`$ where $`\gamma _1\gamma _2\epsilon `$, $`\delta _1\delta _2\epsilon `$. Suppose that $`0<\stackrel{~}{\gamma }_1\stackrel{~}{\gamma }_2\epsilon ^1\gamma _2`$, $`0<\stackrel{~}{\delta }_1\stackrel{~}{\delta }_2\epsilon ^1\delta _2`$ and assume that
$$\mathrm{min}\{\frac{\stackrel{~}{\gamma }_1}{\stackrel{~}{\gamma }_2},\frac{\stackrel{~}{\delta }_1}{\stackrel{~}{\delta }_2}\}\mathrm{max}\{\gamma _2,\delta _2\}.$$
$`3.1`$
Then there are positive constants $`C`$, $`C_1`$ (independent of $`\gamma ,\stackrel{~}{\gamma },\delta ,\stackrel{~}{\delta }`$, $`P`$, $`Q`$) so that for $`\zeta 𝒵_Q^{a_P,b_P}(\stackrel{~}{\gamma }_1,\stackrel{~}{\gamma }_2,\stackrel{~}{\delta }_1,\stackrel{~}{\delta }_2)`$ the function $`C_1^1\zeta `$ belongs to $`𝒵_Q^{a_Q,b_Q}(C\stackrel{~}{\gamma }_1,C\stackrel{~}{\gamma }_2,C\stackrel{~}{\delta }_1,C\stackrel{~}{\delta }_2)`$.
###### Demonstration Proof
Observe that
$$|a_Pa_Q|+|b_Pb_Q|\mathrm{max}\{\gamma _2,\delta _2\}.$$
The relevant geometry is then that by assumption (3.1) the boxes $`B_Q^{a_P,b_P}(\stackrel{~}{\gamma }_1,\stackrel{~}{\gamma }_2,\stackrel{~}{\delta }_1,\stackrel{~}{\delta }_2)`$ and $`B_Q^{a_Q,b_Q}(\stackrel{~}{\gamma }_1,\stackrel{~}{\gamma }_2,\stackrel{~}{\delta }_1,\stackrel{~}{\delta }_2)`$ are contained in fixed dilates of each other. The asserted estimates are easy to check.∎
We shall denote by $`\eta `$ a $`C_0^{\mathrm{}}()`$ function which is supported in $`(1,1)`$ and satisfies $`_n\eta (n)1.`$ Moreover the $`C_0^{\mathrm{}}(^{d1})`$ function $`\chi `$ is defined by $`\chi (x_1,\mathrm{},x_{d1})=_{i=1}^{d1}\eta (x_i)`$. In the proofs of Propositions 3.3-5 we shall use dilates and translates of $`\eta `$ and $`\chi `$ to decompose a suitable cutoff function $`\zeta `$ as
$$\zeta =\underset{X,Z^d}{}\zeta _{XZ}$$
$`3.2`$
the definition of $`\zeta _{XZ}`$ depends on the particular geometry and is given by (3.12), (3.25) and (3.30) below in the three respective cases. We shall then employ orthogonality arguments to estimate the operator norm of $`T_\lambda [\zeta \sigma _{j,k,l}]`$ in terms of the operator norms of
$$T_{XZ}:=T_\lambda [\zeta _{XZ}\sigma _{j,k,l}].$$
This is done by using the Cotlar-Stein Lemma \[25, ch. VII, 2\]. We then have to estimate the kernels of $`T_{XW}^{}T_{\stackrel{~}{X}Z}`$ and $`T_{XZ}T_{Y\stackrel{~}{Z}}^{}`$.
The kernel of $`T_{XW}^{}T_{\stackrel{~}{X}Z}`$ is given by
$$H(w,z)H_{XW\stackrel{~}{X}Z}(w,z)=e^{ı\lambda (\mathrm{\Phi }(x,w)\mathrm{\Phi }(x,z))}\kappa _{XW\stackrel{~}{X}Z}(x,w,z)𝑑x$$
$`3.3`$
where
$$\kappa _{XW\stackrel{~}{X}Z}(x,w,z)=\zeta _{\stackrel{~}{X}Z}(x,z)\overline{\zeta _{XW}(x,w)}\sigma _{j,k,l}(x,z)\overline{\sigma _{j,k,l}(x,w)}.$$
$`3.4`$
The kernel of $`T_{XZ}T_{Y\stackrel{~}{Z}}^{}`$ is given by
$$K(x,y)K_{XZY\stackrel{~}{Z}}(x,y)=e^{ı\lambda (\mathrm{\Phi }(x,z)\mathrm{\Phi }(y,z))}\omega _{XZY\stackrel{~}{Z}}(z,x,y)𝑑z$$
$`3.5`$
with
$$\omega _{XZY\stackrel{~}{Z}}(z,x,y)=\zeta _{XZ}(x,z)\overline{\zeta _{Y\stackrel{~}{Z}}(y,z)}\sigma _{j,k,l}(x,z)\overline{\sigma _{j,k,l}(y,z)}.$$
$`3.6`$
Our localizations $`\zeta _{XZ}`$ will always have the property that the supports of $`\zeta _{XW}`$ and $`\zeta _{\stackrel{~}{X}Z}`$ are disjoint whenever $`|X_i\stackrel{~}{X}_i|3`$ for some $`i\{1,\mathrm{},d\}`$. Moreover the supports of $`\zeta _{XZ}`$ and $`\zeta _{Y\stackrel{~}{Z}}`$ are disjoint whenever $`|Z_i\stackrel{~}{Z}_i|3`$ for some $`i\{1,\mathrm{},d\}`$. This implies that
$$T_{XW}^{}T_{\stackrel{~}{X}Z}=0\text{ if }|X\stackrel{~}{X}|_{\mathrm{}}3,T_{XZ}T_{Y\stackrel{~}{Z}}^{}=0\text{ if }|Z\stackrel{~}{Z}|_{\mathrm{}}3.$$
$`3.7`$
In what follows we shall split variables $`X`$ and $`Z`$ in $`^d`$ as $`X=(X^{},X_d)`$, $`Z=(Z^{},Z_d)`$. The geometric meaning of this splitting depends on the particular situation in Propositions 3.3-5.
The main orthogonality properties will always follow from either the localization properties of the operator in terms of $`h`$, $`V_Lh`$ or $`V_Rh`$, or by an integration by parts with respect to the directions orthogonal to $`a_P`$ or $`b_P`$. To describe this we assume that $`a_P=e_d`$, $`b_P=e_d`$ at a suitable reference point, a situation which we will always be able to achieve by an affine change of variables as described in §2. If $`\mathrm{\Phi }_x^{}(x,w)\mathrm{\Phi }_x^{}(x,z)`$ for all $`x`$ with $`(x,w,z)\text{supp }\kappa _{XW\stackrel{~}{X}Z}`$ then we may integrate by parts with respect to the $`x^{}`$ variables; specifically we have
$$H(w,z)=(i/\lambda )^Ne^{ı\lambda (\mathrm{\Phi }(x,w)\mathrm{\Phi }(x,z)})^N[\kappa _{XW\stackrel{~}{X}Z}](x,w,z)dx$$
$`3.8`$
where the differential operator $``$ is defined by
$$g=\text{div}_x^{}\left(\frac{\mathrm{\Phi }_x^{}(x,z)\mathrm{\Phi }_x^{}(x,w)}{|\mathrm{\Phi }_x^{}(x,z)\mathrm{\Phi }_x^{}(x,w)|^2}g\right).$$
$`3.9`$
Similar formulas hold for the $`z^{}`$ integration by parts for the integral defining $`K(x,y)`$.
We shall give a proof of the estimates (2.17), (2.19) and (2.21), and the proof of (2.18), (2.20) and (2.22) is similar. Here we note that the lower bound on $`|h|`$ in the localization (2.5) is used in the proof of estimate (2.14); however it is not needed for the proof of Propositions 2.3-2.5.
###### Demonstration Remarks on the proof of Proposition 2.2
In order to prove (2.14) it suffices, by Lemma 3.1, to estimate $`𝒜_P(\epsilon 2^l,\epsilon 2^{jl/2},\epsilon 2^l,\epsilon 2^{kl/2})`$ for small $`\epsilon `$.
By an affine change of variable as discussed in (2.9) we may assume that $`P=(0,0)`$ and that $`\mathrm{\Phi }_{z^{}x_d}(0,0)=0`$, $`\mathrm{\Phi }_{x^{}z_d}(0,0)=0`$, thus $`\mathrm{\Phi }_{x^{}z^{}}`$ is close to the identity $`I_{d1}`$ on the support of $`\zeta `$ and the quantities $`|\mathrm{\Phi }_{z^{}x_d}(x,z)|`$ and $`|\mathrm{\Phi }_{x^{}z_d}(x,z)|`$ are bounded by $`A_0\epsilon 2^{kl/2}`$ for $`(x,z)\text{supp }\zeta `$ (recall that $`kj`$). Moreover $`a_P=e_d`$, $`b_P=e_d`$; thus $`\zeta `$ is, up to a constant, a normalized cutoff function associated to a box where $`|x^{}|,|z^{}|\epsilon 2^l`$, $`|x_d|2^{jl/2}`$, $`|z_d|2^{kl/2}`$. This puts us in the situation as in the proof of \[4, (3.6)\]. If $`𝒜_P(\epsilon 2^l,\epsilon 2^{jl/2},\epsilon 2^l,\epsilon 2^{kl/2})`$ does not vanish identically then the function $`|h(x,z)|`$ is comparable to $`2^l`$ on the box $`B_P^{a_P,b_P}(\epsilon 2^l,\epsilon 2^{jl/2},\epsilon 2^l,\epsilon 2^{kl/2})`$. Set $`S=T_\lambda [\zeta \sigma _{j,k,l}]`$. For $`jk`$ (assumed here) the kernel of $`SS^{}`$ can be estimated using integration by parts, and all the details of this argument are provided in .
The estimate (2.15) is more standard, but we sketch the argument for completeness. We may assume that $`(V_L)_P=_{z_d}`$ and $`(V_R)_P=_{x_d}`$ and then “freezing” $`x_d,z_d`$ we may write
$$Sf(x^{},x_d)=_{}S^{x_d,z_d}[f(,z_d)](x^{})𝑑z_d.$$
Each $`S^{x_d,z_d}`$ is an oscillatory integral operator of the form (1.1) in $`^{d1}`$ and the mixed Hessian of the phase function has maximal rank $`d1`$; however the amplitudes have less favorable differentiability properties. Note that each $`(x^{},z^{})`$ differentiation causes a blowup of $`O(2^l)=O(\lambda ^{1/2})`$. These estimates for the amplitudes are analogous to the differentiability properties of symbols of type $`(1/2,1/2)`$, and in this situation the classical bound remains true; one can combine Hörmander’s argument in with almost-orthogonality arguments in the proof of the Calderón-Vaillancourt theorem for pseudo-differential operators . See also for related but somewhat different arguments for Fourier integral operators associated to canonical graphs. Here it follows that the $`L^2`$ operator norm of $`S^{x_d,z_d}`$ is $`O(\lambda ^{(d1)/2})`$ uniformly in $`x_d,z_d`$. ¿From the definition of $`\sigma _{j,k,l}`$ we see that there are intervals $`I`$ and $`J`$ of length $`O(2^{jl/2})`$ and $`O(2^{kl/2})`$, respectively, so that $`S^{x_d,z_d}=0`$ unless $`x_dI`$ and $`z_dJ`$. Thus from applications of Minkowski’s and Cauchy-Schwarz’ inequalities it follows that $`S2^{(j+l/2)/2}2^{(k+l/2)/2}\lambda ^{(d1)/2}`$. (2.16) is proved in the same way.∎
###### Demonstration Proof of Proposition 2.3
Fix $`P`$. By Lemma 3.1 it suffices to estimate $`T_\lambda [\zeta \sigma _{j,k,l}]`$ where $`\zeta `$ belongs to $`𝒵_P^{a_P,b_P}(\epsilon 2^{j+kl},\epsilon 2^{kl/2},\epsilon 2^{j+kl},\epsilon 2^{kl/2})`$, with norm independent of $`P`$.
By an affine change of variable as discussed in (2.9) we may assume that $`P=(0,0)`$ and that $`\mathrm{\Phi }_{z^{}x_d}(0,0)=0`$, $`\mathrm{\Phi }_{x^{}z_d}(0,0)=0`$, hence
$$\begin{array}{c}\mathrm{\Phi }_{x^{}z^{}}I2^d\\ |\mathrm{\Phi }_{z^{}x_d}(x,z)|+|\mathrm{\Phi }_{x^{}z_d}(x,z)|A_0\epsilon 2^{kl/2}\end{array}$$
$`3.10`$$`3.11`$
for $`(x,z)\text{supp }\zeta `$. Moreover $`a_P=e_d`$, $`b_P=e_d`$; thus $`\zeta `$ is, up to a constant, a normalized cutoff function associated to a box where $`|x^{}|,|z^{}|\epsilon 2^{j+kl}`$, $`|x_d|,|z_d|2^{kl/2}`$.
For $`X,Z^d`$ let
$$\begin{array}{c}\zeta _{XZ}(x,z)=\hfill \\ \hfill \zeta (x,z)\chi (2^l\epsilon ^1x^{}X^{})\eta (2^{j+l/2}\epsilon ^1x_dX_d)\chi (2^l\epsilon ^1z^{}Z^{})\eta (2^{k+l/2}\epsilon ^1z_dZ_d)\end{array}$$
$`3.12`$
and let $`T_{XZ}=T_\lambda [\zeta \sigma _{j,k,l}]`$. By (3.10-11) and Lemma 3.2 there are positive constants $`C,C_1`$ so that $`C_1^1\zeta _{XZ}`$ belongs to $`𝒵_Q^{a_Q,b_Q}(C2^l,C2^{jl/2},C2^l,C2^{kl/2})`$. Thus
$$T_{XZ}\underset{P}{sup}𝒜_P(2^l,2^{jl/2},2^l,2^{kl/2})$$
$`3.13`$
and it remains to show almost orthogonality of the pieces $`T_{XZ}`$.
By our localization the orthogonality properties (3.7) are satisfied. Therefore the assertion (2.17) follows from
$`T_{XW}^{}`$ $`T_{\stackrel{~}{X}Z}2^{l(2d1)jk}(\lambda 2^{2l}|W^{}Z^{}|)^N`$ $`3.14`$
$`\text{if }2A_0|W^{}Z^{}||W_dZ_d|\text{ and }|W^{}Z^{}|C_1,`$
$$T_{XW}^{}T_{\stackrel{~}{X}Z}=0,\text{ if }2A_0|W^{}Z^{}|<|W_dZ_d|\text{ and }|W_dZ_d|C_1,$$
$`3.15`$
(for suitable $`C_11`$) and
$`T_{XZ}`$ $`T_{Y\stackrel{~}{Z}}^{}2^{l(2d1)jk}(\lambda 2^{2l}|X^{}Y^{}|)^N,`$ $`3.16`$
$`\text{if }2A_0|X^{}Y^{}||X_dY_d|\text{ and }|X^{}Y^{}|C_1,`$
$$T_{XZ}T_{Y\stackrel{~}{Z}}^{}=0,\text{ if }2A_0|X^{}Y^{}|<|X_dY_d|\text{ and }|X_dY_d|C_1.$$
$`3.17`$
We now show (3.15) and (3.14). The kernel $`H`$ of $`T_{XW}^{}T_{\stackrel{~}{X}Z}`$ is given by (3.3), (3.4). In order to see (3.15) pick points $`(x,w)\text{supp }\zeta _{XW}`$ and $`(x,z)\text{supp }\zeta _{\stackrel{~}{X}Z}`$ and also assume that $`(x,w)`$ and $`(x,z)`$ belong to $`\text{supp }\sigma _{j,k,l}`$ (if there are no two such points then $`T_{XW}^{}T_{\stackrel{~}{X}Z}=0`$). By definition of $`\sigma _{j,k,l}`$ we have
$$|h(x,z)h(x,w)|2^{l+2}.$$
$`3.18`$
Also for all $`(x,\stackrel{~}{z})\text{supp }\zeta `$ we have that
$$|h_{z_d}(x,\stackrel{~}{z})V_Lh(x,\stackrel{~}{z})|A_0\epsilon 2^{kl/2}$$
$`3.19`$
so that $`|h_{z_d}(x,\stackrel{~}{z})|2^{kl/22}`$. Note that $`\epsilon 2^{kl/2}(|W_dZ_d|2)|w_dz_d|2^{kl/2}(|W_dZ_d|+2)\epsilon `$ and $`|w^{}z^{}|C_0(|W^{}Z^{}|+2)2^l\epsilon `$. Therefore
$`|h(x,w)h(x,z)|`$ $`|h_{z_d}(x,z)||w_dz_d|A_0|w^{}z^{}|`$ $`3.20`$
$`2^{kl/24}\epsilon 2^{kl/2}|W_dZ_d|=\epsilon 2^{l2}|W_dZ_d|`$
if $`|W_dZ_d|\mathrm{max}\{2A_0|W^{}Z^{}|_{\mathrm{}},2A_0\}`$ and $`|W_dZ_d|2A_0`$. Observe that (3.18) and (3.20) can hold simultaneously only when $`|W_dZ_d|`$ stays bounded; this implies (3.15).
Now assume that $`|W_dZ_d|2A_0|W^{}Z^{}|_{\mathrm{}}`$, and we show (3.14) if $`|W^{}Z^{}|C_1`$ for sufficiently large $`C_1`$.
We perform integration by parts with respect to the $`x^{}`$ variables in (3.3), using (3.8/3.9). Now in view of (3.10/11) we have
$`|\mathrm{\Phi }_x^{}(x,w)\mathrm{\Phi }_x^{}(x,z)|`$ $`|\mathrm{\Phi }_{x^{}z^{}}(x,z)(w^{}z^{})|A_0\epsilon 2^{kl/2}|w_dz_d|A_0|wz|^2`$ $`3.21`$
$`\epsilon 2^l|W^{}Z^{}|`$
if $`|W^{}Z^{}|C_1`$ for suitable $`C_1`$. Moreover the $`x`$-derivatives of $`\mathrm{\Phi }(x,w)\mathrm{\Phi }(x,z)`$ are $`O(2^l|W^{}Z^{}|)`$, and differentiating the symbol causes a blowup of $`O(2^l)`$ for each differentiation. Thus for $`|W^{}Z^{}|C_1`$
$$|^N(\kappa _{XW\stackrel{~}{X}Z})|(\lambda 2^{2l}|W^{}Z^{}|)^N.$$
Taking into account the $`x`$ support this yields the estimate
$$|H(w,z)|2^{l(d1)}2^{l/2j}(\lambda 2^{2l}|W^{}Z^{}|)^N.$$
By Schur’s test we have to bound $`sup_w|H(w,z)|𝑑z`$ and $`sup_z|H(w,z)|𝑑w`$. Since the integrals are extended over sets of measure $`O(2^{l(d1)}2^{l/2k})`$ we obtain the bound (3.14).
We still have to estimate the kernel $`K`$ given by (3.5), (3.6). Note that $`|h_{x_d}(x,\stackrel{~}{z})V_Rh(x,\stackrel{~}{z})|A_0\epsilon 2^{kl/2}`$ so that $`|h_{x_d}(x,z)|2^{jl/2}`$ (recall that $`jk`$). Thus in place of (3.20) we have
$$|h(x,z)h(y,z)|2^{jl/2}|x_dy_d|A_0|x^{}y^{}|$$
$`3.22`$
and in place of (3.21) we have
$$|\mathrm{\Phi }_z^{}(x,z)\mathrm{\Phi }_z^{}(y,z)||\mathrm{\Phi }_{z^{}x^{}}(y,z)(x^{}y^{})|A_0\epsilon 2^{kl/2}|x_dy_d|A_0|xy|^2.$$
$`3.23`$
Since $`|x_dy_d||X_dY_d|`$ we proceed as before to obtain (3.16) and (3.17).∎
###### Demonstration Proof of Proposition 2.4
We continue to use the same notations as in the previous proof although our localizations are with respect to different (larger) boxes. By Lemma 3.1 it suffices to estimate the operator norm of $`T_\lambda [\zeta \sigma _{j,k,l}]`$ where now $`\zeta 𝒵_P(\epsilon 2^{jl/2},\epsilon 2^{jl/2},\epsilon 2^{kl/2},\epsilon 2^{kl/2})`$. Again we may assume that by an affine change of variable $`P=(0,0)`$ and that $`\mathrm{\Phi }_{z^{}x_d}`$, $`\mathrm{\Phi }_{x^{}z_d}`$ vanish at $`(0,0)`$. It follows that
$$|\mathrm{\Phi }_{z^{}x_d}(x,z)|+|\mathrm{\Phi }_{x^{}z_d}(x,z)|A_0\epsilon 2^{jl/2},(x,z)\text{supp }\zeta ,$$
$`3.24`$
and again $`a_P=e_d`$, $`b_P=e_d`$.
For $`X,Z^d`$ we now define
$$\begin{array}{c}\zeta _{XZ}(x,z)=\zeta (x,z)\times \hfill \\ \hfill \chi (2^{jk+l}\epsilon ^1x^{}X^{})\eta (2^{k+l/2}\epsilon ^1x_dX_d)\chi (2^{jk+l}\epsilon ^1z^{}Z^{})\eta (2^{k+l/2}\epsilon ^1z_dZ_d)\end{array}$$
$`3.25`$
and set $`T_{XZ}=T_\lambda [\zeta _{XZ}\sigma _{j,k,l}]`$. In view of (3.24), Lemma 3.1 and Lemma 3.2
$$T_{XZ}\underset{P}{sup}𝒜_P(2^{j+kl},2^{kl/2},2^{j+kl},2^{kl/2}).$$
To show the orthogonality observe that (3.7) remains valid. Moreover the width of the smaller boxes in the $`z_d`$ direction is comparable to the $`z_d`$-width of the original boxes, namely $`2^{kl/2}`$. This shows that
$$T_{XW}^{}T_{\stackrel{~}{X}Z}=0\text{ if }|W_dZ_d|C_1$$
$`3.26`$
for sufficiently large $`C_1`$.
This estimate is complemented by
$$T_{XW}^{}T_{\stackrel{~}{X}Z}2^{2(j+k)(d1)}2^{2k}2^{l(2d1)}(\lambda 2^{j+k2l}|W^{}Z^{}|)^N,\text{ if }|W^{}Z^{}|C_1,$$
$`3.27`$
for large $`C_1`$.
To see (3.27) we integrate by parts with respect to $`x^{}`$. Our kernel is still given by (3.3), (3.4). To perform the integration by parts we may assume that $`|W_dZ_d|C_1`$ by (3.26). We now see from (3.24) that
$$|\mathrm{\Phi }_x^{}(x,w)\mathrm{\Phi }_x^{}(x,z)||\mathrm{\Phi }_{x^{}z^{}}(x,z)(w^{}z^{})|A_0\epsilon 2^{jl/2}|w_dz_d|A_0|wz|^2$$
but $`|w^{}z^{}||W^{}Z^{}|\epsilon 2^{j+kl}`$, and $`|w_dz_d|2^{kl/2}\epsilon |W_dZ_d|C\epsilon 2^{kl/2}`$. Thus if $`|W^{}Z^{}|`$ is sufficiently large we have the lower bound
$$|\mathrm{\Phi }_x^{}(x,w)\mathrm{\Phi }_x^{}(x,z)|2^{j+kl}|W^{}Z^{}|$$
for $`(x,w,z)\text{supp }\kappa _{XW\stackrel{~}{X}Z}`$. Therefore analyzing $`^N(\kappa _{XW\stackrel{~}{X}Z})`$ as in the proof of Proposition 3.3 we see that
$$|^N(\kappa _{XW\stackrel{~}{X}Z})|(2^{j+k2l}|W^{}Z^{}|)^N.$$
From this we get the pointwise bound
$$|H(w,z)|2^{(j+k2l)(d1)}2^{kl/2}(\lambda 2^{j+k2l}|W^{}Z^{}|)^N.$$
For Schur’s test we have to integrate this in $`x`$ or $`y`$ over a set of measure $`2^{(j+kl)(d1)}2^{kl/2}`$ and we obtain in fact a slightly better estimate than (3.27).
Next, it remains to show that
$`T_{XZ}`$ $`T_{Y\stackrel{~}{Z}}^{}2^{2(j+k)(d1)}2^{2k}2^{l(2d1)}(\lambda 2^{j+k2l}|X^{}Y^{}|)^N,`$ $`3.28`$
$`\text{if }2A_0|X^{}Y^{}||X_dY_d|\text{ and }|X^{}Y^{}|C_1,`$
and
$$T_{XZ}T_{Y\stackrel{~}{Z}}^{}=0\text{ if }2A_0|X^{}Y^{}|<|X_dY_d|\text{ and }|X_dY_d|C_1.$$
$`3.29`$
The proof of these estimates is similar to the proof of the corresponding estimates in Proposition 2.3. The estimate (3.22) continues to hold and the estimate (3.23) is replaced by the weaker estimate
$$|\mathrm{\Phi }_z^{}(x,z)\mathrm{\Phi }_z^{}(y,z)||\mathrm{\Phi }_{z^{}x^{}}(y,z)(x^{}y^{})|A_0\epsilon 2^{jl/2}|x_dy_d|A_0|xy|^2$$
which however still gives the asserted bound since $`|x_dy_d|2^{kl/2}\epsilon |X_dY_d|`$ and $`|x^{}y^{}|2^l|X^{}Y^{}|`$. ∎
###### Demonstration Proof of Proposition 2.5
We may assume that the support of $`\zeta `$ is small (i.e. contained in a ball of radius $`\epsilon `$). By Lemma 3.1 it suffices to estimate the operator norm of $`T_\lambda [\zeta \sigma _{j,k,l}]`$ where now $`\zeta 𝒵_P(\epsilon ,\epsilon ,\epsilon ,\epsilon )`$. By affine changes of variables we may assume that $`P=(0,0)`$ and that $`\mathrm{\Phi }_{z^{}x_d}`$, $`\mathrm{\Phi }_{x^{}z_d}`$ vanish at $`(0,0)`$. Thus
$$|\mathrm{\Phi }_{z^{}x_d}(x,z)|+|\mathrm{\Phi }_{x^{}z_d}(x,z)|A_0\epsilon ,(x,z)\text{supp }\zeta .$$
For $`X,Z^d`$ we now consider $`T_{XZ}=T_\lambda [\zeta \sigma _{j,k,l}]`$ with
$$\begin{array}{c}\zeta _{XZ}(x,z)=\zeta (x,z)\chi (2^{j+l/2}\epsilon ^1x^{}X^{})\times \hfill \\ \hfill \eta (2^{j+l/2}\epsilon ^1x_dX_d)\chi (2^{k+l/2}\epsilon ^1z^{}Z^{})\eta (2^{k+l/2}\epsilon ^1z_dZ_d),\end{array}$$
$`3.30`$
and again $`T_{XZ}sup_P𝒜_P(2^{j+kl},2^{kl/2},2^{j+kl},2^{kl/2}).`$
For the orthogonality of the pieces we now use besides (3.7) the assumptions (2.3). By our choice of $`\epsilon `$ we have that
$$|V_L^2h_{z_d}V_Lh|A_0\epsilon A_1/10$$
and similarly
$$|V_R^2h_{x_d}V_Rh|A_1/10.$$
Thus
$`42^{kl/2}`$ $`|V_Lh(x,w)V_Lh(x,z)|(2A_1)^1|w_dz_d|A_0|w^{}z^{}|`$
$`42^{jl/2}`$ $`|V_Rh(x,z)V_Rh(y,z)|(2A_1)^1|x_dy_d|A_0|x^{}y^{}|.`$
This shows that
$`T_{XW}^{}T_{\stackrel{~}{X}Z}`$ $`=0\text{ if }|W_dZ_d|C`$ $`3.31`$$`3.32`$
$`T_{XZ}T_{Y\stackrel{~}{Z}}^{}`$ $`=0\text{ if }|X_dY_d|C`$
for $`C=10A_0A_1`$.
Now assume that $`|W_dZ_d|C_1`$. Then if $`(x,w,z)\text{supp }\overline{\zeta _{XW}}\zeta _{XZ}`$ we have
$$|\mathrm{\Phi }_x^{}(x,w)\mathrm{\Phi }_x^{}(x,z)||\mathrm{\Phi }_{x^{}z^{}}(x,z)(w^{}z^{})|A_0\epsilon |w_dz_d|A_0|wz|^2$$
but now $`|w^{}z^{}||W^{}Z^{}|2^{kl/2}`$, and $`|w_dz_d|2C2^{kl/2}`$. Thus for large $`|W^{}Z^{}|`$ we have the lower bound
$$|\mathrm{\Phi }_x^{}(x,w)\mathrm{\Phi }_x^{}(x,z)|2^{kl/2}|W^{}Z^{}|$$
and it follows that
$$|^N[\kappa _{XW\stackrel{~}{X}Z}](x,w,z)|(\lambda 2^{k3l/2}|W^{}Z^{}|)^N.$$
Consequently $`|H(w,z)|2^{(jl/2)d}(\lambda 2^{k3l/2}|W^{}Z^{}|)^N.`$ To apply Schur’s test we observe that for fixed $`z`$ the $`w`$ integral is extended over a set of measure $`O(2^{(kl/2)d})`$ (likewise for fixed $`w`$ the $`z`$ integral). We obtain the bound
$$T_{XW}^{}T_{\stackrel{~}{X}Z}2^{(j+kl)d}(\lambda 2^{k3l/2}|W^{}Z^{}|)^N$$
$`3.33`$
if $`|W^{}Z^{}|C^{}`$. By a similar argument
$$T_{XZ}T_{Z\stackrel{~}{Z}}^{}2^{(j+kl)d}(\lambda 2^{j3l/2}|X^{}Y^{}|)^N.$$
$`3.34`$
The asserted estimate (2.21) now follows from combining (3.31-34) and the estimate for the individual pieces.∎
## 4. One-sided type three singularities
In this section we discuss the proof of Theorem 1.2. The reasoning is very close to the one given by the authors in , but the assumptions there are somewhat different. We thus only sketch the proof and refer the reader to , for details of some of the arguments.
First we shall need an extension of Theorem 1.1 to oscillatory integral operators of the form
$$𝒯_\mu f(x)=f(y)e^{ı\mu \psi (x,y,\vartheta )}a(x,y,\vartheta )𝑑\vartheta 𝑑y$$
where the frequency variable $`\vartheta `$ lives in an open set $`\mathrm{\Theta }^N`$ and we assume that $`aC_0^{\mathrm{}}(\mathrm{\Omega }_L\times \mathrm{\Omega }_R\times \mathrm{\Theta })`$. It is assumed that $`\psi `$ is a nondegenerate phase function in the sense of Hörmander (but not necessarily homogeneous), i.e. $`_{\vartheta _i}(_{x,y,\vartheta }\mathrm{\Psi })`$, $`i=1,\mathrm{},N`$ are linearly independent. The canonical relation $`C_\psi T^{}\mathrm{\Omega }_L\times T^{}\mathrm{\Omega }_R`$ is given by
$$C_\psi =\{(x,\psi _x,y,\psi _y):\psi _\vartheta =0\}.$$
###### Lemma 4.1
Suppose that the projections $`\pi _L:C_\psi T^{}\mathrm{\Omega }_L`$, $`\pi _R:C_\psi T^{}\mathrm{\Omega }_R`$ are of type $`2`$. Then $`𝒯_\mu _{L^2L^2}=O(\mu ^{(d+N1)/21/4})`$, $`\mu \mathrm{}`$. This estimate is stable under small perturbations of $`\psi `$ and $`a`$ in the $`C^{\mathrm{}}`$-topology.
The reduction to the situation in Theorem 1.2 involves canonical transformations on $`T^{}\mathrm{\Omega }_L`$ and $`T^{}\mathrm{\Omega }_R`$ and then as in an application of the method of stationary phase to reduce the number of frequency variables (see for details).
The following Lemma deals with phase functions $`\mathrm{\Phi }(x,z)`$ without frequency variables.
###### Lemma 4.2
Let $`\mathrm{\Phi }`$ be a real–valued phase function defined near $`(x^0,z^0)`$ and assume that $`_x(det\mathrm{\Phi }_{xz}(x^0,z^0))0`$, and $`|det\mathrm{\Phi }_{xz}(x^0,z^0)|+|V_Rdet\mathrm{\Phi }_{xz}(x^0,z^0)|c|_x(det\mathrm{\Phi }_{xz}(x^0,z^0))|`$. Let $`M>0`$.
Then, if $`c`$ is sufficiently small, there are neighborhoods $`\mathrm{\Omega }_L^0`$ of $`x_0`$, $`\mathrm{\Omega }_R^0`$ of $`z_0`$, neighborhoods $`𝒰`$ and $`𝒱`$ of $`(x_0,_x\mathrm{\Phi }(x^0,z^0))`$ in $`T^{}\mathrm{\Omega }_L`$, a canonical transformation $`\chi :𝒰𝒱`$, and a unitary operator $`U_\lambda `$, so that the following statements hold if $`\sigma `$ is supported in $`\mathrm{\Omega }_L^0\times \mathrm{\Omega }_R^0`$.
(i) If $`T_\lambda `$ is the integral operator with kernel $`\sigma (x,z)e^{ı\lambda \mathrm{\Phi }(x,z)}`$ then
$$U_\lambda T_\lambda =S_\lambda +R_\lambda $$
where $`S_\lambda `$ is an integral operator with kernel $`\tau (x,z)e^{ı\lambda \mathrm{\Psi }(x,z)}`$ and $`R_\lambda _{L^2L^2}=O(\lambda ^M)`$,
(ii) If $`𝒞_\mathrm{\Phi }=\{(x,\mathrm{\Phi }_x,z,\mathrm{\Phi }_z),(x,z)\text{supp }\sigma \}`$ then for $`𝒞_\mathrm{\Psi }=\{(x,\mathrm{\Psi }_x,z,\mathrm{\Psi }_z),(x,z)\text{supp }\tau \}`$ we have
$$𝒞_\mathrm{\Psi }\{(\chi (x,\xi ),z,\zeta ):(x,\xi ,z,\zeta )𝒞_\mathrm{\Phi }\}.$$
(iii) $`_z(det\mathrm{\Psi }_{xz})0`$ for $`(x,z)\text{supp }\tau `$.
###### Demonstration Proof
This can be extracted from the arguments in §4 of .
###### Demonstration Proof of Theorem 1.2
We work with $`T_\lambda `$ as in (1.1) where $`(x,z)`$ is close to the origin, and the origin lies on the singular surface $`\{(x,z):det\mathrm{\Phi }_{xz}=0\}`$. We may assume, after a change of variable in $`z`$ that
$$\mathrm{\Phi }(0,z)=0,\mathrm{\Phi }_{x^{}z^{}}(0,z)=I,\mathrm{\Phi }_{x^{}z_d}(0,z)=0$$
$`4.1`$
(cf. the proof of Lemma 2.7 in ) and by a change of variable in $`x`$ we may also assume
$$\mathrm{\Phi }_{z^{}x_d}(0,0)=0.$$
We assume that $`\pi _L`$ is of type $`2`$ and that $`_{x,z}(det\mathrm{\Phi }_{xz})0`$ where $`det\mathrm{\Phi }_{xz}`$ vanishes.
If $`\mathrm{\Phi }_{x_dx_dz_d}0`$ or $`\mathrm{\Phi }_{x_dz_dz_d}0`$ or $`\mathrm{\Phi }_{x_dz_dz_dz_d}0`$ then we have a fold or cusp singularity and better results then the one claimed in Theorem 1.2 were proved in , . Therefore assume that $`\mathrm{\Phi }_{x_dx_dz_d}`$, $`\mathrm{\Phi }_{x_dz_dz_d}`$ and $`\mathrm{\Phi }_{x_dz_dz_dz_d}`$ are small. By (4.1) and Lemma 4.2 we may assume the more restrictive assumption that $`_z(det\mathrm{\Phi }_{xz})0`$ which near the origin is equivalent with $`_z(\mathrm{\Phi }_{x_dz_d})0`$, again by (4.1). After a rotation we may assume
$$\mathrm{\Phi }_{x_dz_dz_1}(0,0)0.$$
$`4.2`$
We now consider the operator $`T_\lambda T_\lambda ^{}`$ and estimate it by the slicing technique in (also familiar from the proof of Strichartz estimates). Now $`T_\lambda T_\lambda ^{}f(x)=𝒦^{x_d,y_d}[f(,y_d)](x^{})𝑑y_d`$ where the kernel of $`𝒦^{x_d,y_d}`$ as an integral operator acting on functions in $`^{d1}`$ is given by
$$K^{x_d,y_d}(x^{},y^{})=e^{ı\lambda [\mathrm{\Phi }(x^{},x_d,z)\mathrm{\Phi }(y^{},y_d,z)]}\sigma (x,z)\overline{\sigma (y,z)}𝑑z.$$
The computation in shows that after rescaling the estimation is reduced to showing that two integral operators $`_\mu ^\pm _{\mu ,\gamma ,c}^\pm `$ with kernel
$$H_\mu ^\pm (u,v)=e^{ı\mathrm{\Psi }^\pm (u,v,z;\gamma ,c)}b_{\gamma ,c}(u,v,z)𝑑z$$
are bounded on $`L^2(^{d1})`$ with norm $`O(\mu ^{(d1)1/4})`$, $`\mu 1`$. Here $`b_{\gamma ,c}`$ is $`C_0^{\mathrm{}}`$ and
$$\mathrm{\Psi }^\pm (u,v,z;\gamma ,c)=uv,\mathrm{\Phi }_x^{}(0,c,z)\pm \mathrm{\Phi }_{x_d}^{}(0,c,z)+r_\pm (u,v,z,\gamma ,c)$$
where $`\gamma `$ and $`c`$ are small parameters and $`r_\pm (u,v,z,0,c)0`$. The dependence of $`\mathrm{\Psi }^\pm `$, $`r_\pm `$ and $`b`$ on $`\gamma `$ and $`c`$ is smooth and the bounds have to be uniform for small $`\gamma ,c`$. For this it remains to show that the operators $`_\mu ^\pm `$ are oscillatory integral operators with two-sided type two singularities to which we can apply Lemma 4.1 in $`d1`$ dimensions (with $`d`$ frequency variables). It suffices to check the type two condition at $`\gamma =0`$, $`c=0`$.
The condition (4.2) guarantees that $`\mathrm{\Psi }^\pm `$ is indeed a nondegenerate phase function with critical set
$`\text{Crit}_{\mathrm{\Psi }^\pm }`$ $`=\{(u,v,z):_z\mathrm{\Psi }^\pm =0\}`$
$`=\{(u,v,z):v=u\mathrm{\Phi }^{x^{}z^{}}\mathrm{\Phi }_{z^{}x_d},\mathrm{\Phi }_{x_dz_d}\mathrm{\Phi }_{x_dz^{}}\mathrm{\Phi }^{z^{}x^{}}\mathrm{\Phi }_{x^{}z_d}=0\}`$
where the $`\mathrm{\Phi }`$ derivatives are evaluated at $`(0,z)`$. At $`x=0`$ the second equation becomes $`\mathrm{\Phi }_{x_dz_d}(0,z)=0`$ and by (4.2) we may solve this equation expressing $`z_1`$ as a function $`z_1^\pm `$ of $`\stackrel{~}{z}=(z^{\prime \prime },z_d)`$ with $`z^{\prime \prime }=(z_2,\mathrm{},z_{d1})`$. Set $`G^\pm (\stackrel{~}{z})=\mathrm{\Phi }_x^{}(0,z_1^\pm (\stackrel{~}{z}),\stackrel{~}{z})`$ and $`B^\pm (\stackrel{~}{z})=\mathrm{\Phi }^{x^{}z^{}}(0,z_1^\pm (\stackrel{~}{z}),\stackrel{~}{z})\mathrm{\Phi }_{z^{}x_d}(0,z_1^\pm (\stackrel{~}{z}),\stackrel{~}{z})`$ then the canonical relation for vanishing $`\gamma ,c`$ is given by
$$C_{\mathrm{\Psi }^\pm }=\{(u,G^\pm (\stackrel{~}{z}),uB^\pm (\stackrel{~}{z}),G^\pm (\stackrel{~}{z}))\}$$
which is parametrized by the coordinates $`(u,\stackrel{~}{z})`$. The derivative of the projection to $`T^{}\mathrm{\Omega }_L`$ in these coordinates is given by
$$DG^\pm =\left(\begin{array}{c}\mathrm{\Phi }_{x^{}z_1}\frac{z_1}{\stackrel{~}{z}}+\mathrm{\Phi }_{x^{}\stackrel{~}{z}}\end{array}\right)$$
and by (4.1) we see that its determinant equals $`(1)^{d1}z_1^\pm /z_d`$ and $`\stackrel{~}{V}_L=/_{z_d}`$ is a kernel vector field for the left projection. Moreover $`\stackrel{~}{V}_R=\stackrel{~}{V}_L+_{i=1}^{d1}c_i(z)/_{u_i}`$ so that $`\stackrel{~}{V}_L`$ and $`\stackrel{~}{V}_R`$ coincide when acting on the determinant. By implicit differentiation we see that $`_{z_d}^kz_1^\pm \mathrm{\Phi }_{z_dx_dz_1}^1\mathrm{\Phi }_{x_dz_d^{k+1}}`$ belongs to the ideal generated by $`\mathrm{\Phi }_{x_dz_d^j}`$, $`jk`$. ¿From this one deduces that $`\stackrel{~}{V}_L`$ and $`\stackrel{~}{V}_R`$ are of type $`k1`$ if one of the derivatives $`\mathrm{\Phi }_{x_dz_d^j}`$ , $`jk+1`$, does not vanish. We apply this for $`k=3`$ to conclude the proof.∎
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# Untitled Document
hep-th/0007255 RUNHETC-2000-07
Noncommutative Scalar Solitons at Finite $`\theta `$
Chen-Gang Zhou
Department of Physics and Astronomy
Rutgers University
Piscataway, NJ 08855–0849
czhou@physics.rutgers.edu.
We investigate the behavior of the noncommutative scalar soliton solutions at finite noncommutative scale $`\theta `$. A detailed analysis of the equation of the motion indicates that fewer and fewer soliton solutions exist as $`\theta `$ is decreased and thus the solitonic sector of the theory exhibits an overall hierarchy structure. If the potential is bounded below, there is a finite $`\theta _c`$ below which all the solitons cease to exist even though the noncommutativity is still present. If the potential is not bounded below, for any nonzero $`\theta `$ there is always a soliton solution, which becomes singular only at $`\theta =0`$. The $`\varphi ^4`$ potential is studied in detail and it is found the critical $`(\theta m^2)_c=13.92`$ ($`m^2`$ is the coefficient of the quadratic term in the potential) is universal for all the symmetric $`\varphi ^4`$ potential.
July 2000
1. Introduction
Recently it was found that noncommutative geometry arises naturally in string theory with a constant $`B`$ field background \[1,,2\]. In particular, in the large $`B`$ limit, the string field algebra factors into a direct product with the noncommutative algebra being an independent subalgebra\[3\], resulting in a noncommutative field theory as a decoupled low energy effective description of string theory in this limit. It is highly nonlocal and contains infinitely high order derivative terms, but in a controlled and self-consistent way. Compared to the commutative case, the renormalizability is improved but remains an open question, while nontrivial behavior such as UV/IR mixing adds to the difficulties \[4,,5,,6,,7,,8,,9\]. The phase structure of the noncommutative scalar field theory is analyzed in the $`\varphi ^4`$ case \[10\] in which an unusual phase structure is uncovered.
The soliton sector of the noncommutative scalar field theory also exhibits an intriguingly rich structure. In the limit of large noncommutative parameter $`\theta `$ and ignoring the kinetic term, the solitons can be explicitly constructed via an isomorphism between the noncommutative fields and the operators on a single particle Hilbert space \[11\], and it was found that there are infinitely many solutions as long as the potential has more than one extremum. In spacetime, the soliton interpolates between the pseudo-vacuum at the core and the true vacuum at spatial infinity. This construction has been applied in the description of the unstable D-branes and of tachyonic condensation \[12,,13\], based on the idea that D-branes can be constructed as solitons or lumps in the open string field theory \[14,,15,,16,,17,,18,,19,,20\]. These solitons are expected to disappear in the commutative limit $`\theta =0`$ because Derrick’s theorem states that there is no soliton solution in more than 1+1 dimensions. Precisely how this happens, as $`\theta `$ changes from infinity to zero, will be investigated in this paper.
The general picture we found is as follows. First, all the solutions at $`\theta =\mathrm{}`$ can be extrapolated into the finite $`\theta `$ region, when the contribution from the kinetic term is taken into account. As we will show later, the noncommutative parameter $`\theta `$ becomes an overall multiplicative factor of the scalar potential, after a simple scaling. As $`\theta `$ decreases, the potential seen by the soliton scales down, which finally makes the soliton solution impossible. Each soliton at $`\theta =\mathrm{}`$ has its own critical point in $`\theta `$, and when $`\theta `$ gets smaller this particular solution disappears. In general different solutions at $`\theta =\mathrm{}`$ have different critical points for their own existence, so there is a hierarchy structure controlled by $`\theta `$, and we will see fewer and fewer soliton solutions as $`\theta `$ is decreased.
Depending on whether the potential is bounded below we have qualitatively different results. If the polynomial potential is bounded below and assumed to have its global minimum at the origin, there exists a lowest critical point $`\theta _c`$ which is nonzero. The bound for this critical point can be estimated for each particular potential, in some cases it can even be found precisely, such as in $`\varphi ^4`$ theory, which will be studied in detail in this paper. If instead the potential is not bounded below, such as for the cubic potential, there always exists a certain soliton solution at any nonzero $`\theta `$ and it becomes singular only in the commutative limit $`\theta =0`$. Thus the noncommutative soliton owes its existence not only to the noncommutativity, but also to the dynamics.
The paper is organized as follows. In section 2 we review the construction of the noncommutative scalar soliton at $`\theta =\mathrm{}`$. In particular the isomorphism between the noncommutative algebra and the operator algebra on a single particle Hilbert space is discussed in detail. Ignoring the kinetic energy, the soliton in this limit is easily obtained by solving an algebraic equation. In section 3, we consider the solitons at finite $`\theta `$ for a general polynomial potential, and provide a proof of our result through a qualitative study of the full equation of motion expanded in the projection operator basis. In section 4, we provide a detailed numerical study of the soliton solution at finite $`\theta `$, for the $`\varphi ^4`$ theory. We use a more natural dimensionless parameter $`\theta m^2`$ instead of $`\theta `$, where $`m^2=V^{}(0)`$. The approximation method in \[11\] gives a good estimation of $`(\theta m^2)_c`$. Using a simple scaling argument, we find that the dimensionless critical parameter $`(\theta m^2)_c`$ is the same for all $`\varphi ^4`$ potential with two degenerate vacua. In section 5, we discuss string theoretic implications of our results and give our conclusions.
2. Review of Noncommutative Scalar Solitons
The basis for noncommutative geometry lies in the deformation of the usual commutative product of the smooth functions on a flat space $`R^n`$, into the noncommutative star product. The simplest realization of this deformation is Weyl quantization, when a Poisson bracket structure over the space of smooth functions $`C^{\mathrm{}}(R^n)`$ is determined by a constant bivector field $`\theta ^{ij}\frac{}{x^i}\frac{}{x^j}`$, and the product assumes the form
$$f(x)g(x)=e^{i\theta ^{ij}\frac{}{a^i}\frac{}{b^j}}f(x+a)g(x+b)|_{a=b=0}.$$
Actually this is the familiar quantization procedure in quantum mechanics, as the Poisson bracket is lifted up to the canonical commutation relations, and functions on the phase space become operators on the quantum Hilbert space. So after the deformation quantization, the noncommutative star algebra $`(C^{\mathrm{}}(R^n),)`$ is naturally isomorphic to the operator algebra over a quantum Hilbert space corresponding to a finite number of particles. Under the Weyl prescription, the algebraic isomorphism is
$$f(x)=\frac{1}{(2\pi )^n}d^{2n}k\stackrel{~}{f}(k)e^{i(kx)},$$
$$\widehat{f}(\widehat{x})=\frac{1}{(2\pi )^n}d^{2n}k\stackrel{~}{f}(k)e^{i(k\widehat{x})}.$$
The star product is mapped to the operator product, and the integration of the function over the phase space is equal to the trace of the corresponding operator in the Hilbert space $``$
$$fg\widehat{f}\widehat{g},$$
$$\frac{1}{(2\pi )^n}d^{2n}xf(x)=Tr_{}\widehat{f}.$$
As Hilbert space is separable and naturally equipped with a positive definite inner product, it always has a complete set of orthonormal basis of vectors $`\{|n,n=0,1,2,\mathrm{}\}`$. Then the bounded linear operators on the Hilbert space have a corresponding basis composed of operators $`\{|mn|,m,n=0,1,2,\mathrm{}\}`$. Conversely, the above isomorphism, (2.1) and (2.1), allows us to find the smooth function corresponding to each operator.
We will give the detailed construction in the case that the space is two dimensional, and the generaliation to the higher dimensional case is easy. $``$ now is a single particle Hilbert space. Using (2.1), the complex variable $`\overline{z}`$ and $`z`$ are mapped to the creation operator $`a^{}`$ and the annihilation operator $`a`$ respectively. We choose the simple harmonic oscillator basis, which are eigenstates of the operator $`a^{}a\frac{r^2}{2}`$
$$|n=\frac{(a^{})^n}{\sqrt{n!}}|0,a|n=\sqrt{n}|n1,a^{}|n=\sqrt{n+1}|n+1.$$
A real function $`\varphi ^{}(x)=\varphi (x)`$ corresponds to an hermitian operator, and so can be diagonalized using a unitary operator $`U`$ ($`U^{}U=UU^{}=I`$)
$$\widehat{\varphi }=U(\underset{n=0}{\overset{\mathrm{}}{}}\lambda _n|nn|)U^{},\lambda _nR.$$
The projection operator is easily expressed in the normal ordered form, $`|nn|=\frac{1}{n!}:(a^{})^ne^{a^{}a}a^n:`$ which is proportional to the n-th Laguerre polynomial $`L_n`$ in momentum space
$$\varphi _N^{(n)}(k)=e^{\frac{k^2}{2}}L_n(\frac{k^2}{2}).$$
Actually Weyl ordering and normal ordering are two equivalent isomorphisms from the star algebra to the operators on Hilbert space (Weyl quantization uses the symmetric ordering), and in particular they differ by an integration kernel in momentum space
$$\stackrel{~}{f}_W(k)=\stackrel{~}{f}_N(k)e^{\frac{k^2}{4}}.$$
Now, using the relation (2.1)(2.1) and (2.1), we can easily find that the projection operator corresponds to a radially symmetric function
$$\begin{array}{cc}\hfill \varphi _n(r^2)& =\frac{1}{2\pi }d^2ke^{\frac{k^2}{4}}L_n(k^2/2)e^{ik.x}\hfill \\ & =2(1)^ne^{r^2}L_n(2r^2)\hfill \end{array}.$$
The functions $`\varphi _n(x)`$ have the same properties under the star product as the corresponding projection operators. This will greatly facilitate the construction of the noncommutative solitons.
Consider noncommutative scalar field theory in 2+1 dimensions with noncommutativity only in the spatial dimensions. The soliton is the classical extremum of the energy functional
$$E[\varphi ]=d^2x[_\mu \varphi ^\mu \varphi +V(\varphi )].$$
Upon changing variables to $`(z,\overline{z})`$, and performing a rescaling $`z\frac{z}{\sqrt{\theta }}`$, the star product will be independent of $`\theta `$, and the sole effect of the noncommutative parameter $`\theta `$ will appear as an overall scale factor for the potential
$$E=d^2z(\frac{1}{2}\varphi \overline{}\varphi +\theta V(\varphi )).$$
In the operator representation, the energy functional becomes
$$E(\varphi )=K(\varphi )+U(\varphi ),K(\varphi )=Tr[a,\varphi ][\varphi ,a^+],U(\varphi )=\theta TrV(\varphi ).$$
At $`\theta =\mathrm{}`$, the kinetic energy is much smaller than the potential energy and so can be ignored. The potential $`U(\widehat{\varphi })`$ has a $`U(\mathrm{})`$ symmetry, and the scalar field can be diagonalized as in (2.1). In the operator representation, the potential is a function of the coefficient series $`\{\lambda _n,n=0,1,\mathrm{}\}`$, where the $`\lambda _n`$’s are decoupled from each other
$$E(\varphi )=\theta V(\varphi )=\theta \underset{n=0}{\overset{+\mathrm{}}{}}V(\lambda _n).$$
The classical equation of motion $`E^{}(\varphi )=0`$ is a set of independent algebraic equations, $`V^{}(\lambda _n)=0`$. So the solution at $`\theta =\mathrm{}`$ is a sequence of components $`(\lambda _0,\lambda _1,\lambda _2,\mathrm{})`$ in the projection operator basis, with each $`\lambda _n`$ being an extremum of the potential, and $`\lambda _n=0`$ as $`n\mathrm{}`$, which is the finite energy requirement for the soliton solution. Here we assume the potential has zero vacuum energy at the origin.
Therefore there are infinitely many soliton solutions at $`\theta =\mathrm{}`$, as long as the potential has more than one extremum. Each solution can be regarded as a map from the positive integer to the extrema of the potential. A general solution will spontaneously break the $`U(\mathrm{})`$ group down to a finite unitary subgroup, depending on how many $`\lambda _n`$ are the same in that particular solution. This has been interpreted as describing the decay of the unstable D-brane into multiple lower dimensional D-branes in the string theory\[13\].
3. Noncommutative solitons at finite $`\theta m^2`$: general analysis
In this section we will discuss the noncommutative soliton solutions at finite $`\theta `$ in 2+1 dimensions. We assume the potential has a true vacuum at the origin with value zero. This can always be satisfied by a constant shift of the scalar field, if the highest power term is even with positive coefficient. So for example, $`\varphi ^3`$ potential does not satisfy this condition while $`\varphi ^4`$ potential does. This trivial looking assumption turns out to be essential in understanding the existence of the critical point for the solitons. We will comment on the case when the potential has odd highest power at the end of the section.
The kinetic energy has to be taken into account. It breaks the $`U(\mathrm{})`$ symmetry, so the energy functional contains the unitary matrix $`U`$. Let $`U_{mn}=m|U|n`$ be the matrix element of $`U`$ in SHO basis, the energy functional is
$$E(\{\lambda _n\},\{U_{mn}\})=\underset{n=0}{\overset{\mathrm{}}{}}\lambda _n^2(1+2\underset{m=0}{\overset{\mathrm{}}{}}m|U_{mn}|^2)2\underset{m,n=0}{\overset{\mathrm{}}{}}\lambda _m\lambda _n|A_{mn}|^2+\theta \underset{n=0}{\overset{\mathrm{}}{}}V(\lambda _n),$$
where
$$A_{mn}=\underset{k=1}{\overset{\mathrm{}}{}}\sqrt{k}U_{kn}U_{k1,m}^{}.$$
Only the radially symmetric soliton solutions will be considered, which means the scalar field in the operator representation is diagonalized, $`U=I`$. Actually adding a noncommutative $`U(1)`$ gauge field into the action can restore the $`U(\mathrm{})`$ symmetry, while the scalar field lies in the adjoint representation of this $`U(\mathrm{})`$ group. Then by a proper $`U(\mathrm{})`$ transformation, the radially symmetric form of the scalar field can always be assumed.
Under such an assumption, the energy functional simplifies greatly,
$$E(\{\lambda _n\})=\underset{n=0}{\overset{\mathrm{}}{}}[(2n+1)\lambda _n^22(n+1)\lambda _{n+1}\lambda _n+\theta V(\lambda _n)].$$
The classical equation of motion $`E/\lambda _n=0`$ becomes a set of infinite number of coupled equations
$$\begin{array}{cc}\hfill (n+1)\lambda _{n+1}(2n+1)\lambda _n+n\lambda _{n1}& =\frac{1}{2}\theta V^{}(\lambda _n),n1\hfill \\ \hfill \lambda _1\lambda _0& =\frac{1}{2}\theta V^{}(\lambda _0).\hfill \end{array}$$
In addition we impose the asymptotic boundary condition required by the finiteness of the total energy
$$\lambda _n0,\mathrm{as}n\mathrm{}$$
This is a second order difference equation for which it is hard to find a closed form solution. In general, the difference equation allows more solutions than its corresponding differential equation. In this section we first try a qualitative analysis to find the effect of $`\theta `$ on the solution. In the next section, we will use both numerical and analytical methods to analyze the $`\varphi ^4`$ potential in detail. m Add the set of equations (3.1) up to the N-th to get an equation that is “integrated” once
$$\lambda _{N+1}\lambda _N=\frac{\theta }{2}\frac{1}{N+1}\underset{n=0}{\overset{N}{}}V^{}(\lambda _n).$$
Take the $`N\mathrm{}`$ limit and use asymptotic condition (3.1), we obtain the necessary condition
$$\underset{n=0}{\overset{\mathrm{}}{}}V^{}(\lambda _n)<\mathrm{}.$$
If we regard $`n`$ as the discrete time, the above equation describes a non-autonomous dynamic system. The noncommutative soliton solution is like a particle starting from a large nonzero $`\lambda _0`$ and approaches zero as time $`n`$ goes to infinity. In general it is possible for this particle to go back and forth, but it should ultimately approach zero monotonically as dictated by the asymptotic condition. We will only study this part below, and previous “motion” only shows up as an initial condition of $`\frac{1}{N}_{n=0}^{N1}V^{}(\lambda _n)`$ which is bounded for $`\lambda >0`$.
First let us observe whether the particle approaches zero from below or above. Assume this imaginary particle starts from positive $`\lambda `$ and decreases. If at some time $`N`$, it jumps to $`\lambda _N<0`$ close to zero, then $`_{n=0}^{N1}V^{}(\lambda _n)<0`$. As the origin is a true vacuum, $`V^{}`$ is a line with positive slope near the origin so $`V^{}(\lambda _N)<0`$. Then necessarily $`\lambda _{N+1}<\lambda _N<0`$ from (3.1) and $`\lambda _n`$ would not converge to zero. Similarly, if the particle starts from the negative point, it should approach the origin from below. So we need only consider $`\lambda _n`$ positive only. Then for large $`n`$, $`\lambda _n`$ approach zero from above monotonically, which by (3.1)gives a sharper constraint than (3.1)
$$\underset{n=0}{\overset{\mathrm{}}{}}V^{}(\lambda _n)0.$$
Fig. 1: The general $`\varphi ^4`$ potential $`V(\varphi )`$ and its derivative $`V^{}(\varphi )`$. Notice the absolute value of the minima (the valley) of $`V^{}(\varphi )`$ for $`\varphi >0`$ is smaller than the maxima of $`V^{}(\varphi )`$ (the hump) in this case.
To see what happens as $`\theta `$ decreases, let us take a general $`\varphi ^4`$ potential as an example, as shown in fig. 1, and study the extrapolation of the solution corresponding to the particle’s starting from a nonzero $`\lambda _0>0`$ and jumping to $`\lambda _1<<1`$ and then approach zero afterwards. This soliton solution has the lowest $`\theta _c`$, as will be proved in the next section. (3.1) requires $`V^{}(\lambda _0)<0`$ and $`V^{}(\lambda _1)+V^{}(\lambda _0)<0`$. But the minimum of $`V^{}`$ is bounded for $`\lambda >0`$, so no matter how $`\lambda _0`$ changes, $`\lambda _1=\lambda _0+\theta V^{}(\lambda _0)/2`$ would finally increase with $`\theta `$ decreasing when $`\theta `$ becomes small enough, and the particle climbs up the hump at $`\lambda _1`$. In other words, $`\theta `$ controls how far the imaginary particle can go at each step. In this $`\varphi ^4`$ potential, it can be easily proved that the requirement $`V(x)V(0)`$ always makes the hump larger than the absolute value of the valley. So as $`\theta `$ becomes small enough, it’s impossible for $`\lambda _1`$ to be close enough to the origin to satisfy the necessary condition $`V^{}(\lambda _1)+V^{}(\lambda _0)<0`$. Then $`\lambda _2>\lambda _1`$ and $`\lambda _n`$ wouldn’t converge to zero. Thus this solution can’t exist for such a $`\theta `$.
It seems that the above argument depends on the special property of the $`\varphi ^4`$ potential and presumably would not hold for a more general polynomial potential. But there is a more general argument to establish the existence of the finite critical point, although this may not reflect the actual situation of how the solution disappears. Notice that there is a bound on $`\mathrm{\Delta }\lambda _N\lambda _N\lambda _{N1}`$
$$|\mathrm{\Delta }\lambda _N|=|\frac{\theta }{2}\frac{1}{N}\underset{n=0}{\overset{N1}{}}V^{}(\lambda _n)|<\frac{\theta }{2}|\mathrm{inf}_{x>0}V^{}(x)|.$$
When $`\theta `$ becomes very small, $`\mathrm{\Delta }\lambda `$ becomes very small, and the sum can be approximated by an integral from $`\lambda _0`$ to $`\lambda _N`$. This is exactly the same as using the limit process to find Riemann integral of a function. But this integral $`V(\lambda _0)=_0^{\lambda _0}V^{}(x)𝑑x>0`$ because zero is the true vacuum of the potential at the origin and $`\lambda _0`$ cab not be a minimum, so there must be a critical point $`\theta _c`$ such that the sum becomes zero or even positive at some large N. Then the convergent process stops and the solution disappears.
It is interesting to observe how the difference equation allows a solution to avoid the constraint $`_0^{\lambda _0}V^{}(x)𝑑x>0`$ in the commutative limit and satisfies the constraint (3.1). It is comparable to the definition of the Riemann integral through the limit of a finite sum by using a particular partition of the coordinate region determined by $`\lambda _0>\lambda _1>\lambda _2>\mathrm{}>0`$ and letting the partition be smaller and smaller. When $`\theta `$ is big enough, the partition is coarse and it is possible for the sum to be different from the continuous limit and remain always negative. Decreasing $`\theta `$ is the same as taking the continuous limit to calculate the Riemann integral and it will finally make the solution impossible at some finite $`\theta _c`$. In this sense, the noncommutative scale $`\theta `$ really controls the continuum limit.
For the polynomial potential whose highest power is odd, the above proof breaks down as either $`V^{}(\varphi )\mathrm{}`$ as $`\varphi \mathrm{}`$ or $`V^{}(\varphi )\mathrm{}`$ as $`\varphi \mathrm{}`$. So in principle, it is always possible to find a soliton solution satisfying the required constrains at any nonzero $`\theta `$. Certainly the commutative limit $`\theta =0`$ invalidates the scaling transformation of $`xx/\sqrt{\theta }`$ and so it is a singular critical point of all the soliton solutions in this case.
Our qualitative analysis is also valid for non-polynomial potentials, such as the periodic cosine shaped potential, as long as it satisfies the bounded-below condition.
4. Noncommutative Solitons at finite $`\theta `$: $`\varphi ^4`$ potential
In this section we discuss the $`\varphi ^4`$ potential in detail. First we will use the numerical method to explicit construct these solutions at finite $`\theta `$ and see clearly how the solutions varies with $`\theta `$. Second we will see that the method in \[11\] can be extended to find the lowest critical point $`(\theta m^2)_c`$ explicitly to a very good approximation. Finally using a simple scaling argument, we find that $`(\theta m^2)_c`$ is the same for all the symmetric $`\varphi ^4`$ potentials. Note, we use the dimensionless paremeter $`\theta m^2`$ in this section instead of $`\theta `$, where $`m^2=V^{}(0)`$.
4.1. Numerical Results
First we explain the numerical method briefly. We use the relaxation method normally applied in solving differential equation with two point boundary conditions. The two boundary conditions for (3.1) are
$$\lambda _1=\lambda _2=0,\lambda _n0\mathrm{as}n\mathrm{}.$$
An initial guess for the solution is required as an input. We can estimate the asympotic value of $`\lambda _n`$ by going to the continuum limit and convert the difference equation into a differential equation, and ignore the nonlinear terms
$$\lambda (u)=\frac{2}{\theta m^2}u\frac{d^2\lambda (u)}{du^2.}$$
The solution satisfying the asympotic boundary condition is \[11\]
$$\lambda (n)=An^{\frac{1}{4}}e^{\sqrt{\frac{2n}{\theta m^2}}}$$
which can be used as the initial input to recursively find the true solution of the difference equation.
We use as an illustrative $`\varphi ^4`$ potential
$$V(\varphi )=\frac{1}{4}\varphi ^4\frac{5}{3}\varphi ^3+3\varphi ^2$$
which is shown in fig. 1. It has a local maximum at $`\varphi =2`$ and a local minimum at $`\varphi =3`$ in addition to the global minimum at $`\varphi =0`$.
The results from the numerical analysis are as follows:
1) At $`\theta =\mathrm{}`$ and ignoring the kinetic term, $`\varphi (x)=_n\lambda _n\varphi _n(x)`$ with $`V^{}(\lambda _n)=0`$ is the general soliton solution. There are several changes to the solution after including the kinetic term at finite $`\theta `$. First, those $`\lambda _n`$ which are zero at $`\theta =\mathrm{}`$ become nonzero. They increase when $`\theta `$ decreases, but never become appreciably large. Second, those $`\lambda _i`$ which are nonzero at $`\theta =\mathrm{}`$ will also change such that $`V^{}(\lambda _n)`$ is negative and decreasing when $`\theta `$ decreases. So if $`\lambda _i`$ at $`\theta =\mathrm{}`$ is a local minimum/maximum, it will decrease/increase when $`\theta `$ decreases. A typical example is the extrapolation of the solution $`\varphi =3\varphi _0`$ to $`\theta m^2=600`$, $`\varphi (x)=2.89\varphi _0+0.047\varphi _1+0.0015\varphi _2+0.00006\varphi _3+\mathrm{}`$. Here $`\lambda _0`$ decreases with $`\theta `$ because $`\lambda =3`$ is the local minimum of the illustrative potential given above. Finally, while several $`\lambda _n`$ can take the same nonzero value at $`\theta =\mathrm{}`$, this is not possible at finite $`\theta `$. For example, the solution $`\varphi (x)=3\varphi _1(x)+3\varphi _2(x)`$ at $`\theta =\mathrm{}`$, extrapolated to $`\theta m^2=78`$, is $`\varphi (x)=0.038\varphi _0+2.90\varphi _1+2.70\varphi _2+0.10\varphi _3+0.004\varphi _4\mathrm{}`$, which clearly shows that $`\lambda _1\lambda _2`$.
2) The existence of the critical point $`\theta _c`$ can be seen as follows. The solution at $`\theta =\mathrm{}`$ is characterized by those nonzero $`\lambda _n`$’s. If one of them becomes zero (approximately), then this particular solution can be regarded as being nonexistant. A nonzero $`\lambda _n`$ at $`\theta =\mathrm{}`$ always changes in the direction such that $`V^{}(\lambda _i)`$ decreases when $`\theta `$ decreases. Because $`V^{}(\lambda )`$ is bounded below for $`\lambda 0`$, $`\lambda _n`$ will reach a critical value at a finite $`\theta _c`$. With further decrease in $`\theta `$, this particular $`\lambda _n`$ will jump to a small value which is approximately zero, and this solution will cease to exist.
3) When $`\theta `$ decreases, the nonzero $`\lambda _n`$ with the largest $`n`$ reaches the critical value first. This can be explained considering the argument of section 3. The quantity $`\mathrm{\Delta }\lambda _N\lambda _N\lambda _{N1}`$ should be negatively large enough to make the solution possible. As it is proportional to $`1/N`$, the $`\lambda _N`$ with the larger $`N`$ will fail this criterion first. So the solution $`\varphi =\lambda _0\varphi _0`$ gives the lowest critical point $`\theta _c`$.
4) We emphasize that at the critical point the solitonic solution has no singular behavior because $`\theta _c`$ is finite. The solution looks almost the same as the example given in 1). The only sign of the criticality is the discontiuity of the $`\lambda _n`$, which changes from a finite nonzero value to zero.
4.2. Determination of $`(\theta m^2)_c`$
In this section we will find the method of \[11\] useful to explicitly express the critical point $`(\theta m^2)_c`$ within a good approximation. In particular for all the symmetric $`\varphi ^4`$ potential, $`(\theta m^2)_c`$ is the same, which can be proved using a scaling argument.
The lowest critical point $`(\theta m^2)_c`$ corresponds to the extrapolation of the solution $`\varphi (x)=\lambda _0\varphi _0(x)`$. The numerical solution indicates that $`\lambda _n`$ is extremely small for $`n0`$, so ignoring the nonliear terms is a good approximation
$$(n+1)\lambda _{n+1}(2n+1)\lambda _nn\lambda _{n1}=\frac{1}{2}\theta m^2\varphi .$$
Going back to the coordinate space representation, and noticing that $`\lambda _0>>0`$ is equivalent to adding a source term proportional to $`\varphi _0(x)`$
$$(\frac{1}{\theta m^2}^2+1)\varphi (x)=A\varphi _0(x),$$
except that it has a different boundary condition at $`n=0`$
$$\lambda _0=\frac{2}{\theta m^2}(\lambda _1\lambda _0)+A.$$
Compatibility with the boundary condition of the original equation of motion determines A.
Equation (4.1) is solved using properties of Laguerre polynomials $`\varphi _n`$,
$$\lambda _n=e_0^+\mathrm{}\frac{e^x}{1+\frac{2}{\theta m^2}x}L_n(x).$$
Define the function F by
$$F(a)=_0^+\mathrm{}\frac{e^x}{1+2ax},$$
as shown in fig.2 , and $`\lambda _0=F(\frac{1}{\theta m^2})`$. The two boundary conditions (4.1) and (4.1) should agree, which gives an equation for the scale factor A
$$m^2A(F(\frac{1}{\theta m^2})1)=V^{}(AF)=m^2(AF)a(AF)^2+b(AF)^3,$$
where we assume a general form of the $`\varphi ^4`$ potential, $`V^{}(\varphi )=m^2\varphi a\varphi ^2+b\varphi ^3`$. The existence of a real solution for A requires $`(\frac{1}{\theta m^2})\frac{4m^2b}{a^2}`$. $`F`$ is a monotonic function, so the equality determines the critical point $`(\theta m^2)_c`$
$$F(\frac{1}{(\theta m^2)_c})=\frac{4m^2b}{a^2}.$$
At this critical point, $`\lambda _0`$ is equal to
$$\lambda _c=(eF)_c=\frac{a}{2b.}$$
These two expressions agree with the numerical results.
Fig. 2: Function F(a) (4.1) which determines the coefficient $`\lambda _0=F(\frac{1}{\theta m^2})`$ of the single $`\delta `$ function like solution.
We have assumed that the potential has a global minimum at the origin, so $`V(\varphi )0`$. It sets a lower bound $`4m^2b/a^28/9`$. By (4.1), it sets a lower bound for $`F_c`$, which in turn determines a lower bound $`(\theta m^2)_c14.374`$. It is saturated exactly by the symmetric $`\varphi ^4`$ potential. Numerical analysis gives the exact result to be $`(\theta m^2)_c=13.92`$.
Let’s study the case of the symmetric $`\varphi ^4`$ potential in more detail. The effect of the potential on $`(\theta m^2)_c`$ enters through its derivative $`V^{}`$, as seen in the equation of motion in the projection operator basis, equation (3.1). Assuming one of the degenerate vacuum is at the origin, the symmetric $`\varphi ^4`$ potential is characterized by the zeros of its derivative, assumed to be at 0, 1/a, 2/a. So in general $`V^{}(\lambda )=\frac{m^2}{2}\lambda (a\lambda 1)(a\lambda 2)`$, and the variation of $`a`$ and $`m^2`$ gives all the symmeric $`\varphi ^4`$ potential. Writing out (3.1) explicitly as
$$(n+1)\lambda _{n+1}(2n+1)\lambda _n+n\lambda _{n1}=\frac{1}{2}\theta \frac{m^2}{2}\lambda _n(a\lambda _n1)(a\lambda _n2).$$
Under scaling transformation $`\lambda b\lambda `$, it becomes
$$(n+1)\lambda _{n+1}(2n+1)\lambda _n+n\lambda _{n1}=\frac{1}{2}\theta \frac{m^2}{2}\lambda _n(ab\lambda _n1)(ab\lambda _n2).$$
Effectively it transforms the moduli $`a`$ by the scaling factor $`b`$, which can be absorbed into $`m^2`$. This scaling of the variable will not affect the existence of the solution, and $`\theta m^2`$ remains invariant under the transformation, so the critical point $`(\theta m^2)_c=13.92`$ for the existence of the nontrivial solution is the same for all the symmetric $`\varphi ^4`$ potential.
5. Discussion
From the analysis in this paper, we find that noncommutative soliton sector exhibits a nontrivial hierarchy structure controlled by the noncommutative scale $`\theta `$. For quite general potentials which have a global minimum, the $`\theta _c`$ is finite and determined by the details of the potential. This indicates that the noncommutative geometry should not be considered as merely a passive kinematic background, but may have the similar dynamic content comparable to the potential, as shown by its effect on the existence of the soliton solutions.
The noncommutative soliton solutions can be interpreted as describing the tachyon condensation on the unstable brane at large B field background \[12,,13\]. It results in co-dimension-two branes obtained from the decay of the original brane. A general noncommutative soliton solution would disappear at finite $`\theta `$. But notice that the string field theory algebra factors into a direct product only in the limit $`B\mathrm{}`$, and those string degrees of freedom other than those involving only the center of mass coordinates have to be taken into account at finite $`\theta `$. This may change the picture in the string theory content.
The soliton solution at finite $`\theta `$ in general breaks the $`U(\mathrm{})`$ symmetry completely. At $`\theta =\mathrm{}`$ and ignoring the kinetic energy, the level N soliton solution, such as $`\varphi =_{n=0}^N\lambda _n\varphi _n`$ with all $`\lambda _n`$ equal, break the $`U(\mathrm{})`$ down to $`U(N)\times U(\mathrm{}N)`$ \[13\]. This soliton solution describes N coincident D(n-2)-branes with $`U(N)`$ symmetry. Inclusion of the kinetic energy term brings $`1/\theta ^1`$ corrections and all the $`\lambda _n`$ are different. It breaks the residual $`U(N)`$ symmetry completely into $`U(1)^N`$. So the $`U(N)`$ symmetry for this solution is at most approximate, and it seems that it should be interpreted as describing N branes that are not coincident. But this conclusion may not be true if we consider the additional terms in the action coming from the open string field theory, because at finite $`\theta `$ it is necessary to include those string degree of freedom other than those accounted for by the noncommutative field theory. In particular, noncommutative Yang-Mills theory is not enough and the analysis of Dirac-Born-Infeld action may give a different result which will be worth exploring.
Also notice that the open superstring tachyon potential for type II superstring theory is exactly of the symmetric $`\varphi ^4`$ shape, which follows from the reflection symmetry of the potential, although the details of the potential are unknown. But we find the critical point $`(\theta m^2)_c=13.92`$ for the existence of the noncommutative soliton is ignorant of the exact form of the $`\varphi ^4`$ potential, and so is charecteristic of type-II in the B-field background. We haste to add that this conclusion may be changed after inclusion of the open string field theory degrees of freedom.
It would be interesting to study the stability of these soliton solutions, and their effects on the quantum structure of the whole theory. We hope to explore these issues in the future.
Note added:
We notice that the recent two papers \[21,,22\] discussed the noncommutative solitons in the context of noncommutative scalar theory coupled with the noncommutative gauge field. These new solutions involve nontrivial gauge field configutation, which may answer some of the puzzles put forward in the last section of the paper concerning the finite $`\theta `$ behavior and are worth exploring.
Acknowledgements
The author cordially thanks Michael Douglas for invaluable advice, discussions and encouragement, without which this work would not have been possible.
References
relax A.Connes, M.R.Douglas and A.Schwarz,” Noncommutative Geometry and Matrix Theory:compactification on Tori”, JHEP 9802 (1998) 003, hep-th/9711162. relax N.Seiberg and E.Witten, ” String Theory and Noncommutative Geometry”,JHEP 9909 (1999) 032, hep-th/9908142. relax E. Witten, “Noncommutative Tachyons and String Field Theory”, hep-th/0006071. relax T.Filk, “Divergencies in a Field Theory on Quantum Space”, Phys.Lett.B376 (1996) 53. relax D.Bigatti and L.Susskind, “Magnetic Fields, Branes and Noncommutative Geometry”, hep-th/9908056. relax I.Chepelev and R.Roiban, “Renormalization of Quantum Field Theories on Noncommutative R\**d. I: Scalars”, hep-th/9911098. relax S.Minwalla, M.V.Raamsdonk and N.Seiberg, “Noncommutative Perturbative Dynamics”, hep-th/9912072. relax M.V.Raamsdonk and N.Seiberg, “Comments on Noncommutative Perturbative Dynamics”, hep-th/0002186. relax H.O.Girotti, M.Gomes, V.O.Rivelles and A.J.da Silva, “A Consistent Noncommutative Field Fheory: the Wess-Zumino Model”, hep-th/0005272. relax S.S.Gubser, S.L.Sondhi, “Phase Structure of Noncommutative Scalar Field Theories”, hep-th/0006119. relax R.Gopakumar, S.Minwalla and A.Strominger, ”Noncommutative Solitons”, hep-th/0003160. relax K.Dasgupta, S.Mukhi, G. Rajesh, “Noncommutative Tachyons”, hep-th/0005006. relax J.A.Harvey, P.kraus, F.Larsen, E.J.Martinec, “D-branes and Strings ss Noncommutative Solitons”, hep-th/0005031. relax A.Sen, “Descent relations among bosonic D-branes”, Int.J.Mod.Phys.A14(1999) 4061, hep-th/9902105. relax A.Sen, “Universality of the Tachyon Potential”, JHEP 9912 (1999) 027, hep-th/9911116. relax A. Sen and B.Zwiebach, “Tachyon Condensation in String Field Theory”, hep-th/9912249. relax J.A.Harvey and P.Kraus, “D-branes as Lumps in Bosonic Open String Field Theory”,JHEP 0004:012, 2000 hep-th/0002117. relax R.de Mello Koch, A.Jevicki, M.Mihailescu and R.Tatar, “Lumps and P-branes in Open String Field Theory”, hep-th/0003031. relax V.A.Kostelecky and S.Samuel, “On a Nonperturbative Vacuum for the Open Bosonic String”, Nucl.Phys.B336(1990) 263. relax N.Berkovits, A.Sen and B.Zwiebach, “Tachyon Condensation in Superstring Field Theory”, hep-th/0002211. relax C. Sochichiu, “Noncommutative Tachyonic Solitions. Intercation with Gaube Field”, hep-th/0007217. relax R.Gopakumar, S.Minwalla and A.Strominger, ”Symmetry Restoration and Tachyon Condensation in Open String Theory”, hep-th/0007226.
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# Acknowledgments
## Acknowledgments
We thank Jochen Bartels and Genya Levin for useful comments and Larry McLerran for useful correspondance. This research has been partially supported by the EU Framework TMR programme, contract FMRX-CT98-0194, Deutsche Forschungsgemeinschaft and the Polish Commitee for Scientific Research grants Nos. KBN 2P03B 120 19, 2P03B 051 19.
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# Stellar Companions and the Age of HD 141569 and its Circumstellar Disk
## 1 Introduction
Stars with far-infrared excesses generated in circumstellar disks were discovered in abundance by the IRAS mission. Since circumstellar material disappears over time due to processes such as stellar winds, radiation pressure, and accretion onto stars and planetesimals, most infrared excess stars are young, $``$10<sup>8</sup> yr (Habing et al., 1999; Spangler et al., 1999; Silverstone, 2000). To develop an understanding of the evolution of circumstellar material and its relationship to the formation of planets, detailed studies of individual disks must be combined with information about stellar ages. However, stars which are nearing the main sequence are difficult to age because they lack traditional indicators such as Lithium and x-ray emission and they move through the color-magnitude diagram very quickly.
The Herbig AeBe stars (HAEBEs) are thought to be in this transition phase, intermediate in evolution between protostars and stars on the zero age main sequence (ZAMS). The classical definition of HAEBEs, i.e. that they are of spectral type A or earlier, are found in clouds, and show emission lines and reflection nebulae (Herbig, 1960) almost certainly selects for such young objects. In addition, more recent studies (The, de Winter & Perez, 1994) use the presence of near or far-infrared excess to select HAEBEs. Thus, objects such as HD 141569, a B9.5 Ve star with H$`\alpha `$ in emission and 12–100 µm excess, which are not associated with any cloud or reflection nebula, fall into the HAEBE class. There have been attempts to date HAEBE stars within this general picture based on the strength of their infrared excesses (Hillenbrand et al., 1992), and authors have proposed an evolutionary sequence from the embedded HAEBEs to the isolated $`\beta `$ Pictoris or Vega-like main sequence stars (Malfait, Bogaert, & Waelkens, 1998, e.g.).
The presence of lower mass companions to these objects represents an independent way of measuring their age, assuming they and the companions are coeval. Lower mass companions may not yet be on the ZAMS and their ages can be estimated from theoretical pre-main sequence evolutionary tracks. For early-type stars near the ZAMS which also have circumstellar dust, this is probably the most accurate way to determine the disk age. The well-studied disk stars, $`\beta `$ Pic, Fomalhaut, and HR 4796A, were all dated by their association with lower mass co-moving companions (Barrado y Navascués et al., 1999, 1997; Stauffer et al., 1995).
HD 141569 is joining the ranks of well-studied dusty disk stars with the discovery of its large and morphologically complicated debris disk (Silverstone et al., 1998; Weinberger et al., 1999; Augereau et al., 1999; Fisher et al., 2000). In 1938, Rossiter identified HD 141569 as a member of a potential triple system with a second star 7.<sup>′′</sup>5 away (B) and a third star 1.<sup>′′</sup>5 from it (C) (Rossiter, 1943). In a study of HAEBEs, Gahm, Ahlin & Lindroos (1983) measured the spectral type of HD141569B as G0 V, and Lindroos (1985) concluded, based on its magnitude and color, that it was a background star. However, Gahm, Ahlin & Lindroos (1983) noted the presence of “peculiar” emission lines in the spectrum of HD 141569B, including Ca H and K and H$`\beta `$, which are often associated with young stars.
The parallactic distance to HD 141569A was measured by the Hipparcos mission as 99$`\pm `$10 pc, which combined with its visual magnitude of 6.8, makes it underluminous for its spectral type, just as are other young A-type infrared excess stars such as $`\beta `$ Pic, 49 Ceti and HR 4796A (Jura et al., 1998; Lowrance et al., 2000) (see also Figure 6).
Given the apparent proximity of two stars to HD 141569A, and given that the disk around HD 141569A extends to half of the projected distance between HD 141569A and B, this system could be dynamically interesting if the three stars are physically associated. In this paper, we show, via comparison to previous astrometry plus new near infrared imaging and visual spectroscopy, that HD 141569 A/B/C form a triple system, and we estimate their ages.
## 2 Observations and Data Analysis
### 2.1 HST and Ground-based Imaging
On 1998 September 27, Short integration time images of the HD 141569 system were obtained with the NICMOS camera 2 on the Hubble Space Telescope in order to acquire the primary star for coronagraphic observations. Two simple ACCUM mode images of 0.342 s each were taken with the F171M filter ($`\lambda _{eff}`$=1.721, FWHM=0.071) and ten dark frames of the same integration time were taken two days earlier for calibration of those acquisition images. The two source images were used for cosmic ray rejection by taking the minimum of the two values for every pixel. The median of the dark images was subtracted from the science image to remove the effects of detector shading and bias. The resulting image was divided by an on-orbit flat field frame taken in the same filter. The best available photometric calibration was applied in which 0 mag=948 Jy and 1 adu s<sup>-1</sup> = 1.071$`\times `$10<sup>-5</sup> Jy.
Near-infrared J ($`\lambda _c`$=1.27, $`\mathrm{\Delta }\lambda `$=0.25), H ($`\lambda _c`$=1.65, $`\mathrm{\Delta }\lambda `$=0.32), K ($`\lambda _c`$=2.20, $`\mathrm{\Delta }\lambda `$=0.40), and L<sub>s</sub>-band ($`\lambda `$=3.45, $`\mathrm{\Delta }\lambda `$=0.57) images were taken with the Hale 200-inch Telescope on 1999 May 25 and 28. The infrared camera had a pixel scale of 0.<sup>′′</sup>125 per pixel and a full field of view of 32$`{}_{}{}^{\prime \prime }\times `$ 32<sup>′′</sup>. The night of 25 May was cloudy, but the flux ratios of B and C to star A were measured at J, H, and K bands. Short integration time, 0.374 s, images were taken in which all three stars appear in every full-field frame. The seeing was 0.<sup>′′</sup>82 at K band, so stars B and C were easily resolved. So as not to saturate HD 141569A at J and H bands, the chopping secondary was used to smear the light in a direction perpendicular to the position angle (PA), 301.9° between components B and C.
The night of 28 May was photometric and measurements of HD 141569A at all four wavelengths were preceded by measurements of the photometric standard HD 129655. For both target and standard, an 8<sup>′′</sup> square subframe of the full array was employed to allow fast readout and prevent the images from saturating the bright primary at J and H or the thermal background at L<sub>s</sub>. At each filter, 100 integrations of 0.07 s were coadded for a total on-source integration time of 7 s. One 0.54 s image of HD141569 B and C was also obtained at L<sub>s</sub>.
Sky frames of the same integration time were obtained after every set of exposures on the stars. After sky subtraction, the images were flat fielded and corrected for hot pixels. Aperture photometry was then performed on the final images.
### 2.2 Spectroscopy from the W. M. Keck Observatory
Resolution $``$5000 spectra of stars B and C were taken with the LRIS (Oke et al., 1995) instrument on 1999 February 25 covering the spectral range of 6250$``$7550Å. A long slit of width 0.<sup>′′</sup>7 was placed at a PA of 301.9 deg to obtain simultaneous spectra of both B and C, and two integrations of 120 s each were made. The amplifier bias was estimated from the CCD over-read areas and subtracted from the spectra. A spectral flat field was made from halogen lamp exposures taken immediately after the spectra and was divided into the spectra. The seeing was 1.<sup>′′</sup>1, so B and C were not completely spatially resolved; their spectra were deblended by fitting the spatial dimension at every spectral element with two 1-D Gaussians plus a linear background. The two observed spectra of each star were then averaged.
High resolution spectra of all three stars were obtained with HIRES (Vogt et al., 1994) on 1999 July 19. The C1 decker was employed, which gives a slit 0.<sup>′′</sup>86 wide and 7<sup>′′</sup> long, yielding R=45,000 over 21 orders, from 5420-7880Å with some gaps. The small separation of B and of C required attention to scattered light. Hence, the image rotator was used to align the slit so that it was perpendicular to the vector separating the stars. Because of the good (0.<sup>′′</sup>8) seeing, small slit width, and care in position of the slit, we believe that the individual spectra of B and C are uncontaminated by scattered light from the other components. A radial velocity standard, HR 6056 (-19.9 km s<sup>-1</sup>), was observed immediately following the spectroscopy of HD 141569 A/B/C. The spectra were extracted using the MAKEE package written by T. Barlow; however, equivalent widths and radial velocities were measured using the SPLOT package in IRAF.
## 3 Results
### 3.1 Relative Motion
A summary of measurements of relative separation and orientation of HD 141569 A/B/C is given in Table 1 and Figure 1. The first point, from 1938 (Rossiter, 1955), averages four measurements of B–C and three measurements of A–B made in 1938 and 1943. Although no uncertainties were given for the measured separations and PAs, we used the scatter in the measurements to estimate them. The 1998 data are from our NICMOS acquisition images and the 1999 data from our Palomar images. The uncertainties for the Palomar positions were found from the standard deviation of many independent measurements as described in §3.4. From 1938 to 1998, the motions of the three stars are negligible to within the uncertainties, and therefore they all have the same proper motion.
In 1995-1996, Pirzkal, Spillar, & Dyck (1997) imaged HD 141569 A/B/C with shift-and-add and measured the A–B separation as 6.<sup>′′</sup>8 at a position angle of 312 and A–C as 8<sup>′′</sup>.0 with PA=314. No uncertainties were provided in that paper, and these results are inconsistent with the others.
### 3.2 Spectral Types and Features
The LRIS spectra of stars B and C are shown in Figure 2. Notable features are H$`\alpha `$ in emission, Li I 6708Å in absorption and the TiO bands at $``$7000Å which all appear in both spectra. The spectra have not been divided by a spectral standard, so global slopes are not meaningful. Figure 3 presents a portion of the LRIS spectrum of B compared to spectral standard stars in Kirkpatrick, Henry, & McCarthy (1991). Figure 4 does the same for C. We assign a spectral type of M2 V to B and M4 V to C based on the depth of the TiO features.
The equivalent widths of H$`\alpha `$ and Li I are presented in Table 2 and the H$`\alpha `$ line profiles for B and C are shown in Figure 5, both from the HIRES spectra. Stars B and C have double peaked lines with nearly the same shape and separations of their peaks of 1Å. H$`\alpha `$ in star A is much broader, double peaked with a peak to peak separation of 5.3Å (242 km s<sup>-1</sup>), and has a stronger blue than red peak; all of which are consistent with previous measurements of H$`\alpha `$ by several authors (Andrillat, Jaschek, & Jaschek, 1990; Zuckerman, 1994; Dunkin, Barlow, & Ryan, 1997).
### 3.3 Radial Velocity
From a cross-correlation of the high resolution B and C spectra over each of the 21 spectral orders, B is moving away from us faster than C by 0.9$`\pm `$0.4 km s<sup>-1</sup>. The uncertainty is the standard deviation in the mean of the 21 cross-correlations. The radial velocity of B, from cross correlation between its spectrum and that of the standard, is –1.5 $`\pm `$ 0.6 km s<sup>-1</sup>. We did not determine the radial velocity of A because it had no lines suitable for comparison with the M0.5III spectral-type radial velocity standard. From the literature its radial velocity is –6 $`\pm `$ 5 km s<sup>-1</sup> (Frisch, 1987; Dunkin, Barlow, & Ryan, 1997).
### 3.4 Photometry
In each photometric (i.e. 28 May) Palomar image, the magnitude of HD 141569A was measured in a 8<sup>′′</sup> diameter aperture. In every short non-photometric (i.e. 25 May) frame containing all three stars, A was used as a point spread function (PSF) for fitting the locations of B and C and flux ratios B/A and C/A, via minimization of the chi-square. All 40 images taken at J, and 20 images taken in each of the H and K bands were fit independently. Once the flux ratios were determined, the magnitudes of B and C were calculated from the photometry of star A. The results are summarized in Table 3. The uncertainties reported are purely statistical and contain the statistical uncertainty in the magnitude of A and the standard deviation of the independent determinations of the flux ratios to B and C. In the photometric L<sub>s</sub> image, the total flux from both B and C was measured in an 8<sup>′′</sup> diameter aperture centered between the stars and then star B was used as a template PSF to fit the flux ratio of these two close stars.
## 4 Discussion
### 4.1 Companionship and Associations
The proper motion of component A has been measured by Hipparcos as $``$16.86 $`\pm `$ 0.98 mas yr<sup>-1</sup> in right ascension and $``$21.11 $`\pm `$ 0.71 mas yr<sup>-1</sup> in declination, so over the 60 years between Rossiter’s measurements and our own, it should have moved 1.<sup>′′</sup>01 west and 1.<sup>′′</sup>27 south. If B is a background star with negligible proper motion, the PA of the A–B pair would have gone from 310.5° in 1938 to 323.5° in 1998 due to the proper motion of A (Figure 1). A change of this magnitude is ruled out by the measurements. Furthermore, as can be seen from Table 2, all three stars have consistent radial velocities.
If stars B and C were ordinary main sequence M-dwarfs, their photometric distance would be $``$35 pc. From the Hipparcos catalog, the average proper motion of the 217 stars at this distance which are within 30 deg of HD 141569A is 240 mas/yr, which is nearly nine times larger than that of HD 141569A itself. The relative separations of A–B and A–C would have most likely changed by $``$15<sup>′′</sup> over 60 years if B and C were foreground stars whereas in fact the separations have remained nearly constant.
The measured changes in separation and position angle are consistent with what would be expected from orbital motion. Star B is at least 760 AU from A. If it is orbiting A, the orbital period of B is $``$13,500 yr and in 60 years it would move at most 1.6 deg in PA (depending on the inclination of its orbit). Star C is at least 140 AU from B. If C is orbiting B, its period is $``$2100 yr and in 60 yr, it would move at most 10 deg in PA. The relative radial velocities of B and C would be $`<`$2 km s<sup>-1</sup>. The actual changes in PA (see Table 1) and the measured relative radial velocities (see Table 2) are well within these constraints and imply that the stars could be orbiting one another.
Since A, B and C have common proper motions and common radial velocities, they are with high probability physical companions. Whether the three stars are actually bound, however, cannot be determined. As noted in Weinberger et al. (1999), if the companion stars are in the plane of the disk, the physical separation of A–B is 990 AU and that of B–C is 190 AU. Then, the ratio of the semi-major axes of the wider to the closer pair would be only $``$5.2. This is not expected to be a stable triple system (Eggleton & Kiseleva, 1995), although such young stars may not yet have had time to become unbound.
The presence of three young stars unassociated with a star-forming cloud begs the question of how they formed. HD 141569 is located in projection near a complex of high latitude molecular cloud cores, MBM 34-39 associated with the dark clouds L169, L183-4, and L134 (together called L134N). The last of these has 100µm emission contours which actually encompass HD 141569 (Sahu et al., 1998). The excess B–V color of HD 141569A, 0.095 mag, implies a reddening A<sub>v</sub>=0.3 mag, so the star cannot lie deep within the cloud. Sahu et al. (1998) concluded, based on the strength of interstellar absorption lines toward HD 141569A, that it lies behind part of L134N/MBM 37. However, these clouds have LSR radial velocities of 0.8–3.2 km s<sup>-1</sup> (Magnani, Blitz, & Mundy, 1985), compared to $``$20 km s<sup>-1</sup> for HD 141569A, which suggests that they are not co-moving with the star. The clouds are also quiescent and compact with internal temperatures of 3-12 K (Snell et al., 1981; Clemens & Barvainis, 1988), which suggests that they are not undergoing protostellar collapse or being heated by nearby stars.
The discovery of other young stars dispersed across the sky has prompted speculation about fast cloud dispersal mechanisms and runaway stars (Feigelson, 1996). One way to address the question of star formation in this region is to look for other nearby stars which may have formed at the same epoch as HD 141569. As a first attempt to search for members of such an association, we have queried the Hipparcos catalog using the approximate characteristics of the TW Hya association (Webb et al., 1999), a radius of 19 pc and a proper motion dispersion of 7 km s<sup>-1</sup>, which at the distance to HD 141569 correspond to a circle of radius 10° and a proper motion dispersion of $``$9 mas yr<sup>-1</sup> in each direction. In addition, we require the Hipparcos parallaxes to agree with that of HD 141569A to within their uncertainties. This search produces 14 stars. In comparison, a search of 20 other fields at the same absolute galactic latitude and a range of galactic longitudes produces an average of 7 $`\pm `$ 5 stars using the same search parameters. The data suggest an overdensity of stars near HD 141569. We note that at the distance of 100 pc, the Hipparcos catalog is highly incomplete for stars fainter than V=9 mag (later than F).
Two of the 14 stars near HD 141569 are spectral type A and lie below or along the ZAMS (Figure 6). This location in the HR diagram is populated by young A type stars (Lowrance et al., 2000) including well known disk sources such as $`\beta `$ Pic and HR 4796A as well as HD 141569 itself.
We have computed the space motion of HD 141569A following Johnson & Soderblom (1987) and find U,V,W = -3.0, -13.0, and -3.0 km s<sup>-1</sup>. This is within twice the velocity dispersion in the space motions of the local star formation associations such as $`\eta `$ Chameleon, TW Hya, and Tucanae (Zuckerman & Webb, 2000). So, the HD 141569 system may have formed as part of a larger episode of star formation near the Sun.
### 4.2 Age
The pre-main sequence nature of stars B and C is confirmed by the presence of Lithium in absorption. Since M-type stars should be fully convective, Lithium in the spectra indicates that these stars are not yet hot enough in their cores to burn it or have had insufficient time to burn all of it. The equivalent widths of 0.5Å for B and C is very close to the boundary of 0.54Å set by Martin (1998) for separating weak-line T Tauri stars (WTTS) from post-T Tauri stars. This sets an upper limit on their age of 10 Myr (Martin et al., 1994). The H$`\alpha `$ equivalent widths ($`<`$10Å) also argue for classifying these stars as WTTS as opposed to classical TTS. The strength of the double peaked H$`\alpha `$ emission is within the distribution of chromospherically active main-sequence M dwarfs (Stauffer & Hartmann, 1986). The central reversal of the line indicates a high level of chromospheric activity, which is generally associated with a variety of factors including youth and rotation. Finally, Lindroos (1985) classified the combined B/C spectrum as “peculiar” because he detected hydrogen and calcium emission lines, and pre-main-sequence stars often have such emission lines.
The A/B/C system falls within the 90% confidence error circle of a ROSAT sky survey point-source detection. The PSPC count rate of 0.093 s<sup>-1</sup> converts to an x-ray luminosity of 8.72 $`\times `$ 10<sup>29</sup> erg s<sup>-1</sup> assuming that it comes from stars at 100 pc with T-Tauri-like x-ray spectra (Neuhauser et al., 1995). It is likely that the x-ray emission comes from the later-type stars, B and C rather than from the primary star. We estimate the bolometric luminosities of stars B and C as log(L/L)=$``$0.62 and $``$0.95 respectively from their spectral types and J-band magnitudes using bolometric corrections from Hartigan, Strom, & Strom (1994). This makes their total x-ray to bolometric luminosity ratio log(L<sub>x</sub>/L<sub>bol</sub>)=$``$3.2. For pre-main sequence stars, the ratio of x-ray to bolometric luminosity increases with stellar age, and the ratio for stars B/C is typical for that of stars in clusters of age $`<`$20 Myr. (Kastner et al., 1997).
Finally, we estimate the age of stars B and C on the basis of their location on theoretical pre-main sequence evolutionary tracks. Here, a major source of uncertainty is the not-well understood effective temperature scale for M-dwarfs (Allard et al., 1997). Following the spectral-type to temperature calibration of Luhman & Rieke (1998), we can assign an effective temperature for B (M2 V) and C (M4 V) of 3500$`\pm `$85 K and 3200$`\pm `$85 K, respectively, where the uncertainties correspond to 1/2 spectral sub-class. We note that the temperature scale of Kirkpatrick, et al. (1993) gives a temperature 175 K hotter for M4 V.
Shown in Figure 7 are tracks by Baraffe et al. (1998) with the effective temperatures and absolute J-band magnitudes of stars B and C plotted. The plotted magnitudes assume that stars B and C are at the same distance as A, 100 pc, and the uncertainties in the magnitudes correspond to the uncertainty in the Hipparcos parallax to A. Within the uncertainties, stars B and C appear to be the same age of 2 to 8 Myr. Star A is indistinguishable from the ZAMS and is not shown on the figure. The tracks indicate masses for stars B and C of 0.5 and 0.25 M, respectively. Tracks by other authors give the same basic result. For comparison, Figure 7 also shows HR 4796B with an effective temperature of a M2.5 star from Luhman & Rieke (1998). This companion to HR 4796A, another star with a well studied circumstellar disk, was assigned an age of 8$`\pm `$2 Myr by Stauffer et al. (1995) using tracks from D’Antona & Mazzitelli (1994). The age given by the Baraffe et al. (1998) tracks is consistent with that determination and shows that HD 141569 B/C are about a factor of two younger than HR 4796B.
Based on all of the above arguments, we estimate an age for the HD 141569 system of 5 $`\pm `$ 3 Myr.
### 4.3 Dynamics
The disk around HD 141569A has two features which suggest dynamical sculpting (Weinberger et al., 1999). First, the density of scatterers is as high at 360 AU as 200 AU from the star. If the companions are out of the plane of the disk, they could excite significant vertical velocities in the disk dust.
Second, there is a dip in surface brightness, or a “gap” in the disk at a radius of 250 AU with a width of 60 AU. No point source is seen in the gap to a limit of F110W=20.3 mag. The gap is as circular as the disk and must be cleared continually to remove particles drawn through it by radiation pressure and Poynting-Robertson drag. If the companion stars orbit each other and the primary and are in the plane of the disk, their center of mass is 1053 AU from the primary. The 2–9:1 Lindblad resonances between the orbital period of the companion center of mass and orbiting particles in the disk are shown with the dashed ellipses on Figure 1. There is no obvious agreement between the resonances and the structure observed in the disk, even if the companions are assumed to be somewhat out of the plane and so the resonances shifted. The resonances closest to the gap lie at 243 AU (9:1) and 263 AU (8:1) from the star. It is not clear, however, why only these high order resonances and not any of the others would mold the disk. If the companions were in very eccentric orbits and currently near apastron, the locations of the resonances might change dramatically with time. However, the circularity of the gap implies a stable dynamical influence over long time periods.
## 5 Conclusions
On the basis of common proper motion, common radial velocity and the low probability of a chance superposition of three young stars away from any known star-forming cloud, we conclude that HD 141569 A/B/C form a physical association which may or may not be bound. The age of the stars as determined from spectroscopic features, x-ray emission, and placement on pre-main sequence tracks is 5 Myr. At least two other stars, HD 141693 and HD 140574, may be members of this common proper motion group.
HD 141569 has a now well studied disk with grains present from 25 – 500 AU of the star. The inferred size of the grains, $`<`$5 $`\mu `$m (Fisher et al., 2000), means that radiation pressure blows them away in at least an order of magnitude less time than our derived stellar age. Thus, they must be continuously regenerated, probably through collisions of larger bodies. The timescales for the formation of planetesimal cores is 10<sup>4</sup> \- 10<sup>5</sup> yr (Wetherill, 1980) which is easily consistent with the age of the system. From measurements of CO around HD 141569, the remnant mass of H<sub>2</sub> is estimated as 20 to 460 M<sub>Earth</sub> (Zuckerman, Forveille, & Kastner, 1995), so any gas giants present must have formed very quickly.
The morphology of the disk, including a gap at 250 AU, indicates the presence of dynamical sculpting. Resonant interactions between the companions and the disk do not appear to account for the structure, and it would be hard given current models of planet formation to generate Jupiter sized bodies in a disk of such young age at distances so far from the central star (Boss, 1998). The cause of the structure in the disk thus remains unknown.
We thank G. Neugebauer for donating 200-inch Telescope time to this project. We thank B. Schaefer and R. Quick for making the LRIS observations at the W. M. Keck Observatory, which is operated as a scientific partnership between the California Institute of Technology, the University of California, and the National Aeronautics and Space Administration and was made possible by the generous financial support of the W. M. Keck Foundation. This work is supported in part by NASA grant NAG 5-3042, and based on observations with the NASA/ESA Hubble Space Telescope, obtained at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc. under NASA contract NAS5-26555. We also thank M. Jura for many helpful discussions.
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# Contents
## 1 Introduction
One of the corner-stones of all recent developments in string theory has been the exact microscopic description of D-branes provided by J. Polchinski . D-branes are Ramond-Ramond charged objects defined as hypersurfaces on which open strings can end. From the low-energy point of view, the D-branes instead appear as classical solutions of the supergravity field equations which preserve a fraction of the original set of supercharges, and hence are BPS saturated. The fact that type II supergravity possesses classical BPS solutions was known well before Polchinski’s paper (see for instance ), but only after the stringy interpretation their fundamental importance has been fully appreciated.
More recently, after a series of papers by A. Sen , a lot of attention has been devoted to the study of non-BPS D-branes (for reviews see ). The main motivations are the following: $`i)`$ stable non-BPS D-branes are crucial in testing some non-perturbative string dualities without relying on supersymmetry arguments ; $`ii)`$ the existence of non-BPS D-branes could be used to define or find duality relations in a non-supersymmetric context ; $`iii)`$ non-BPS D-branes hopefully can play a crucial role in describing non-perturbative properties of non-supersymmetric gauge theories, similarly to what happened with the BPS D-branes in connection with supersymmetric Yang-Mills theory.
Despite the big amount of knowledge that has been accumulated, non-BPS D-branes have still to be completely understood. There are two main types of D-branes which do not saturate the BPS bound: those which are unstable due to the existence of tachyons on their world-volume, and those which are instead stable and free of tachyons. For example, the D$`p`$ branes of type IIA with $`p`$ odd are of the first kind, while the D-particle of type I is of the second kind. In both cases, their microscopic description is fairly well under control, for instance, in terms of world-volume effective actions or boundary states <sup>1</sup><sup>1</sup>1For a review of the boundary state formalism and its applications, see ., but very little is known about the nature or even the existence of the classical geometry associated to them. Actually, in the case of an unstable non-BPS brane one should first of all specify what is the meaning of a classical solution, but even in the simpler case of the stable non-BPS branes a discussion about the conditions which guarantee the consistency of the classical geometry with the microscopic string description is lacking.
From the effective field theory point of view, the problem of finding the classical solution corresponding to a given brane configuration is always well defined, because it amounts to solve the inhomogeneous field equations of the supergravity theory in the presence of a source term represented by the brane effective action, which has a delta-function singularity at the position of the brane. Equivalently, the same problem can be addressed by solving the homogeneous field equations and then imposing that the solution has the asymptotic behaviour prescribed by the boundary state description of the brane. As we shall see, this procedure clearly shows how the integration constants appearing in the solution are related to the physical parameters of the brane configuration, typically its tension and charges. The existence and the nature of possible singularities of the solution therefore depend on such constants, and the physical requirements that a classical solution should fulfill (for instance, the absence of naked singularities) constrain the acceptable range of their values.
From the microscopic point of view, however, the supergravity action describes only the effective dynamics of the model at low energies, and thus one has to be sure that the constraints imposed on the brane parameters by the existence of a meaningful classical solution of the field equations are compatible with the approximations that lead to such an action. As is well known, this may happen only when the brane tension is very large, i.e. $`M_p\mathrm{}`$. The validity of the no-force condition, which allows to construct a superposition of an arbitrary large number of D-branes, is therefore the necessary condition for the consistency of a classical solution.
In this paper we address these questions in general and discuss in particular the case of the non-BPS D-particle in six dimensions arising from the compactification of the type II string theory on a K3 manifold at the orbifold point. This configuration, which can be easily described using the boundary state formalism, is stable because the orbifold projection removes the open string tachyons; moreover, at a particular value of the volume of the compact space it satisfies a no-force condition at one loop . This system seems therefore to possess all the required features to produce a non-trivial classical geometry whose leading behavior at large distances has recently been found in . However, differently from what expected, we will give evidences that this does not happen. In particular, we will find that the geometry corresponding to a stack of such non-interacting non-BPS branes displays pathological features which make the configuration unacceptable. This is not in contradiction with the result of . In fact, one-loop calculations can extend their validity to the supergravity regime only when some preserved supersymmetries cancel higher order corrections. In the case of a non-BPS configuration clearly this is not guaranteed and our result is indeed an evidence that the no-force found at first order is lost at two-loop level. It is interesting to observe that the singularities we encounter in our solution manifest themselves as divergences in the metric tensor which make the gravitational force to become “repulsive” at small distances. This kind of singularities are known in the literature as repulsons , and have been recently considered in string theory in where a mechanism for their resolution has been proposed. It would be interesting to investigate whether a similar mechanism can be applied also in our case.
The content of this paper is the following: in Section 2 we consider the unstable non-BPS D-branes of type II in ten dimensions, and discuss the limits and the validity of the corresponding classical geometry. In Section 3 we compactify the type II string on the orbifold $`T_4/𝐙_2`$ and write the six-dimensional low-energy effective action which describes the model in the field theory limit. Even though it is known that this action is that of the (1,1) supergravity coupled to 20 $`U(1)`$ vector multiplets, for our purposes we find more convenient to recover its explicit form using the S-duality which relates our model to the heterotic theory compactified on $`T^4`$. In Section 4, by exploiting the precise knowledge of the action obtained with the duality map, and of the couplings between the fields and the D-brane given by the boundary state, we write down the equations of motion and solve them iteratively in the effective open string coupling $`gN`$. Contrarily to what expected, the perturbative series may be resummed to obtain the exact solution with the asymptotic behaviour described by the boundary state. As already stressed, the solution presents pathologies which signal the impossibility of having a macroscopic stable configuration of $`N`$ coincident non-BPS D-particles. It may be possible that more complicated bound states made up of stable non-BPS D-branes can resolve these singularities and lead to a consistent classical solution. This issue is discussed in Section 5, where following , we solve the system of the homogeneous field equations in full generality and comment on possible different settings in which a macroscopic solution can exist and the singularity be resolved. Finally, in Appendix A and B we present some technical details and explicit calculations.
## 2 Non-BPS D-branes in type II theories
By now it is well-known that type II string theories in ten dimensions possess non-BPS D$`p`$-branes (with $`p`$ odd in IIA and $`p`$ even in IIB) <sup>2</sup><sup>2</sup>2Throughout this section, we always take $`p<7`$. which are unstable due to the existence of open string tachyons on their world-volumes. If we give a vanishing v.e.v. to such tachyons, these non-BPS D-branes are easily described by a boundary state $`|Dp`$ which has only a NS-NS component
$$|Dp=\sqrt{2}|Bp_{\mathrm{NS}\mathrm{NS}}$$
(2.1)
where $`|Bp_{\mathrm{NS}\mathrm{NS}}`$ is the GSO projected boundary state in the NS-NS sector whose explicit expression can be found for example in . The factor of $`\sqrt{2}`$ in (2.1) is required by the open-closed string consistency and implies that a non-BPS D$`p`$-brane of type IIA (or B) is heavier than the corresponding BPS D$`p`$-brane of type IIB (or A). In fact, from (2.1) one can see that the tension is
$$\widehat{M}_p=\sqrt{2}\frac{T_p}{\kappa _{10}g}$$
(2.2)
where $`T_p=\sqrt{\pi }\left(2\pi \sqrt{\alpha ^{}}\right)^{3p}`$ is the factor that appears in the normalization of the boundary state, $`\kappa _{10}=8\pi ^{7/2}\alpha _{}^{}{}_{}{}^{2}`$ is the gravitational coupling constant in ten dimensions, and $`g`$ is the string coupling constant.
By projecting $`|Dp`$ onto perturbative closed string states, one can easily see that the only massless bulk fields emitted by the non-BPS D-branes are the graviton $`G_{\mu \nu }`$ and the dilaton $`\phi `$, and that their couplings are described by the DBI action (in the Einstein frame) <sup>3</sup><sup>3</sup>3 We label the world-volume directions of the brane with indices $`\alpha ,\beta ,\mathrm{}=0,\mathrm{},p`$ and the transverse directions by indices $`i,j,\mathrm{}=p+1,\mathrm{},9`$.
$$S_{\mathrm{boundary}}=\widehat{M}_pd^{p+1}x\mathrm{e}^{\frac{p3}{4}\phi }\sqrt{detG_{\alpha \beta }}.$$
(2.3)
The graviton and the dilaton can in principle propagate in the entire ten dimensional space-time where their dynamics is governed by the following bulk action (again in the Einstein frame)
$$S_{\mathrm{bulk}}=\frac{1}{2\kappa _{10}^2}d^{10}x\sqrt{detG}\left((G)\frac{1}{2}_\mu \phi ^\mu \phi \right)$$
(2.4)
which is a consistent truncation of the type IIA (or B) supergravity action containing only those fields emitted by the non-BPS D-brane.
Following the procedure described in and using the explicit form of the boundary state (2.1), one can find the metric and dilaton profiles at large distances from the brane. These turn out to be given by
$`G_{\alpha \beta }`$ $``$ $`\left(1+{\displaystyle \frac{p7}{8}}{\displaystyle \frac{\widehat{Q}_p}{r^{7p}}}+\mathrm{}\right)\eta _{\alpha \beta }`$
$`G_{ij}`$ $``$ $`\left(1+{\displaystyle \frac{p+1}{8}}{\displaystyle \frac{\widehat{Q}_p}{r^{7p}}}+\mathrm{}\right)\delta _{ij}`$ (2.5)
$`\phi `$ $``$ $`{\displaystyle \frac{3p}{4}}{\displaystyle \frac{\widehat{Q}_p}{r^{7p}}}+\mathrm{}`$
where $`r`$ is the distance in transverse space and
$$\widehat{Q}_p=\frac{2\widehat{M}_p\kappa _{10}^2g^2}{(7p)\mathrm{\Omega }_{8p}}$$
(2.6)
with $`\mathrm{\Omega }_q=2\pi ^{\frac{1}{2}(q+1)}/\mathrm{\Gamma }\left(\frac{1}{2}(q+1)\right)`$ being the area of a unit $`q`$-sphere. We remark that the expressions (2.5) are the same as those of the usual BPS D$`p`$-branes in the Einstein frame (except for the different value of the tension appearing in $`\widehat{Q}_p`$). It is also worth pointing out that the terms in (2.5) proportional to $`\widehat{Q}_p`$ can be obtained by evaluating the 1-point diagrams for the graviton and the dilaton (see Figure 1) in which the couplings with the brane are read from the boundary action (2.3).
At this point it is natural to ask whether these non-BPS D-branes can yield a non-trivial classical geometry in ten dimensions, just like the BPS D-branes do. In other words, one can ask whether the metric in (2.5) can be interpreted not only as a small deformation of the flat Minkowski space-time due to the emission of a graviton, but also as the asymptotic behavior of a non-trivial space-time geometry. One way to answer this question is to compute higher order terms in $`\widehat{Q}_p`$, both for the metric and the dilaton profiles, and eventually re-sum the perturbative series. Despite its conceptual simplicity, this procedure clearly requires calculations which become more and more daunting as one proceeds in the perturbative expansion. However, there is another (and more efficient) way to answer the question, namely one can write the field equations of the metric and the dilaton that follow from the bulk action (2.4), and then look for a solution with the asymptotic behavior (2.5). In our case these equations are simply
$$_\mu \left(\sqrt{detG}G^{\mu \nu }_\nu \phi \right)=0$$
(2.7)
for the dilaton, and
$$R_{\mu \nu }\frac{1}{2}G_{\mu \nu }\frac{1}{2}\left(_\mu \phi _\nu \phi \frac{1}{2}G_{\mu \nu }_\rho \phi ^\rho \phi \right)=0$$
(2.8)
for the metric. If we require Poincaré invariance in the world-volume and rotational invariance in the transverse space, we can use the following Ansatz
$`ds^2`$ $`=`$ $`B^2(r)\eta _{\alpha \beta }dx^\alpha dx^\beta +F^2(r)\delta _{ij}dx^idx^j`$
$`\phi `$ $`=`$ $`\phi (r)`$ (2.9)
and then solve for the functions $`B(r)`$, $`F(r)`$ and $`\phi (r)`$. Using this Ansatz, the field equations (2.7) and (2.8) become
$`\phi ^{\prime \prime }+\left(\xi ^{}+{\displaystyle \frac{8p}{r}}\right)\phi ^{}=0`$
$`(\mathrm{log}B)^{\prime \prime }+\left(\xi ^{}+{\displaystyle \frac{8p}{r}}\right)(\mathrm{log}B)^{}=0`$
$`(\mathrm{log}F)^{\prime \prime }+\left(\xi ^{}+{\displaystyle \frac{8p}{r}}\right)(\mathrm{log}F)^{}+{\displaystyle \frac{\xi ^{}}{r}}=0`$ (2.10)
$`\xi ^{\prime \prime }+(\mathrm{log}F)^{\prime \prime }\left(\xi ^{}{\displaystyle \frac{8p}{r}}\right)(\mathrm{log}F)^{}`$
$`+(p+1)[(\mathrm{log}B)^{}]^2+(7p)[(\mathrm{log}F)^{}]^2+(\phi ^{})^2=0`$
where $`{}_{}{}^{}d/dr`$, and
$$\xi (p+1)\mathrm{log}B+(7p)\mathrm{log}F.$$
(2.11)
The general solution of these equations can be easily deduced from the analysis of <sup>4</sup><sup>4</sup>4See also for a related discussion. and depends on several integration constants which can be uniquely fixed by imposing the asymptotic behavior (2.5) dictated by the boundary state. If we introduce the harmonic functions
$$f_\pm (r)=1\pm x\frac{\widehat{Q}_p}{r^{7p}},$$
(2.12)
the non-BPS D-brane solution can be written in a rather simple form and reads
$`B^2(r)`$ $`=`$ $`\left({\displaystyle \frac{f_{}(r)}{f_+(r)}}\right)^\lambda `$
$`F^2(r)`$ $`=`$ $`f_{}(r)^\mu _{}f_+(r)^{\mu _+}`$ (2.13)
$`\mathrm{e}^{\phi (r)}`$ $`=`$ $`\left({\displaystyle \frac{f_{}(r)}{f_+(r)}}\right)^\nu `$
where
$`\lambda ={\displaystyle \frac{7p}{16x}},\mu _\pm `$ $`=`$ $`{\displaystyle \frac{2}{7p}}\pm {\displaystyle \frac{p+1}{16x}},\nu ={\displaystyle \frac{p3}{8x}},`$ (2.14)
$`x`$ $`=`$ $`\sqrt{{\displaystyle \frac{7p}{8(8p)}}}.`$
It is not difficult to realize that the metric described by (2.13) possesses a curvature singularity at $`r_p(x\widehat{Q}_p)^{1/(7p)}`$ . Thus, the solution (2.13) is meaningful only in the “physical” region $`r>r_p`$.
The problem that we want to address now is the consistency of the geometrical description (2.13) with the microscopic string interpretation of the non-BPS D-branes. In other words we want to check whether the classical solution is consistent with the approximations that lead to the action from which it descends. In this respect, we observe that (2.4) is a valid effective action only when curvature effects are small with respect to the string scale so that higher derivative terms can be consistently neglected in the Lagrangian. In our case this happens when $`M_p`$ is large. However, as we have seen at the beginning of this section, the microscopic interpretation implies that the tension $`M_p`$ is not a free parameter, and the only way to make it large is to consider a superposition of $`N`$ D-branes and then take the limit $`N\mathrm{}`$. But this is possible only if the D-branes do not interact with each other, i.e. if they satisfy a no-force condition.
Unfortunately, the unstable non-BPS D-branes in ten dimensions do not enjoy this property. To see this, let us compute the interaction energy $`\mathrm{\Gamma }`$ between two D-branes, which, from the closed string point of view, is simply given by
$$\mathrm{\Gamma }=Dp|𝒫|Dp$$
(2.15)
where $`𝒫`$ is the closed string propagator. Using standard techniques, it is easy to see that $`\mathrm{\Gamma }`$ is not vanishing; moreover, by taking the field theory limit, one may find that
$$\mathrm{\Gamma }|_{\alpha ^{}0}=\widehat{M}_pV_{p+1}\frac{\widehat{Q}_p}{r^{7p}}$$
(2.16)
where $`V_{p+1}`$ is the (infinite) world-volume of the D$`p`$-brane. Eq. (2.16) explicitly shows that there exists a non-vanishing force between two non-BPS D-branes: in fact the attraction due to the exchange of gravitons and dilatons is not compensated by any repulsion because these branes do not carry any charge. Thus, according to our previous discussion, we can conclude that, since it does not satisfy the no-force condition, the classical solution (2.13) is acceptable (for $`r>r_p`$) as long as we do not require a microscopic string interpretation of the underlying theory.
One may wonder whether these conclusions may change by taking into account the presence of tachyons on the brane world-volume. As is well-known, these fields have non-trivial consequences on the open-string dynamics and modify the structure of the boundary action . Even if these effects are taken into account, and consequently the form of the solution is changed (see for example for a recent discussion on this point), the no-force condition still cannot be satisfied. Furthermore, if we appeal to the existence of tachyonic modes, a more fundamental question arises, namely to what extent a classical geometry can be associated to an unstable system.
In conclusion we see that the two essential requirements for the existence of a consistent geometrical description of a D-brane are its stability and the validity of the no-force condition. As we have seen, these two properties are not satisfied by the non-BPS D-branes considered so far. However, it turns out that in suitable orbifold compactifications of type II theories, there exist non-BPS D-branes which are stable, i.e. tachyon free, and do not interact pairwise, at one-loop, . These are therefore the natural candidates to be considered for a classical supergravity interpretation consistent with a microscopic string description. The study of these branes will be the subject of the remaining part of this paper.
## 3 Low-energy actions
In this section we provide the necessary ingredients to analyze the geometry associated to the stable non-BPS D-branes in six dimensions. These are non-perturbative configurations of the type II string compactified on $`T_4/𝐙_2`$ orbifolds, which have been extensively studied using the boundary state formalism . Here we will focus on the simplest case, namely the stable non-BPS D-particle. From its boundary state description it is easy to realize that such a particle is a source for a graviton, a dilaton, four scalars and a vector potential in six dimensions. Our goal is to determine the classical configuration for these fields and study its consistency. Therefore, in subsection 3.1 we first derive the bulk action that governs the dynamics of the fields emitted by the D-particle. As we will see, this is a consistent truncation of the $`D=6`$ supergravity action. Later, in subsection 3.2 we will derive the boundary action that describes the couplings of the bulk fields with the D-particle.
### 3.1 Bulk theory
The theory we consider is the non-chiral supergravity in six dimensions with sixteen supercharges. To describe it, we start from the ten dimensional type IIA string compactified on a torus $`T_4`$ and orbifolded with a discrete parity $`𝐙_2`$ generated by the reflection $`_4`$ of the four compact directions (labeled by indices $`a,b=6,\mathrm{},9`$). Equivalently, we could start from the type IIB string compactified on the orbifold generated by $`_4(1)^{F_L}`$, where $`(1)^{F_L}`$ is the operator that changes the sign of all R-R and R-NS states. These two theories are related to each other by a $`T`$-duality along one of the compact directions, and thus yield the same low-energy Lagrangian. Both orbifolds have been extensively studied in the literature (see for example ) and it is well known that their massless spectrum is described by the six dimensional $`(1,1)`$ supergravity coupled to 20 $`U(1)`$ vector multiplets. The action of this theory is commonly written in the following compact form
$`S`$ $``$ $`{\displaystyle }d^6x\sqrt{detg}\mathrm{e}^{2\phi }[(g)+4_\mu \phi ^\mu \phi {\displaystyle \frac{1}{12}}H_{\mu \nu \rho }H^{\mu \nu \rho }`$
$`{\displaystyle \frac{1}{4}}\left(M^1\right)_{IJ}F_{\mu \nu }^IF^{J\mu \nu }+{\displaystyle \frac{1}{8}}\mathrm{Tr}\left(_\mu M^\mu M^1\right)+\mathrm{}],`$
where $`g_{\mu \nu }`$ is the string frame metric, $`\phi `$ is the dilaton, $`H_{\mu \nu \rho }`$ is the field strength of the NS-NS two-form, and $`M`$ is the matrix parameterizing the coset manifold of the scalar fields. In our case, this coset is $`SO(4,20)/SO(4)\times SO(20)`$, which has the right dimension to accommodate the 80 scalars of 20 vector multiplets. Finally, $`F_{\mu \nu }^I`$ contains the field strengths of all the $`U(1)`$’s present in the spectrum and transforms as a vector under $`T`$-duality <sup>5</sup><sup>5</sup>5Note that, in general, a $`T`$-duality transformation mixes fields which are in different multiplets..
The action (3.1) explicitly displays the full $`T`$-duality invariance of the theory. However, for our purposes this form is too general and blurs a few crucial details. First of all, since our aim is to study the theory in the region of its moduli space corresponding to the orbifold $`T_4/𝐙_2`$, we are not interested in having a manifestly $`T`$-dual invariant formulation. Indeed, the expectation values of the scalars change under $`T`$-duality, and so does the shape of the compact space. Secondly, we need to know the precise normalizations of the various terms in the action, and in particular their dependence on the moduli. In fact, from the microscopic description we know that the stability of the non-BPS branes crucially depends on the radii of the compact space . Moreover, if we want to construct from these branes a macroscopic object, the shape of the internal $`T_4`$ must be further restricted by fixing all radii to some critical value . For these reasons, we need to write the supergravity action (3.1) in a form capable to make more explicit its relation with the orbifold construction. This means that we must break the $`SO(4,20)/SO(4)\times SO(20)`$ invariance and select, out of the 80 scalars, 4 fields that describe the characteristic lengths of the internal space. As will become clear later, we will also need to know the relationship between the fields appearing in the supergravity action and their string counterparts which are associated to the massless vertex operators of the orbifold conformal field theory.
It would be interesting to perform these steps directly within the supergravity context, but this turns out to be a rather non-trivial task. Therefore, we take a different route and exploit the $`S`$-duality that relates our model to the heterotic string compactified on a torus $`T^4`$. In this way, we can carry out the reduction from ten to six dimensions on the heterotic side by using the standard machinery of toroidal compactification, and then translate the result in the type II theory by using the duality map. Notice that this short-cut is possible because this $`S`$-duality in six dimensions acts trivially on the moduli spaces of the two theories. This means that all couplings involving the 80 scalars and their dependence on the internal radii do not change in going from the heterotic to the type II string.
We now sketch the derivation of the type II effective action starting from the heterotic string compactified on a 4-torus whose low-energy action (in the string frame) is
$`S_\mathrm{h}`$ $`=`$ $`{\displaystyle \frac{(2\pi \sqrt{\alpha ^{}})^4}{2\kappa _{10}^2}}{\displaystyle d^6x\sqrt{detg^\mathrm{h}}\mathrm{e}^{2\phi ^\mathrm{h}}\left[(g^\mathrm{h})+4_\mu \phi ^\mathrm{h}^\mu \phi ^\mathrm{h}+\frac{1}{4}_\mu \varphi _a^\mathrm{h}^\mu (\varphi _a^\mathrm{h})^1\right]}`$ (3.2)
$``$ $`{\displaystyle \frac{(2\pi \sqrt{\alpha ^{}})^4}{4g_{10}^2}}{\displaystyle d^6x\sqrt{detg^\mathrm{h}}\mathrm{e}^{2\phi ^\mathrm{h}}F_{\mu \nu }^IF^{I\mu \nu }}+\mathrm{}`$
where the gravitational and the gauge couplings are related in the usual way $`\kappa _{10}/g_{10}=\sqrt{\alpha ^{}}/2`$. Note that (3.2) is only a subset of the whole action coming from the toroidal compactification since most of the original fields have been put to zero. We will discuss the validity of this truncation after having translated the action (3.2) in the type II language; its motivations will become clearer when the boundary action related to the non-BPS brane is discussed (that is, when the source term is taken into account). Here, we just notice that in (3.2) only the scalars related to dilatations of the compact dimensions have been explicitly written: since we consider a compactification where the torus is just a product of four orthogonal circles, these scalars are simply the four diagonal components of the ten-dimensional metric in the internal space, i.e. $`g_{aa}^\mathrm{h}\varphi _a^\mathrm{h}`$ with $`a=6,\mathrm{},9`$, and the dilaton $`\phi ^\mathrm{h}`$. The v.e.v.’s of these scalar fields are given by
$$\varphi _a^\mathrm{h}=\frac{(R_a^\mathrm{h})^2}{\alpha ^{}}\text{ and}\mathrm{e}^{\phi ^\mathrm{h}}=\frac{\alpha ^{}}{V_\mathrm{h}^{1/2}}g^{},$$
(3.3)
where $`g^{}`$ is the heterotic string coupling constant, and $`V_\mathrm{h}_{a=6}^9R_a^\mathrm{h}`$. In writing (3.2) we have also assumed that the gauge group is broken to $`U(1)^{16}`$ by suitable Wilson lines, and $`F_{\mu \nu }^I`$ denotes the surviving field strengths ($`I=1,\mathrm{},16`$).
The correspondence between the heterotic and the type II theories can be established by means of the following chain of dualities
$$\text{Het. }T^4\begin{array}{c}S\\ \end{array}\text{Type I }T^4\begin{array}{c}4T\\ \end{array}\text{IIB }\frac{T^4}{\mathrm{\Omega }_4}\begin{array}{c}S\\ \end{array}\text{IIB }\frac{T^4}{(1)^{F_L}_4}(\begin{array}{c}T\\ \end{array}\text{IIA }\frac{T^4}{_4}).$$
(3.4)
The last step is not really necessary for our purposes, but it is useful to have it in mind. In fact, for some practical calculations the type IIB picture is easier, whereas the geometrical interpretation is clearer for the orbifold of type IIA which is a singular limit of a smooth K3 manifold. In Figure 2 we briefly summarize how the bosonic fields of this theory emerge from three different points of view.
By following the various steps of (3.4), one can see how the parameters of the different theories are related to each other . For instance, the radii $`R_a^\mathrm{B}`$ and the string coupling constant $`g`$ in the type IIB orbifold are related to the corresponding quantities of the original heterotic theory as
$$R_a^\mathrm{B}=\frac{\sqrt{2}V_h^{\mathrm{\hspace{0.17em}1}/2}}{R_a^\mathrm{h}},g=\frac{\sqrt{2}V_h}{\alpha _{}^{}{}_{}{}^{2}}\frac{1}{g^{}}.$$
(3.5)
By performing a further $`T`$-duality in one of the four compact directions (say $`x^9`$) we can reach the type IIA orbifold, for which we have
$$R_a^\mathrm{A}=\frac{\sqrt{2}V_h^{\mathrm{\hspace{0.17em}1}/2}}{R_a^\mathrm{h}}\text{for}a9,R_9^\mathrm{A}=\frac{R_9^\mathrm{h}\alpha ^{}}{\sqrt{2}V_h^{\mathrm{\hspace{0.17em}1}/2}},g=\frac{R_9^\mathrm{h}V_h^{\mathrm{\hspace{0.17em}1}/2}}{\alpha _{}^{}{}_{}{}^{3/2}}\frac{1}{g^{}}.$$
(3.6)
The numerical coefficient in these relations have been fixed by checking that the masses of BPS objects take the expected values after a duality transformation. This is the same derivation used in ; however, here we do not follow their conventions and our results are slightly different. We keep the string length fixed, i.e. $`\alpha _\mathrm{h}^{}=\alpha _\mathrm{B}^{}=\alpha _\mathrm{A}^{}\alpha ^{}`$, and define the dilaton v.e.v. in the orbifold compactification as
$$\mathrm{e}^{\phi ^\mathrm{B}}=\frac{\sqrt{2}\alpha ^{}}{V_\mathrm{B}^{1/2}}g,$$
(3.7)
$`V_\mathrm{B}=_{a=6}^9R_a^\mathrm{B}`$, and similarly for the IIA case. The factor of $`\sqrt{2}`$ in the above definition is quite natural. In fact, as usual, the dilaton v.e.v. in a compactified theory contains the volume of the compact space. In the toroidal case one has $`\text{Vol}V`$, while in the orbifold the $`𝐙_2`$ identification halves the “physical” volume of the internal space: $`\text{Vol}V/2`$.
Recalling that the radii are related to the v.e.v. of the four scalar fields $`\varphi _a`$, and using (3.3) and (3.7), we can lift the above duality maps to the field level and obtain
$$\phi ^{\mathrm{A},\mathrm{B}}=\phi ^\mathrm{h},\varphi _a^{\mathrm{A},\mathrm{B}}=2\frac{\sqrt{\underset{b=6}{\overset{9}{}}\varphi _b^\mathrm{h}}}{\varphi _a^\mathrm{h}}.$$
(3.8)
Finally, by exploiting the invariance of the metric in the Einstein frame under $`S`$-duality, one can find the usual relation between the string-frame metrics:
$$g_{\mu \nu }^{\mathrm{A},\mathrm{B}}=\mathrm{e}^{2\phi ^\mathrm{h}}g_{\mu \nu }^\mathrm{h}.$$
(3.9)
Equipped with this machinery, we are ready to perform the $`S`$-duality on the heterotic action (3.2), and rewrite it in terms of IIB quantities. Using Eq.s (3.3), (3.8) and (3.9), we get
$`S_B`$ $`=`$ $`{\displaystyle \frac{(2\pi \sqrt{\alpha ^{}})^4}{2\kappa _{10}^2}}{\displaystyle d^6x\sqrt{detg^\mathrm{B}}\mathrm{e}^{2\phi ^\mathrm{B}}\left[(g_\mathrm{B})+4_\mu \phi ^\mathrm{B}^\mu \phi ^\mathrm{B}+\frac{1}{4}_\mu \varphi _a^\mathrm{B}^\mu (\varphi _a^\mathrm{B})^1\right]}`$ (3.10)
$``$ $`{\displaystyle \frac{(2\pi \sqrt{\alpha ^{}})^4}{4g_{10}^2}}{\displaystyle d^6x\sqrt{detg^\mathrm{B}}F_{\mu \nu }^IF^{I\mu \nu }}+\mathrm{}.`$
It is not difficult to see that this action is consistent with the perturbative string amplitudes that can be calculated in the IIB orbifold. However, in order to do this comparison one has first to rewrite the Lagrangian (3.10) in the Einstein frame by rescaling the metric $`g_{\mu \nu }=\mathrm{e}^{\stackrel{~}{\phi }}G_{\mu \nu }`$. Here $`\stackrel{~}{\phi }=(\phi \phi _{\mathrm{}})`$, where $`\phi _{\mathrm{}}`$ is the constant value of the dilaton at spatial infinity, which in our case is simply the v.e.v. defined in (3.7). Another rescaling is usually done on the gauge fields. In fact, in type II theory these are taken to be dimensionless regardless of the number of indices they carry, while on the heterotic side they have canonical dimensions. Thus, we introduce $`\stackrel{~}{F}=2\sqrt{\alpha ^{}}gF`$. Finally, for later convenience, we write $`\varphi _a^\mathrm{B}=\varphi _a^\mathrm{B}\mathrm{e}^{2\stackrel{~}{\eta }_a^\mathrm{B}}`$. In terms of these rescaled fields, the action (3.10) becomes
$`S_B`$ $`=`$ $`{\displaystyle \frac{1}{2\kappa _{orb}^2}}{\displaystyle }d^6x\sqrt{detG_\mathrm{B}}[(G_\mathrm{B})_\mu \stackrel{~}{\phi }^\mathrm{B}^\mu \stackrel{~}{\phi }^\mathrm{B}_\mu \stackrel{~}{\eta }_a^\mathrm{B}^\mu \stackrel{~}{\eta }_a^\mathrm{B}`$ (3.11)
$`{\displaystyle \frac{1}{4}}\left({\displaystyle \frac{2\kappa _{orb}^2}{g_{orb}^2}}\right)\mathrm{e}^{\stackrel{~}{\phi }^\mathrm{B}}\stackrel{~}{F}_{\mu \nu }^I\stackrel{~}{F}^{I\mu \nu }],`$
where
$$\kappa _{orb}^2=\frac{2\kappa _{10}^2g^2}{(2\pi )^4V_\mathrm{B}},\text{and}\frac{2\kappa _{orb}^2}{g_{orb}^2}=\frac{\alpha _{}^{}{}_{}{}^{2}}{4V_\mathrm{B}}.$$
(3.12)
Contrary to what happens in the heterotic theory, here the gravitational and gauge couplings have a different dependence on the radii of the compact space. This fact can be naturally understood by comparing the tree-level string amplitudes with the vertices derived from the field theory Lagrangian. As usual, the moduli dependence of the couplings is directly related to the “nature” of strings involved in the amplitudes. Since in the heterotic theory both the gauge and the gravitational fields are made out of the same kind of closed strings, it is natural that all couplings have a uniform dependence on the moduli. On the contrary, in the type II setup, the gauge fields we are looking at come from the twisted R-R sector. In this case the mode expansion of the string coordinates does contain momentum along the compact directions. Because of this, the twisted and untwisted vacua are differently normalized
$${}_{U}{}^{}p_1,n|p_2,m_{U}^{}\delta ^{(6)}(p_1p_2)\delta _{nm}V,{}_{T}{}^{}p_1,0|p_2,0_{T}^{}\delta ^{(6)}(p_1p_2),$$
and this difference reflects on the volume dependence of the gauge and gravitational couplings (3.12). It is interesting to remark that the gauge kinetic term in (3.11) becomes canonically normalized only for a particular value of the radii of the internal space, namely at $`R_a=\sqrt{\alpha ^{}/2}`$. As we will see in the next paragraph, this value plays a privileged role also from a different point of view.
We finally comment on the consistency of the truncation we did in deriving the effective action (3.11). As we already mentioned, we have considered only those massless fields which couple directly to the non-BPS D-particle we want to study, and switched off all other fields. One may ask whether this truncation is consistent with the equations of motion of the complete theory . In particular, problems can arise if in the full Lagrangian there are interaction terms which are only linear in one of the fields here disregarded, for instance the twisted NS-NS scalars, call them $`\xi `$. In this case, the equations of motion for $`\xi `$ will contain a term which is not automatically vanishing in our approximation, since it is independent of the field itself and a contradiction may arise. However, it can be checked with perturbative arguments that these terms cannot be present in the complete Lagrangian. For instance, the twisted NS-NS scalars $`\xi `$ are described by vertex operators containing a left and a right spin field of the internal space. Thus, they have a non-zero $`M`$-point amplitude, only if one of the other external vertices also contains these spin fields. However, this is not the case for the vertices corresponding to the fields we were considering. This shows that the above mentioned problems cannot arise with the twisted NS-NS scalars and that we can safely set them to zero. Similar world-sheet considerations can be done also for the R-R fields we switched off. For the untwisted scalars, like the compact off-diagonal entries of ten dimensional metric $`g_{ab}`$, it is easier to look directly at the part of the action where they appear. In fact, as well known, the terms containing only untwisted fields can be derived by means of the toroidal compactification from the original ten dimensional description: in the usual calculation, one has simply to put to zero all the fields which are odd under the orbifold operation. In this way one can check that the untwisted scalars appear at least quadratically in the action, and thus our truncation is consistent.
### 3.2 Boundary action
From the supergravity point of view, the boundary action is seen as the source term which one must add to the bulk action in order to describe a D-brane configuration. In string theory, this source can be efficiently represented by a boundary state and, in particular, by its overlaps with the massless closed string states. The boundary states $`|Dp`$ that describe the non-BPS D$`p`$-branes of the type II orbifolds, have been studied in detail in the literature . Here we just recall the main features that will be employed in the following section.
A first important point is that $`|Dp`$ is non-trivial only in the NS-NS untwisted and R-R twisted sectors of the theory . In particular, focusing on the non-BPS D-particle present in the type IIB$`/𝐙_2`$ orbifold, one has
$$|D0=|B0_{\mathrm{NS}\mathrm{NS},U}+|B0_{\mathrm{R}\mathrm{R},T_I},$$
(3.13)
where the index $`I=1,\mathrm{},16`$ in the twisted part indicates on which orbifold plane the D-brane is placed. The explicit form of the coherent states $`|B0`$ in (3.13) and their overlaps with perturbative closed string states have been studied in . This paper, however, uses fields and vertex operators that are essentially written in the framework of the original string theory in ten dimensions. But, if one wants to make contact with the six dimensional bulk theory discussed in the previous subsection, one must use more appropriate fields. These can be easily obtained by observing, for example, that the vertex operators that describe the six-dimensional graviton and dilaton have the same structure of their ten-dimensional analogues, but contain only oscillators with indices in the non-compact directions. The internal part of the ten-dimensional vertices describes instead the six dimensional scalars. Keeping this in mind, it is not difficult to find the relation between the canonically normalized fields $`h_{\mu \nu }^{}`$ and $`\phi ^{}`$ used in and the canonically normalized fields $`\widehat{h}_{\mu \nu }`$, $`\widehat{\eta }_a`$, $`\widehat{\phi }`$ to be used in six dimensions. This relation is
$`\widehat{h}_{\mu \nu }`$ $`=`$ $`h_{\mu \nu }^{}\eta _{\mu \nu }\left[{\displaystyle \frac{\sqrt{2}}{2}}\phi ^{}{\displaystyle \frac{1}{4}}{\displaystyle \underset{a}{}}h_{aa}^{}\right],`$
$`\widehat{\eta }_a`$ $`=`$ $`h_{aa}^{}{\displaystyle \frac{1}{2\sqrt{2}}}\phi ^{},`$ (3.14)
$`\widehat{\phi }`$ $`=`$ $`{\displaystyle \frac{3\sqrt{2}}{2}}\phi ^{}{\displaystyle \frac{1}{2}}{\displaystyle \underset{a}{}}h_{aa}^{}.`$
Writing the vertex operators associated to the hatted field in (3.2) and to the (canonically normalized) twisted R-R potential $`\widehat{A}_\mu ^I`$, we can compute the overlaps with the boundary state (3.13) and find the couplings between the bulk fields and the D-particle. Using the results of and the redefinitions (3.2), we get
$`D0|\widehat{h}_{\mu \nu }`$ $`=`$ $`M_0V_1\kappa _{orb}\widehat{h}_{00},D0|\widehat{\phi }={\displaystyle \frac{M_0V_1}{2}}\kappa _{orb}\widehat{\phi },`$ (3.15)
$`D0|\widehat{\eta }_a`$ $`=`$ $`{\displaystyle \frac{M_0V_1}{2}}\kappa _{orb}\widehat{\eta }_a,D0|\widehat{A}_\mu ^I=\sqrt{{\displaystyle \frac{8V_\mathrm{B}}{\alpha _{}^{}{}_{}{}^{2}}}}M_0V_1\kappa _{orb}\widehat{A}_0^I,`$
where
$$M_0=\frac{\sqrt{2}T_0}{(2\pi )^2\kappa _{orb}V_\mathrm{B}^{1/2}}=\frac{1}{\sqrt{\alpha ^{}}g}$$
(3.16)
is the mass of the non-BPS D-particle and $`V_1`$ is the (infinite) length of its world-line. The overlaps (3.15) represent the one-point functions of the bulk fields encoded in the boundary action. To write it we find convenient to use the same notation of the previous subsection, and not to work any more with the canonically normalized hatted fields. The relation between the latter and the fields appearing in the bulk Lagrangian (3.11) is simply given by
$$G_{\mu \nu }^\mathrm{B}=\eta _{\mu \nu }+2\kappa _{orb}\widehat{h}_{\mu \nu },\widehat{\phi }=\frac{\stackrel{~}{\phi }^\mathrm{B}}{\kappa _{orb}},\widehat{\eta }_a=\frac{\eta _a^\mathrm{B}}{\kappa _{orb}},\widehat{A}_\mu ^I=\frac{\stackrel{~}{A}_\mu ^I}{g_{orb}},$$
(3.17)
where $`\kappa _{orb}`$ and $`g_{orb}`$ are defined in (3.12). Then, it is easy to realize that the overlaps (3.15) are consistent with the following boundary action
$$S_{\mathrm{boundary}}=M_0𝑑\tau \mathrm{e}^{\frac{1}{2}\stackrel{~}{\phi }^\mathrm{B}\frac{1}{2}_a\stackrel{~}{\eta }_a^\mathrm{B}}\sqrt{G_{00}^\mathrm{B}}+M_0𝑑\tau \stackrel{~}{A}_0^I.$$
(3.18)
Of course, this action does not describe the complete world-volume dynamics of the non-BPS D-particle. In fact, in deriving it we considered only the trivial configuration for the fields related to open strings and we switched off all non-linear couplings with closed strings<sup>6</sup><sup>6</sup>6For instance, anomalous couplings, similar to those of usual D-branes, are present also for non-BPS branes . However, for our purposes the action (3.18) will be sufficient. Note that its untwisted part can be obtained also from the action (2.3) of the unstable non-BPS branes discussed in Section 2: one has just to perform a toroidal compactification and remove all fields that are odd under the orbifold projections. On the other hand, the twisted part of (3.18) accounts for the minimal coupling of a charged particle with its gauge field. Notice that the strength we found for this coupling is consistent with the $`S`$-dual interpretation of the D-particle. In fact, on the heterotic side this particle corresponds to a perturbative massive state with charge $`q_\mathrm{h}=2`$ under one of the 16 unbroken $`U(1)`$’s. Of course, this $`U(1)`$ charge does not change in the duality map and, thus, translating the heterotic result in type II language, one should find agreement with the overlap (3.15). This is indeed the case; in fact, taking into account the relation between the type II gauge fields $`\stackrel{~}{A}`$ and those of the heterotic theory introduced after Eq. (3.10), we can see that the gauge charge $`q_\mathrm{h}=2`$ becomes exactly the one that we read from (3.18).
So far we have discussed the boundary action for a single non-BPS brane. However, in order to form a macroscopic object one should consider a superposition of many branes. Only in this case, in fact, one can hope that the source creates a smooth classical geometry where it is possible to neglect string and loop corrections. Of course, the dynamics of many coincident branes is quite complicated and can radically change the couplings previously derived for a single object. For BPS configurations, this is not the case because the D-branes do not interact with each other. Thus the effect of the superposition of $`N`$ D-branes is simply to multiply by $`N`$ the strength of all couplings. In a non-supersymmetric setup, instead, the branes in general interact among them in a non-trivial way. However, as shown in , the non-BPS D-particles of the IIB$`/𝐙_2`$ orbifold enjoy two fundamental properties that make them similar to the usual BPS D-branes. In fact, the orbifold projection always kills the tachyon zero-mode and, if $`R_a^\mathrm{B}\sqrt{\alpha ^{}/2}`$ the winding excitations of the tachyon, which survive the projection, have a positive mass<sup>2</sup>. Thus, the non-BPS D-particles are stable. Moreover, if $`R_a^\mathrm{B}=\sqrt{\alpha ^{}/2}`$ (which is also the value where the gauge kinetic term in (3.11) becomes canonically normalized), the “would-be” tachyons become massless and an accidental Bose-Fermi degeneracy appears in the spectrum. Thus, at the critical radii the force between two non-BPS D-particles vanishes at one-loop , i.e.
$$\mathrm{\Gamma }=D0|𝒫|D0=0.$$
(3.19)
For this reason, it is natural to conjecture that, in this particular case, it is possible to describe $`N`$ non-BPS D-particles simply by taking the naïve sum of $`N`$ boundary states previously introduced to describe a single object, that is
$$|D0,N=N|D0.$$
(3.20)
We want to stress that the no-force condition (3.19) is clearly a necessary ingredient for this simplification to hold, but it is not sufficient to really prove the validity of the assumption (3.20). In fact, Eq. (3.19) proves the vanishing of the interaction only at one-loop (from the open-string point of view), and does not guarantee that a similar result occurs at higher loops. In other words, from Eq. (3.19) we can see that the interactions between $`N`$ D-particles vanish at the leading order in $`N`$, while we know that the supergravity description is reliable in the opposite regime, $`N\mathrm{}`$. Thus, we take Eq. (3.20) as a working hypothesis and, in the next section, we will check whether this leads to acceptable space-time configurations for the metric and the other fields.
## 4 The non-BPS D-particle solution
In the previous section we have shown that the action describing the dynamics of the fields emitted by a non-BPS D-particle in six dimensions, is given by the sum of Eq.s (3.11) and (3.18) which we rewrite here in a simplified notation
$`S`$ $`=`$ $`{\displaystyle \frac{1}{2\kappa _{orb}^2}}{\displaystyle d^6x\sqrt{detG}\left[(G)_\mu \phi ^\mu \phi _\mu \eta _a^\mu \eta _a\frac{1}{4}\mathrm{e}^\phi F_{\mu \nu }F^{\mu \nu }\right]}`$ (4.1)
$`M{\displaystyle 𝑑\tau \mathrm{e}^{\frac{1}{2}\phi \frac{1}{2}_a\eta _a}\sqrt{G_{00}}}+M{\displaystyle 𝑑\tau A_0}.`$
Notice that here we have fixed the compact volume to its critical value $`V_c=\alpha _{}^{}{}_{}{}^{\mathrm{\hspace{0.17em}2}}/4`$ where the no-force condition holds at one loop, and, according to our working hypothesis (3.20), we have put $`M=NM_0`$. The field equations that follow from this action are <sup>7</sup><sup>7</sup>7We use the static gauge $`X^0=\tau `$.
$$\frac{1}{\sqrt{detG}}_\mu \left(\sqrt{detG}G^{\mu \nu }_\nu \phi \right)\frac{1}{8}\mathrm{e}^\phi F_{\mu \nu }F^{\mu \nu }=\frac{1}{2}T(x)\delta ^5(\stackrel{}{x})$$
(4.2)
for the dilaton,
$$\frac{1}{\sqrt{detG}}_\mu \left(\sqrt{detG}G^{\mu \nu }_\nu \eta _a\right)=\frac{1}{2}T(x)\delta ^5(\stackrel{}{x})$$
(4.3)
for the 4 scalar fields,
$$_\mu \left(\sqrt{detG}\mathrm{e}^\phi F^{\mu \nu }\right)=2M\kappa _{orb}^2G_{\mathrm{\hspace{0.33em}\hspace{0.17em}0}}^\nu \delta ^5(\stackrel{}{x})$$
(4.4)
for the gauge field, and
$$R_{\mu \nu }\frac{1}{2}G_{\mu \nu }\left(_\mu \phi _\nu \phi \frac{1}{2}G_{\mu \nu }_\rho \phi ^\rho \phi \right)\left(_\mu \eta _a_\nu \eta _a\frac{1}{2}G_{\mu \nu }_\rho \eta _a^\rho \eta _a\right)$$
$$\frac{1}{2}\mathrm{e}^\phi \left(F_{\mu \rho }F_\nu ^\rho \frac{1}{4}G_{\mu \nu }F_{\rho \sigma }F^{\rho \sigma }\right)=G_{00}G_\mu ^{\mathrm{\hspace{0.33em}\hspace{0.17em}0}}G_\nu ^{\mathrm{\hspace{0.33em}\hspace{0.17em}0}}T(x)\delta ^5(\stackrel{}{x})$$
(4.5)
for the metric, where
$$T(x)=M\kappa _{orb}^2\mathrm{e}^{\frac{1}{2}\phi \frac{1}{2}_a\eta _a}\frac{\sqrt{G_{00}}}{\sqrt{detG}}$$
(4.6)
Our task is to find a solution to these equations describing a static and spherically symmetric non-BPS D-particle in which the fields depend only on the distance in transverse space, $`r`$. There are several ways to reach this goal. A first possibility is to build up the solution iteratively via a perturbative approach by expressing the various fields as series in powers of $`1/r^3`$ with arbitrary coefficients (recall that the usual $`1/r^{Dp3}`$ dependence of a $`p`$-brane in $`D`$ dimensions reduces in the present case to $`1/r^3`$). Inserting this Ansatz in the coupled system (4.2)–(4.5), one can determine all coefficients by solving the equations order by order in $`1/r^3`$. For example, up to third order in $`1/r^3`$ the solution looks like
$`\phi `$ $``$ $`{\displaystyle \frac{1}{4}}{\displaystyle \frac{Q}{r^3}}{\displaystyle \frac{1}{8}}\left({\displaystyle \frac{Q}{r^3}}\right)^2+{\displaystyle \frac{1}{24}}\left({\displaystyle \frac{Q}{r^3}}\right)^3+\mathrm{}`$ (4.7)
$`\eta _a`$ $``$ $`{\displaystyle \frac{1}{4}}{\displaystyle \frac{Q}{r^3}}+\mathrm{}`$ (4.8)
$`A_0`$ $``$ $`{\displaystyle \frac{Q}{r^3}}+{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{Q}{r^3}}\right)^2{\displaystyle \frac{1}{3}}\left({\displaystyle \frac{Q}{r^3}}\right)^3+\mathrm{}`$ (4.9)
$`G_{00}`$ $``$ $`1+{\displaystyle \frac{3}{4}}{\displaystyle \frac{Q}{r^3}}{\displaystyle \frac{21}{32}}\left({\displaystyle \frac{Q}{r^3}}\right)^2+{\displaystyle \frac{61}{128}}\left({\displaystyle \frac{Q}{r^3}}\right)^3+\mathrm{}`$ (4.10)
$`G_{ij}`$ $``$ $`\delta _{ij}\left[1+{\displaystyle \frac{1}{4}}{\displaystyle \frac{Q}{r^3}}{\displaystyle \frac{3}{32}}\left({\displaystyle \frac{Q}{r^3}}\right)^2+{\displaystyle \frac{5}{384}}\left({\displaystyle \frac{Q}{r^3}}\right)^3+\mathrm{}\right],`$ (4.11)
where
$$Q=\frac{2M\kappa _{orb}^2}{3\mathrm{\Omega }_4}Ng\alpha _{}^{}{}_{}{}^{\mathrm{\hspace{0.17em}3}/2}$$
(4.12)
and the indices $`i,j=1,\mathrm{},5`$ label the transverse directions.
This same result can be obtained via an alternative approach based on the use of the boundary state, which, as shown in , allows to find the asymptotic behavior of the various fields at large distance from the source. For the non-BPS D-particle, this method has been recently used in where the leading terms proportional to $`Q/r^3`$ have been obtained, and generalized for any $`p`$ in (provided the redefinitions (3.2) are taken into account). Actually, one can do more. Remembering that $`Q`$ is proportional to the ’t Hooft coupling $`\lambda Ng`$ (see Eq. (4.12)), the above expansions for large $`r`$ can also be interpreted as expansions in $`\lambda `$. Since in string theory different powers of $`\lambda `$ characterize open string diagrams of different topologies, the various terms in (4.7)-(4.11) can be associated to the one-point functions of massless bulk fields evaluated on world-sheets with an increasing number of boundaries. Specifically, the terms linear in $`Q`$ arise from one-point function on a disk diagram, the terms proportional to $`Q^2`$ from one-point functions on an annulus, and so on. However, for the purpose of finding the classical solution, it is not really necessary to perform such calculations in string theory, but it is sufficient to do them directly in the low-energy field theory described by the action (4.1). Here one has simply to compute (in configuration space) diagrams like the ones represented in Figure 3, where the couplings with the sources are determined by the boundary part of the effective action and the interaction vertices from its bulk part.
Despite its conceptual simplicity, this diagrammatic method requires calculations which become more and more cumbersome as one proceeds in the perturbative expansion. Nevertheless, it is useful because it clarifies the origin and the meaning of the various terms. In Appendix A we present the detailed calculations of the diagrams that contribute to the classical solution up to $`Q^2`$.
Let us now return to the perturbative expansions (4.7)-(4.11) of the D-particle solution. Differently from what expected , it turns out that it is possible to re-sum these series and present the fields in a closed form. Indeed we find
$`\phi `$ $`=`$ $`{\displaystyle \frac{1}{4}}\mathrm{ln}\left[1+\mathrm{sin}\left({\displaystyle \frac{Q}{r^3}}\right)\right]`$ (4.13)
$`\eta _a`$ $`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle \frac{Q}{r^3}}`$ (4.14)
$`A_0`$ $`=`$ $`1+{\displaystyle \frac{\mathrm{cos}\left(\frac{Q}{r^3}\right)}{1+\mathrm{sin}\left(\frac{Q}{r^3}\right)}}`$ (4.15)
$`G_{00}`$ $`=`$ $`\left[1+\mathrm{sin}\left({\displaystyle \frac{Q}{r^3}}\right)\right]^{3/4}`$ (4.16)
$`G_{ij}`$ $`=`$ $`\delta _{ij}\left[1+\mathrm{sin}\left({\displaystyle \frac{Q}{r^3}}\right)\right]^{1/4}.`$ (4.17)
It is not difficult, but rather tedious, to check that (4.13)-(4.17) is indeed a solution of the differential equations (4.2)-(4.5). In Appendix B we will provide all details for deriving the above expressions directly from the field equations (4.2)-(4.5), and present also their extension to the case of a generic $`p`$-brane.
We now discuss the properties of the solution (4.13)-(4.17). First of all, we observe that it is well-defined only for $`rQ^{1/3}`$. In fact, in the region $`r<Q^{1/3}`$ the dilaton, the gauge field and $`G_{00}`$ have branch cut singularities at
$`r_n=\left[{\displaystyle \frac{Q}{(3/2+2n)\pi }}\right]^{1/3}\mathrm{for}n=0,1,2,\mathrm{}.`$ (4.18)
Moreover, the scalar curvature $``$ diverges at these singular points. Clearly, when the curvature is big, the classical supergravity description is not any more reliable; moreover, since the singularities at $`r=r_n`$ are naked, the entire solution is unacceptable according to the cosmic censorship conjecture. This result seems to indicate that the prediction made in , about the possibility of having a classical description for stable non-BPS D-branes, does not hold, at least in the case that we have considered here. Let us be more specific about this fact. From the supergravity point of view, $`Q`$ is a free parameter, and solutions with different values of $`Q`$ are all on equal footing. Hence they are all unacceptable for the reason we have explained above. However, if we appeal to the underlying string theory, some crucial differences come into play. First of all, $`Q`$ is not any more a free parameter since it is related to the fundamental quantities of the microscopic theory as shown in (4.12). Moreover, in a string context one may expect a priori that the classical supergravity description can break down at distances of the order of the string scale, where the massless closed string states cease to be good probes for the geometry. Thus, at $`r\sqrt{\alpha ^{}}`$ stringy effects must be taken into account, and the entire supergravity approximation must be reconsidered. In other words, if the naked singularities of a classical solution are at distances smaller or equal to the string scale, no contradictions arise, and the supergravity solution can be accepted at larger distances.
This is what happens for a single D-particle, i.e. $`N=1`$. In fact, since the first singularity is at $`r_0=(2Q/3\pi )^{1/3}<Q^{1/3}`$, and, for $`N=1`$, $`Q^{1/3}g^{1/3}\sqrt{\alpha ^{}}`$, the region where the classical solution starts to have problems is inside the region in which stringy corrections are relevant. Thus, Eq.s (4.13)-(4.17) represent a valid classical solution associated to a stable non-BPS D-particle in six dimensions at distances much larger than the string scale. However, if we also want to justify the classical approximation and give a reason for disregarding loop corrections in the bulk, we must also to take the limit $`g0`$. In this case the fields produced by the D-brane become just small fluctuations around the trivial background and the only relevant terms are the leading ones in the $`Q/r^3`$ expansion.
Things are different for $`N1`$. In fact, as is clear from (4.12), if we increase the value of $`N`$, the parameter $`Q`$ becomes macroscopic (i.e. $`Q^{1/3}\sqrt{\alpha ^{}}`$), and the solution (4.13)-(4.17) exhibits naked singularities also in a region which is not affected by stringy fuzziness effects and where the classical approximation is reliable. Therefore, according to the cosmic censorship conjecture, the solution (4.13)-(4.17) must be rejected, and its source, namely a stack of many non-interacting D-particles, must be regarded as non physical.
We would like to stress that our conclusions are not in contradiction with the result of about the vanishing of the force between two non-BPS D-particles at critical radius. In fact, the result of is exact in $`\alpha ^{}`$ but perturbative in the ’t Hooft coupling $`\lambda Ng`$, and is due to a cancellation occurring at one loop because of an accidental Bose-Fermi degeneracy of the open string spectrum. The classical solution (4.13)-(4.17) is instead valid in a very different regime, since it is perturbative in $`\alpha ^{}`$ but exact in $`\lambda `$. Therefore, our result should be compared with that of only in the limit $`\lambda 0`$, and if we do this, we too find a vanishing force at the first order in $`\lambda `$. This can be easily seen by inserting the classical solution (4.13)-(4.17) into the boundary part of action (4.1). Expanding at first order in $`Q`$, and subtracting the vacuum energy, we find
$`S_{\mathrm{boundary}}`$ $`=`$ $`M{\displaystyle 𝑑\tau \mathrm{e}^{\frac{1}{2}\phi \frac{1}{2}_a\eta _a}\sqrt{G_{00}}}+M{\displaystyle 𝑑\tau A_0}`$ (4.19)
$``$ $`M{\displaystyle 𝑑\tau \frac{Q}{r^3}\left(\frac{1}{8}\frac{1}{2}\frac{3}{8}+1\right)}=0.`$
A similar calculation shows however, that the no-force condition is not satisfied at the next-to-leading order. It would be interesting to derive this result also from an open string computation at two loops (see for a recent discussion on this issue). In this case the non-trivial dynamics of the open strings living on the non-BPS D-brane becomes relevant, and due to the lack of supersymmetry one does not expect any cancellation to occur.
We conclude this section with a final observation. The fact that the stable non-BPS D-particles do not satisfy a no-force condition at all orders in $`\lambda `$, is also suggested by the behavior of $`G_{00}`$ around $`rQ^{1/3}`$. One can check that, before reaching the first singularity at $`r_0`$, the derivative of $`G_{00}`$ changes sign in $`r=r_G>r_0`$, see Figure 4.
As is well known, this fact indicates that something strange is happening: indeed, the gravitational force changes its sign at $`r=r_G`$, and the singularity located at $`r=r_0`$ (where $`G_{00}\mathrm{}`$) can be called repulson , because close to it, massive particles repel each other. These naked singularities have been recently studied in the context of string theory in where a mechanism for their resolution has been proposed. Since our configuration exhibits properties similar to those discussed in , it would be interesting to see whether the same kind of mechanism works also in our case and resolves the singularity we have found.
## 5 General solution and discussion
As we have stressed several times, the explicit form of the solution (4.13)-(4.17) crucially depends on the detailed knowledge of the source term, and in particular on the hypothesis (3.20). In this last section, we relax this assumption and consider a more general solution of the field equations (4.2)-(4.5). In fact, on general grounds, one may expect that there exist more complicated configurations of non-BPS D-particles which are stable and do not display any pathological behavior in the corresponding classical geometry. For example, one can think of a non-trivial bound state of non-BPS D-particles described by some complicated boundary state $`|B`$ and, correspondingly, by a boundary action which could differ from (3.18) even in its functional dependence on the fields.
To explore these possibilities, we therefore study the general solution of the differential equations (4.2)-(4.5) under the minimal amount of requirements. We simply ask that the solution describe an asymptotically flat geometry, be spherically symmetric in the transverse directions and also that all scalar fields have vanishing v.e.v. at infinity. Instead, we do not enforce any specific behavior on the leading terms in the large distance expansion. In fact, as we have explicitly seen in the previous section, these are directly related to the specific microscopic structure of the source. Our only hypotheses about it are therefore that it couples to the graviton, the dilaton, the scalars $`\eta _a`$ and to one twisted R-R gauge potential. As we have shown in Section 3, this set of fields defines a consistent truncation of the full six-dimensional supergravity theory.
Under these assumptions, we now study the most general solution of the field equations (4.2)-(4.5), after removing the source terms in the right hand sides. The resulting homogeneous equations can be analyzed by generalizing the methods of and the general solution can be written in a closed form in terms of elementary functions. It will depends on some integration constants (two for each equation); half of them are fixed by the general requirements discussed above, and the remaining ones are free parameters which can be associated to the microscopic structure of the source. In Appendix B we will provide the details to solve explicitely the homogeneous field equations (4.2)-(4.5); here we simply write the result. To do this, it is first convenient to introduce the functions
$`f_\pm (r)`$ $`=`$ $`1\pm x{\displaystyle \frac{Q}{r^3}}`$ (5.1)
$`X(r)`$ $`=`$ $`\alpha +\beta \mathrm{ln}\left({\displaystyle \frac{f_{}(r)}{f_+(r)}}\right)`$
where $`Q`$ is defined in (4.12), and $`\alpha ,\beta `$ and $`x`$ are constants. Then, the general solution gets the following form
$`e^{\eta _a}`$ $`=`$ $`\left({\displaystyle \frac{f_{}(r)}{f_+(r)}}\right)^\delta `$ (5.2)
$`e^{2\phi }`$ $`=`$ $`\left({\displaystyle \frac{\mathrm{cosh}X(r)+\gamma \mathrm{sinh}X(r)}{\mathrm{cosh}\alpha +\gamma \mathrm{sinh}\alpha }}\right)\left({\displaystyle \frac{f_{}(r)}{f_+(r)}}\right)^{\frac{3}{4}ϵ}`$ (5.3)
$`A_0`$ $`=`$ $`\sqrt{2(\gamma ^21)}\left({\displaystyle \frac{\mathrm{sinh}X(r)\left(\mathrm{cosh}\alpha +\gamma \mathrm{sinh}\alpha \right)}{\mathrm{cosh}X(r)+\gamma \mathrm{sinh}X(r)}}\mathrm{sinh}\alpha \right)`$ (5.4)
$`G_{00}`$ $`=`$ $`\left({\displaystyle \frac{\mathrm{cosh}X(r)+\gamma \mathrm{sinh}X(r)}{\mathrm{cosh}\alpha +\gamma \mathrm{sinh}\alpha }}\right)^{\frac{3}{2}}\left({\displaystyle \frac{f_{}(r)}{f_+(r)}}\right)^{\frac{3}{8}ϵ}`$ (5.5)
$`G_{ij}`$ $`=`$ $`\delta _{ij}\left({\displaystyle \frac{\mathrm{cosh}X(r)+\gamma \mathrm{sinh}X(r)}{\mathrm{cosh}\alpha +\gamma \mathrm{sinh}\alpha }}\right)^{\frac{1}{2}}\left(f_{}(r)\right)^{\frac{2}{3}\frac{1}{8}ϵ}\left(f_+(r)\right)^{\frac{2}{3}+\frac{1}{8}ϵ}`$ (5.6)
where
$$ϵ=\pm \frac{4}{3}\sqrt{43\beta ^212\delta ^2}.$$
(5.7)
As anticipated, this solution depends on some integration constants, namely $`\alpha `$, $`\beta `$, $`\gamma `$, $`\delta ,`$ and $`x`$, which can be fixed by specifying the form of the source term. For example, if we impose the boundary conditions corresponding to $`N`$ coincident and non-interacting D-particles, which was the physical situation considered in Section 4, namely if we require that at large distance the fields behave as
$`\phi {\displaystyle \frac{1}{4}}{\displaystyle \frac{Q}{r^3}}+\mathrm{},\eta _a{\displaystyle \frac{1}{4}}{\displaystyle \frac{Q}{r^3}}+\mathrm{},A_0{\displaystyle \frac{Q}{r^3}}+\mathrm{},`$ (5.8)
$`G_{00}1+{\displaystyle \frac{3}{4}}{\displaystyle \frac{Q}{r^3}}+\mathrm{},G_{ij}\delta _{ij}\left(1+{\displaystyle \frac{1}{4}}{\displaystyle \frac{Q}{r^3}}+\mathrm{}\right),`$
one finds that the constants $`\beta `$, $`\gamma `$, $`\delta ,`$ and $`x`$ must be chosen as
$`x\beta =\pm {\displaystyle \frac{\mathrm{i}}{4}},\gamma ={\displaystyle \frac{\mathrm{sinh}2\alpha \pm \mathrm{i}}{\mathrm{cosh}2\alpha }},x\delta ={\displaystyle \frac{1}{8}}\mathrm{and}x0`$ (5.9)
Inserting these values into (5.2)-(5.6), one can easily check that the solution given in Eq.s (4.13)-(4.17) is recovered and the $`\alpha `$-dependence drops out. Notice, in particular, that even if the vanishing value for $`x`$ renders both $`f_{}`$ and $`f_+`$ trivial, the function $`X(r)`$ and the scalars $`\eta _a`$ do not become constant since the products $`x\beta `$ and $`x\delta `$ are non zero. On the other hand, the fact that $`x\beta `$ is purely imaginary, makes all hyperbolic functions become periodic.
The structure of the general solution (5.2)-(5.6) clearly indicates that different choices of the integration constants do not necessarily yield a classical geometry with a pathological behavior. However, when one regards the supergravity as the low-energy description of string theory, one should ask which of all possible choices in the classical context have some physical interpretation from the string viewpoint. A first very natural possibility is to consider an orbifold compactification with an internal volume $`VV_c`$ (recall that at $`V_c=\alpha _{}^{}{}_{}{}^{2}/4`$ the non-BPS D-particle becomes “extremal”, in some sense). In fact, provided that $`V>V_c`$, the D-particle remains a stable configuration , even if it does not exhibit any more the Fermi-Bose degeneracy which, at first order, was responsible for the no-force condition. However, since this condition does not hold at higher orders, it is not necessary to focus on the critical volume any more, and one can hope that the departure from extremality will eventually lead to develop an event horizon which, hiding all singularities in a casual disconnected region from the physical space, would make the solution (5.2)-(5.6) free of singularities.
Another possibility could be to consider a stable bound state obtained by displacing along the transverse directions the stack of D-branes in a sphere of radius $`rQ^{1/3}`$. This kind of mechanism for resolving the repulsive singularities which we have found in our solution, is similar to the one advocated in , where these problems have been discussed in great detail. In our non-BPS situation, this option, however, deserves further investigation. Of course, one could also imagine more exotic possibilities that rely on very different settings, with different source terms and more bulk fields that couple to them. These configurations would need a different truncation of the original $`(1,1)`$ supergravity theory discussed in Section 3, and are clearly not accomplished by the solution (5.2)-(5.6).
The highly non-trivial role that stable non-BPS D-branes could play in deeper understanding of non-perturbative dualities in string theory and eventually on non-supersymmetric versions of the AdS/CFT correspondence, clearly makes quite challenging to find some positive answers to these open problems.
Acknowledgments
We would like to thank L. Andrianopoli, M. Billò, L. Gallot, A. Liccardo, I. Pesando and M. Trigiante for very useful discussions. M.F., A.L. and R.R. thank NORDITA, and M.B. the Physics Institute of the University of Neuchâtel for kind hospitality. M.B. acknowledges support by INFN and R.R. by the Fond National Suisse.
Appendix A
In this appendix we present the diagrammatic calculations of the leading and next-to-leading terms in the large distance expansion of the fields emitted by a non-BPS D-particle. To do this, we first rewrite the bulk part of the action (4.1) in terms of canonical normalized fields (see Eq. (3.17)) and get
$$S_{\mathrm{bulk}}=d^6x\sqrt{detG}\left[\frac{1}{2\kappa _{orb}^2}(G)\frac{1}{2}_\mu \widehat{\phi }^\mu \widehat{\phi }\frac{1}{2}_\mu \widehat{\eta }_a^\mu \widehat{\eta }_a\frac{1}{4}\mathrm{e}^{\kappa _{orb}\widehat{\phi }}\widehat{F}^{\mathrm{\hspace{0.17em}2}}\right]$$
(A.1)
with $`G_{\mu \nu }=\eta _{\mu \nu }+2\kappa _{orb}\widehat{h}_{\mu \nu }`$. Expanding (A.1) in $`\kappa _{orb}`$, we get
$$S_{\mathrm{bulk}}=S_0+\kappa _{orb}S_I+𝒪(\kappa _{orb}^2)$$
(A.2)
where $`S_0`$ is the free action and
$$S_I=S_{\phi \phi h}+S_{AAh}+S_{\eta \eta h}+S_{hhh}.$$
(A.3)
The four terms in $`S_I`$ describe respectively the interaction of a graviton with two dilatons, of a graviton with two gauge fields, of a graviton with two scalars, and the coupling among three gravitons.
The interactions of the bulk fields with the D-brane are encoded in the boundary part of the action (4.1), which, at the linearized level, is
$$S_{\mathrm{boundary}}=d^6x\left(J^{\mu \nu }\widehat{h}_{\mu \nu }+J\widehat{\phi }+J_a\widehat{\eta }_a+J^\mu \widehat{A}_\mu \right)$$
(A.4)
where the currents are
$`J_{\mu \nu }(x)=\kappa _{orb}M\eta _{\mu 0}\eta _{\nu 0}\delta ^5(\stackrel{}{x}),J(x)=J_a(x)={\displaystyle \frac{\kappa _{orb}M}{2}}\delta ^5\left(\stackrel{}{x}\right),`$
$`J_\mu (x)=\sqrt{2}\kappa _{orb}M\eta _{\mu 0}\delta ^5(\stackrel{}{x}).`$ (A.5)
The quadratic (and higher order) terms of the boundary action will not be needed in our calculations since they give rise only to tadpole diagrams which vanish in dimensional regularization.
In this theory, the one-point function of a generic bulk field $`\widehat{\mathrm{\Psi }}(x)`$ is given by
$$\left[D\widehat{h}D\widehat{\phi }D\widehat{A}D\widehat{\eta }\right]\widehat{\mathrm{\Psi }}(x)\mathrm{e}^{\mathrm{i}S_0}\mathrm{e}^{\mathrm{i}S_I}\mathrm{e}^{\mathrm{i}S_{\mathrm{boundary}}}\widehat{\mathrm{\Psi }}(x)\mathrm{e}^{\mathrm{i}S_I}\mathrm{e}^{\mathrm{i}S_{\mathrm{boundary}}}.$$
(A.6)
Expanding the two exponentials, we generate a perturbative series whose various terms correspond to diagrams that contain a different number of bulk and boundary interactions. The first two diagrams in this series are represented in Figure 3 and describe, respectively, the leading and next-to-leading terms in the large distance expansion of the classical bulk fields. In particular, the leading contribution is obtained from (A.6) by neglecting $`S_I`$ and expanding at first order the exponential containing $`S_{\mathrm{boundary}}`$. Applying this procedure to the dilaton, we find that the leading contribution to its one-point function is
$$\widehat{\phi }^{(1)}(x)=\mathrm{i}d^6y<\widehat{\phi }(x)J(y)\widehat{\phi }(y)>=\frac{\kappa _{orb}M}{2}\frac{d^5k}{(2\pi )^5}\frac{\mathrm{e}^{\mathrm{i}kx}}{k^2}.$$
(A.7)
By using the following expression for the Fourier transform
$$\frac{d^dk}{(2\pi )^d}\frac{\mathrm{e}^{\mathrm{i}kx}}{k^{2\alpha }}=\frac{2\alpha }{2^{2\alpha }\pi ^{d/2}}\frac{\mathrm{\Gamma }\left(\frac{d}{2}+1\alpha \right)}{\mathrm{\Gamma }\left(1+\alpha \right)}\frac{1}{(d2\alpha )}\frac{1}{|x|^{d2\alpha }},$$
(A.8)
we easily see that (A.7) becomes
$$\widehat{\phi }^{(1)}(x)=\frac{1}{\kappa _{orb}}\frac{1}{4}\frac{Q}{r^3}$$
(A.9)
where $`Q`$ is the parameter defined in Eq. (4.12). This is precisely the leading term at large distance of the dilaton produced by the non-BPS D-particle (see Eq. (4.7)). In a similar manner we can compute the asymptotic behavior of the other bulk fields and find complete agreement with the results reported in Eq.s (4.8)-(4.11). We now compute the next-to-leading order of the one-point function (A.6). This is obtained by expanding the exponential of $`S_I`$ at first order and the exponential of $`S_{\mathrm{boundary}}`$ at second order.
Applying this procedure to the dilaton, we find two contributions corresponding to the diagrams in Figure 5, so that we can write
$$\widehat{\phi }^{(2)}(x)=A_\phi ^{\phi AA}(x)+A_\phi ^{\phi hh}(x).$$
(A.10)
The first term, due to the coupling of a dilaton with two gravitons (see Figure 5a), is equal to
$`A_\phi ^{\phi hh}(x)`$ $`=`$ $`\mathrm{i}^3{\displaystyle d^6yd^6zd^6u}`$
$`\widehat{\phi }(x)\widehat{\phi }(y)J(y)_\mu \widehat{\phi }(z)_\nu \widehat{\phi }(z)\left(\widehat{h}^{\mu \nu }(z){\displaystyle \frac{1}{2}}\widehat{h}_\tau ^\tau (z)\eta ^{\mu \nu }\right)\widehat{h}_{\rho \sigma }(u)J^{\rho \sigma }(u).`$
By performing all contractions and using the explicit expressions for the propagators, it is not difficult to see that $`A_\phi ^{\phi hh}(x)=0`$. Moreover, it is interesting to notice that this result holds in any space-time dimension. The second term in (A.10) corresponds to the diagram of Figure 5b that is given by
$`A_\phi ^{\phi AA}(x)`$ $`=`$ $`{\displaystyle \frac{\mathrm{i}\kappa _{orb}}{8}}{\displaystyle d^6yd^6zd^6u\widehat{\phi }(x)\widehat{\phi }(y)\widehat{F}^{\mathrm{\hspace{0.17em}2}}(y)\widehat{A}_\mu (z)J^\mu (z)\widehat{A}_\nu (u)J^\nu (u)}`$ (A.12)
$`=`$ $`\kappa _{orb}^3M^2{\displaystyle \frac{d^5k}{(2\pi )^5}\frac{\mathrm{e}^{\mathrm{i}kx}}{k^2}\frac{d^5p}{(2\pi )^5}\frac{p(p+k)}{p^2(k+p)^2}}.`$
The second integral in (A.12) can be easily evaluated with standard techniques. For the sake of generality we give the result of the previous integral for an arbitrary value $`d`$ of the number of transverse directions. One gets
$$\frac{d^dp}{(2\pi )^d}\frac{p(p+k)}{k^2p^2(k+p)^2}=\frac{1}{2}\frac{1}{(4\pi )^{d/2}}B(\frac{d}{2}1,\frac{d}{2}1)\mathrm{\Gamma }\left(2\frac{d}{2}\right)(k^2)^{d/22}.$$
(A.13)
Inserting this result in (A.12) and using (A.8) for $`\alpha =1/2`$ and $`d=5`$, we finally get
$$\widehat{\phi }^{(2)}(x)=\frac{1}{\kappa _{orb}}\frac{1}{8}\left(\frac{Q}{r^3}\right)^2$$
(A.14)
which agrees with the next-to-leading term of Eq. (4.7).
The same calculations that lead to $`A_\phi ^{\phi hh}(x)=0`$, also imply that $`A_\eta ^{\eta hh}(x)=0`$. In fact, the gravitational couplings of the scalar fields $`\widehat{\eta }_a`$ are the same as those of the dilaton, and therefore also this diagram does not contribute. On the other hand, since $`\widehat{\eta }_a`$ does not couple to any other bulk field, from the vanishing of $`A_\eta ^{\eta hh}(x)`$ we can deduce that $`\widehat{\eta }_a`$ does not receive any correction at the next-to-leading order. Actually, also the higher orders for these fields are vanishing and thus the leading term for $`r\mathrm{}`$ already gives the exact result (see Eq. (4.14)).
With this same method we can compute the next-to-leading term for the gauge field $`\widehat{A}_0`$. This is the sum of two terms which arise from the bulk interaction of two gauge fields with a dilaton and a graviton respectively, namely
$$\widehat{A}_{0}^{}{}_{}{}^{(2)}(x)=A_A^{A\phi A}(x)+A_A^{AhA}(x).$$
(A.15)
Finally, the next-to-leading term of the graviton is produced by the bulk interactions involving three gravitons, two dilatons and one graviton, two gauge fields and one graviton, and two scalars and one graviton, that is
$$\widehat{h}_{\mu \nu }^{(2)}(x)=A_{\mu \nu }^{hhh}(x)+A_{\mu \nu }^{h\phi \phi }(x)+A_{\mu \nu }^{hAA}(x)+A_{\mu \nu }^{h\eta \eta }(x).$$
(A.16)
All terms in (A.15) and (A.16) can be computed following the procedure outlined before, and after some lengthy algebra one gets precisely the next-to-leading behavior of the gauge field and the metric written in Eq.s (4.9)-(4.11).
Appendix B
In this appendix we explicitly derive the non-BPS D-particle solution (4.13)-(4.17), and the most general one presented in Eq.s (5.2)-(5.6). For the sake of generality we start from a $`D`$-dimensional action containing the metric, the dilaton, the scalars $`\eta _a`$ and a $`(p+1)`$-form R-R potential with $`p<D3`$. The case of the non-BPS D-particle, considered in Section 2, can be obtained by taking in all our equations $`p=0`$ and $`D=6`$. However, our equations can also be used to derive the non-BPS solution in $`D=10`$ discussed in Section 2. Actually they can also be used for the usual BPS D-branes in ten dimensions. We start from the following action
$$S=S_{\mathrm{bulk}}+S_{\mathrm{boundary}}$$
(B.1)
where
$$S_{\mathrm{bulk}}=\frac{1}{2\kappa _{orb}^2}d^Dx\sqrt{detG}\left[(G)_\mu \phi ^\mu \phi _\mu \eta _a^\mu \eta _a\frac{1}{2(p+2)!}\mathrm{e}^{a\phi }F_{p+2}^2\right]$$
(B.2)
and
$$S_{\mathrm{boundary}}=Md^{p+1}\xi \mathrm{e}^{\frac{a}{2}\phi \frac{1}{2}_a\eta _a}\sqrt{detG_{\alpha \beta }}+MA_{p+1}$$
(B.3)
where $`\kappa _{orb}`$ has been defined in (3.12), while $`M=NM_p`$ with
$$M_p=\frac{\sqrt{2}T_p}{(2\pi )^2\kappa _{orb}V^{1/2}}.$$
(B.4)
As mentioned above, we treat simultaneously the cases of non-BPS branes in both $`D=6`$ and $`D=10`$. This can be done by taking the constant $`a`$ to be given by
$$a=\frac{D42p}{\sqrt{D2}}.$$
(B.5)
Clearly, if $`D=10`$ in (B.4) we have to put $`\kappa _{10}`$ in place of $`\kappa _{orb}`$ and delete the factor of $`(2\pi )^2V^{1/2}`$. Moreover, in the ten dimensional case there are no scalars $`\eta _a`$ and in the case of the non-BPS branes there is no R-R field. Finally for the BPS branes there is no factor of $`\sqrt{2}`$ in the brane tension (B.4).
By varying the action (B.1), we get the equations of motion for the various fields. In particular, we have
$$\frac{1}{\sqrt{detG}}_\mu \left(\sqrt{detG}G^{\mu \nu }_\nu \phi \right)\frac{a}{4}\mathrm{e}^{a\phi }\frac{1}{(p+2)!}F_{p+2}^2=\frac{a}{2}T(x)\delta ^d(\stackrel{}{x})$$
(B.6)
for the dilaton,
$$\frac{1}{\sqrt{detG}}_\mu \left(\sqrt{detG}G^{\mu \nu }_\nu \eta _a\right)=\frac{1}{2}T(x)\delta ^d(x)$$
(B.7)
for the scalars and
$$_{\mu _1}\left(\sqrt{detG}\mathrm{e}^{a\phi }G^{\mu _1\nu _1}\mathrm{}G^{\mu _{p+2}\nu _{p+2}}\frac{F_{\nu _1\mathrm{}\nu _{p+2}}}{(p+1)!}\right)=2M\kappa _{orb}^2G_0^{\mu _2}\mathrm{}G_p^{\mu _{p+2}}\delta ^d(\stackrel{}{x})$$
(B.8)
for the R-R field. Finally, the Einstein equations for the metric can be written in a simple form by first evaluating their trace, and then plugging it back into the original equations, obtaining
$`_{\mu \nu }{\displaystyle \frac{\mathrm{e}^{a\phi }}{2(p+2)!}}\left[(p+2)F_{\mu \mu _2\mathrm{}\mu _{p+2}}F_\nu ^{\mu _2\mathrm{}\mu _{p+2}}G_{\mu \nu }{\displaystyle \frac{(1+p)}{4}}F_{p+2}^2\right]`$ (B.9)
$`_\mu \phi _\nu \phi _\mu \eta _a_\nu \eta _a=T(x)\left(G_{\mu \alpha }G_{\nu \beta }G^{\alpha \beta }{\displaystyle \frac{p+1}{4}}G_{\mu \nu }\right)\delta ^d(\stackrel{}{x})`$
where
$$T(x)=M\kappa _{orb}^2\mathrm{e}^{\frac{a}{2}\phi \frac{1}{2}_a\eta _a}\frac{\sqrt{detG_{\alpha \beta }}}{\sqrt{detG}}.$$
(B.10)
We now solve the previous equations using the following Ansatz for the metric
$$ds^2=B^2(r)\eta _{\alpha \beta }dx^\alpha dx^\beta +F^2(r)\delta _{ij}dx^idx^j$$
(B.11)
where $`\alpha ,\beta =0,\mathrm{},p`$ and $`i,j=p+1,\mathrm{},dDp1`$, and assuming that all other fields are functions only of the radial coordinate $`r`$. Under these assumptions, the dilaton equation (B.6) becomes
$`{\displaystyle \frac{1}{r^{d1}}}\left(r^{d1}B^{p+1}F^{d2}\phi ^{}\right)^{}+{\displaystyle \frac{a}{4}}\mathrm{e}^{a\phi }F^{d2}B^{p1}\left(A_{01\mathrm{}p}^{}\right)^2`$
$`={\displaystyle \frac{a}{2}}B^{p+1}F^dT(x)\delta ^d(\stackrel{}{x}),`$ (B.12)
the scalar equation (B.7) becomes
$$\frac{1}{r^{d1}}\left(r^{d1}B^{p+1}F^{d2}\eta _a^{}\right)^{}=\frac{1}{2}B^{p+1}F^dT(x)\delta ^d(\stackrel{}{x}),$$
(B.13)
while R-R field equation (B.8) becomes
$$\frac{1}{r^{d1}}\left(r^{d1}B^{p1}F^{d2}\mathrm{e}^{a\phi }A_{01\mathrm{}p}^{}\right)^{}=2M\kappa _{orb}^2\delta ^d(\stackrel{}{x})$$
(B.14)
where $`{}_{}{}^{}d/dr`$. Finally, from the Einstein equations (5) we get
$`F^2\{\xi ^{\prime \prime }\left(\mathrm{log}F\right)^{\prime \prime }{\displaystyle \frac{d1}{r}}\left(\mathrm{log}F\right)^{}(p+1)\left[\left(\mathrm{log}B\right)^{}\right]^2+\xi ^{}\left(\mathrm{log}F\right)^{}`$
$`(d2)\left[\left(\mathrm{log}F\right)^{}\right]^2\}F^2(\phi ^{})^2F^2{\displaystyle }_a(\eta _a^{})^2`$
$`{\displaystyle \frac{\mathrm{e}^{a\phi }}{2}}{\displaystyle \frac{d2}{D2}}F^2B^{2(p+1)}\left(A_{01\mathrm{}p}^{}\right)^2={\displaystyle \frac{p+1}{D2}}T(x)\delta ^d(\stackrel{}{x})`$ (B.15)
for the components $`R_r^r`$,
$`F^2\left\{(\mathrm{log}B)^{\prime \prime }(\mathrm{log}B)^{}\left[\xi ^{}+{\displaystyle \frac{d1}{r}}\right]\right\}`$
$`{\displaystyle \frac{\mathrm{e}^{a\phi }}{2}}\left({\displaystyle \frac{d2}{D2}}\right)F^2B^{2(p+1)}\left(A_{01\mathrm{}p}^{}\right)^2={\displaystyle \frac{d2}{D2}}T(x)\delta ^d(\stackrel{}{x})`$
(B.16)
for the components $`R_\alpha ^\alpha `$, and
$`F^2\left[(\mathrm{log}F)^{\prime \prime }{\displaystyle \frac{d1}{r}}(\mathrm{log}F)^{}(\mathrm{log}F)^{}\xi ^{}{\displaystyle \frac{\xi ^{}}{r}}\right]`$
$`{\displaystyle \frac{\mathrm{e}^{a\phi }}{2}}{\displaystyle \frac{p+1}{D2}}B^{2p2}\left(A_{01\mathrm{}p}^{}\right)^2={\displaystyle \frac{p+1}{D2}}T(x)\delta ^d(\stackrel{}{x})`$ (B.17)
for the components $`R_{\overline{\alpha }}^{\overline{\alpha }}`$ where the index $`\overline{\alpha }`$ corresponds to the angular variables. In these equations we have introduced the function
$$\xi =(p+1)\mathrm{log}B+(d2)\mathrm{log}F.$$
(B.18)
Following for example , we now multiply Eq. (B.16) by a factor of $`(p+1)`$ and Eq. (B.17) by a factor of $`(d2)`$, and then sum the two expressions. In this way we see that the function $`\xi `$ obeys a simple differential equation, namely
$$\left[r^{2d3}\left(\mathrm{e}^\xi \right)^{}\right]^{}=0.$$
(B.19)
This is the Laplace equation in $`2d1`$ dimensions and its most general solution can be written as
$$e^\xi =\widehat{C}+C\left(\frac{Q_p}{r^{d2}}\right)^2$$
(B.20)
where $`\widehat{C}`$ and $`C`$ are arbitrary constants, and, for later convenience, we have introduced the dimensionful quantity
$$Q_p=\frac{2\kappa _{orb}^2M_p}{(d2)\mathrm{\Omega }_{d1}}.$$
(B.21)
In order to have an asymptotically flat metric, we must choose $`\widehat{C}=1`$, and thus we can write
$$\mathrm{e}^\xi B^{p+1}F^{d2}=f_{}(r)f_+(r)$$
(B.22)
with
$$f_\pm (r)=1\pm x\frac{Q_p}{r^{d2}};x^2=C.$$
(B.23)
Inserting Eq. (B.22) into Eq.s (B.12)-(B.17), we get
$$\frac{\mathrm{e}^\xi }{r^{d1}}\left(r^{d1}\mathrm{e}^\xi \phi ^{}\right)^{}+\frac{a}{4}\mathrm{e}^{a\phi }B^{2(p+1)}\left(A_{01\mathrm{}p}^{}\right)^2=\frac{a}{2}F^2T(x)\delta ^d(\stackrel{}{x})$$
(B.24)
for the dilaton,
$`{\displaystyle \frac{\mathrm{e}^\xi }{r^{d1}}}\left(r^{d1}\mathrm{e}^\xi \eta _a^{}\right)^{}={\displaystyle \frac{1}{2}}F^2T(x)\delta ^d(\stackrel{}{x})`$ (B.25)
for the scalars $`\eta _a`$,
$`{\displaystyle \frac{1}{r^{d1}}}\left(r^{d1}\mathrm{e}^{a\phi }B^{(p+1)}F^{d2}A_{01\mathrm{}p}^{}\right)^{}=2M\kappa _{orb}^2\delta ^d(\stackrel{}{x})`$ (B.26)
for the R-R field, while Eq. (B.16) can be rewritten as
$`{\displaystyle \frac{e^\xi }{r^{d1}}}\left[r^{d1}\mathrm{e}^\xi (\mathrm{log}B)^{}\right]^{}{\displaystyle \frac{\mathrm{e}^{a\phi }}{2}}{\displaystyle \frac{d2}{D2}}B^{2(1+p)}\left(A_{01\mathrm{}p}^{}\right)^2`$
$`={\displaystyle \frac{d2}{D2}}F^2T(x)\delta ^d(\stackrel{}{x}).`$ (B.27)
Multiplying Eq. (B.24) by $`2(d2)/(D2)`$ and Eq. (B.27) by $`a`$, and then summing the resulting expressions we get:
$$\left(r^{d1}\mathrm{e}^\xi Y^{}\right)^{}=0$$
(B.28)
where
$$Y\frac{D2}{d2}\mathrm{log}B+\frac{2}{a}\phi .$$
(B.29)
The solution of this equation is
$$\mathrm{e}^YB^{(D2)/(d2)}\mathrm{e}^{2\phi /a}=\left(\frac{f_{}(r)}{f_+(r)}\right)^ϵ$$
(B.30)
where $`ϵ`$ is an arbitrary integration constant <sup>8</sup><sup>8</sup>8Here and in the following, when $`a=0`$ (e.g. $`p=1`$ in $`D=6`$ and $`p=3`$ in $`D=10`$) the equations are ill-defined. However, their general solutions are valid also in these cases.. Actually the most general solution of Eq. (B.28) admits an additional arbitrary constant which, however, we have fixed by requiring that $`Y`$ vanish for $`r\mathrm{}`$. Using Eq. (B.17) in Eq. (B.15), and expressing $`B`$ and $`F`$ in terms of $`\xi ,Y`$ and $`\phi `$ we can rewrite Eq. (B.15) as follows
$$\frac{\mathrm{e}^{a\phi }}{2}B^{2(p+1)}\left(A_{01\mathrm{}p}^{}\right)^2=\frac{4}{D2}\left[\frac{a^2(D2)}{4}+(p+1)(d2)\right]\left(\frac{\phi ^{}}{a}\right)^2+$$
$$+\underset{a}{}\left(\eta _a^{}\right)^2+\xi ^{\prime \prime }\frac{1}{d2}\left(\xi ^{}\right)^2\frac{\xi ^{}}{r}+\frac{(d2)(p+1)}{D2}\left[\left(Y^{}\right)^2Y^{}\left(\frac{4\phi ^{}}{a}\right)\right].$$
(B.31)
Let us start by examining the case of a non-BPS D-brane in $`D=6`$. When we use Eq. (B.31) in the dilaton equation (B.24), the latter becomes
$$\frac{\mathrm{e}^{\xi 2\phi /(1p)+(d2)(p+1)Y/2}}{r^{d1}}\left[r^{d1}\mathrm{e}^{\xi (d2)(p+1)Y/2}\left(\mathrm{e}^{\frac{2}{1p}\phi }\right)^{}\right]^{}+\underset{a}{}\left(\eta _a^{}\right)^2$$
$$\frac{d1}{d2}(\xi ^{})^22(d1)\frac{\xi ^{}}{r}+\frac{(d2)(p+1)}{4}\left(Y^{}\right)^2=F^2T(x)\delta ^{(5p)}(\stackrel{}{x}).$$
(B.32)
Using Eq.s (B.22) and (B.30) with $`D=6`$ for the functions $`\xi `$ and $`Y`$, and the fact that the scalar fields $`\eta _a`$ satisfy the same homogeneous equation as $`Y`$ and therefore are given by
$$\mathrm{e}^{\eta _a}=\left(\frac{f_{}(r)}{f_+(r)}\right)^\delta ,$$
(B.33)
one can see that the dilaton equation (B.32) becomes
$`{\displaystyle \frac{\mathrm{e}^{\xi 2\phi /(1p)+(d2)(p+1)Y/2}}{r^{d1}}}\left[r^{d1}\mathrm{e}^{\xi (d2)(p+1)Y/2}\left(\mathrm{e}^{\frac{2}{1p}\phi }\right)^{}\right]^{}+C{\displaystyle \frac{d2}{r^2}}\left({\displaystyle \frac{Q_p}{r^{d2}}}\right)^2\mathrm{e}^{2\xi }`$
$`\times [16(d2)\delta ^2+4(d1)(d2)^2(p+1)ϵ^2]=F^2T(x)\delta ^{(5p)}(\stackrel{}{x}).`$
(B.34)
This must be considered together with the equation (B.26) for the R-R field, which becomes
$$\frac{1}{r^{d1}}\left(r^{d1}\mathrm{e}^{\xi (p+1)(3p)Y/2+\frac{4}{1p}\phi }A_{01\mathrm{}p}^{}\right)=2M\kappa _{orb}^2\delta ^{(5p)}(\stackrel{}{x}).$$
(B.35)
Finally, using Eq.s (B.29) and (B.28), Eq. (B.27) can be rewritten as follows
$$\frac{\mathrm{e}^\xi }{r^{d1}}\left[r^{d1}\mathrm{e}^\xi \left(\frac{2\phi ^{}}{1p}\right)\right]^{}+\frac{1}{2}\mathrm{e}^{4\phi /(1p)(p+1)(d2)Y/2}\left(A_{01\mathrm{}p}^{}\right)^2=T(x)\delta ^{(5p)}(\stackrel{}{x}).$$
(B.36)
In the following we want to find the most general solution of Eq.s (B.31) and (B.34)-(B.36) excluding the origin where the boundary action is located and corresponding to vanishing values of $`\phi `$ and $`A_{01\mathrm{}p}`$ for $`r\mathrm{}`$. Under these conditions, we find that (B.34) is solved by
$$\mathrm{e}^{2\phi /(1p)}=\left(\frac{\mathrm{cosh}X(r)+\gamma \mathrm{sinh}X(r)}{\mathrm{cosh}\alpha +\gamma \mathrm{sinh}\alpha }\right)\left(\frac{f_{}(r)}{f_+(r)}\right)^{(d2)(1+p)ϵ/4}$$
(B.37)
where
$$X(r)=\alpha +\beta \mathrm{log}\frac{f_{}(r)}{f_+(r)},$$
(B.38)
provided that the following relation is satisfied
$`\mathrm{\hspace{0.17em}4}(d2)(x\beta )^216(d2)(x\delta )^2+4(d1)x^2`$
$`+(d2)^2(p+1)(xϵ)^2\left({\displaystyle \frac{(d2)(p+1)}{4}}1\right)=0.`$ (B.39)
Inserting the solution for the dilaton into Eq. (B.35) and neglecting again the source term, we find that the R-R field is given by
$$A_{01\mathrm{}p}=\sqrt{2(\gamma ^21)}\left[\frac{\mathrm{sinh}X(r)(\mathrm{cosh}\alpha +\gamma \mathrm{sinh}\alpha )}{\mathrm{cosh}X(r)+\gamma \mathrm{sinh}X(r)}\mathrm{sinh}\alpha \right]$$
(B.40)
where the overall constant has been determined in terms of $`\gamma `$ through Eq. (B.36).
Finally from Eq.s (B.22), (B.30) and (B.37), one can find the explicit expressions for the components of the metric:
$$B^2=\left(\frac{\mathrm{cosh}X(r)+\gamma \mathrm{sinh}X(r)}{\mathrm{cosh}\alpha +\gamma \mathrm{sinh}\alpha }\right)^{(d2)/2}\left(\frac{f_{}(r)}{f_+(r)}\right)^{ϵ(d2)/2[1(d2)(p+1)/4]}$$
(B.41)
and
$`F^2`$ $`=`$ $`\left({\displaystyle \frac{\mathrm{cosh}X(r)+\gamma \mathrm{sinh}X(r)}{\mathrm{cosh}\alpha +\gamma \mathrm{sinh}\alpha }}\right)^{(p+1)/2}\left(f_{}(r)f_+(r)\right)^{\frac{2}{(d2)}}`$ (B.42)
$`\times \left({\displaystyle \frac{f_{}(r)}{f_+(r)}}\right)^{ϵ(p+1)/2[1(d2)(p+1)/4]}.`$
Eq.s (B.33), (B.37) and (B.40)-(B.42) represent the most general solution of the field equations derived from the action (B.2) which describe a static, spherically symmetric configuration, with asymptotically flat geometry and vanishing v.e.v.’s at infinity for the gauge and scalar fields. The solution depends on five arbitrary parameters $`\alpha `$, $`\beta `$, $`\gamma `$, $`\delta `$ and $`x`$ ($`ϵ`$ is in fact determined in terms of the others through Eq. (B.39)). Setting $`p=0`$ one recovers the solution presented in Section 5 (see Eq.s (5.2)-(5.6)).
In writing this solution, we have not used the precise form of the source terms, or equivalently we have not imposed that the behavior of the various fields at infinity be consistent with what follows from the boundary state, which is given by
$$\phi \frac{1p}{4}\frac{Q_p}{r^{d2}},\eta _a\frac{1}{4}\frac{Q_p}{r^{d2}},A_{01\mathrm{}p}\frac{Q_p}{r^{d2}}$$
(B.43)
and
$$G_{\alpha \beta }\eta _{\alpha \beta }(1\frac{d2}{4}\frac{Q_p}{r^{d2}}),G_{ij}\delta _{ij}(1+\frac{p+1}{4}\frac{Q_p}{r^{d2}}).$$
(B.44)
If we impose that our general solution behaves for large $`r`$ as required by the previous conditions, we must choose the integration constants as follows
$$x\beta =\pm \frac{\mathrm{i}}{4},\gamma =\frac{\mathrm{sinh}2\alpha \pm \mathrm{i}}{\mathrm{cosh}2\alpha },x\delta =\frac{1}{8}\mathrm{and}x0,$$
(B.45)
with $`\alpha `$ arbitrary. For this choice, it is not difficult to see that the $`\alpha `$ dependence drops out and the solution is given by
$`\eta _a`$ $`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle \frac{Q_p}{r^{3p}}}`$ (B.46)
$`\mathrm{e}^{2\phi }`$ $`=`$ $`\left(1+\mathrm{sin}{\displaystyle \frac{Q_p}{r^{3p}}}\right)^{\frac{1p}{2}}`$ (B.47)
$`F^2`$ $`=`$ $`\left(1+\mathrm{sin}{\displaystyle \frac{Q_p}{r^{3p}}}\right)^{\frac{1+p}{4}}`$ (B.48)
$`B^2`$ $`=`$ $`\left(1+\mathrm{sin}{\displaystyle \frac{Q_p}{r^{3p}}}\right)^{\frac{p3}{4}}`$ (B.49)
$`A_{01\mathrm{}p}`$ $`=`$ $`1+{\displaystyle \frac{\mathrm{cos}\frac{Q_p}{r^{3p}}}{1+\mathrm{sin}\frac{Q_p}{r^{3p}}}}.`$ (B.50)
For $`p=0`$ this is precisely the solution (4.13)-(4.17).
In the final part of this appendix we use the general equations we have derived to find the solution corresponding to the non-BPS branes in $`D=10`$ discussed in Section 2. In this case we have to switch off the scalar fields $`\eta _a`$ and the R-R field. Keeping this in mind, the dilaton equation (B.24) becomes
$$\frac{\mathrm{e}^\xi }{r^{d1}}\left(r^{d1}\mathrm{e}^\xi \phi ^{}\right)^{}=\frac{a}{2}F^2T(x)\delta ^{(9p)}(x),$$
(B.51)
whereas the metric equations (B.27) and (B.17) become respectively
$$\frac{\mathrm{e}^\xi }{r^{d1}}\left[r^{d1}\mathrm{e}^\xi (\mathrm{log}B)^{}\right]^{}=\frac{7p}{8}F^2T(x)\delta ^{(9p)}(x)$$
(B.52)
and
$$\frac{\mathrm{e}^\xi }{r^{d1}}\left[r^{d1}\mathrm{e}^\xi (\mathrm{log}F)^{}\right]^{}+\frac{\xi ^{}}{r}=\frac{p+1}{8}F^2T(x)\delta ^{(9p)}(x).$$
(B.53)
Finally, Eq. (B.31) becomes
$$\xi ^{\prime \prime }+\frac{(\xi ^{})^2}{d2}+\frac{\xi ^{}}{r}\frac{(p+1)(d2)}{8}\left[(Y^{})^24Y^{}\frac{\phi ^{}}{a}\right]=2\left(\frac{2\phi ^{}}{a}\right)^2.$$
(B.54)
If we use Eq.s (B.18) and (B.29), we can easily see that they coincide with (2.10), provided that the boundary term is omitted.
Neglecting for a moment the origin, where the boundary term is located, the most general solution of Eq.s (B.51) and (B.52) is given by
$$\mathrm{e}^\phi =\left(\frac{f_{}(r)}{f_+(r)}\right)^\nu ,B^2=\left(\frac{f_{}(r)}{f_+(r)}\right)^\lambda $$
(B.55)
where $`\lambda `$ and $`\nu `$ are constants to be determined and $`f_{}`$ and $`f_+`$ are given in Eq. (B.23) with the substitution of $`Q_p`$ with $`\widehat{Q}_p`$ (the ten dimensional non-BPS D-brane charge defined in Eq. (2.6)). From the previous equations and Eq. (B.30) we get
$$ϵ=\frac{4}{7p}\lambda +\frac{2}{a}\nu .$$
(B.56)
Inserting in Eq. (B.54) the Eq.s (B.22), (B.30) and the first equation in (B.55) one gets
$$8\left(\frac{8p}{7p}\right)(p+1)(7p)\left(ϵ2\frac{\nu }{a}\right)^2=8\nu ^2.$$
(B.57)
Finally imposing that the solution matches also the boundary term we get
$$x=\frac{p3}{8\sqrt{2}\nu }=\frac{7p}{16\lambda }.$$
(B.58)
Eq.s (B.57) and (B.58) imply that
$$ϵ=0,x=\sqrt{\frac{7p}{8(8p)}},\lambda =\frac{7p}{16x},\nu =\frac{p3}{8\sqrt{2}x}.$$
(B.59)
Taking into account that the kinetic term for the dilaton in Eq.s (2.4) and (B.2) have a factor $`2`$ of difference in the normalization we see that Eq.s (B.59) reproduce the solution for the ten dimensional non-BPS D-branes given in Eq.s (2.13) and (2.14).
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# Phase transition in linear sigma model and disoriented chiral condensate
## ACKNOWLEDGMENTS
The author would like to thank K, Rajagopal and C. Gale and J. Randrup for discussions. He also gratefully acknowledge the kind hospitality of Centre for Theoretical Physics, MIT and McGill University where part of the work was done.
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# Critical exponents of the Gross-Neveu model from the effective average action
## Abstract
The phase transition of the Gross-Neveu model with $`N`$ fermions is investigated by means of a non-perturbative evolution equation for the scale dependence of the effective average action. The critical exponents and scaling amplitudes are calculated for various values of $`N`$ in $`d=3`$. It is also explicitely verified that the Neveu-Yukawa model belongs to the same universality class as the Gross-Neveu model.
HD-THEP-00-30
The Gross-Neveu (GN) model is one of the simplest models for interacting fermions. Nevertheless, in three dimensions our quantitative understanding beyond some universal characteristics of the phase transition has remined rather incomplete. The universality class of the GN-model in dimensions between 2 and 4 has been argued to be the same as the Neveu-Yukawa (NY) model in $`4ϵ`$ dimensions . Both the large $`N`$ and $`ϵ`$ expansion indicate that a second order phase transition takes place for some critical value of the coupling constant if the number of fermion species $`N`$ is larger than one . The anomalous dimensions have been calculated up to the third order in the $`1/N`$ expansion , while some critical exponents have been computed to the order $`1/N`$ in the phase with spontaneous symmetry breaking (SSB) . In this letter we find the second order phase transition and calculate the critical exponents employing an analytical method based on nonperturbative flow equations for scale dependent effective couplings. We directly obtain results for arbitrary dimension and without a restriction to large $`N`$. Despite the presence of massless fermions we are able to investigate the symmetric phase. Due to the fermion fluctuations the infrared physics is not trivial in the NY-language and requires a careful discussion of the critical exponents. Beyond the universal critical behaviour our method gives a description for arbitrary values of the GN-coupling away from the critical point. In particular, we compute the non-universal critical amplitudes.
The running couplings parameterize the effective average action $`\mathrm{\Gamma }_k`$ which is a type of coarse grained free energy. It includes the effects of the quantum fluctuations with momenta larger than an infrared cutoff $`k`$. In the limit where the average scale $`k`$ tends to zero $`\mathrm{\Gamma }_k`$ becomes therefore the usual effective action, i.e. the generating functional of $`1PI`$ Green functions. In the limit $`k\mathrm{}`$ it approaches the classical action. In a theory with bosons and fermions the scale dependence of $`\mathrm{\Gamma }_k`$ can be described by an exact nonperturbative evolution equation
$`{\displaystyle \frac{}{t}}\mathrm{\Gamma }_k[\varphi ,\psi ]`$ $`=`$ $`{\displaystyle \frac{1}{2}}\text{Tr}\{(\mathrm{\Gamma }_k^{(2)}[\varphi ,\psi ]+_k)_B^1{\displaystyle \frac{}{t}}R_{kB}`$ (1)
$``$ $`(\mathrm{\Gamma }_k^{(2)}[\varphi ,\psi ]+_k)_F^1{\displaystyle \frac{}{t}}R_{kF}\}`$ (2)
where $`t=\mathrm{ln}(k/\mathrm{\Lambda })`$ with $`\mathrm{\Lambda }`$ some suitable high momentum scale. The trace represents a momentum integration as well as a summation over internal indices and $`\mathrm{\Gamma }_k^{(2)}`$ is the exact inverse propagator given by the matrix of second functional derivatives of the action with respect to bosonic and fermionic field variables. The infrared cutoff
$$_k(q,q^{})=\left(\begin{array}{ccc}R_{kB}& 0& 0\\ 0& 0& R_{kF}\\ 0& R_{kF}& 0\end{array}\right)(2\pi )^d\delta ^d(qq^{})$$
is parameterized by the bosonic and fermionic cutoff functions $`R_{kB}(q)=q^2Z_{\sigma ,k}r_B(q),R_{kF}(q)=i\mathit{}Z_{\psi ,k}r_F(q)`$. We choose
$$R_{kB}=\frac{Z_{\sigma ,k}q^2}{e^{\frac{q^2}{k^2}}1};R_{kF}=iZ_{\psi ,k}\mathit{}\left(\frac{1}{\sqrt{1e^{\frac{q^2}{k^2}}}}1\right)$$
(3)
where $`Z_{\sigma ,k},Z_{\psi ,k}`$ are wave function renormalizations. The momentum integration in Eq. (2) is both infrared and ultraviolet finite. Equation (2) is an exact but complicated functional differential equation wich can only be solved approximately by truncating the most general form of $`\mathrm{\Gamma }_k`$. Once a suitable nonperturbative truncation is found the flow equation can be integrated from some short distance scale $`\mathrm{\Lambda }`$, where $`\mathrm{\Gamma }_\mathrm{\Lambda }`$ can be taken as the classical action, to $`k0`$ thus solving the model approximately.
The GN model is described in terms of a $`O(N)`$ symmetric action for a set of $`N`$ massless Dirac fermions. The classical Euclidean action is given by
$$S_{GN}=d^dx[\overline{\psi }_i(x)(/+\sigma (x))\psi _i(x)+\frac{1}{2\overline{G}}\sigma ^2(x)].$$
(4)
(Here and in the following we distinguish with a bar the dimensionful couplings.) The (pseudo)-scalar $`\sigma (x)`$ is an auxiliary non dynamical field which can be integrated out from the partition function, leading to the replacement $`\sigma (x)\overline{G}\overline{\psi }(x)\psi (x)`$. Its vacuum expectation value $`\sigma _0`$ is proportional to the fermion condensate, $`\sigma _0=\overline{G}<\overline{\psi }\psi >`$. The model is asymptotically free and perturbatively renormalizable in 2 dimensions, hence it exhibits a non-trivial fixed point in $`d=2+ϵ`$. It is $`1/N`$ renormalizable in $`2<d<4`$.
The NY model whose classical action is
$`S_{NY}`$ $`=`$ $`{\displaystyle }d^dx[\overline{\psi }_i(x)(/+\overline{h}\sigma (x))\psi _i(x)+{\displaystyle \frac{1}{2}}(_\mu \sigma (x))^2`$ (5)
$`+`$ $`{\displaystyle \frac{m^2}{2}}\sigma ^2(x)+{\displaystyle \frac{\overline{g}}{4!}}\sigma ^4(x)]`$ (6)
has a Gaussian fixed point in $`d=4`$ where it is perturbatively renormalizable and a non-trivial fixed point in $`d=4ϵ`$. Both models have, in even dimensions, a discrete chiral symmetry which prevents from the addition of a fermion mass term, while in odd dimensions a mass term is forbidden by space parity. Performing a large $`N`$ analysis the universal properties of the two models are argued to be the same in $`2<d<4`$ :in such limit the two models are equivalent in the scaling region if we rescale $`\overline{h}\sigma `$ to $`\sigma `$ and set $`\overline{G}=\overline{h}^2/m^2`$.
We consider a truncation of the effective action $`\mathrm{\Gamma }_k`$ which contains a potential for the scalar field and a Yukawa term. In momentum space it reads ($`𝑑q=d^dq/(2\pi )^d`$)
$`\mathrm{\Gamma }_k[\sigma ,\psi ,\overline{\psi }]={\displaystyle }d^dxU_k(\sigma )+{\displaystyle }dq[{\displaystyle \frac{Z_{\sigma ,k}}{2}}\sigma (q)q^2\sigma (q)`$ (7)
$`Z_{\psi ,k}\overline{\psi }_i(q)i\mathit{}\overline{\psi }^i(q){\displaystyle }dp\overline{h}_k\overline{\psi }_i(q)\sigma (p)\psi ^i(qp)].`$ (8)
The scalar potential is assumed to be a function of the invariant $`\sigma ^2(x)`$ and we make the further simplification
$`U_k(\sigma )`$ $`=`$ $`{\displaystyle \frac{m_k^2}{2}}(\sigma ^2(x)\sigma _{0k}^2)+{\displaystyle \frac{\overline{g}_k}{4!}}(\sigma ^2(x)\sigma _{0k}^2)^2`$ (9)
$`+`$ $`{\displaystyle \frac{\overline{b}_k}{6!}}(\sigma ^2(x)\sigma _{0k}^2)^3.`$ (10)
The symmetric regime is characterized by the minimum being at $`\sigma _{0k}^2=0`$. In the SSB regime a $`k`$-dependent minimum $`\sigma _{0k}^20`$ develops, whereas $`m_k^2=0`$.
Inserting Eqs. (8), (10) into (2) we obtain a set of evolution or renormalization group equations (RGE) for the effective parameters of the theory in the two regimes. Integrating the RGE between some high momentum scale $`\mathrm{\Lambda }`$ and $`k=0`$ we will find the phase transition point and extract the critical exponents of the theory. We find it convenient to work with dimensionless quantities $`h_k^2=Z_\sigma ^1Z_\psi ^2k^{d4}\overline{h}_k^2,g_k=Z_\sigma ^2k^{d4}\overline{g}_k,b_k=Z_\sigma ^3k^{2d6}\overline{b}_k,e_k=Z_\sigma ^1k^2m_k^2,\stackrel{~}{\rho }=\frac{1}{2}Z_\sigma k^{2d}\sigma ^2,\kappa _k=\frac{1}{2}Z_\sigma k^{2d}\sigma _{0k}^2,u_k=U_kk^d`$ and we use $`u_k^{}=\frac{u_k}{\stackrel{~}{\rho }}`$ etc.
The evolution equation for the potential obtains from (2) by evaluating $`\mathrm{\Gamma }_k^{(2)}`$ in the truncation (8) for a constant background scalar field. We find
$`{\displaystyle \frac{_tU_k(\sigma )}{k^d}}`$ $`=`$ $`v_d{\displaystyle _0^{\mathrm{}}}dyy^{d/2}\{{\displaystyle \frac{\eta _\sigma r_B2y\dot{r}_B}{u_k^{}+2\stackrel{~}{\rho }u_k^{\prime \prime }+y(1+r_B)}}`$ (11)
$`+`$ $`2N^{}{\displaystyle \frac{(\eta _\psi r_F+2y\dot{r}_F)(1+r_F)}{2h_k^2\stackrel{~}{\rho }+y(1+r_F)^2}}\}\zeta _k(\stackrel{~}{\rho }).`$ (12)
Here we have introduced the notation $`N^{}=2^{\gamma /2}N`$ with $`2^{\gamma /2}`$ the dimension of the $`\gamma `$ matrices and $`y=q^2/k^2,\dot{r}=\frac{r}{y}`$. Also we have defined $`v_d^1=2^{d+1}\pi ^{d/2}\mathrm{\Gamma }(d/2)`$. In the SSB regime the evolution equations for the parameters $`\kappa _k,g_k`$ and $`b_k`$ are then obtained as
$`_t\kappa _k`$ $`=`$ $`(2d\eta _\sigma )\kappa _k{\displaystyle \frac{3}{g_k}}\left[_{\stackrel{~}{\rho }}\zeta _k\right]_{\kappa _k}`$ (13)
$`_tg_k`$ $`=`$ $`(2\eta _\sigma +d4)g_k+3\left[_{\stackrel{~}{\rho }}^2\zeta _k\right]_{\kappa _k}+{\displaystyle \frac{1}{5}}b_k_t\kappa _k`$ (14)
$`_tb_k`$ $`=`$ $`(3\eta _\sigma +2d6)b_k+15\left[_{\stackrel{~}{\rho }}^3\zeta _k\right]_{\kappa _k}.`$ (15)
In the symmetric regime $`\kappa _k=0`$ so we replace eq. (13) by
$$_te_k=(\eta _\sigma 2)e_k+\left[_{\stackrel{~}{\rho }}\zeta _k\right]_0.$$
(16)
The anomalous dimensions $`\eta _\sigma `$ and $`\eta _\psi `$ are defined as
$$\eta _\sigma (k)=_t\mathrm{ln}Z_{\sigma ,k},\eta _\psi (k)=_t\mathrm{ln}Z_{\psi ,k}.$$
(17)
The wave function renormalizations $`Z`$ parametrize the momentum dependence of the propagators at zero momentum and $`\sigma =\sigma _{0k}`$. One finds
$`\eta _\sigma (k)`$ $`=`$ $`_\alpha \{{\displaystyle \frac{v_d}{d}}{\displaystyle }dyy^{d/2}[({\displaystyle \frac{2}{15}}b_k\kappa _k+g_k)^2\kappa _k\dot{H}(y,\alpha )^2`$ (18)
$``$ $`2Nh_k^2(2^\gamma y\dot{G}(y,\alpha )^22h_k^2\kappa _k\dot{F}(y,\alpha )^2)]\}_{\alpha =0}`$ (19)
$`\eta _\psi (k)`$ $`=`$ $`4h_k^2_\alpha \left\{{\displaystyle \frac{v_d}{d}}{\displaystyle 𝑑yy^{d/2}\dot{H}(y,\alpha )G(y,\alpha )}\right\}_{\alpha =0}`$ (20)
with
$`H(y,\alpha )`$ $`=`$ $`{\displaystyle \frac{1}{e_k+\frac{2}{3}g_k\kappa _k+y[(1+r_B)\alpha (\eta _\sigma r_B+2y\dot{r}_B)]}}`$ (21)
$`G(y,\alpha )`$ $`=`$ $`{\displaystyle \frac{1+r_F\alpha (\eta _\psi r_F+2y\dot{r}_F)}{y[1+r_F\alpha (\eta _\psi r_F+2y\dot{r}_F)]^2+2h_k^2\kappa _k}}`$ (22)
$`F(y,\alpha )`$ $`=`$ $`{\displaystyle \frac{1}{y[1+r_F\alpha (\eta _\psi r_F+2y\dot{r}_F)]^2+2h_k^2\kappa _k}}`$ (23)
Finally, the evolution equation for the Yukawa coupling obtains from taking derivatives of eq. (2) with respect to $`\overline{\psi },\psi `$ and $`\sigma `$:
$`_th_k^2`$ $`=`$ $`(2\eta _\psi +\eta _\sigma +d4)h_k^2`$ (24)
$``$ $`4h_k^4v_d{\displaystyle _0^{\mathrm{}}}dyy^{d/2}[(\eta _\sigma r_B+2y\dot{r}_B)H^2(y,0)F(y,0)`$ (25)
$``$ $`2(\eta _\psi r_F+2y\dot{r}_F){\displaystyle \frac{G^2(y,0)H(y,0)}{1+r_F}}].`$ (26)
The equations (19),(20) and (26) are valid in both regimes provided we set $`\kappa _k`$ and $`e_k`$ appropriately.
If we expand this set of equations in the coupling constants we recover the one-loop results obtained in the $`4ϵ`$ expansion for the NY model
$`_tg`$ $`=`$ $`ϵg+{\displaystyle \frac{1}{8\pi ^2}}\left({\displaystyle \frac{3}{2}}g^2+4Ngh^224Nh^4\right)`$ (27)
$`_th^2`$ $`=`$ $`ϵh^2+{\displaystyle \frac{h^4}{8\pi ^2}}(2N+3)`$ (28)
$`\eta _\sigma `$ $`=`$ $`{\displaystyle \frac{Nh^2}{4\pi ^2}},\eta _\psi ={\displaystyle \frac{h^2}{16\pi ^2}}`$ (29)
Moreover, after identifying the running coupling constant of the GN model as $`G=h_k^2/e_k`$ we also recover the one-loop result obtained in the $`2+ϵ`$ expansion for the GN model
$$_tG=(d2)G(N^{}2)\frac{G^2}{2\pi }+O(G^3).$$
(30)
We numerically evolve the flow equations (13)-(20) from a large momentum scale $`\mathrm{\Lambda }`$ to $`k0`$. The initial values of the parameters are chosen in such a way that $`\mathrm{\Gamma }_\mathrm{\Lambda }=S_{GN}`$: $`Z_{\sigma \mathrm{\Lambda }}0,Z_{\psi \mathrm{\Lambda }}=1,\overline{h}_\mathrm{\Lambda }^2=\mathrm{\Lambda },\overline{g}_\mathrm{\Lambda }=0,\overline{b}_\mathrm{\Lambda }=0.`$ Then $`e_\mathrm{\Lambda }=\left(Z_{\sigma \mathrm{\Lambda }}G_\mathrm{\Lambda }\right)^1`$ is the only free parameter of the theory and plays the role of the temperature. This value has to be tuned in order to be near the second order phase transition. The value corresponding to the critical temperature $`T_c`$ is denoted by $`e_{\mathrm{\Lambda }cr}`$. For $`Z_{\sigma \mathrm{\Lambda }}=10^{10}`$, $`d=3`$ and $`N=3`$ we find in our truncation $`e_{\mathrm{\Lambda }cr}=1.87808521201610^9`$ and the critical coupling in the GN model is $`\overline{G}_{\mathrm{\Lambda }cr}\mathrm{\Lambda }=5.32`$.
The relevant parameter for the deviation form $`T_c`$ is $`\delta e=e_\mathrm{\Lambda }e_{\mathrm{\Lambda }cr}=H(TT_c)`$ with constant $`H`$. In Fig. 1 we show the behaviour of the dimensionless couplings as functions of $`k`$ for $`N=3`$. As can be seen, they all reach a constant value in the symmetric regime ($`e_k>0`$, $`\kappa _k=0`$) corresponding to a nontrivial fixed point with vanishing beta functions. For $`\delta e<0`$ the symmetry is broken, as expected: the coupling $`e_k`$ goes to zero and the order parameter $`\sigma _{0k}`$ acquires a non zero value.
In the following, whenever numerical results are reported, we fix $`d=3`$ and $`N=3`$. However, in Table 1. we summarize the results obtained for $`N=2,4,12`$.
In the symmetric (high $`T`$) phase the fermions are massless. Their fluctuations induce a nontrivial dependence of $`Z_\sigma `$ and the renormalized scalar mass $`m_R^2(k)=Z_\sigma ^1(k)m^2(k)`$ on the scale $`k`$ even away from the phase transition. This contrasts with the standard situation where the running of $`m_R(k)`$ essentially stops in the symmetric phase once $`k`$ becomes much smaller than $`m_R`$. The issue of critical exponents in a situation with two different infrared cutoffs $`k`$ and $`m_R`$ is therefore more complex than usual. In a standard situation we would define the exponents $`\gamma `$ and $`\nu `$ by following the temperature dependence of the unrenormalized and renormalized mass, $`m^2(k)`$ and $`m_R^2(k)`$ for $`k0`$ . Here we define the renormalized mass at some fixed small ratio $`k/m_R`$ by
$$\overline{m}_R^2=m_R^2(k_c)m_{Rcr}^2(k_c),k_c=r_c\overline{m}_R$$
(31)
with $`m_{Rcr}^2(k)=ek^2`$ on the critical trajectory. (In the numerical simulations we will fix the ratio $`r_c`$ to be equal to 0.01). This mass corresponds to the only relevant parameter characterizing the critical behaviour. It is directly related to the deviation from the critical temperature $`TT_c`$ or $`\delta e`$. In the following we use the arguments $`\delta e`$ or $`\overline{m}_R`$ interchangeably. We also define the inverse susceptibility or unrenormalized mass by
$$\overline{m}^2=\overline{m}_R^2Z_\sigma (k_c,\overline{m}_R).$$
(32)
Correspondingly, the critical exponents $`\nu `$ and $`\gamma `$ are defined for fixed $`r_c`$ and we find
$`\nu `$ $`=`$ $`{\displaystyle \frac{1}{2}}\underset{\delta e0}{lim}{\displaystyle \frac{\mathrm{ln}\overline{m}_R^2(\delta e)}{\mathrm{ln}\delta e}}=1.041`$ (33)
$`\gamma `$ $`=`$ $`\underset{\delta e0}{lim}{\displaystyle \frac{\mathrm{ln}\overline{m}^2(\delta e)}{\mathrm{ln}\delta e}}=1.323.`$ (34)
From the definition (32) one has the relation
$$2\nu =\gamma \frac{\mathrm{ln}Z_\sigma (k_c,\overline{m}_R)}{\mathrm{ln}\delta e}.$$
(35)
A typical form of $`Z_\sigma `$ is
$$Z_\sigma Z_0\left(\frac{\overline{m}_R^2+k^2}{\mathrm{\Lambda }^2}\right)^{\frac{1}{2}\overline{\eta }_\sigma }\left(\frac{k^2}{\overline{m}_R^2+k^2}\right)^{\frac{1}{2}\eta _2}$$
(36)
and we conclude
$$\frac{\mathrm{ln}Z_\sigma (k_c,\overline{m}_R)}{\mathrm{ln}\delta e}=\overline{\eta }_\sigma \nu ,\gamma =\nu (2\overline{\eta }_\sigma ).$$
(37)
This is the usual index relation. We find $`\overline{\eta }_\sigma =0.729`$ and the index relation (37) yields $`\gamma =1.323`$, consistent with the value of $`\gamma `$ computed directly. The index $`\eta _2`$, which vanishes in the standard situation, determines the dependence of $`\overline{m}_R`$ on $`r_c`$.
For a more detailed understanding of the scale dependence we consider next the running of the renormalized mass and unrenormalized mass with $`k`$. We define
$`\widehat{\nu }(k,\delta e)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{\mathrm{ln}[m_R^2(k,\delta e)m_{Rcr}^2(k)]}{t}}|_{\delta e}`$ (38)
$`\widehat{\gamma }(k,\delta e)`$ $`=`$ $`{\displaystyle \frac{\mathrm{ln}[m^2(k,\delta e)\stackrel{~}{m}_{cr}^2(k,\delta e)]}{t}}|_{\delta e}`$ (39)
with $`\stackrel{~}{m}_{cr}^2(k,\delta e)=m_{Rcr}^2(k)Z_\sigma (k,\delta e).`$ The relation $`m_R^2=m^2/Z_\sigma `$ implies the index relation
$$\widehat{\gamma }=2\widehat{\nu }\eta _\sigma $$
(40)
which differs from the usual relation $`\gamma =\nu (2\eta _\sigma )`$. For both $`k`$ and $`\overline{m}_R`$ sufficiently small and $`k>>\overline{m}_R`$ the indices $`\widehat{\nu }`$, $`\widehat{\gamma }`$, $`\eta _\sigma `$ and $`\eta _\psi `$ approach constant values independent of $`k`$ and $`\overline{m}_R`$ and we find
$$\widehat{\nu }=0.502,\widehat{\gamma }=0.295,\eta _\sigma =0.710,\eta _\psi =0.040.$$
(41)
These values agree well with Eq. (40). In the opposite regime, $`k<<\overline{m}_R`$, the running of the renormalized mass is only due to the anomalous dimension $`\eta _\sigma `$ which is now different from the value (41). We find $`\widehat{\nu }=0.500,\widehat{\gamma }=0.000,\eta _\sigma =1.000,\eta _\psi =0.000.`$ Again, these values agree well with Eq. (40) and the expectation $`\widehat{\gamma }=0`$. The nontrivial exponents $`\widehat{\nu }`$, $`\eta _\sigma `$ in the NY-language correspond to the absence of renormalization effects for $`\overline{G}_k`$ for $`k0`$ in the GN language. We note that for fixed $`\delta e`$ the renormalized scalar mass (which corresponds to the inverse correlation length) scales as $`m_Rk`$ for $`m_R<<k`$ and $`m_R\sqrt{k}`$ for $`m_R>>k`$. The value of $`\eta _\sigma =1`$ for $`k<<m_R`$, corresponds to $`\eta _2`$ in eq. (36).
We next turn to the dependence on the deviation $`\delta e`$ at fixed $`k`$. We define the exponents as
$`\stackrel{~}{\nu }(k,\delta e)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{\mathrm{ln}[m_R^2(k,\delta e)m_{Rcr}^2(k)]}{\mathrm{ln}\delta e}}|_k`$ (42)
$`\stackrel{~}{\gamma }(k,\delta e)`$ $`=`$ $`{\displaystyle \frac{\mathrm{ln}[m^2(k,\delta e)\stackrel{~}{m}_{cr}^2(k,\delta e)]}{\mathrm{ln}\delta e}}|_k.`$ (43)
In the limit $`\delta e0`$ we find $`\stackrel{~}{\nu }=0.520,\stackrel{~}{\gamma }=1.326`$. The indices (38),(39) and (42),(43) are related to the definition of the critical exponents (33), (34) by
$`\nu `$ $`=`$ $`\underset{\delta e0}{lim}(\stackrel{~}{\nu }(k_c,\delta e)+\widehat{\nu }(k_c,\delta e)\nu )`$ (44)
$`\gamma `$ $`=`$ $`\underset{\delta e0}{lim}(\stackrel{~}{\gamma }(k_c,\delta e)+\widehat{\gamma }(k_c,\delta e)\nu )`$ (45)
Here $`k_c(\delta e)`$ is given by evaluating (31) at fixed $`r_c`$. Equations (44) and (45) yield respectively $`\nu =1.041`$ and $`\gamma =1.326`$, consistent with the results of (33) and (34).
We also have computed the (nonuniversal) critical amplitudes which describe the dependence of $`\overline{m}_R`$ and $`\overline{m}`$ on the coupling $`G_\mathrm{\Lambda }`$ of the GN model. Observing $`\delta e/e_{\mathrm{\Lambda }cr}=G_{\mathrm{\Lambda }cr}\delta (1/G_\mathrm{\Lambda })`$ we obtain, for small deviations from criticality,
$$\overline{m}_R=A_\nu \left|\frac{\delta G_\mathrm{\Lambda }}{G_{\mathrm{\Lambda }cr}}\right|^\nu ,\overline{m}^2=A_\gamma \left|\frac{\delta G_\mathrm{\Lambda }}{G_{\mathrm{\Lambda }cr}}\right|^\gamma .$$
(46)
Our numerical values of the amplitudes $`A_\nu ,A_\gamma `$ are $`A_\nu /\mathrm{\Lambda }=0.016,A_\gamma /\mathrm{\Lambda }^2=0.212`$.
We finally have investigated the low temperature phase with spontaneous symmetry breaking. Here the running of $`\sigma _0`$ stops for small $`k`$ and the complications of the symmetric phase are absent. The critical exponent $`\beta `$ is defined with $`\sigma _0=lim_{k0}\sigma _{0k}`$
$$\beta =\frac{1}{2}\underset{\delta e0}{lim}\frac{\mathrm{ln}\sigma _0^2}{\mathrm{ln}\delta e}$$
(47)
such that
$$\sigma _0=A_\beta \left|\frac{\delta G_\mathrm{\Lambda }}{G_{\mathrm{\Lambda }cr}}\right|^\beta .$$
(48)
We find $`A_\beta /\sqrt{\mathrm{\Lambda }}=0.008`$, $`\beta =0.903`$, in good agreement with the scaling relation $`\beta =\frac{\nu }{2}(d2+\eta _\sigma )`$. In Fig. 2 we plot the condensate $`\sigma _0`$ as a function of $`G_\mathrm{\Lambda }`$.
In the Table we report our results for different values of $`N`$. For all $`N2`$ the existence of a second order phase transition is confirmed by our analysis. As can be checked, the scaling relations are well verified. To compare with existing results obtained in the $`1/N`$ expansion, let us fix $`N=12`$. In the critical exponents have been calculated to the order $`1/N`$:
$`\nu `$ $`=`$ $`1+{\displaystyle \frac{8}{3N\pi ^2}}=1.022,\gamma =1+{\displaystyle \frac{8}{N\pi ^2}}=1.068,`$ (49)
$`\beta `$ $`=`$ $`1+O({\displaystyle \frac{1}{N^2}}),\eta _\sigma =1{\displaystyle \frac{16}{3N\pi ^2}}=0.955.`$ (50)
In the same paper Montecarlo simulations for $`N12`$ are also reported. Conformal techniques have been used to calculate the anomalous dimensions to $`O(1/N^3)`$ and yield, for $`N=12`$, $`\eta _\psi =0.013,\eta _\sigma =0.913`$. However,such techniques rely on being exactly at the critical point and hence cannot be used to calculate the other exponents.
| N | 2 | 3 | 4 | 12 |
| --- | --- | --- | --- | --- |
| $`\nu `$ | 0.961 | 1.041 | 1.010 | 1.023 |
| $`\gamma `$ | 1.384 | 1.323 | 1.228 | 1.075 |
| $`\nu (2\overline{\eta }_\sigma )`$ | 1.403 | 1.323 | 1.230 | 1.075 |
| $`\beta `$ | 0.745 | 0.903 | 0.910 | 0.998 |
| $`\frac{\nu }{2}(1+\eta _\sigma )`$ | 0.750 | 0.890 | 0.903 | 0.991 |
| $`A_\nu /\mathrm{\Lambda }`$ | 0.007 | 0.016 | 0.009 | 0.014 |
| $`A_\gamma /\mathrm{\Lambda }^2`$ | 0.042 | 0.212 | 0.233 | 0.968 |
| $`A_\beta /\sqrt{\mathrm{\Lambda }}`$ | 0.007 | 0.008 | 0.005 | 0.007 |
| $`\eta _\sigma `$ | 0.561 | 0.710 | 0.789 | 0.936 |
| $`\eta _\psi `$ | 0.066 | 0.040 | 0.027 | 0.007 |
| $`\overline{\eta }_\sigma `$ | 0.541 | 0.729 | 0.765 | 0.971 |
| $`G_{\mathrm{\Lambda }cr}`$ | 9.989 | 5.325 | 3.613 | 1.006 |
TABLE 1. Critical exponents and amplitudes for different values of N
The case $`N=1`$ appears to be different from $`N>1`$. We find a phase transition. For small $`G_\mathrm{\Lambda }`$ (large $`e_\mathrm{\Lambda }`$) eq. (30) is valid ($`N^{}=2`$) and $`G_k`$ scales according to its canonical dimension. The model is in the symmetric phase. For $`G_\mathrm{\Lambda }>G_{\mathrm{\Lambda }cr}`$, $`G_{\mathrm{\Lambda }cr}=19.416`$ the mass term at the origin of the potential becomes negative, indicating spontaneous symmetry breaking. We find no scaling solution, neither for $`e_k0`$ nor for $`\kappa _k0`$. This may suggest a first order transition.
In conclusion, a simple truncation of the exact flow equation for the effective average action gives a consistent picture for a second order phase transition for the GN model with $`N2`$ in three dimensions. We have computed critical exponents and amplitudes and we relate directly physical observables like the correlation length or the order parameter to the value of the coupling $`G`$. By choosing different initial conditions we have also explicitely verified that the Neveu-Yukawa model belongs to the same universality class as the Gross-Neveu model. The universal exponents are independent of the initial parameters, but not quantities like $`\sigma _0`$, the renormalized mass and the corresponding amplitudes.
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# References
The award of the 1999 Nobel Prize for physics to G. ’t Hooft and M. Veltman, and the success of the predictions of their formulation of the renormalized non-Abelian quantum loop corrections for the Standard Model of the electroweak interactions in confrontation with data of LEP experiments, underscores the need to continue to test this theory at the quantum loop level in the gauge boson sector itself. This emphasizes the importance of the on-going precision studies of the processes $`e^+e^{}W^+W^{}+n(\gamma )4f+n(\gamma )`$ at LEP2 energies , as well as the importance of the planned future higher energy studies of such processes in LC physics programs . We need to stress that hadron colliders also have considerable reach into this physics and we hope to come back to their roles elsewhere .
In what follows, we present precision predictions for the event selections (ES) of the LEP2 MC Workshop , for the processes $`e^+e^{}W^+W^{}+n(\gamma )4f+n(\gamma )`$, based on our new exact $`𝒪(\alpha )_{prod}`$ YFS-exponentiated LL $`𝒪(\alpha ^2)`$ FSR leading-pole approximation (LPA) formulation, as it is realized in the MC program YFSWW3-1.14 , in combination with all four-fermion processes MC event generator KoralW-1.42 so that the respective four-fermion background processes are taken into account in a gauge-invariant way. Indeed, gauge invariance is a crucial aspect of our work and we stress that we maintain it throughout our calculations. Here, FSR denotes final-state radiation and LL denotes leading-log as usual.
Recently, the authors in Refs. have also presented MC program results for the processes $`e^+e^{}W^+W^{}+n(\gamma )4f+n(\gamma ),n=0,1`$, in combination with the complete background processes that feature the exact LPA $`𝒪(\alpha )`$ correction. Thus, we will compare our results, where possible, with those in Refs. in an effort to check the over-all precision of our work. As we argue below, the two sets of results should agree at a level below $`0.5\%`$ on observables such as the total cross section.
More specifically, in YFSWW3-1.13 , the leading-pole approximation (LPA) is used to develop a fully gauge-invariant YFS-exponentiated calculation of the signal process $`e^+e^{}W^+W^{}+n(\gamma )4f+n(\gamma )`$, which features the exact $`𝒪(\alpha )`$ electroweak correction to the production process and the $`𝒪(\alpha ^2)`$ LL corrections to the final-state decay processes. The issue is how to combine this calculation with that of KoralW-1.42 in Ref. for the corresponding complete Born-level cross section with YFS-exponentiated initial-state $`𝒪(\alpha ^3)`$ LL corrections. In this connection, we point out that the LPA enjoys some freedom in its actual realization, just as does the LL approximation in the precise definition of the big log L, without spoiling its gauge invariance. This can already be seen from the book of Eden et al. , wherein it is stressed that the analyticity of the S-matrix applies to the scalar form factors themselves in an invariant Feynman amplitude, without any reference to the respective external wave functions and kinematical (spinor) covariants. The classic example illustrated in Ref. is that of pion–nucleon scattering, with the amplitude
$$=\overline{u}(p_2)[A(s,t)+B(s,t)(\overline{)}q_1+\overline{)}q_2)]u(p_1),$$
(1)
where the $`p_i`$ are the nucleon 4-momenta, the $`q_i`$ are the pion 4-momenta, $`u(p)`$ is the usual Dirac wave function of the nucleon, and the invariant scalar functions $`A(s,t)`$ and $`B(s,t)`$ of the Mandelstam invariants $`s=(p_1+q_1)^2,t=(q_2q_1)^2`$ realize the analytic properties of the S-matrix themselves in the complex $`s`$ and $`t`$ planes. This means that, whenever we have spinning particles, we may focus on the analogs of $`A`$ and $`B`$ in eq. (1) in isolating the respective analytic properties of the corresponding S-matrix elements. We note that Stuart has emphasized this point in connection with the production and decay of $`Z`$-pairs in $`e^+e^{}`$ annihilation and in connection with the production and decay of single $`W`$’s in $`e^+e^{}`$ annihilation. What this means is that, in formulating the Laurent expansion of the S-matrix about its poles to isolate the dominant leading-pole term (the LPA is then realized by dropping all but this leading term), we may focus on only $`A`$ and $`B`$, or we may insist that in evaluating the residues of the poles in the S-matrix the wave functions and kinematical covariants are also evaluated at the pole positions. When we focus only on the analogs of $`A`$ and $`B`$ in formulating the LPA, we shall refer to the result simply as the LPA<sub>a</sub>; when we also evaluate the wave functions and kinematical covariants at the pole positions in isolating the poles in the analogs of $`A`$ and $`B`$ for the LPA, we shall refer to the respective result as the LPA<sub>b</sub>. As Stuart stressed as well, both the LPA<sub>a</sub> and the LPA<sub>b</sub> are fully gauge-invariant.
For the process under discussion, a general representation is
$$=\underset{j}{}\mathrm{}_jA_j\left(\{q_kq_l\}\right),$$
(2)
where $`\{\mathrm{}_j\}`$ are a complete set of kinematical covariants which carry the same transformation properties as does $``$, and the Lorentz scalars $`\{q_kq_l\}`$ are a complete set of Lorentz scalar invariants for the external 4-momenta of $``$. In the LPA<sub>a</sub>, we make a Laurent expansion of the $`A_j`$ and retain only their leading poles, without touching the $`\{\mathrm{}_j\}`$; in the LPA<sub>b</sub>, we also evaluate the $`\mathrm{}_j`$ at the position of the respective leading poles. Evidently, in the latter case, we must make an analytic continuation of the phase-space point originally associated with the $`\{\mathrm{}_j\}`$ to a corresponding such point for the respective pole positions. See Ref. for an illustration of such a continuation in the context of the YFS-exponentiated exact $`𝒪(\alpha )`$ calculation for the production process in $`e^+e^{}W^+W^{}+n(\gamma )4f+n(\gamma )`$, and Refs. for a similar illustration in the context of the $`𝒪(\alpha )`$ correction to $`e^+e^{}W^+W^{}4f`$. Having isolated the appropriate realization of the LPA at the level of $``$, it must still be decided whether to treat the phase space used to integrate the cross section exactly or approximately to match what was done for the $`\{\mathrm{}_j\}`$ in the case of the LPA<sub>b</sub>. In all of our work, we stress that we always treat the exact phase space, both in the LPA<sub>a</sub> and in the LPA<sub>b</sub>.
In the context of YFS exponentiation, we realize the LPA as follows, as was briefly described already in Ref. . Taking the respective 4-fermion plus $`n`$-photon process kinematics to be as given by (here, $`d\tau _{n+4}`$ is the respective phase space differential with the appropriate normalization):
$$\begin{array}{cc}& e^{}(p_1)+e^+(p_2)f_1(r_1)+\overline{f}_2(r_2)+f_1^{}(r_1^{})+\overline{f}_2^{}(r_2^{})+\gamma (k_1),\mathrm{},\gamma (k_n)\hfill \\ & \sigma _n=\frac{1}{flux}𝑑\tau _{n+4}(p_1+p_2;r_1,r_2,r_1^{},r_2^{},k_1,\mathrm{},k_n)\hfill \\ & \underset{Ferm.Spin}{}\underset{Phot.Spin}{}|_{4f}^{(n)}(p_1,p_2,r_1,r_2,r_1^{},r_2^{},k_1,\mathrm{},k_n)|^2,\hfill \end{array}$$
(3)
and that of the corresponding $`W^+W^{}`$ production and decay process to be as given by
$$\begin{array}{cc}& e^{}(p_1)+e^+(p_2)W^{}(q_1)+W^+(q_2),\hfill \\ & W^{}(q_1)f_1(r_1)+\overline{f}_2(r_2),W^+(q_2)f_1^{}(r_1^{})+\overline{f}_2^{}(r_2^{}),\hfill \\ & \sigma _n=\frac{1}{flux}𝑑\tau _{n+4}(p_1+p_2;r_1,r_2,r_1^{},r_2^{},k_1,\mathrm{},k_n)\hfill \\ & \underset{Ferm.Spin}{}\underset{Phot.Spin}{}|_{LPA}^{(n)}(p_1,p_2,r_1,r_2,r_1^{},r_2^{},k_1,\mathrm{},k_n)|^2\hfill \end{array}$$
(4)
in the context of YFS exponentiation , we proceed according to Refs.
$$\begin{array}{cc}& _{4f}^{(n)}(p_1,p_2,r_1,r_2,r_1^{},r_2^{},k_1,\mathrm{},k_n)\genfrac{}{}{0pt}{}{LPA}{=>}_{LPA}^{(n)}(p_1,p_2,r_1,r_2,r_1^{},r_2^{},k_1,\mathrm{},k_n)\hfill \\ & =\underset{Phot.Partitions}{}_{Prod}^{(a),\lambda _1\lambda _2}(p_1,p_2,q_1,q_2,k_1,\mathrm{},k_a)\hfill \\ & \times \frac{1}{D(q_1)}_{Dec_1,\lambda _1}^{(ba)}(q_1,r_1,r_2,k_{a+1},\mathrm{},k_b)\times \frac{1}{D(q_2)}_{Dec_2,\lambda _2}^{(nb)}(q_2,r_1^{},r_2^{},k_{b+1},\mathrm{},k_n),\hfill \\ & D(q_i)=q_i^2M^2,M^2=(M_W^2i\mathrm{\Gamma }_WM_W)(1\mathrm{\Gamma }_W^2/M_W^2+𝒪(\alpha ^3)),\hfill \\ & q_1=r_1+r_2+k_{a+1}+\mathrm{}+k_b;q_2=r_1^{}+r_2^{}+k_{b+1}+\mathrm{}+k_n,\hfill \end{array}$$
(5)
so that $`M^2`$ is the pole in the complex $`q^2`$ plane when $`q`$ is the respective $`W`$ 4-momentum, and $`M_W`$ and $`\mathrm{\Gamma }_W`$ are the on-shell scheme mass and width, respectively. The residues in (5) are all defined at $`q_i^2=M^2`$ with a prescription according to whether we have LPA<sub>a</sub> or LPA<sub>b</sub>, so that (5) is our YFS generalization of the formula in eq. (12) in the first paper in Ref. :
$$^{(n)}=\underset{\lambda _1,\lambda _2}{}\mathrm{\Pi }_{\lambda _1,\lambda _2}(M_1,M_2)\frac{\mathrm{\Delta }_{\lambda _1}^+(M_1)}{D_1}\frac{\mathrm{\Delta }_{\lambda _2}^{}(M_2)}{D_2},n=0,1,$$
(6)
where $`D_i=D(q_i)`$ and $`M_i^2=M^2`$. We stress that, unlike what is true of the formula in eq. (12) in the first paper in Ref. and in eq. (6) here, in eq. (5) $`n`$ is arbitrary. The sum over “photon partitions” is over all $`10^n`$ possible attachments of $`n`$ photons to the six external fermion lines and the two $`W^\pm `$ lines (one for the $`W`$ production and one for the $`W`$ decay, respectively). We make the further approximation that $`M_i^2=M_W^2`$ in the residues in (5), always maintaining gauge invariance, as explained. Equations (3) and (4) in Ref. then give us, in the presence of renormalization-group-improved perturbation theory, for the representation
$$_{LPA}^{(n)}(p_1,p_2,r_1,r_2,r_1^{},r_2^{},k_1,\mathrm{},k_n)=\underset{j=0}{\overset{\mathrm{}}{}}_j^{(n)}(p_1,p_2,r_1,r_2,r_1^{},r_2^{},k_1,\mathrm{},k_n),$$
(7)
where $`_j^{(n)}`$ is the $`j`$-th virtual photon loop contribution to the residues in $`_{LPA}^{(n)}`$, the identifications
$$_j^{(n)}(p_1,p_2,r_1,r_2,r_1^{},r_2^{},k_1,\mathrm{},k_n)=\underset{r=0}{\overset{j}{}}𝔪_{jr}^{(n)}\frac{(\alpha B^{})^r}{r!},$$
(8)
where $`B^{}`$ is now, for the LPA<sub>b</sub> case to be definite, the on-shell virtual YFS infrared function, which reduces to that given in eqs. (8) and (9) in Ref. when we restrict our attention to the production process, and $`\alpha `$ is indeed $`\alpha (0)`$ when it multiplies $`B^{}`$ here. Let us keep this limit of $`B^{}`$ in mind, as we focus on the gauge invariance of YFSWW3-1.11 in Ref. , which treats the radiation in the production process, and on that of YFSWW3-1.13 and YFSWW3-1.14, in which the radiation from the decay processes is also treated. In the LPA<sub>a</sub> case, the corresponding $`B^{}`$ function is off-shell. Let us discuss first the LPA<sub>b</sub> case and comment later on how the corresponding results for the LPA<sub>a</sub> case are obtained.
Here, since the $`SU(2)_L\times U(1)`$ Ward–Takahashi identities require (see eq. (47) in Ref. )
$$k^\mu M_\mu ^\gamma =0,k^\mu M_\mu ^Z=i\sqrt{\mu _Z}M^\chi ,k^\mu M_\mu ^{W^\pm }=\pm \sqrt{\mu _W}M^{\varphi ^\pm },$$
(9)
for $`\mu _V`$ denoting the squared $`V`$ boson mass (so that, for $`V=W`$, $`\mu _W=M^2`$), we find that $`B^{}`$ is $`SU(2)_L\times U(1)`$-invariant from the equations in (9) and our result eq. (8) in Ref. . From eq. (8) it then follows that the infrared residuals $`𝔪_{jr}^{(n)}`$ are also $`SU(2)_L\times U(1)`$-invariant. Here, $`\chi `$ and $`\varphi ^\pm `$ are the usual unphysical Higgs fields in our general renormalizable gauges and we use the notation of Ref. , so that $`M_\mu ^Z`$ is their respective amplitude for the emission of a $`Z`$ of Lorentz index $`\mu `$ and 4-momentum $`k`$, and $`M^\chi `$ is their corresponding amplitude for the emission of a $`\chi `$ with the same 4-momentum, etc.
Introducing eq. (8) into (7) gives
$$_{LPA_b}^{(n)}(p_1,p_2,r_1,r_2,r_1^{},r_2^{},k_1,\mathrm{},k_n)=e^{\alpha B^{}}\underset{j=0}{\overset{\mathrm{}}{}}𝔪_j^{(n)}(p_1,p_2,r_1,r_2,r_1^{},r_2^{},k_1,\mathrm{},k_n).$$
(10)
Equation (2.13) in Ref. and eq. (7) in Ref. then give our $`n`$-photon differential cross section, for $`P=p_1+p_2`$, $`\stackrel{}{P}=0`$, as
$$d\sigma _{LPA_b}^n=e^{2\mathrm{}\alpha B^{}}\frac{1}{n!}\underset{j=1}{\overset{n}{}}\frac{d^3k}{k_j^0}\delta ^{(4)}\left(PR\underset{j}{}k_j\right)\left|\underset{n^{}=0}{\overset{\mathrm{}}{}}𝔪_n^{}^{(n)}\right|^2\frac{d^3r_1}{r_1^0}\frac{d^3r_2}{r_2^0}\frac{d^3r_1^{}}{r_{}^{}{}_{1}{}^{0}}\frac{d^3r_2^{}}{r_{}^{}{}_{2}{}^{0}},$$
(11)
where we note that, when we only focus on the production process in eq. (11), $`R`$ is the produced $`WW`$ intermediate state; $`R=r_1+r_2+r_1^{}+r_2^{}`$. Using the second theorem of the YFS program (eq. (2.15) in ), we get
$$\begin{array}{cc}\hfill \left|\underset{n^{}=0}{\overset{\mathrm{}}{}}𝔪_n^{}^{(n)}\right|^2& =\stackrel{~}{S}(k_1)\mathrm{}\stackrel{~}{S}(k_n)\overline{\beta }_0+\underset{i=1}{\overset{n}{}}\stackrel{~}{S}(k_1)\mathrm{}\stackrel{~}{S}(k_{i1})\stackrel{~}{S}(k_{i+1})\mathrm{}\stackrel{~}{S}(k_n)\overline{\beta }_1(k_i)\hfill \\ & +\mathrm{}+\underset{i=1}{\overset{n}{}}\stackrel{~}{S}(k_i)\overline{\beta }_{n1}(k_1,\mathrm{},k_{i1},k_{i+1},\mathrm{},k_n)+\overline{\beta }_n(k_1,\mathrm{},k_n),\hfill \end{array}$$
(12)
where the real emission function $`\stackrel{~}{S}(k)`$ is given by $`\stackrel{~}{S}_{Prod}(k)`$, the real emission infrared function in eq. (8) in Ref. for on-shell $`W^{}s`$, when we only focus on the emission from the production process as we did in Refs. . Since, in general,
$$\stackrel{~}{S}(k)=\stackrel{~}{S}_{Prod}+\stackrel{~}{S}_{Dec_1}+\stackrel{~}{S}_{Dec_2}+\stackrel{~}{S}_{Int},$$
(13)
with
$$\begin{array}{cc}\hfill \stackrel{~}{S}_{Prod}(k)& =\frac{\alpha }{4\pi ^2}[(\frac{p_1}{kp_1}\frac{p_2}{kp_2})^2+(\frac{𝒜q_1}{k𝒜q_1}\frac{𝒜q_2}{k𝒜q_2})^2\hfill \\ & +\left(\frac{p_1}{kp_1}\frac{𝒜q_1}{k𝒜q_1}\right)^2+\left(\frac{p_2}{kp_2}\frac{𝒜q_2}{k𝒜q_2}\right)^2\hfill \\ & (\frac{p_1}{kp_1}\frac{𝒜q_2}{k𝒜q_2})^2(\frac{p_2}{kp_2}\frac{𝒜q_1}{k𝒜q_1})^2\left]\right|_{(𝒜q_i)^2=M_W^2},\hfill \end{array}$$
(14)
$$\begin{array}{cc}\hfill \stackrel{~}{S}_{Dec_1}(k)& =\frac{\alpha }{4\pi ^2}[Q_1Q_2(\frac{r_1}{kr_1}\frac{r_2}{kr_2})^2Q_1Q_W(\frac{r_1}{kr_1}\frac{𝒜q_1}{k𝒜q_1})^2\hfill \\ & +Q_2Q_W(\frac{r_2}{kr_2}\frac{𝒜q_2}{k𝒜q_2})^2\left]\right|_{(𝒜q_1)^2=M_W^2},\hfill \end{array}$$
(15)
$$\begin{array}{cc}\hfill \stackrel{~}{S}_{Dec_2}(k)& =\frac{\alpha }{4\pi ^2}[Q_1^{}Q_2^{}(\frac{r_1^{}}{kr_1^{}}\frac{r_2^{}}{kr_2^{}})^2+Q_1^{}Q_W(\frac{r_1^{}}{kr_1^{}}\frac{𝒜q_2}{k𝒜q_2})^2\hfill \\ & Q_2^{}Q_W(\frac{r_2^{}}{kr_2^{}}\frac{𝒜q_2}{k𝒜q_2})^2\left]\right|_{(𝒜q_2)^2=M_W^2},\hfill \end{array}$$
(16)
$$\begin{array}{cc}\hfill \stackrel{~}{S}_{Int}(k)& =\frac{\alpha }{4\pi ^2}(\frac{p_1}{kp_1}\frac{p_2}{kp_2}+Q_1\frac{r_1}{kr_1}Q_2\frac{r_2}{kr_2}\hfill \\ & +Q_1^{}\frac{r_1^{}}{kr_1^{}}Q_2^{}\frac{r_2^{}}{kr_2^{}})^2\stackrel{~}{S}_{Prod}\stackrel{~}{S}_{Dec_1}\stackrel{~}{S}_{Dec_2}\hfill \end{array}$$
(17)
is really composed of the scalar product of emission currents $`\{j_b^\mu (k)\}`$ with $`k_\mu j_b^\mu (k)=0`$, eq. (13) is also $`SU(2)_L\times U(1)`$-invariant. Here
$$𝒜q_i\text{analytical continuation of }q_i\text{to the point}q_i^2=M_W^2.$$
(18)
This analytical continuation, already described in Ref. , does not spoil the gauge invariance, as we see from Eqs. (13,18). It follows that the hard-photon residuals $`\{\overline{\beta }_n\}`$ are also $`SU(2)_L\times U(1)`$-invariant.
Substitutung (12) into (11) we finally get the $`SU(2)_L\times U(1)`$-invariant expression, which is the fundamental formula of our calculation,
$$\begin{array}{cc}& d\sigma _{LPA_b}=e^{2\mathrm{}\alpha B^{}+2\alpha \stackrel{~}{B}}\frac{1}{(4\pi )^4}d^4ye^{iy(p_1+p_2q_1q_2)+D}\hfill \\ & \times \left[\overline{\beta }_0+\underset{n=1}{\overset{\mathrm{}}{}}\frac{d^3k_j}{k_j^0}e^{iyk_j}\overline{\beta }_n(k_1,\mathrm{},k_n)\right]\frac{d^3r_1}{\overline{r}_1^0}\frac{d^3r_2}{r_2^0}\frac{d^3r_1^{}}{\overline{r^{}}_1^0}\frac{d^3r_2^{}}{\overline{r^{}}_2^0},\hfill \end{array}$$
(19)
where we have defined
$$D=\frac{d^3k}{k_0}\stackrel{~}{S}\left[e^{iyk}\theta (K_{max}|\stackrel{}{k}|)\right],2\alpha \stackrel{~}{B}=\frac{d^3k}{k_0}\theta (K_{max}|\stackrel{}{k}|)\stackrel{~}{S}(k).$$
(20)
This shows that the parameter $`K_{max}\sqrt{s}`$ is a dummy parameter on which $`d\sigma `$ does not depend. In (19), when the complete value of $`\stackrel{~}{S}(k)`$ is used, then all $`W^\pm `$ radiative effects are contained in the respective $`\overline{\beta }_n`$ residuals, in accordance with the YFS theory in Ref. , as the non-zero widths of the $`W`$’s prevent any IR singularities when a $`W`$ radiates a photon in (4). In our work in YFSWW3, as we indicate below, we make the approximation of dropping all interference effects between the production and decay stages and between the two decay stages of (4). This means that we drop the $`\stackrel{~}{S}_{Int}(k)`$ in $`\stackrel{~}{S}(k)`$ in (13) and in (19), so that the YFS theory then determines the corresponding forms of the YFS functions $`\overline{\beta }_n`$, $`B^{}`$ and $`D`$ as also having the respective interferences dropped. This approximation, which resums a certain class of large $`W`$ radiative effects, corresponds to the YFS exponentiation of the $`W`$ production and decay radiation in the LPA, neglecting of all interferences between the production and decay stages and between the two decay processes.
We will now comment further on our use of (19). Since the residuals on the RHS (right-hand side) of (5) are on-shell amplitudes in the $`SU(2)_L\times U(1)`$ theory, for both the production and the decay process, it follows that they satisfy the renormalization group equations for the $`SU(2)_L\times U(1)`$ theory, as explained in Ref. :
$$\left(\mu \frac{}{\mu }+\beta _j(\{g_{iR}\})\frac{}{g_{jR}}\gamma _{\mathrm{\Theta }_j}(\{g_{iR}\})m_{jR}\frac{}{m_{jR}}\gamma _\mathrm{\Gamma }(\{g_{iR}\})\right)\mathrm{\Gamma }=0$$
(21)
for $`\mathrm{\Gamma }=_{Prod}^{(n)\lambda _1,\lambda _2},_{Dec,\lambda _i}^{(n^{})}`$, where $`\mu `$ is the arbitrary renormalization point, $`\{g_{iR}\}`$ are the respective renormalized $`SU(2)_L\times U(1)`$ couplings, and $`\{m_{iR}\}`$ are the corresponding renormalized mass parameters, etc., as defined in Ref. . It follows that any two schemes for computing $`_{Prod}^{(n)\lambda _1,\lambda _2},_{Dec,\lambda _i}^{(n^{})}`$ are related by a finite renormalization-group transformation. Thus, the complex pole scheme (CPS) and the fermion loop scheme (FLS) values of $`_{Prod}^{(n)\lambda _1,\lambda _2},_{Dec_i,\lambda _i}^{(n^{})}`$ are related by such a finite renormalization-group transformation. Specifically, if we denote CPS and FLS values of any quantity $`A`$ by $`A(\mathrm{CPS})`$ and $`A(\mathrm{FLS})`$, respectively, then we have the identity
$$Z_{\mathrm{\Gamma }(CPS)}^1\mathrm{\Gamma }(\mathrm{CPS})=\mathrm{\Gamma }_{un}=Z_{\mathrm{\Gamma }(\mathrm{FLS})}^1\mathrm{\Gamma }(\mathrm{FLS}),$$
(22)
where $`\mathrm{\Gamma }_{un}`$ is the respective unrenormalized value of $`\mathrm{\Gamma }`$ and the $`Z_{\mathrm{\Gamma }(\mathrm{R})},\mathrm{R}=\mathrm{CPS},\mathrm{FLS}`$ are the respective field renormalization constants. It therefore follows that we have the finite renormlization group transformation
$$\mathrm{\Gamma }(\mathrm{CPS})=Z_{\mathrm{\Gamma }(\mathrm{CPS})}Z_{\mathrm{\Gamma }(\mathrm{FLS})}^1\mathrm{\Gamma }(\mathrm{FLS})$$
(23)
connecting the FLS and CPS schemes. Of course, as the usual implementation of the FLS scheme omits a gauge invariant set of contributions to the heavy vector boson widths, for example, in checking the the result in (23) we must omit the corresponding contributions on both sides of the equation for consistency, as is usualy the case when we compare renormalzation group improved quantities. Indeed, as we take the normalization points for the two schemes to be the complex pole position $`M^2=\mu _W`$, and as we evaluate $`_{Prod}^{(n)\lambda _1,\lambda _2},_{Dec_i,\lambda _i}^{(n^{})}`$ at this normalization point, the only difference in using the FLS instead of the CPS will be the approximate treatment in the FLS of the actual values of $`_{Prod}^{(n)\lambda _1,\lambda _2},_{Dec_i,\lambda _i}^{(n^{})},\text{and}\mu _W`$ for the pole position, as already shown in eq. (8) of the first paper in Ref. : thus, for example, if we keep only fermion loops in $`𝒪(\alpha )`$, only the lowest-order width appears in eq. (6) of the first paper in Ref. ; thus if one would express the resulting $`\mu _W=M^2`$ in terms of the respective on-shell mass and width, only the lowest-order part of $`\mathrm{\Gamma }_W`$ would be given properly. Similarly, as both schemes normalize at $`\mu _W`$, the difference in the residues $`_{Prod}^{(n)\lambda _1,\lambda _2},_{Dec_i,\lambda _i}^{(n^{})}`$ is that in the FLS only the fermion-loop contributions are retained, whereas the CPS keeps all loops. Evidently, we may extend our $`𝒪(\alpha )`$ calculation in the FLS by using the complete value of $`\mu _W`$ and including all the one-loop corrections and attendant $`𝒪(\alpha )`$ real corrections from Ref. , as we did in Ref. . We conclude that, after adding in the entire $`𝒪(\alpha )`$ correction from Ref. , our LPA exact $`𝒪(\alpha )_{prod}`$ YFS-exponentiated calculation arrives at the same amplitudes, independent of whether we started with the FLS or the CPS.
All of the above results extend directly to the calculation when we use LPA<sub>a</sub> amplitudes, as these are also gauge-invariant by the gauge-invariance of our leading poles in the S-matrix. Thus, the only change we must make is that the respective residues must be calculated in the LPA<sub>a</sub> rather than in the LPA<sub>b</sub>; for example, in eqs. (14-17), $`q_i`$ would be used instead of $`𝒜q_i`$, etc. We have done this, as we further illustrate in the following.
Let us now comment on the issue of the pure FSR YFS exponentiation for the decay processes treated in the LPA. We proceed in analogy to what is done in Ref. for the MC YFS3 for the respective FSR. Specifically, for both decay residue amplitudes $`_{Dec_i,\lambda _i}^{(n^{})}`$, we may have contributions to the respective hard-photon residuals $`\overline{\beta }_n`$ due to emission for the final-state decay processes; in these we follow the procedure, described in Refs. and already illustrated in Ref. , for including these contributions, using the same YFS methods as we used above for radiative effects to the initial, intermediate and final states. Here, we shall neglect all interference effects between the production and decay processes as we explained above; this is analogous neglecting all interference effects between the initial and final states in the $`\overline{\beta }_n`$ in Refs. . We note that it is possible to retain these interference effects, as we have illustrated in the exact $`𝒪(\alpha )`$ YFS-exponentiated BHWIDE MC in Ref. for wide-angle Bhabha scattering and, more recently, using the new CEEX exponentiation theory in Ref. , to all orders in $`\alpha `$ in the new KK MC for the 2 fermion processes from the $`\tau `$ threshold to $`1`$ TeV. In this way we see that the use of eq. (19) to include exponentiation of the FSR is fully realizable by Monte Carlo methods we already tested. We stress that, for the same reasons as we gave for the exponentiation of the complete process, these FSR contributions to the $`\overline{\beta }_n`$ are fully gauge-invariant.
In the current version of YFSWW3, version 1.14, we also drop the $`\stackrel{~}{S}_{Dec_i}`$ terms in $`\stackrel{~}{S}`$ in eq. (13) and the corresponding terms in the functions $`\overline{\beta }_n`$, $`B^{}`$ and $`D`$, and include FSR using the program PHOTOS , which gives us a LL $`𝒪(\alpha ^2)`$ realization of the FSR in which finite $`p_T`$ effects are represented as they are in the $`𝒪(\alpha )`$ soft-photon limit. This LL implementation of FSR is fully gauge-invariant. The ratio of BRs is then used to obtain the $`𝒪(\alpha )`$ correction in the normalization associated with the $`𝒪(\alpha )`$ correction to the decay processes themselves. Evidently, these ratios of BRs are also gauge-invariant. As we illustrate below, for the corresponding non-factorizable corrections we use an efficient approximation in terms of the so-called screened Coulomb ansatz , which has been shown to be in good agreement with the exact calculations for singly inclusive distributions . This ansatz is gauge-invariant.
We also point out that the current version 1.14 differs from version 1.13 in Ref. in that it uses a different renormalization scheme. Specifically, the scheme used in version 1.13 is the so-called $`G_\mu `$ of Ref. , in which the weak-scale coupling $`\alpha _{G_F}`$ is used for all terms in the virtual correction, except those that are infrared-singular, which are given the coupling $`\alpha \alpha (0)`$. In the renormalization-group-improved YFS theory, as formulated in Ref. , all the terms in the amplitude that involve corrections, in which the emitted photon of 4-momentum $`k`$ has $`k^20`$, should have the coupling strength corresponding to $`\alpha (0)`$ – not just those that are IR singular. We therefore have introduced into YFSWW3 this requirement of the renormalization-group-improved YFS theory to arrive at version 1.14. We refer to this scheme as our scheme $`(A)`$. According to the renormalization-group-improved YFS theory, it gives a better representation of the higher orders effects than does the $`G_\mu `$ scheme of Ref. . We stress that this scheme $`(A)`$ is also gauge-invariant. The main effect of this change in renormalization scheme between versions 1.13 and 1.14 is to change the normalization of version 1.14 by $`0.3\%`$ to $`0.4\%`$ with respect to that of version 1.13 .
The generic size of the resulting shift in the YFSWW3 prediction, which we just quoted, can be understood by isolating the well-known soft plus virtual LL ISR correction to the process at hand, that has in $`𝒪(\alpha )`$ the expression
$$\delta _{ISR,LL}^{v+s}=\beta \mathrm{ln}k_0+\frac{\alpha }{\pi }\left(\frac{3}{2}L+\frac{\pi ^2}{3}2\right),$$
(24)
where $`\beta \frac{2\alpha }{\pi }(L1)`$$`L=\mathrm{ln}(s/m_e^2)`$, and $`k_0`$ is a dummy soft cut-off that cancels out of the cross section as usual. In the $`G_\mu `$ scheme of Refs. , which is used in YFSWW3-1.13, only the part $`\beta \mathrm{ln}k_0+(\alpha /\pi )(\pi ^2/3)`$ of $`\delta _{ISR,LL}^{v+s}`$ has the coupling $`\alpha (0)`$ and the remaining part of $`\delta _{ISR,LL}^{\mathrm{v}+\mathrm{s}}`$ has the coupling $`\alpha _{G_\mu }\alpha (0)/(10.0371)`$. The renormalization-group-improved YFS theory implies, however, that $`\alpha (0)`$ should be used for all the terms in $`\delta _{ISR,LL}^{v+s}`$. This is done in YFSWW3-1.14 and results in the normalization shift $`\left((\alpha (0)\alpha _{G_\mu })/\pi \right)(1.5L2)`$, whichis $`0.33\%`$ at 200 GeV. This explains most of the change in the normalization of YFSWW3-1.14 with respect to YFSWW3-1.13. Moreover, it does not contradict the expected total precision tag of either version of YFSWW3 at their respective stages of testing. We stress that, according to the renormalization-group equation, version 1.14 is an improvement over version 1.13 in that it better represents the true effect of the respective higher-order corrections.
For the purposes of cross checking with ourselves and with Ref. , we also created a second scheme, scheme ($`B`$), for the realization of the renormalization in YFSWW3-1.14. In this scheme, we put the entire $`𝒪(\alpha )`$ correction from Refs. at the coupling strength $`\alpha =\alpha (0)`$. Since the pure NL hard $`𝒪(\alpha )`$ correction is only $`0.006`$ at 200 GeV, scheme $`(B)`$ differs in the normalization from scheme $`(A)`$ by $`(\alpha (0)/\alpha _{G_\mu }1)(0.006)0.0002`$, which is well below the $`0.5\%`$ precision tag regime of interest for LEP2. Thus, scheme $`(B)`$, which is used in Ref. , is a perfectly acceptable scheme for LEP2 applications. It gives us a useful reference point from which to interpret our comparison with the results of RacoonWW from Ref. , which we discuss below.
Having presented our gauge-invariant calculation as it is realized in the MC YFSWW3-1.14, we will now turn to illustrating it in the context of LEP2 applications. Specifically, we always have in mind that one will combine the cross section from YFSWW3 with that from KoralW-1.42 MC, which is capable to calculate the non-resonant background contribution in a gauge-invariant way. We can do this in two ways, which we will now briefly describe and refer the reader to Refs. for the more detailed discussion. In the first way, we start with LPA<sub>a</sub> and we denote the corresponding cross section from YFSWW3-1.14 as $`\sigma (Y_a)`$. It is corrected for the missing background contribution by adding to it a correction $`\mathrm{\Delta }\sigma (K)`$ from KoralW-1.42 to form
$$\sigma _{Y/K}=\sigma (Y_a)+\mathrm{\Delta }\sigma (K),$$
(25)
where $`\mathrm{\Delta }\sigma (K)`$ is defined by
$$\mathrm{\Delta }\sigma (K)=\sigma (K_1)\sigma (K_3).$$
(26)
Here, the cross section $`\sigma (K_1)`$ is the complete 4-fermion result from KoralW-1.42 with all background diagrams and with the YFS-exponentiated $`𝒪(\alpha ^3)`$ LL ISR and the cross section $`\sigma (K_3)`$ is the restricted CC03 Born-level result from KoralW-1.42 – again with the YFS-exponentiated $`𝒪(\alpha ^3)`$ LL ISR. The result in eq. (25) is then accurate to $`𝒪(\frac{\alpha }{\pi }\frac{\mathrm{\Gamma }_W}{M_W})`$. Alternatively, one may start with the cross section for LPA<sub>a,b</sub> in YFSWW3-1.14, which we refer to as $`\sigma (Y_a)`$ and $`\sigma (Y_b)`$ correspondingly, and isolate the respective YFS-exponentiated $`𝒪(\alpha )`$ correction, $`\mathrm{\Delta }\sigma (Y)`$, which is missing from the cross section $`\sigma (K_1)`$ as
$$\mathrm{\Delta }\sigma (Y_j)=\sigma (Y_j)\sigma (Y_4),$$
(27)
where $`\sigma (Y_4)`$ is the corresponding cross section from YFSWW3-1.14, with the non-leading (NL) non-ISR $`𝒪(\alpha )`$ corrections to $`\overline{\beta }_n`$, $`n=0,1`$, switched off. Then the cross section
$$\sigma _{K/Y}=\sigma (K_1)+\mathrm{\Delta }\sigma (Y_j)$$
(28)
has the accuracy of $`𝒪(\frac{\alpha }{\pi }\frac{\mathrm{\Gamma }_W}{M_W})`$. We have checked that the results $`\sigma _{Y/K}`$ and $`\sigma _{K/Y}`$ are numerically equivalent at the 0.1% level of interest to us here. In the following we only show results from the former. For completeness, we also note that we sometimes identify $`\sigma (Y_1)=\sigma (Y_a),\sigma (Y_2)=\sigma (Y_b),\sigma (Y_3)=\sigma (K_3)`$, and $`\sigma (K_2)`$ is to be identified as the cross section from KoralW-1.42 with the restricted on-pole CC03 Born-level matrix element with YFS-exponentiated $`𝒪(\alpha ^3)`$ LL ISR. This latter cross section is a future option of KoralW . It would allow further combinations of YFSWW3 and KoralW with the desired $`𝒪(\frac{\alpha }{\pi }\frac{\mathrm{\Gamma }_W}{M_W})`$ accuracy. Such combinations would be of use in cross checks of our work.
We now illustrate our precision predictions using $`\sigma _{Y/K}`$. We have checked that the correction $`\mathrm{\Delta }\sigma (K)`$ is small, $`0.1\%`$ for CMS energies $`200`$ GeV. This is summarized in Tables 1 and 2, in which we compare the size of the correction $`\mathrm{\Delta }\sigma (K)`$, labeled $`\delta _{4f}`$, with the size of the respective NL non-ISR $`𝒪(\alpha )`$ correction for the 4-lepton, 2-lepton–2-quark, and 4-quark final states, with and without the cuts of Ref. at 200 GeV.
Thus, in what follows, we shall ignore $`\mathrm{\Delta }\sigma (K)`$, as our ultimate precision tag objective, $`<0.5\%`$, does not require that we keep it. It will be analyzed in more detail elsewhere . Further, for the cross section $`\sigma (Y_a)`$ we have already presented, for version 1.13, in Figs. 1-8 of Ref. , for the $`c\overline{s}\mathrm{}\overline{\nu }_{\mathrm{}}`$, $`\mathrm{}=e^{},\mu ^{}`$ final states, the $`W^{+,}`$ angular distributions in the $`e^+e^{}`$ CMS system, the $`W^{+,}`$ mass distributions, the distributions of the final-state lepton energies in the LAB frame ($`e^+e^{}`$ CMS frame), and the final-state lepton angular distributions in the $`W^{}`$ rest frame their corrections (relative to the Born-level). The main effect on these differential distributions of the improved normalization of version 1.14 is to shift the normalization, as we discussed above. Thus, we do not repeat their presentation here. We refer the reader to the results in Ref. for an investigation into the size of the EW=NL and FSR effects in the cases listed above insofar as YFSWW3-1.14 is concerned with the understanding that the shapes of the distributions apply directly to version 1.14, but that the normalization of the EW correction should be reduced by $`0.3\%`$ to $`0.4\%`$. In general, we found in Ref. that, depending on the experimental cuts and acceptances, both the FSR and the EW corrections were important in precision studies of these distributions; this conclusion still holds for version 1.14, of course. For example, in the lepton decay angle distribution, for the BARE acceptance (the final charged lepton is not combined with any photons), where both the FSR and the EW correction modulate the distribution, whereas for the CALO acceptance of Ref. (all photons within $`5^{}`$ of the final-state charged lepton are combined with it) the FSR effect is almost nil whereas the EW correction effect remains at the level of $`2.0\%`$. Here, we focus on the total CMS photon energy distribution (Fig. 1a), and the CMS photon angular distribution (Fig. 1b). We show these results both for the BARE and CALO acceptances, as defined in the 4 fermion Section of the Proceedings of the LEP2 MC Workshop .
In Fig. 1a, we see that the total photon energy distributions are different for the BARE and CALO cases but that the NL non-ISR correction does not affect them strongly. In Fig. 1b, we see that, for both the BARE and CALO cases, the NL correction does affect the photon angular distribution away from the beam directions, as we expect. Note that this is the NL correction implied by the YFS exponentiation of our exact $`𝒪(\alpha )_{prod}`$ correction. Finally, in Fig. 2, we show the effect, in the $`W`$ mass and angular distributions, of using the screened Coulomb correction from Ref. , as against the usual Coulomb correction from Ref. . The effect we see is a 5 MeV shift in the peak position, associated with the difference between the screened and usual Coulomb corrections; we see almost no effect, as expected, associated with this difference on the $`e^+e^{}`$ CMS W angular distribution. Since we calculate the finite $`p_T`$ $`n(\gamma )`$ corrections to these distributions, these results are new. Indeed, in Ref. it is shown that the results from RacoonWW and YFSWW3 (Best) for the distribution in Fig. 1a differ by $`20\%`$ and, as we have the dominant $`𝒪(\alpha ^2)`$ LL corrections to this distribution whereas in Fig. 20 in Ref. the RacoonWW result only has the exact $`𝒪(\alpha )`$ Born result for the hard photon observable, we expect that most of this discrepancy would be removed if the dominant $`𝒪(\alpha ^2)`$ LL corrections were included in the RacoonWW results. This has recently been confirmed in Ref. , where the authors of RacoonWW show that, when they include the latter corrections in their predictions for the $`\mathrm{cos}\theta _\gamma `$ distribution in Fig. 1a, the discrepancy is reduced to the level of $`5\%`$. In summary, from the results in Ref. and those presented here, we see that the FSR and EW corrections are necessary for a precision study of the distributions in the $`W`$-pair production and decay process at LEP2 energies.
We have made a detailed comparison between our results and those from Ref. based on the program RacoonWW in the context of the LEP2 MC Workshop . A complete unpublished preliminary description of the respective results of this comparison has appeared in Ref. . Here, we focus on the normalization comparison of the two calculations at LEP2 energies. We show in Table 3 the comparison of the RacoonWW and YFSWW3-1.14 results for the cross sections as indicated, without cuts at 200 GeV (we have looked at the lower energies 184 and 189 GeV and the comparisons there are similar, if not better). In Table 4, we show the analogous comparisons with the LEP2 MC Workshop cuts as described in Ref. .
We see that for all channels considered, the two sets of results agree to the level of 0.3%. This gives a total precision estimate of 0.4% for the theoretical uncertainty on the 200 GeV CMS energy $`WW`$ signal cross-section normalization when allowance is made for further possible uncertainties in the higher-order radiative corrections and the implementation of the LPA . This is a significant improvement over the originally quoted $`2\%`$ for this uncertainty when the NL non-ISR $`𝒪(\alpha )`$ corrections are not taken into account . An effort to further reduce this 0.4% is in progress.
Finally, with an eye toward the LC projects, we have made simulations using YFSWW3-1.14 for a CMS energy of 500 GeV. We show our results in Table 5 for the total cross section without cuts; here we again compare them to the corresponding ones from RacoonWW . The NL corrections are significant in these results. Precision studies at LC energies must take these effects into account. As expected, the percentage difference between YFSWW3-1.14 and RacoonWW remains below 0.5% at 500 GeV CMS energy and is somewhat larger than at 200 GeV CMS energy.
In summary, we have presented two recipes for combining YFSWW3 and KoralW-1.42 to arrive at a gauge-invariant calculation of the $`WW`$-pair production and decay in which the YFS-exponentiated exact $`𝒪(\alpha )_{prod}`$ corrections are taken into account as well as the $`𝒪(\alpha ^2)`$ LL FSR and YFS-exponentiated $`𝒪(\alpha ^3)`$ LL ISR correction to the background processes. We have illustrated our calculation with several sample MC results and we have compared our results on the cross-section normalization with those on Refs. at 200 GeV. In this way, new precision tag of 0.4% has been established for this normalization, which represents a considerable improvement over the original result of $`2\%`$ when NL non-ISR corrections are not taken into account.
Acknowledgments
Two of us (S.J. and B.F.L.W.) acknowledge the kind hospitality of Prof. A. De Rújula and the CERN Theory Division while this work was being completed. Three of us (B.F.L.W., W.P. and S.J.) acknowledge the support of Prof. D. Schlatter and of the ALEPH, DELPHI, L3 and OPAL Collaborations in the final stages of this work. S.J. is thankful for the kind support of the DESY Directorate and Z.W. acknowledges the support of the L3 Group of ETH Zurich during the time this work was performed. All of us thank the members of the LEP2 MC Workshop for valuable interactions and stimulation during the course of this work. The authors especially thank Profs. A. Denner, S. Dittmaier and F. Jegerlehner and Drs. M. Roth and D. Wackeroth for useful discussions and interactions.
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# 1 Introduction
## 1 Introduction
The 1/N expansion of matrix–valued field theories is probably the most important nonperturbative and nonnumerical theoretical tool presently available in the study of such models as non–Abelian gauge theories and two–dimensional quantum gravity. The first major result concerning the large–N limit of matrix theories was due to ’t Hooft, who made the crucial observation that in the 1/N expansion of continuum gauge theories, the set of Feynman diagrams contributing to any given order admits a simple topological interpretation .
Unfortunately, our knowledge of exact solution for the large–N limit is limited to a small number of few–matrix systems, and it becomes even smaller if we restrict ourselves to the case of unitary matrix fields, which is relevant to the problem of lattice QCD. After the solution of Gross and Witten of the single link problem and its generalizations, exact results were obtained only for the external field problem and a few toy models (L=3,4 chiral chains). Therefore extending the number of solved few–matrix systems is an important progress from different points of view: first from pure theoretical reason. Second by noticing that any few–matrix system have a reinterpretation, via double scaling limit, as some different kind of matter coupled to two–dimensional quantum gravity. And third because every few–matrix system involving unitary matrices can be reinterpreted as the generating functional for a class of integrals over unitary groups, and these integrals in turn are the essential missing ingredient in the context of a complete algorithmization of the strong coupling expansion of many interesting models .
One of the interesting classes of the few–matrix models is the simplicial chiral model (SCM). In this model, each unitary matrix interacts in a fully symmetric way with all other matrices, by preserving the global chiral invariance. The resulting system can be described as a chiral model on $`(d1)`$–dimensional simplex. A simplex is formed by connecting $`d`$ vertices by $`d(d1)/2`$ links. In ref., where the large–N saddle–point equations for density function $`\rho (z)`$ have been found, the authors have introduced a single auxiliary variable $`A`$ (a complex matrix) to decouple the matrix interaction, and by performing the single–link external field integral, the saddle–point equation in strong and weak regions have been found. In $`d=2`$, $`d=4`$, and $`d\mathrm{}`$, the saddle–point equation has been solved analytically and the phase transition of the model has been specified. In the action of SCM in terms of the matrix $`A`$, this matrix appears as a $`\mathrm{Tr}(AA^{})`$ term.
On the other hand, it is well known that the pure 2–dimensional Yang–Mills theory (YM<sub>2</sub>) can be represented by the Lagrangian $`i\mathrm{Tr}(BF)+\mathrm{Tr}(B^2)`$, in which $`F`$ is the usual field–strength tensor and $`B`$ is an auxiliary pseudo–scalar field. Many of the main properties of YM<sub>2</sub> does not change if one considers the generalized two–dimensional YM<sub>2</sub> (gYM<sub>2</sub>), which can be found by replacing $`\mathrm{Tr}(B^2)`$ term in YM<sub>2</sub> action by an arbitrary class function $`f(B)`$. The partition function of these theories have been fully studied in different contexts in refs. and the phase structure of some of the specific examples has been studied in refs. and . In all the studied cases, it is seen that the models have third–order phase transition, which is the same behavior as YM<sub>2</sub>. The crucial point in this area is that such investigation for continuum gYM<sub>2</sub> is very complicated and there is no general result for the phase transition of an arbitrary gYM<sub>2</sub> theory.
In ref. , the procedure used to obtain the gYM<sub>2</sub> from YM<sub>2</sub>, i.e. $`\mathrm{Tr}(B^2)f(B)`$, has been used in SCM to introduce the generalized simplicial chiral model (gSCM). That is, the $`\mathrm{Tr}(AA^{})`$ term in SCM has been replaced by an arbitrary class function $`V(AA^{})`$. The large–N limit of the model has been discussed and the saddle–point equations in strong and weak regimes have been found. Note that as the SCM at $`d=2`$ is the discrete version of YM<sub>2</sub> , we can consider the gSCM, at $`d=2`$, as the equivalent matrix theory of gYM<sub>2</sub>. The phase structure of the model for $`d=2`$ and $`V=\mathrm{Tr}(AA^{})^n`$ ($`n=2,3,4`$) is also obtained in .
In this paper we want to complete our investigation of $`d=2`$ gSCM and show that all $`V=\mathrm{Tr}(AA^{})^n`$ models (with $`n`$ an arbitrary positive integer) have third–order phase transition. We think that it is an important result as it indirectly proves the equivalence of all $`\mathrm{Tr}(B^n)`$ gYM<sub>2</sub> theories with YM<sub>2</sub>, from the phase structure point of view. The plane of the paper is as follows. In section 2, we bring a brief review of SCM and gSCM and also the saddle–point equations in weak and strong regimes. In section 3, we focus on $`d=0`$ case and show that the calculation of the partition function of the model, described in terms of $`A`$ fields, leads to the trivial result which obtained from the main action. This is a consistency check of our procedure. Finally in section 4 we investigate the phase structure of $`V=\mathrm{Tr}(AA^{})^n`$ gSCM in $`d=2`$ and show that all of these models have third–order phase transition. We also discuss the variation of the discontinuity of $`\beta ^2C(\beta ,N)/\beta `$ with $`n`$ at the critical point ($`C`$ is the specific heat and $`\beta =(g^2N)^1`$, where $`g`$ is the coupling constant).
## 2 The gSCM
If we assign a U($`N`$) matrix to each vertex of a $`(d1)`$–dimensional simplex, then the partition function of simplicial chiral models is defined by
$$Z_d(\beta ,N)=\underset{i=1}{\overset{d}{}}dU_i\mathrm{exp}\{N\beta \underset{i=1}{\overset{d}{}}\underset{j=i+1}{\overset{d}{}}\mathrm{Tr}(U_iU_j^{}+U_i^{}U_j)\}.$$
(1)
To find this partition function, the authors of have introduced a single auxiliary variable $`A`$ to decouple the matrix interactions. The resulting partition function is
$$Z_d=\stackrel{~}{Z}_d/\stackrel{~}{Z}_0,$$
(2)
where
$$\stackrel{~}{Z}_d=\underset{i=1}{\overset{d}{}}dU_idA\mathrm{exp}\{N\beta \mathrm{Tr}AA^{}+N\beta \mathrm{Tr}A\underset{i}{}U_i^{}+N\beta \mathrm{Tr}A^{}\underset{i}{}U_iN^2\beta d\}.$$
(3)
Performing the single–link integral over the matrices $`U_i`$
$$e^{NW(BB^{})}𝑑U\mathrm{exp}[N\mathrm{Tr}(B^{}U+U^{}B)],$$
(4)
we obtain
$$\stackrel{~}{Z}_d=𝑑A\mathrm{exp}\{N\beta \mathrm{Tr}AA^{}+NdW(\beta ^2AA^{})N^2\beta d\}.$$
(5)
The gSCM is defined through the action
$$Z_{d,V}=\stackrel{~}{Z}_{d,V}/\stackrel{~}{Z}_{0,V},$$
(6)
in which
$$\stackrel{~}{Z}_{d,V}=𝑑A\mathrm{exp}\{N\beta V(AA^{})+NdW(\beta ^2AA^{})N^2\beta d\},$$
(7)
where $`V(AA^{})`$ is an arbitrary class function of $`AA^{}`$:
$$V(AA^{})=\underset{n=1}{}a_n\mathrm{Tr}(AA^{})^n.$$
(8)
In special case $`\alpha _n=\delta _{n,1}`$, the gSCM reduces to SCM. If we denote the eigenvalues of the Hermitian semi–positive–definite matrix $`\beta ^2AA^{}`$ by $`x_i`$’s, and at the large–N limit, which we are interested in, using the expression $`W(x_i)`$ in the strong and weak regimes , it can be shown that the saddle–point equation for the eigenvalue density function $`\rho (z)`$ ($`z_i=\sqrt{x_i+c}`$) is :
$$z\underset{k=1}{}\frac{ka_k}{\beta ^{2k1}}(z^2c)^{k1}d=\frac{1}{2}𝒫_a^b𝑑z^{}\rho (z^{})\left(\frac{2}{zz^{}}\frac{d2}{z+z^{}}\right).$$
(9)
In the above equation $`𝒫`$ indicates the principal value of integral, and the parameters $`a`$ and $`b`$ must be determined dynamically. The normalization condition of $`\rho (z)`$ is
$$_a^b\rho (z^{})𝑑z^{}=1.$$
(10)
In the weak coupling regime ($`\beta >\beta _c`$), $`c_w=0`$ and the density function $`\rho (z)`$ must satisfy
$$_a^b𝑑z^{}\frac{\rho (z^{})}{z^{}}2.$$
(11)
In the strong coupling regime $`(\beta <\beta _c)`$, $`c_s=a^2`$ and $`\rho (z)`$ must satisfy
$$_a^b𝑑z^{}\frac{\rho (z^{})}{z^{}}=2.$$
(12)
Using the standard method of solving the integral equation , it can be shown that the density function $`\rho (z)`$ in weak regime satisfies the following equation
$`\rho _w(z)={\displaystyle \frac{\sqrt{(bz)(za)}}{\pi }}`$ $`\{`$ $`{\displaystyle \underset{m,p,q}{}}{\displaystyle \frac{ma_m}{\beta ^{2m1}}}C_pC_{2mpq2}z^qa^pb^{2mpq2}`$ (14)
$`{\displaystyle \frac{d2}{2}}{\displaystyle _a^b}{\displaystyle \frac{dy}{y+z}}{\displaystyle \frac{\rho _w(y)}{\sqrt{(b+y)(y+a)}}}\},\mathrm{for}\beta >\beta _c,`$
where $`C_n=(2n1)!!/(2^nn!)`$, and in strong regime satisfies
$`\rho _s(z)={\displaystyle \frac{z}{\pi }}\sqrt{{\displaystyle \frac{bz}{za}}}`$ $`\{`$ $`{\displaystyle \underset{n,m,p,q}{}}{\displaystyle \frac{na_n}{\beta ^{2n1}}}(a^2)^p\left({\displaystyle \genfrac{}{}{0pt}{}{n1}{p}}\right)B_mC_{2n2pmq2}z^qa^mb^{2n2pmq2}`$ (17)
$`+{\displaystyle \frac{d2}{2}}{\displaystyle _a^b}{\displaystyle \frac{dy}{y+z}}\sqrt{{\displaystyle \frac{y+a}{y+b}}}{\displaystyle \frac{\rho _s(y)}{y}}\},\mathrm{for}\beta <\beta _c,`$
where $`B_m=(2m3)!!/(2^mm!)`$ ($`B_01`$) . Also by investigating the behavior of the integrals at $`z\mathrm{}`$, it can be shown that the parameters $`a`$ and $`b`$ in $`\beta >\beta _c`$ must be determined from the following equations
$$\underset{n,m}{}\frac{na_n}{\beta ^{2n1}}C_mC_{2nm1}a^mb^{2nm1}2=0,$$
(18)
and
$$\underset{n,m}{}\frac{na_n}{\beta ^{2n1}}C_mC_{2nm}a^mb^{2nm}(a+b)=1,$$
(19)
and in strong regime, $`\beta <\beta _c`$, from the equations
$$\underset{n,m,p}{}\frac{na_n}{\beta ^{2n1}}(a^2)^p\left(\genfrac{}{}{0pt}{}{n1}{p}\right)B_mC_{2n2pm1}a^mb^{2n2pm1}+2=0,$$
(20)
and
$$\underset{n,m,p}{}\frac{na_n}{\beta ^{2n1}}(a^2)^p\left(\genfrac{}{}{0pt}{}{n1}{p}\right)B_mC_{2n2pm}a^mb^{2n2pm}+ab=1.$$
(21)
Finally if we denote the internal energy per unit link by $`U_{d,V}`$, we have by definition
$$\frac{d(d1)}{2}U_{d,V}=\frac{1}{2N^2}\frac{}{\beta }\mathrm{ln}Z_{d,V}(\beta ,N).$$
(22)
After some calculation, it can be shown that at large–N limit, the internal energies in weak and strong regimes for $`V=\mathrm{Tr}(AA^{})^n`$ are:
$$d(d1)U_{d,n}^{(w)}=\frac{2n1}{\beta ^{2n}}_a^bz^{2n}\rho _w(z)𝑑zd+(\frac{1}{n}2)\frac{1}{\beta },$$
(23)
and
$$d(d1)U_{d,n}^{(s)}=\frac{2n1}{\beta ^{2n}}_a^b(z^2a^2)^n\rho _s(z)𝑑zd+(\frac{1}{n}2)\frac{1}{\beta }.$$
(24)
These equations can be used to deduce the order of phase transition of the models.
## 3 $`d=0`$ gSCM with $`V=\mathrm{Tr}(AA^{})^n`$
It is clear from the definition of $`Z_{d,V}`$ in eq.(6) that $`Z_{0,V}=1`$, however $`\stackrel{~}{Z}_{0,V}`$ is not trivial. To see this, let us focus on $`V=\mathrm{Tr}(AA^{})^n`$ from now on. It is not difficult to show that $`\stackrel{~}{Z}_{0,n}`$ from eq.(7) is
$$\stackrel{~}{Z}_{0,n}\mathrm{exp}[\frac{N^2(2n1)}{2}\mathrm{ln}\beta ].$$
(25)
Now to check the procedure introduced in the last section for studying the theory in the large–N limit, it is instructive to reproduce this result by using the density function $`\rho (z)`$ at $`d=0`$.
At $`d=0`$, the saddle–point equation (9) in weak regime ($`c=0`$), for the case $`a_k=\delta _{k,n}`$, becomes
$$\frac{n}{2\beta ^{2n1}}z^{2n2}=𝒫_a^b\frac{\rho _w(z^{})}{z^2z^2}dz^{},\beta >\beta _c$$
(26)
and in the strong regime ($`c=a^2`$) becomes
$$\frac{n}{2\beta ^{2n1}}(z^2a^2)^{n1}=𝒫_a^b\frac{\rho _s(z^{})}{z^2z^2}dz^{},\beta <\beta _c.$$
(27)
Let us focus on $`\beta <\beta _c`$ case. To make the integral equation (27) more conventional, we use the change of variable $`z\lambda =z^2`$, with density function $`\rho _s(\lambda )`$ satisfies
$$\rho _s(\lambda )d\lambda =\rho _s(z^{})dz^{}.$$
(28)
So
$$\frac{n}{2\beta ^{2n1}}(\lambda a^2)^{n1}=𝒫_{a^2}^{b^2}\frac{\rho _s(\xi )}{\lambda \xi }d\xi ,\beta <\beta _c.$$
(29)
Now consider the function $`H_s(\lambda )`$ in complex–$`\lambda `$ plane
$$H_s(\lambda )=_{a^2}^{b^2}\frac{\rho _s(\xi )}{\lambda \xi }𝑑\xi .$$
(30)
This function is analytic on the entire complex plane except for a cut on the positive real axis in the interval $`[a^2,b^2]`$. Then one has
$$H_s(\lambda \pm iϵ)=R_s(\lambda )i\pi \rho _s(\lambda ),a^2\lambda b^2,$$
(31)
where $`R_s(\lambda )`$ is, from eq.(29),
$$R_s(\lambda )=\frac{n}{2\beta ^{2n1}}(\lambda a^2)^{n1}.$$
(32)
Using the standard method of solving the integral equations , one can show that the expression
$$H_s(\lambda )=\frac{1}{2\pi i}\sqrt{\frac{\lambda b^2}{\lambda a^2}}_c\sqrt{\frac{\xi a^2}{\xi b^2}}\frac{R_s(\xi )}{\lambda \xi }𝑑\xi ,$$
(33)
has the correct analytical behavior in $`\beta <\beta _c`$ region (see for more details). Here the contour $`c`$ is a contour encircling the cut $`[a^2,b^2]`$ and excluding $`\lambda `$. Deforming $`c`$ to a contour around the point $`\lambda `$ and the contour $`c_{\mathrm{}}`$ (a contour at infinity), one finds
$$H_s(\lambda )=R_s(\lambda )+\frac{1}{2\pi i}\sqrt{\frac{\lambda b^2}{\lambda a^2}}_c_{\mathrm{}}\sqrt{\frac{\xi a^2}{\xi b^2}}\frac{R_s(\xi )}{\lambda \xi }𝑑\xi .$$
(34)
Inserting (32) in (34), it can be easily shown that $`\rho _s(\lambda )`$ is (using (31)):
$$\rho _s^{(n)}(\lambda )=\frac{n}{2\pi \beta ^{2n1}}\sqrt{\frac{b^2\lambda }{\lambda a^2}}\underset{m,p,s=0}{}(1)^s\left(\genfrac{}{}{0pt}{}{n1}{s}\right)B_pC_{npms1}\lambda ^ma^{2(p+s)}b^{2(npms1)}.$$
(35)
At $`n=1`$, where gSCM reduces to SCM, it can be shown that (35) is equal
$$\rho _s^{n=1}(\lambda )=\frac{1}{2\pi \beta }\sqrt{\frac{b^2\lambda }{\lambda a^2}},$$
(36)
which is the same as one calculated in . To find the parameters $`a`$ and $`b`$ in eq.(35), we note that at $`|\lambda |\mathrm{}`$, eqs.(30) and (10) imply $`H(\lambda )1/\lambda `$ or $`\sqrt{(\lambda a^2)/(\lambda b^2)}H(\lambda )1/\lambda `$. Therefore we can expand $`\sqrt{(\lambda a^2)/(\lambda b^2)}(R_s(\lambda )i\pi \rho _s(\lambda ))`$ and take the coefficients of $`\lambda ^0`$ and $`1/\lambda `$ equal to 0 and 1, respectively. It can be see that the coefficient of $`\lambda ^0`$ is identically zero, and the second condition yields to (for $`\beta <\beta _c`$)
$$\frac{n}{2\beta ^{2n1}}\underset{p,s=0}{}(1)^s\left(\genfrac{}{}{0pt}{}{n1}{s}\right)B_pC_{nsp}a^{2(p+s)}b^{2(nps)}=1.$$
(37)
In $`n=1`$ case, eq.(37) reduces to $`b^2a^2=4\beta `$, which is one obtained in . The same calculation for the density function $`\rho _w(\lambda )`$ yields to
$`\rho _w^{(n)}(\lambda )=\{\begin{array}{cc}0n=1\hfill & \\ \frac{n}{2\pi \beta ^{2n1}}\sqrt{(b^2\lambda )(\lambda a^2)}_{p,s=0}C_pC_{nps2}\lambda ^sa^{2p}b^{2(nps2)}n>1\hfill & \end{array}`$ (40)
and the following equations which specify the parameters $`a`$ and $`b`$ in $`\beta >\beta _c`$,
$$\frac{n}{2\beta ^{2n1}}\underset{p=0}{}C_pC_{np1}a^{2p}b^{2(np1)}=0,$$
(41)
$$\frac{n}{2\beta ^{2n1}}\underset{p=0}{}C_pC_{np}a^{2p}b^{2(np)}=1.$$
(42)
Now let us focus on critical point $`\beta =\beta _c`$ in which $`a=a_c`$ and $`b=b_c`$ must satisfy eqs.(37), (41) and (42) simultanously. It is not difficult to see at $`\beta =\beta _c`$, except the first term of eqs. (37) and (42) which is equal, all the other terms are not the same. The only unique solution of this inconsistency is
$$a_c=0.$$
(43)
Inserting (43) in eqs.(41) and (42)(or (37)), results
$$\frac{b_c^{2n2}}{\beta _c^{2n1}}=0,$$
(44)
$$\frac{n}{2\beta _c^{2n1}}C_nb_c^{2n}=1.$$
(45)
The solution of these equations are
$`b_c`$ $``$ $`\mathrm{},`$
$`\beta _c`$ $``$ $`\mathrm{},`$ (46)
such that eq.(45) holds. Now in the weak–coupling region, where $`\rho _w1/\beta ^{2n1}`$, $`\beta `$ is always greater than $`\beta _c=\mathrm{}`$, so $`\rho _w0`$. Therefore there is only one regime at $`d=0`$, i.e. the strong–coupling regime.
To check eq.(25), it is easier first to calculate the contribution of $`\stackrel{~}{Z}_{0,n}`$ to internal energy $`U_{0,n}^{(s)}`$ and then find $`\stackrel{~}{Z}_{0,n}`$ from it. If we denote this contribution by $`\stackrel{~}{U}_{0,n}^{(s)}`$, it is equal to (see eq.(24))
$`\stackrel{~}{U}_{0,n}^{(s)}`$ $`=`$ $`{\displaystyle \frac{2n1}{2\beta ^{2n}}}{\displaystyle _a^b}(z^2a^2)^n\rho _s^{(n)}(z)𝑑z`$ (47)
$`=`$ $`{\displaystyle \frac{2n1}{2\beta ^{2n}}}{\displaystyle _{a^2}^{b^2}}(\lambda a^2)^n\rho _s^{(n)}(\lambda )𝑑\lambda ,`$
in which $`\rho _s^{(n)}(\lambda )`$ is given by eq.(35) and the relation between $`a`$ and $`b`$ can be obtained from (37). To see the consistency of the results, it is sufficient to check some specific examples. In $`n=2`$ , (47) reduces to
$$\stackrel{~}{U}_{0,2}^{(s)}=\frac{27}{256\beta ^7}(b^2a^2)^4,$$
(48)
and (37) reduces to
$$\frac{3}{8}(b^2a^2)^2=\beta ^3,$$
(49)
therefore
$$\stackrel{~}{U}_{0,2}^{(s)}=\frac{3}{4\beta }=\frac{1}{2N^2}\frac{}{\beta }\mathrm{ln}\stackrel{~}{Z}_{0,2}.$$
(50)
This equation gives $`\stackrel{~}{Z}_{0,2}\mathrm{exp}(\frac{3N^2}{2}\mathrm{ln}\beta )`$, which is the same behavior as eq.(25). It can be also shown that for $`n=3,4`$, and $`5`$, these two methods have the same result. As another reason for the fact that only the strong–coupling regime exists in $`d=0`$, it can be seen that the $`\stackrel{~}{Z}_{0,n}`$ calculated from $`\rho _w(\lambda )`$ (eq.(40)) does not coincide with eq.(25).
## 4 Phase structure of gSCM at $`d=2`$
As mentioned earlier, the gSCM can be solved analytically only in $`d=0`$ and $`d=2`$ cases, which the latter is an important case because of its equivalence to gYM<sub>2</sub> (and in special YM<sub>2</sub>). So in this section we want to study the gSCM at $`d=2`$ for $`V=\mathrm{Tr}(AA^{})^n`$.
It is clear from eqs.(14) and (17) that the density function are known in the weak and strong regimes at $`d=2`$, as the second term in the right–hand sides of these equations vanishes. It can also be shown that at $`d=2`$, $`\rho _w(0)`$ must be equal to zero at critical point $`\beta =\beta _c`$. In this way, one can obtain the precise value of $`a`$, $`b`$, and $`\beta `$ at critical point as following :
$`a_c`$ $`=`$ $`0,`$
$`b_c`$ $`=`$ $`{\displaystyle \frac{2n}{2n1}},`$
$`\beta _c`$ $`=`$ $`{\displaystyle \frac{2n}{2n1}}[{\displaystyle \frac{n(4n3)!!}{2^{2n}(2n1)!}}]^{\frac{1}{2n1}}.`$ (51)
To study the phase structure of this model, it is necessary to calculate the internal energies in both weak and strong regimes. We know the functional form of the density function in these two regimes, but there are two unknown parameters $`a`$ and $`b`$ in each of these densities which must be obtained from the eqs.(18) and (19) for weak regime and from eqs.(20) and (21) for strong regime. Unfortunately these equations are too complicated to be solved exactly, but the crucial point is that as we want to study the phase structure of the model, it is sufficient to look at the solutions near the critical point. Therefore we expand the equations (18) and (19) around $`a_c=0`$ and $`b_c=2n/(2n1)`$, and find $`a_w`$ and $`b_w`$ in terms of $`\alpha =\beta \beta _c`$ up to second order. After a lengthy calculation one finds
$`a_w`$ $`=`$ $`{\displaystyle \frac{4n3}{2(2n1)\beta _c}}\alpha {\displaystyle \frac{4n3}{2n(4n5)\beta _c^2}}\alpha ^2+\mathrm{}`$
$`b_w`$ $`=`$ $`{\displaystyle \frac{2n}{2n1}}+{\displaystyle \frac{4n1}{2(2n1)\beta _c}}\alpha {\displaystyle \frac{1}{8n\beta _c^2}}\alpha ^2+\mathrm{}`$ (52)
Now inserting $`\rho _w(z)`$ from eq.(14), for $`a_m=\delta _{n,m}`$ and $`d=2`$, into eq.(23), we arrive at
$$U_{2,n}^{(w)}=\frac{n(2n1)}{2\pi \beta ^{4n1}}\underset{p,q,r=0}{}B_rC_pC_{2npq2}a^{p+r}b^{4npq2}I(2n+qr)1+(\frac{1}{2n}1)\frac{1}{\beta },$$
(53)
in which
$$I(k)=\frac{\sqrt{\pi }b^{2+k}\mathrm{\Gamma }(3/2+k)}{2\mathrm{\Gamma }(3+k)}\frac{a^{\frac{3}{2}+k}\sqrt{b}}{\frac{3}{2}+k}{}_{2}{}^{}F_{1}^{}[\frac{1}{2},\frac{3}{2}+k,\frac{5}{2}+k,\frac{a}{b}],$$
(54)
where $`{}_{2}{}^{}F_{1}^{}[a,b,c,x]`$ is the hypergeometric function. Keeping the terms up to second order of $`a_w^2`$, eq.(53) becomes:
$$U_{2,n}^{(w)}=K(T_1b_w^{4n}\frac{1}{2}T_2b_w^{4n1}a_w+T_3b_w^{4n2}a_w^2+\mathrm{})1+(\frac{1}{2n}1)\frac{1}{\beta },$$
(55)
where
$$K=\frac{n(2n1)}{4\sqrt{\pi }\beta ^{4n1}},$$
$$T_1=\underset{q=0}{}C_{2nq2}M(n,q,0),$$
$$T_2=C_{2n2}M(n,0,1),$$
$$T_3=\underset{q=0}{}[\frac{3}{8}C_{2nq4}M(n,q,0)\frac{1}{4}C_{2nq3}M(n,q,1)\frac{1}{8}C_{2nq2}M(n,q,2)],$$
(56)
in which
$$M(n,q,r)=\frac{\mathrm{\Gamma }(3/2+2n+qr)}{\mathrm{\Gamma }(3+2n+qr)}.$$
(57)
In the strong regime, the same calculation leads to
$`a_s`$ $`=`$ $`{\displaystyle \frac{4n3}{2(2n1)\beta _c}}\alpha {\displaystyle \frac{(n1)(4n3)}{2n(4n5)\beta _c^2}}\alpha ^2+\mathrm{},`$
$`b_s`$ $`=`$ $`{\displaystyle \frac{2n}{2n1}}+{\displaystyle \frac{4n1}{2(2n1)\beta _c}}\alpha +{\displaystyle \frac{n1}{2n\beta _c^2}}\alpha ^2+\mathrm{},`$ (58)
for the parameters $`a`$ and $`b`$, and
$$U_{2,n}^{(s)}=K(T_1b_s^{4n}+\frac{1}{2}T_2b_s^{4n1}a_s+S_3b_s^{4n2}a_s^2+\mathrm{})1+(\frac{1}{2n}1)\frac{1}{\beta },$$
(59)
for internal energy. $`S_3`$ in eq.(59) is
$$S_3=\underset{q=0}{}[(\frac{7}{8}n)C_{2nq4}M(n,q,0)\frac{1}{4}C_{2nq3}M(n,q,1)+(\frac{3}{8}n)C_{2nq2}M(n,q,2)].$$
(60)
Now subtracting the internal energies in two regions (eqs.(59) and (55)) and using eqs.(4) and (4), one obtains
$$U_{2,n}^{(s)}U_{2,n}^{(w)}=f(n)(\beta \beta _c)^2+\mathrm{},$$
(61)
where
$$f(n)=\frac{(2n1)(4n3)b_c^{4n1}}{32\sqrt{\pi }\beta _c^{4n+1}}\{4nT_1\frac{4n3}{2(4n5)}T_2\frac{4n3}{2}T_4\},$$
(62)
in which
$$T_4=C_{2n2}M(n,0,2)+C_{2n3}M(n,0,1)+2\underset{q=0}{}C_{2nq4}M(n,q,0).$$
(63)
Eq.(61) shows that all the $`V=\mathrm{Tr}(AA^{})^n`$ gSCM has a third order phase transition near the critical point $`\beta =\beta _c`$. Note that $`f(n)>0`$ for all values of $`n`$, as we expected from the definitions of strong and weak regimes, and also the value of $`f(n)`$ for $`n=2,3,`$ and $`4`$ coincides with those in ref..
As a final comment, it is interesting to compare the behavior of the solutions with respect to $`n`$. The only physical comparable quantity, in the case of third order phase transition, is the difference
$$\frac{C_s}{T}\frac{C_w}{T}=\beta ^2\frac{}{\beta }(\beta ^2\frac{}{\beta }(U^{(s)}U^{(w)})),$$
(64)
where $`C`$ is the specific heat of the theory and $`T=1/\beta `$ is temperature. Using (61), it is found
$$\beta ^2\frac{}{\beta }(\mathrm{\Delta }C)_n=2\beta _c^4f(n).$$
(65)
In Fig. (1), where we have plotted $`\beta ^2\frac{}{\beta }(\mathrm{\Delta }C)_n`$ vs $`n`$, one can see a saturating value 2 for this difference.
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# References
BRX-TH-475
BOW-PH-118
HUTP-00/A021
Seiberg-Witten curves for elliptic models <sup>1</sup><sup>1</sup>1Based on a talk by S. G. Naculich at the Workshop on Strings, Duality and Geometry,
aaa University of Montreal, March 2000.
Isabel P. Ennes<sup>2</sup><sup>2</sup>2Research supported by the DOE under grant DE–FG02–92ER40706.<sup>,a</sup>, Carlos Lozano<sup>2</sup><sup>2</sup>footnotemark: 2<sup>,a</sup>, Stephen G. Naculich<sup>3</sup><sup>3</sup>3Research supported in part by the National Science Foundation under grant no. PHY94-07194 through the
aaa ITP Scholars Program.<sup>,b</sup>, Howard J. Schnitzer<sup>4</sup><sup>4</sup>4Permanent address. $`^,`$ <sup>5</sup><sup>5</sup>5Research supported in part by the DOE under grant DE–FG02–92ER40706.
aaa naculich@bowdoin.edu; ennes,lozano,schnitzer@brandeis.edu
<sup>,a,c</sup>
<sup>a,4</sup>Martin Fisher School of Physics
Brandeis University, Waltham, MA 02454
<sup>b</sup>Department of Physics
Bowdoin College, Brunswick, ME 04011
<sup>c</sup>Lyman Laboratory of Physics
Harvard University, Cambridge, MA 02138
Abstract
> Four-dimensional $`𝒩`$=2 gauge theories may be obtained from configurations of D-branes in type IIA string theory. Unitary gauge theories with two-index representations, and orthogonal and symplectic gauge theories, are constructed from configurations containing orientifold planes. Models with two orientifold planes imply a compact dimension, and correspond to elliptic models. Lifting these configurations to M-theory allows one to derive the Seiberg-Witten curves for these gauge theories. We describe how the Seiberg-Witten curves, necessarily of infinite order, are obtained for these elliptic models. These curves are used to calculate the instanton expansion of the prepotential; we explicitly find the one-instanton prepotential for all the elliptic models considered.
1. Introduction
In the Seiberg-Witten approach to four-dimensional $`𝒩=2`$ supersymmetric gauge theories , one begins by identifying an algebraic curve and meromorphic differential specific to the gauge group and matter content of the theory. One then calculates the periods of this differential, and integrates the result to obtain the exact low-energy prepotential for the gauge theory. The perturbative and instanton contributions to the prepotential may then be compared with results obtained directly from the microscopic Lagrangian of the gauge theory. General methods were presented in Refs. for computing the prepotential for gauge theories with hyperelliptic curves. For theories with non-hyperelliptic curves, a systematic approximation scheme for calculating the instanton expansion of the prepotential was developed in Refs. - and is described in the talk by H. J. Schnitzer at this Workshop .
M-theory provides a systematic means of deriving Seiberg-Witten curves . One identifies a type IIA brane configuration that gives rise to the four-dimensional gauge theory of interest, and then lifts this to an M5 brane configuration; the world-volume of the M5 brane contains the Seiberg-Witten curve as a factor.
Our goal in this talk is to explain how to derive the one-instanton prepotentials for the class of elliptic models with a simple gauge group . In sect. 2, we describe the construction of type IIA brane configurations that correspond to these four-dimensional $`𝒩=2`$ gauge theories. We explain in sect. 3 how the lift to M-theory may be used to obtain the Seiberg-Witten curves for these theories. In sect. 4, the quartic truncation of the SW curve is used to obtain explicit expressions for the one-instanton prepotentials for each of these theories.
2. Type IIA brane configurations and four-dimensional gauge theory
We begin by briefly reviewing type IIA brane configurations associated with various four-dimensional $`𝒩=2`$ gauge theories; see Ref. for a review with references. A typical brane configuration, shown in fig. 1, contains a number of NS 5-branes, extended in the 012345 directions, located at the same point in the 789 directions, and having distinct values of $`x_6`$. The horizontal direction in the figure corresponds to $`x_6`$, the vertical direction to $`v=x_4+ix_5`$, with the remaining directions suppressed. The NS 5-branes are connected by D4-branes, extended in the 01236 directions, but of finite length along $`x_6`$.
The brane configuration in fig. 1 gives rise to an $`𝒩=2`$ $`\mathrm{SU}(N_1)\times \mathrm{SU}(N_2)`$ gauge theory with a matter hypermultiplet in the bifundamental representation . The first two NS 5-branes are connected by $`N_1`$ D4-branes; strings extending between the latter give rise to the adjoint vector multiplet of the gauge group $`\mathrm{SU}(N_1)`$. Strings extending between the $`N_2`$ D4-branes connecting the last two NS 5-branes yield the adjoint vector multiplet of the gauge group $`\mathrm{SU}(N_2)`$. Finally, strings extending between the D4 branes connecting the first two NS 5-branes and the D4 branes connecting the last two NS 5-branes give rise to the hypermultiplet in the $`(\overline{\text{ }},\text{ })`$ representation of $`\mathrm{SU}(N_1)\times \mathrm{SU}(N_2)`$.
Figure 2 contains an orientifold 6-plane extending in the 0123789 directions, and intersecting the central NS 5-brane. The orientifold 6-plane can have either $`+4`$ or $`4`$ units of 6-brane charge, and is designated O$`6^+`$ or O$`6^{}`$ respectively. The O6 plane identifies the points $`(x_6,v)(x_6,vm)`$ in the directions transverse to it. In terms of the gauge theory, the orientifold identifies the two factors of the gauge group $`\mathrm{SU}(N)\times \mathrm{SU}(N)`$, and projects the bifundamental representation onto either the symmetric representation (O$`6^+`$) or the antisymmetric representation (O$`6^{}`$) .
In fig. 3, the orientifold plane is located midway between the NS 5-branes, and serves to project out some of the states of the adjoint representation of the $`\mathrm{SU}(2N)`$ gauge group, leaving either $`\mathrm{SO}(2N)`$ for O$`6^+`$ or $`\mathrm{Sp}(2N)`$ for O$`6^{}`$ .
Most asymptotically-free $`𝒩=2`$ gauge theories can be obtained from a variant of the brane configurations just described, either with or without an orientifold plane. However, $`\mathrm{SU}(N)`$ gauge theory with two antisymmetric hypermultiplets apparently requires a configuration with at least two O$`6^{}`$ planes, as described in Ref. and in the talk of H. J. Schnitzer at this Workshop . These two O$`6^{}`$ planes, moreover, generate an infinite number of O$`6^{}`$ planes and NS 5-branes, equally spaced in the $`x_6`$ direction. The corresponding SW curve would be of infinite order.
Alternatively, we may observe that a pair of reflections through different points generates a translation, so the brane configuration must in fact be periodic in the $`x_6`$ direction. A second compactified direction, $`x_{10}`$, emerges in M-theory, so that the lifted M5 configuration lives on a torus. We are thus naturally led to a discussion of elliptic models . The infinite order SW curve may be regarded as an elliptic curve written on the covering space (see, e.g., ref. ).
We now describe the class of elliptic brane configurations, containing a pair of O$`6^\pm `$ planes, that give rise to four-dimensional gauge theories with simple gauge groups and vanishing beta function . The latter condition requires that the total 6-brane charge vanish, so that configurations with two O$`6^{}`$ planes also contain four D6 branes (plus mirrors) parallel to the O6 planes; configurations with O$`6^+`$ and O$`6^{}`$ require no D6 branes.
Figures 4 and 5 show only the unit cell of the periodic configurations; the left-and right-most NS 5-branes are to be identified (with a possible shift in the $`v`$ direction). The corresponding gauge theories may be identified using the rules described above. The configurations in fig. 4 contain two NS 5-branes per unit cell. Figure 4(a) corresponds to $`\mathrm{SU}(N)`$ gauge theory with two antisymmetric hypermultiplets and four fundamental hypermultiplets, whose degrees of freedom arise from strings stretched between the D4- and D6-branes. Figure 4(b) corresponds to $`\mathrm{SU}(N)`$ gauge theory with an antisymmetric and a symmetric hypermultiplet.
The configurations in fig. 5 contain only one NS 5-brane per unit cell. Figure 5(a) corresponds to $`\mathrm{Sp}(2N)`$ gauge theory with an antisymmetric hypermultiplet and four fundamental hypermultiplets. Figure 5(b) corresponds to $`\mathrm{Sp}(2N)`$ gauge theory with an adjoint hypermultiplet (O$`6^+`$ on the NS 5-brane), or $`\mathrm{SO}(2N)`$ gauge theory with an adjoint hypermultiplet (O$`6^{}`$ on the NS 5-brane).
Finally, a periodic configuration without orientifold planes and with one NS 5-brane per unit cell corresponds to $`\mathrm{SU}(N)`$ gauge theory with an adjoint hypermultiplet .
3. M theory and Seiberg-Witten curves
In the strong-coupling limit, type IIA string theory goes over to eleven-dimensional M-theory with an additional periodic coordinate $`x_{10}`$ (with period $`R`$); the brane configurations described in the previous section are “lifted” to M5-brane configurations . The M5-brane world-volume is $`\mathrm{IR}^4\times \mathrm{\Sigma }`$ where $`\mathrm{IR}^4`$ spans the 0123 directions, and $`\mathrm{\Sigma }`$ is a two-dimensional submanifold of $`Q\text{ }\mathrm{C}^2=(v,t)`$, where $`v=x_4+ix_5`$ and $`t=\mathrm{exp}\left[(x_6+ix_{10})/R\right]`$. The M5-brane is located at a point in the remaining 789 directions. $`\mathrm{\Sigma }Q`$ can be written as an algebraic curve, which is none other than the Seiberg-Witten curve of the corresponding four-dimensional gauge theory.
The M5-brane curve $`\mathrm{\Sigma }`$ corresponding to the IIA configuration shown in fig. 1 is
$`t^3{\displaystyle \underset{i=1}{\overset{N_1}{}}}(va_i)t^2+\mathrm{\Lambda }^{2N_1N_2}{\displaystyle \underset{i=1}{\overset{N_2}{}}}(vb_i)t\mathrm{\Lambda }^{3N_1}=0.`$ (3.1)
The features of this curve can be understood directly in terms of the classical IIA picture. Holding $`v`$ fixed, eq. (3.1) has three solutions for $`t`$; these correspond to the positions of the NS 5-branes. The coefficients of the various powers of $`t`$ vanish at the positions of the D4-branes between adjacent NS 5-branes. Since there are no D4 branes to the left and right of the NS 5-branes, the first and last terms of the curve have constant coefficients. The curve (3.1) is indeed the SW curve of the $`\mathrm{SU}(N_1)\times \mathrm{SU}(N_2)`$ gauge theory associated with this IIA configuration.
The M-theory geometry corresponding to a type IIA configuration involving an orientifold plane is more complicated. For a single O$`6^{}`$ plane, as in fig. 2, the M-theory background is an Atiyah-Hitchin space . This may be described in terms of a submanifold $`\stackrel{~}{Q}`$ of $`\text{ }\mathrm{C}^3=(v,t_L,t_R)`$. Far from the orientifold plane, $`\stackrel{~}{Q}`$ is given by
$`t_Lt_R^1={\displaystyle \frac{\mathrm{\Lambda }^{2N+4}}{\left(v+\frac{1}{2}m\right)^4}}`$ (3.2)
and is invariant under
$`vvm,t_Lt_R^1.`$ (3.3)
In the region far to the left of the orientifold plane ($`x_6\mathrm{}`$), the variable $`t_Lt`$, whereas in the region far to the right of the orientifold plane ($`x_6\mathrm{}`$), the variable $`t_Rt`$.
The M5-brane configuration corresponding to the type IIA configuration shown in fig. 2 is given by $`\mathrm{IR}^4\times \mathrm{\Sigma }`$, where now $`\mathrm{\Sigma }`$ is an algebraic curve embedded in $`\stackrel{~}{Q}`$,
$`t_L^3+{\displaystyle \underset{i=1}{\overset{N}{}}}(va_i)t_L^2+A(v)t_L+B(v)=0.`$ (3.4)
Since $`t_L`$ corresponds to $`t`$ to the left of the orientifold plane, the coefficients of the first two (but not the last two) terms correspond to the positions of the D4-branes in that region of fig. 2. The curve $`\mathrm{\Sigma }`$ must be invariant under (3.3), and so may be rewritten as
$`B(vm)t_R^3+A(vm)t_R^2+{\displaystyle \underset{i=1}{\overset{N}{}}}(vma_i)t_R+1=0`$ (3.5)
where now the coefficients of the last two terms correspond to the positions of the D4-branes to the right of the orientifold plane in fig. 2. Using eq. (3.2), we may rewrite eq. (3.5) in terms of $`t_L`$; equating the result with eq. (3.4), we finally obtain
$`t_L^3+{\displaystyle \underset{i=1}{\overset{N}{}}}(va_i)t_L^2+{\displaystyle \frac{\mathrm{\Lambda }^{N+2}\underset{i=1}{\overset{N}{}}(va_im)}{(v+\frac{1}{2}m)^2}}t_L+{\displaystyle \frac{\mathrm{\Lambda }^{3N+6}}{(v+\frac{1}{2}m)^6}}=0.`$ (3.6)
We have used eq. (3.2), which is valid only far from the orientifold; consequently, eq. (3.6) only gives the leading terms (in powers of $`\mathrm{\Lambda }`$) of the curve. The subleading terms may be determined by a more careful consideration of the Atiyah-Hitchin space $`\stackrel{~}{Q}`$ . Including the subleading terms, and defining $`y=t_L/(v+\frac{1}{2}m)^2`$, we obtain the curve
$`y^3`$ $`+`$ $`y^2\left[(v+\frac{1}{2}m)^2{\displaystyle \underset{i=1}{\overset{N}{}}}(va_i)+3\mathrm{\Lambda }^{N+2}\right]`$ (3.7)
$`+`$ $`y\mathrm{\Lambda }^{N+2}\left[(v+\frac{1}{2}m)^2{\displaystyle \underset{i=1}{\overset{N}{}}}(va_im)+3\mathrm{\Lambda }^{N+2}\right]+\mathrm{\Lambda }^{3N+6}=0,`$
which is the form of the curve given by Landsteiner and Lopez . From this curve, the one-instanton propotential may be calculated .
After these warm-ups, we turn to the calculation of the SW curves for the elliptic IIA configurations shown in figs. 4 and 5. For concreteness, we focus on the configuration in fig. 4(b), which gives rise to the four-dimensional $`𝒩=2`$ supersymmetric $`\mathrm{SU}(N)`$ gauge theory with hypermultiplets in the symmetric and antisymmetric representations, but the other cases are analogous.
Figure 6 shows a piece of the IIA configuration on the covering space. The displacements of the O6 planes in the $`v`$ direction correspond to the masses of the hypermultiplets; here $`m_1`$ ($`m_2`$) is the mass of the antisymmetric (symmetric) hypermultiplet.
We would like to derive the M5-brane curve $`\mathrm{IR}^4\times \mathrm{\Sigma }`$ that corresponds to the configuration in fig. 6. First we must describe the M-theory geometry in which $`\mathrm{\Sigma }`$ is embedded. Because of the presence of orientifold planes, we introduce an infinite set of “local” variables $`t_n`$. In the region between any pair of orientifold planes, the corresponding variable $`t_n`$ shown in fig. 6 corresponds to $`t=\mathrm{exp}\left[(x_6+ix_{10})/R\right]`$. The two variables $`t_0`$ and $`t_1`$ adjacent to the O$`6^{}`$ plane at $`v=\frac{1}{2}m_1`$ are related, far from the plane, by
$`t_0t_1^1={\displaystyle \frac{q^{1/2}}{\left(v+\frac{1}{2}m_1\right)^4}}`$ (3.8)
as in eq. (3.2), where for elliptic models the parameter $`q`$ replaces the scale factor $`\mathrm{\Lambda }`$. The two variables $`t_0`$ and $`t_1`$ adjacent to the O$`6^+`$ plane at $`v=\frac{1}{2}m_2`$ are related, far from the plane, by
$`t_1t_0^1=q^{1/2}\left(v+\frac{1}{2}m_2\right)^4.`$ (3.9)
The pairs of variables adjacent to each of the other O6 planes are related by analogous “transition functions.”
Let us write the (leading terms of the) infinite order curve $`\mathrm{\Sigma }`$ as
$`\mathrm{}+qP_2(v)t_0^2+q^{1/4}P_1(v)t_0+P_0(v)+q^{1/4}P_1(v)t_0^1+qP_2(v)t_0^2+\mathrm{}=0`$ (3.10)
in terms of the local variable $`t_0`$. In the central region of fig. 6, $`t_0`$ corresponds to $`t`$, so the coefficient $`P_0(v)`$ must vanish at the positions $`v=a_i`$ of the $`N`$ D4 branes in that region, i.e. $`P_0(v)=_{i=1}^N(va_i)`$. To determine the other coefficients, $`P_n(v)`$, we must use the invariance of the curve induced by the orientifold planes, together with the “transition functions” (3.8) and (3.9). First, requiring that the curve be invariant under
$`vvm_1,t_0t_1^1`$ (3.11)
we may rewrite eq. (3.10) as
$`\mathrm{}`$ $`+q^{1/4}P_1(vm_1)t_1^1+P_0(vm_1)`$
$`+q^{1/4}P_1(vm_1)t_1+qP_2(vm_1)t_1^2+\mathrm{}=0.`$
We then use eq. (3.8) to obtain
$`\mathrm{}`$ $`+q^{1/4}(v+\frac{1}{2}m_1)^6P_2(vm_1)t_0+(v+\frac{1}{2}m_1)^2P_1(vm_1)`$
$`+q^{1/4}(v+\frac{1}{2}m_1)^2P_0(vm_1)t_0^1+q(v+\frac{1}{2}m_1)^6P_1(vm_1)t_0^2+\mathrm{}=0`$
Equating this with eq. (3.10), we find
$`P_1(v)`$ $`=`$ $`(v+\frac{1}{2}m_1)^2P_0(vm_1),`$
$`P_2(v)`$ $`=`$ $`(v+\frac{1}{2}m_1)^6P_1(vm_1),`$
$`\mathrm{}`$
Next, we require that the curve be invariant under
$`vvm_2,t_0t_1^1.`$ (3.15)
This, together with eq. (3.9), generates a set of relations similar to eq. (LABEL:relations). The two reflections (3.11) and (3.15) generate the entire invariance group, and so are sufficient to give us (the leading terms of) the entire set of coefficients
$`\mathrm{}`$
$`P_2(v)`$ $`=`$ $`(v+\frac{1}{2}m_2)^6(v+m_2\frac{1}{2}m_1)^2{\displaystyle \underset{i=1}{\overset{N}{}}}(va_i+m_2m_1),`$
$`P_1(v)`$ $`=`$ $`(v+\frac{1}{2}m_2)^2{\displaystyle \underset{i=1}{\overset{N}{}}}(va_im_2),`$
$`P_0(v)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}(va_i),`$ (3.16)
$`P_1(v)`$ $`=`$ $`(v+\frac{1}{2}m_1)^2{\displaystyle \underset{i=1}{\overset{N}{}}}(va_im_1),`$
$`P_2(v)`$ $`=`$ $`(v+\frac{1}{2}m_1)^6(v+m_1\frac{1}{2}m_2)^2{\displaystyle \underset{i=1}{\overset{N}{}}}(va_i+m_1m_2),`$
$`\mathrm{}`$
This is equivalent to the curve given in ref. , up to overall multiplication by $`F(v)`$ and redefinition $`t_0=G(v)t`$, where $`F(v)`$ and $`G(v)`$ are rational functions of $`v`$ and the hypermultiplet masses. The prepotential derived from these curves is independent of $`F(v)`$ and $`G(v)`$.
In Refs. , SW curves are given for all the other elliptic models discussed in the previous section.
In the limit $`m_1m_2`$ (i.e., vanishing global mass), there are no subleading terms, and the curve (3.10) for $`\mathrm{SU}(N)+\text{ }+\text{ }`$ reduces, upon change of variable $`t=t_0(v+\frac{1}{2}m)^2`$, to
$`0`$ $`=`$ $`{\displaystyle \underset{n\mathrm{even}}{}}q^{n^2/4}t^n{\displaystyle \underset{i=1}{\overset{N}{}}}(va_i)+{\displaystyle \underset{n\mathrm{odd}}{}}q^{n^2/4}t^n{\displaystyle \underset{i=1}{\overset{N}{}}}(va_im)`$ (3.17)
$`=`$ $`\theta _3\left({\displaystyle \frac{z}{\omega _1}}|2\tau \right){\displaystyle \underset{i=1}{\overset{N}{}}}(va_i)+\theta _2\left({\displaystyle \frac{z}{\omega _1}}|2\tau \right){\displaystyle \underset{i=1}{\overset{N}{}}}(va_im),`$
where $`q=\mathrm{exp}(2\pi i\tau )`$, $`t=\mathrm{exp}(i\pi z/\omega _1)`$, and $`\theta _2(\nu |\tau )`$, $`\theta _3(\nu |\tau )`$ are Jacobi theta functions. In ref. , we have shown that eq. (3.17) is equivalent to the curve for this theory given by Uranga .
4. One-instanton prepotential
Although we have obtained an infinite order curve for the $`\mathrm{SU}(N)+\text{ }+\text{ }`$ theory, the one-instanton $`(𝒪(q))`$ prepotential for this theory may be extracted from the quartic truncation of this curve consisting of just those five terms shown explicitly in eq. (LABEL:coeffs). Define
$`S(v)={\displaystyle \frac{P_1(v)P_1(v)}{P_0(v)^2}}.`$ (4.1)
For this theory, $`S(v)`$ has quadratic poles at $`v=a_k`$ and $`v=\frac{1}{2}m_1`$. At these poles, we define the residue functions $`S_k(v)`$ and $`S_{m_1}(v)`$ by
$`S(v)={\displaystyle \frac{S_k(v)}{(va_k)^2}}={\displaystyle \frac{S_{m_1}(v)}{(v+\frac{1}{2}m_1)^2}}.`$ (4.2)
It may be shown that the one-instanton prepotential is given by
$`2\pi i_{1\mathrm{inst}}={\displaystyle \underset{k=1}{\overset{N}{}}}S_k(a_k)2S_{m_1}(\frac{1}{2}m_1).`$ (4.3)
Although $`S(v)`$ and therefore eq. (4.3) depend explicitly only on three of the five coefficients in eq. (LABEL:coeffs), the entire quartic truncation (including the first subleading term) is necessary for the consistency of the calculation to $`𝒪(q)`$ .
In Table 1 below, we list the expressions for $`S(v)`$ for all the elliptic models described in sect. 2 . The one-instanton prepotential for each of these theories is then given in terms of the residue functions defined in eq. (4.2). For $`\mathrm{SU}(N)+`$ adjoint
$`2\pi i_{1\mathrm{inst}}={\displaystyle \underset{k=1}{\overset{N}{}}}S_k(a_k).`$ (4.4)
For $`\mathrm{SU}(N)+\text{ }+\text{ }`$, $`\mathrm{SO}(2N)+`$ adjoint, and $`\mathrm{SO}(2N+1)+`$ adjoint,
$`2\pi i_{1\mathrm{inst}}={\displaystyle \underset{k=1}{\overset{N}{}}}S_k(a_k)2S_{m_1}(\frac{1}{2}m_1),`$ (4.5)
where $`m_1`$ is the mass of the antisymmetric or adjoint hypermultiplet. For $`\mathrm{SU}(N)+2\text{ }+4\text{ }`$,
$`2\pi i_{1\mathrm{inst}}={\displaystyle \underset{k=1}{\overset{N}{}}}S_k(a_k)2S_{m_1}(\frac{1}{2}m_1)2S_{m_2}(\frac{1}{2}m_2),`$ (4.6)
where $`m_1`$ and $`m_2`$ are the masses of the antisymmetric hypermultiplets. For $`\mathrm{Sp}(2N)+`$ adjoint, and $`\mathrm{Sp}(2N)+\text{ }+4\text{ }`$,
$`2\pi i_{1\mathrm{inst}}=2[\overline{S}_0(0)]^{1/2}`$ (4.7)
where
$`S(v)={\displaystyle \frac{\overline{S}_0(v)}{v^4}}`$ (4.8)
defines the residue function at the quartic pole at $`v=0`$. All of these results have been subjected to a wide variety of consistency checks, as described in ref. .
| $`𝒩=2`$ gauge theory | $`S(v)`$ |
| --- | --- |
| $`\mathrm{SU}(N)+\mathrm{\hspace{0.17em}2}\text{}(m_1,m_2)+\mathrm{\hspace{0.17em}4}\text{}(M_j)`$ | $`\frac{_{i=1}^N(v+a_i+m_1)_{i=1}^N(v+a_i+m_2)_{j=1}^4(v+M_j)}{(v+{\scriptscriptstyle \frac{1}{2}}m_1)^2(v+{\scriptscriptstyle \frac{1}{2}}m_2)^2_{i=1}^N(va_i)^2}`$ |
| $`\mathrm{SU}(N)+\text{}(m_1)+\text{}(m_2)`$ | $`\frac{(v+{\scriptscriptstyle \frac{1}{2}}m_2)^2_{i=1}^N(v+a_i+m_1)_{i=1}^N(v+a_i+m_2)}{(v+{\scriptscriptstyle \frac{1}{2}}m_1)^2_{i=1}^N(va_i)^2}`$ |
| $`\mathrm{SU}(N)+\mathrm{adjoint}(m)`$ | $`\frac{_{i=1}^N[(va_i)^2m^2]}{_{i=1}^N(va_i)^2}`$ |
| $`\mathrm{SO}(2N)+\mathrm{adjoint}(m)`$ | $`\frac{v^4_{i=1}^N[(vm)^2a_i^2]_{i=1}^N[(v+m)^2a_i^2]}{(v+{\scriptscriptstyle \frac{1}{2}}m)^2(v{\scriptscriptstyle \frac{1}{2}}m)^2_{i=1}^N(v^2a_i^2)^2}`$ |
| $`\mathrm{SO}(2N+1)+\mathrm{adjoint}(m)`$ | $`\frac{v^2(v+m)(vm)_{i=1}^N[(vm)^2a_i^2]_{i=1}^N[(v+m)^2a_i^2]}{(v+{\scriptscriptstyle \frac{1}{2}}m)^2(v{\scriptscriptstyle \frac{1}{2}}m)^2_{i=1}^N(v^2a_i^2)^2}`$ |
| $`\mathrm{Sp}(2N)+\mathrm{adjoint}(m)`$ | $`\frac{(v+{\scriptscriptstyle \frac{1}{2}}m)^2(v{\scriptscriptstyle \frac{1}{2}}m)^2_{i=1}^N[(vm)^2a_i^2]_{i=1}^N[(v+m)^2a_i^2]}{v^4_{i=1}^N(v^2a_i^2)^2}`$ |
| $`\mathrm{Sp}(2N)+\text{}(m)+\mathrm{\hspace{0.17em}4}\text{}(M_j)`$ | $`\frac{_{i=1}^N[(vm)^2a_i^2]_{i=1}^N[(v+m)^2a_i^2]_{j=1}^4(v^2M_j^2)}{v^4(v+{\scriptscriptstyle \frac{1}{2}}m)^2(v{\scriptscriptstyle \frac{1}{2}}m)^2_{i=1}^N(v^2a_i^2)^2}`$ |
Table 1
Acknowledgement
We are grateful to the organizers of this Workshop, E. D’Hoker, D.H. Phong, and S.T. Yau, for the opportunity to present this work in such a pleasant and stimulating environment.
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# Nonequilibrium perturbation theory for spin-1/2 fields
## I Introduction
Many interesting physical problems, arising, for example, in the study of the early universe and in relativistic heavy-ion collisions require an understanding of the evolution with time of highly-excited states of a quantum field theory. The properties of high-temperature states in thermal equilibrium have been studied for a long time and much is known about them (see, for example ). For the most part, attempts to study the nonequilibrium properties of systems that evolve with time have been based on the assumption that this evolution can be adequately represented as a sequence of near-equilibrium states. While such an assumption may be justified in some cases, our understanding of the true nonequilibrium dynamics is at present very incomplete. One route towards a more complete understanding that has been pursued by several groups is through the study of the $`N\mathrm{}`$ limit of $`N`$-component scalar field theories (see and references cited in these papers). The major advantage of this limit is that it is a gaussian field theory for which the path integral can be evaluated exactly, yielding closed-form evolution equations suitable for numerical solution. However, it is also a theory devoid of scattering processes, which in general lead to important dissipative effects (and the same is true of the related Hartree and one-loop approximations). Moreover, it seems to be extremely difficult to extend such calculations beyond leading order in $`1/N`$, which is essential for approaching any description of more realistic systems. In this paper, we focus on the most obvious alternative of extracting as much information as possible from perturbation theory.
A serious limitation of standard perturbation theory is that, being an expansion about a non-interacting theory, it too is devoid of scattering at lowest order. This means, in particular, that the occupation numbers of single-particle modes which appear in propagators are fixed at some initial values and do not reflect the evolving nonequilibrium state. Low-order calculations therefore become essentially meaningless unless they are restricted to time intervals much shorter than a typical relaxation time. To improve this situation, one should reformulate perturbation theory so as to describe the nonequilibrium state in terms of its own quasiparticle excitations. These excitations have a nonzero thermal width, which in part also characterizes the rate of relaxation of their occupation numbers in response to a changing environment. To put this idea into practice, it is necessary to construct a lowest-order approximation to the interacting theory in which at least some of the dissipative effects of interactions are resummed. Methods for achieving this in the case of both real and complex scalar theories have been described in and incorporated in a comprehensive perturbative approach to the nonequilibrium dynamics of phase transitions in . The purpose of this paper is to investigate how the same idea might be implemented for spin-$`\frac{1}{2}`$ fields. The nonequilibrium dynamics of spinor fields turns out to be quite complicated. In contrast to scalar fields, their propagators appear to have a structure that is not simply a time-dependent generalization of the one that applies in thermal equilibrium; it is sufficiently complicated that we have not been able to explore it in full generality.
We begin in Section II by reviewing briefly the resummation of 2-point functions for real scalar fields. In Section III, we derive some general properties of the full spinor 2-point functions which serve as a guide to the construction of an effective quasiparticle action, under the simplifying assumption that the latter will be CP-invariant. The quasiparticle action is constructed in Section IV in terms of several undetermined functions of time and spatial momentum that characterize the quasiparticle dispersion relation, thermal width and occupation numbers. These functions appear in a counterterm which is added to the free part of the action and subtracted from the interaction part, and will subsequently be determined self-consistently by requiring the counterterm to cancel part of the higher-order corrections to the self energy. The real- and imaginary-time quasiparticle propagators corresponding to this effective action are derived in Sections V and VI respectively and a self-consistent criterion for determining the quasiparticle parameters is implemented in the context of a simple model in Section VII. For illustrative purposes, we introduce supplementary approximations that allow them to be evaluated in closed form. These approximations correspond to a weakly interacting system close to equilibrium, and for this special situation we find, reassuringly, that the evolution of occupation numbers is described by a Boltzmann equation. Finally, in Section VIII, we summarize our principal conclusions and comment on their relation to some other approaches to non-equilibrium field theory.
## II Dissipative perturbation theory for scalar fields
Consider the usual $`\lambda \varphi ^4`$ theory, defined by the Lagrangian density
$$=\frac{1}{2}_\mu \varphi ^\mu \varphi \frac{1}{2}m^2(t)\varphi ^2\frac{1}{4!}\lambda \varphi ^4,$$
(1)
and suppose that an initial state of thermal equilibrium with inverse temperature $`\beta `$ is set up at time $`t=0`$. In this model, the time-dependent mass $`m(t)`$, which arises, for example, in the case of a scalar field theory in a Robertson-Walker spacetime, drives the subsequent state away from equilibrium. Then the closed-time-path formalism (developed for general time-dependent situations in ) yields a path integral weighted by the action
$$I_\mathrm{C}(\varphi _1,\varphi _2,\varphi _3)=\mathrm{d}^3x\left[_0^{\mathrm{}}dt(\varphi _1)_0^{\mathrm{}}dt(\varphi _2)+i_0^\beta d\tau _\mathrm{E}(\varphi _3)\right],$$
(2)
where the path-integration variables $`\varphi _1`$, $`\varphi _2`$ and $`\varphi _3`$ live on a closed contour C in the complex time plane. This contour runs along the real axis from $`t=0`$ to a final time $`t_f`$, returns along the real axis to $`t=0`$, and finally descends along the imaginary axis to $`t=i\beta `$. Here, we have taken the limit $`t_f\mathrm{}`$. The Euclidean action $`_\mathrm{E}`$ (which uses $`m(0)`$) represents the initial density matrix. In this theory, there is a $`3\times 3`$ matrix of 2-point functions $`G_{ab}(x,x^{})`$, with $`a,b=1,2,3`$, but our attention will focus mainly on the real-time functions, with $`a,b=1,2`$. For the real-time part of the action, we wish to construct a lowest-order version $`I_{\mathrm{C0}}(\varphi _1,\varphi _2)=\frac{1}{2}\mathrm{d}^4x\varphi _a𝒟_{ab}\varphi _b`$, where, after a spatial Fourier transform, the differential operator $`𝒟`$ is
$$𝒟_k(t,_t)=\left(\begin{array}{cc}_t^2+k^2+m^2(t)& 0\\ \\ 0& [_t^2+k^2+m^2(t)]\end{array}\right)_k(t,_t).$$
(3)
The counterterm $`\frac{1}{2}d^4x\varphi _a_{ab}\varphi _b`$ is added to $`I_{\mathrm{C0}}`$ and subtracted from the interaction $`I_{\mathrm{C}\mathrm{int}}=I_\mathrm{C}I_{\mathrm{C0}}`$, so as to leave the whole theory unchanged. A choice of $``$ is a choice of the approximate theory about which we perturb and is, of course, equivalent to a choice of $`𝒟`$. Subject to several constraints (discussed in , and generalized below for spinors), the most general choice for $`𝒟`$ is
$$𝒟_k(t,_t)=\left(\begin{array}{cc}[_t^2+\beta _k(t)i\alpha _k(t)]& [\gamma _k(t)_t+\frac{1}{2}\dot{\gamma }_k(t)+i\alpha _k(t)]\\ \\ [\gamma _k(t)_t\frac{1}{2}\dot{\gamma }_k(t)+i\alpha _k(t)]& [_t^2\beta _k(t)i\alpha _k(t)]\end{array}\right),$$
(4)
where $`\alpha _k(t)`$, $`\beta _k(t)`$ and $`\gamma _k(t)`$ are real functions yet to be determined. Of course, the counterterm $``$ can be read from (3) and (4).
The $`2\times 2`$ matrix of quasiparticle propagators $`g_k(t,t^{})`$ is the solution (subject to suitable boundary conditions) of
$$𝒟_k(t,_t)g_k(t,t^{})=g_k(t,t^{})𝒟_k(t^{},\stackrel{}{}_t^{})=i\delta (tt^{}).$$
(5)
Suppressing the spatial momentum $`k`$, this solution can be written in terms of a single complex function $`h(t,t^{})`$ as
$$g_{ab}(t,t^{})=h_b(t,t^{})\theta (tt^{})+h_a(t^{},t)\theta (t^{}t),$$
(6)
where $`h_1=h`$ and $`h_2=h^{}`$. The function $`h`$ is
$$h(t,t^{})=\frac{1}{2}\mathrm{exp}\left(\frac{1}{2}_t^{}^t\gamma (t^{\prime \prime })𝑑t^{\prime \prime }\right)\left\{[1+N(t^{})]f^{(+)}(t)f^{()}(t^{})+[1+N^{}(t^{})]f^{()}(t)f^{(+)}(t^{})\right\},$$
(7)
with $`f^{(\pm )}(t)=[2\mathrm{\Omega }(t)]^{1/2}\mathrm{exp}\left(i_0^t\mathrm{\Omega }(t^{\prime \prime })𝑑t^{\prime \prime }\right)`$. We see that one of the undetermined functions, $`\gamma _k(t)`$, can be interpreted as a quasiparticle width. The quasiparticle energy $`\mathrm{\Omega }_k(t)`$ is a solution of
$$\frac{1}{2}\frac{\ddot{\mathrm{\Omega }}_k(t)}{\mathrm{\Omega }_k(t)}\frac{3}{4}\frac{\dot{\mathrm{\Omega }}_k^2(t)}{\mathrm{\Omega }_k^2(t)}+\mathrm{\Omega }_k^2(t)=\beta _k(t)\frac{1}{4}\gamma _k^2(t).$$
(8)
Finally, the function $`N_k(t)`$, which we hope to interpret in terms of time-dependent occupation numbers, is a solution of
$$\left[_t+\gamma _k(t)+2i\mathrm{\Omega }_k(t)\frac{\dot{\mathrm{\Omega }}_k(t)}{\mathrm{\Omega }_k(t)}\right]\left[_t+\gamma _k(t)\right]N_k(t)=2i\alpha _k(t).$$
(9)
To give substance to this scheme, a prescription is needed for determining the three functions $`\alpha _k(t)`$, $`\beta _k(t)`$ and $`\gamma _k(t)`$ introduced in (4). To this end, define the $`2\times 2`$ self energy matrix $`\mathrm{\Sigma }_k(t,t^{})`$ by
$$G_k(t,t^{})=g_k(t,t^{})+idt^{\prime \prime }dt^{\prime \prime \prime }g_k(t,t^{\prime \prime })\mathrm{\Sigma }_k(t^{\prime \prime },t^{\prime \prime \prime })G_k(t^{\prime \prime \prime },t^{}).$$
(10)
This self energy has contributions from the counterterm $``$ and from loop diagrams:
$$\mathrm{\Sigma }_k(t,t^{})=_k(t,_t)\delta (tt^{})+\mathrm{\Sigma }_k^{\mathrm{loop}}(t,t^{}).$$
(11)
The general strategy is to optimize $`g_k(t,t^{})`$ as an approximation to the full two-point functions $`G_k(t,t^{})`$ by arranging for $``$ to cancel some part of $`\mathrm{\Sigma }^{\mathrm{loop}}`$. Clearly, since $`\mathrm{\Sigma }^{\mathrm{loop}}`$ is non-local in time, only a partial cancellation can be achieved. Various prescriptions might be possible; perhaps the most obvious is the following. Express $`\mathrm{\Sigma }_k(t,t^{})`$ in terms of the average time $`\overline{t}=\frac{1}{2}(t+t^{})`$ and the difference $`(tt^{})`$ and Fourier transform on $`(tt^{})`$. The components of $`_k(t,_t)`$ contain at most one time derivative, so the self energy can be decomposed into contributions that are even and odd in the frequency:
$$\mathrm{\Sigma }_k(\overline{t},\omega )=_k^{(1)}(\overline{t})+_k^{(2)}(\overline{t})\omega +\mathrm{\Sigma }_k^{(1)\mathrm{loop}}(\overline{t},\omega ^2)+\mathrm{\Sigma }_k^{(2)\mathrm{loop}}(\overline{t},\omega ^2)\omega .$$
(12)
Generalized gap equations to be solved for $`\alpha _k(t)`$, $`\beta _k(t)`$ and $`\gamma _k(t)`$ can now be obtained by requiring
$`_k^{(1)}(\overline{t})`$ $`=`$ $`\mathrm{\Sigma }_k^{(1)\mathrm{loop}}(\overline{t},\mathrm{\Omega }_k^2(\overline{t}))`$ (13)
$`_k^{(2)}(\overline{t})`$ $`=`$ $`\mathrm{\Sigma }_k^{(2)\mathrm{loop}}(\overline{t},\mathrm{\Omega }_k^2(\overline{t})),`$ (14)
which amounts to an on-shell renormalization prescription.
These gap equations provide exact implicit definitions of $`\alpha _k(t)`$, $`\beta _k(t)`$ and $`\gamma _k(t)`$, but they cannot, of course, be exactly solved. If the perturbative expansions for $`\mathrm{\Sigma }_k^{(1)\mathrm{loop}}`$ and $`\mathrm{\Sigma }_k^{(2)\mathrm{loop}}`$ are truncated at some finite order, one obtains concrete expressions for them in terms of the propagators $`g_k(t,t^{})`$. These truncated gap equations, together with equation (8) for the quasiparticle energy and (9) for the function $`N_k(t)`$ form a closed system that one might try to solve numerically. It is to some extent illuminating to establish a connection with kinetic theory through some further approximations. Suppose that the gap equations are truncated at two-loop order – the lowest order that yields a nonzero quasiparticle width $`\gamma _k(t)`$. Then, assuming sufficiently weak coupling and sufficiently slow time evolution, propagators inside the loop diagrams can be approximated by taking $`_t^{}^t\mathrm{\Omega }_k(t^{\prime \prime })dt^{\prime \prime }\mathrm{\Omega }_k(\overline{t})(tt^{})`$ and the limit $`\gamma _k(t)0`$. Then, with quasiparticle occupation numbers $`n_k(t)`$ defined by
$$N_k(t)=\frac{1+2n_k(t)}{1i\gamma _k(t)/2\mathrm{\Omega }_k(t)}$$
(15)
a time-derivative expansion of (9), yields the Boltzmann-like equation
$`_tn_k(t)`$ $``$ $`{\displaystyle \frac{\lambda ^2}{32(2\pi )^5}}{\displaystyle \mathrm{d}^3k_1\mathrm{d}^3k_2\mathrm{d}^3k_3\frac{\delta \left(\mathrm{\Omega }_1+\mathrm{\Omega }_2\mathrm{\Omega }_3\mathrm{\Omega }_k\right)\delta \left(𝐤_1+𝐤_2𝐤_3𝐤\right)}{\mathrm{\Omega }_1\mathrm{\Omega }_2\mathrm{\Omega }_3\mathrm{\Omega }_k}}`$ (17)
$`\times \left[n_1n_2(1+n_3)(1+n_k)(1+n_1)(1+n_2)n_3n_k\right].`$
In the following sections, we investigate how this resummation scheme might be extended to spin-$`\frac{1}{2}`$ fields.
## III Exact properties of spinor 2-point functions
To be concrete, we consider a system defined by the Lagrangian density
$$=\overline{\psi }\left[\gamma ^\mu _\mu m(t)\right]\psi +\mathrm{\Delta }$$
(18)
where $`\mathrm{\Delta }`$ represents the coupling of the spinor $`\psi `$ to other fields. For the purposes of this work, we suppose once more that the system is driven away from thermal equilibrium by the time-dependent mass $`m(t)`$ (and possibly by other time-dependent parameters in $`\mathrm{\Delta }`$), but remains spatially homogeneous. As in the scalar case, spinor field theory in a flat Robertson-Walker universe can be represented as a Minkowski-space theory with time-dependent mass. If the state at an initial time that we shall call $`t=0`$ is one of thermal equilibrium, then standard methods (described, for example, in Ref. ) serve to derive the generating functional
$$Z(\xi ,\overline{\xi })=\underset{a=1}{\overset{3}{}}[d\overline{\psi }_ad\psi _a]\mathrm{exp}[iI_\mathrm{c}(\overline{\psi },\psi )+i\mathrm{\Xi }_\mathrm{c}(\overline{\psi },\psi ,\overline{\xi },\xi )],$$
(19)
which generalizes that described in Section II for a scalar field. In this case, the source term is
$$\mathrm{\Xi }_\mathrm{c}(\overline{\psi },\psi ,\overline{\xi },\xi )=_0^{t_f}𝑑t\left[\overline{\xi }_1(t)\psi _1(t)+\overline{\psi }_1(t)\xi _1(t)+\overline{\xi }_2(t)\psi _2(t)+\overline{\psi }_2(t)\xi _2(t)\right]+_0^\beta 𝑑\tau \left[\overline{\xi }_3(\tau )\psi _3(\tau )+\overline{\psi }_3(t)\xi _3(\tau )\right]$$
(20)
and we do not indicate explicitly the other fields that may appear in $`\mathrm{\Delta }`$. If the initial state is characterized by a temperature $`\beta ^1`$ and chemical potential $`\mu `$, then the path-integration variables at the ends of the closed time path obey the boundary conditions $`\psi _1(0)=e^{\beta \mu }\psi _3(\beta )`$ and $`\overline{\psi }_1(0)=e^{\beta \mu }\overline{\psi }_3(\beta )`$, which are inherited by the Green’s functions.
As before, we are particularly concerned with the real-time 2-point functions
$$𝒮_{\alpha \beta }^{(ab)}(𝒙,t;𝒙^{},t^{})=\frac{\delta }{\delta \overline{\xi }_{a\alpha }(𝒙,t)}\frac{\delta }{\delta \xi _{b\beta }(𝒙^{},t^{})}\mathrm{ln}Z|_{\overline{\xi }=\xi =0}$$
(21)
for $`a,b=1,2`$. In terms of field operators, they are
$$𝒮_{\alpha \beta }^{(ab)}(𝒙,t;𝐱^{},t^{})=\left(\begin{array}{cc}𝓣[\psi _\alpha (𝒙,t)\overline{\psi }_\beta (𝒙^{},t^{})]_\mu & \overline{\psi }_\beta (𝒙^{},t^{})\psi _\alpha (𝒙,t)_\mu \\ \\ \psi _\alpha (𝒙,t)\overline{\psi }_\beta (𝒙^{},t^{})_\mu & \overline{𝓣}[\psi _\alpha (𝒙,t)\overline{\psi }_\beta (𝒙^{},t^{})]_\mu \end{array}\right),$$
(22)
where $`\alpha `$ and $`\beta `$ are spinor indices, while $`𝓣`$ and $`\overline{𝓣}`$ denote time ordering and anti-time ordering respectively. In the presence of a chemical potential, the expectation values are given by
$$A_\mu =\frac{\mathrm{Tr}\left[e^{\beta (\widehat{H}\mu \widehat{N})}A\right]}{\mathrm{Tr}\left[e^{\beta (\widehat{H}\mu \widehat{N})}\right]},$$
(23)
where $`\widehat{H}`$ is the Hamiltonian at the initial time and $`\widehat{N}`$ is the particle number.
We hope to construct a perturbation theory in which the lowest-order propagators are partially resummed versions of these full 2-point functions, and begin by establishing some properties of the full functions that our approximate ones ought to share. Expecting that correlations should decay, very roughly as $`e^{\lambda |tt^{}|}`$, over large time intervals, we write the Wightman function $`𝒮_{\alpha \beta }^>(𝒙,t;𝒙^{},t^{};\mu )=\psi _\alpha (𝒙,t)\overline{\psi }_\beta (𝒙^{},t^{})_\mu `$ as
$$𝒮_{\alpha \beta }^>(𝒙,t;𝒙^{},t^{};\mu )=_{\alpha \beta }(𝒙,t;𝒙^{},t^{})\theta (tt^{})+𝒦_{\alpha \beta }(𝒙,t;𝒙^{},t^{})\theta (t^{}t).$$
(24)
Using $`\overline{\psi }_\beta =\psi _\gamma ^{}\gamma _{\gamma \beta }^0`$, it is easy to see that
$$𝒮^>(𝒙,t;𝒙^{},t^{};\mu )=\gamma ^0𝒮^>(𝒙^{},t^{};𝒙,t;\mu )\gamma ^0$$
(25)
and hence that
$$𝒦(𝒙,t;𝒙^{},t^{};\mu )=\overline{}(𝒙^{},t^{};𝒙,t;\mu ),$$
(26)
where, for any Dirac matrix, we define $`\overline{M}=\gamma ^0M^{}\gamma ^0`$. It would be helpful if the second Wightman function $`𝒮_{\alpha \beta }^<(𝒙,t;𝒙^{},t^{};\mu )=\overline{\psi }_\beta (𝐱^{},t^{})\psi _\alpha (𝐱,t)_\mu `$ could be expressed in terms of the same matrix $`(𝒙,t;𝒙^{},t^{};\mu )`$. This can, in fact, be done in a CP-invariant theory. If the CP transformation of an operator $`A`$ is implemented by a unitary operator $`U_{\mathrm{CP}}`$, so that $`A^{\mathrm{CP}}=U_{\mathrm{CP}}^1AU_{\mathrm{CP}}`$, then a CP-invariant theory has $`\widehat{H}^{\mathrm{CP}}=\widehat{H}`$ and $`\widehat{N}^{\mathrm{CP}}=\widehat{N}`$, and we see from (23) that
$$A_\mu =A^{\mathrm{CP}}_\mu .$$
(27)
For a Dirac spinor, we have $`\psi ^{\mathrm{CP}}(𝒙,t)=\gamma _0C\overline{\psi }^T(𝒙,t)`$, where $`C`$ is the charge conjugation matrix and <sup>T</sup> indicates the transpose. It follows from this that
$`𝒮_{\alpha \beta }^<(𝒙,t;𝒙^{},t^{};\mu )`$ $`=`$ $`\left[C^1\gamma _0𝒮_{\alpha \beta }^>(𝒙^{},t^{};𝒙,t;\mu )\gamma ^0C\right]^T`$ (28)
$`=`$ $`\stackrel{~}{\overline{}}(𝒙,t;𝒙^{},t^{};\mu )\theta (tt^{})+\stackrel{~}{}(𝒙^{},t^{};𝒙,t;\mu )\theta (t^{}t),`$ (29)
where, for a matrix-valued function of the chemical potential, we define $`\stackrel{~}{M}(\mu )=\left[C^1\gamma ^0M(\mu )\gamma ^0C\right]^T`$. It is simple to check that $`\overline{\overline{M}}(\mu )=\stackrel{~}{\stackrel{~}{M}}(\mu )=M(\mu )`$ and that $`\stackrel{~}{\overline{M}}(\mu )=\overline{\stackrel{~}{M}}(\mu )`$. With these definitions, the matrix of real-time 2-point functions for a fermion in a CP-invariant theory can be expressed, after a spatial Fourier transformation, as
$$𝒮^{(ab)}(t,t^{};𝒌)=\left(\begin{array}{cc}(t,t^{};𝒌)& \stackrel{~}{\overline{}}(t,t^{};𝒌)\\ (t,t^{};𝒌)& \stackrel{~}{\overline{}}(t,t^{};𝒌)\end{array}\right)\theta (tt^{})+\left(\begin{array}{cc}\stackrel{~}{}(t^{},t;𝒌)& \stackrel{~}{}(t^{},t;𝒌)\\ \overline{}(t^{},t;𝒌)& \overline{}(t^{},t;𝒌)\end{array}\right)\theta (t^{}t),$$
(30)
with
$$(t,t^{};𝒌)=d^3xe^{i𝒌𝒙}(𝒙,t;\mathrm{𝟎},t^{}).$$
(31)
Equivalently, defining $`^{(1)}(t,t^{};𝒌)=(t,t^{};𝒌)`$ and $`^{(2)}(t,t^{};𝒌)=\stackrel{~}{\overline{}}(t,t^{};𝒌)`$, we can write
$$𝒮^{(ab)}(t,t^{})=^{(b)}(t,t^{};𝒌)\theta (tt^{})+\stackrel{~}{}^{(a)}(t^{},t;𝒌)\theta (t^{}t).$$
(32)
We shall demand of our perturbative propagators that they have the structure shown here (as, indeed, do the propagators of standard perturbation theory). This does not mean that our resummation can be applied only in the context of a CP-invariant theory; it does mean, though, that any CP-violating effects will not be resummed. It is worth pointing out that a relation similar to (29) can be obtained by assuming C invariance rather than CP invariance. This would be equally usable, but phenomenologically a little more restrictive.
The structure expressed by (30) or (32) implies two symmetries that will be useful to us. They are
$$\stackrel{~}{𝒮}^{(ab)}(t,t^{};𝒌)=𝒮^{(ba)}(t^{},t;𝒌)$$
(33)
and
$$\left(\begin{array}{cc}\overline{𝒮}^{(11)}(t,t^{};𝒌)& \overline{𝒮}^{(12)}(t,t^{};𝒌)\\ \overline{𝒮}^{(21)}(t,t^{};𝒌)& \overline{𝒮}^{(22)}(t,t^{};𝒌)\end{array}\right)=\left(\begin{array}{cc}𝒮^{(22)}(t^{},t;𝒌)& 𝒮^{(12)}(t^{},t;𝒌)\\ 𝒮^{(21)}(t^{},t;𝒌)& 𝒮^{(11)}(t^{},t;𝒌)\end{array}\right).$$
(34)
The first of these generalizes to the full $`3\times 3`$ matrix of real- and imaginary-time 2-point functions, but the second makes sense only for the real-time functions.
Finally, we shall need two pieces of information concerning the values of these functions at equal times. The functions $`𝒮^{(12)}`$ and $`𝒮^{(21)}`$ have unique values at $`t=t^{}`$, which implies
$$\overline{}(t,t;𝒌)=(t,t;𝒌).$$
(35)
On the other hand, the time-ordered function $`𝒮^{(11)}`$ has a discontinuity at $`t=t^{}`$, which reproduces the equal-time anticommutator
$$\underset{t^{}t0}{lim}𝒮_{\alpha \beta }^{(11)}(𝒙,t;𝒙^{},t^{})\underset{t^{}t+0}{lim}𝒮_{\alpha \beta }^{(11)}(𝒙,t;𝒙^{},t^{})=\{\psi _\alpha (𝒙,t),\overline{\psi }_\beta (𝒙^{},t)\}_\mu =\gamma _{\alpha \beta }^0\delta (𝒙𝒙^{}),$$
(36)
and this implies
$$(t,t;𝒌)\stackrel{~}{}(t,t;𝒌)=\gamma ^0.$$
(37)
## IV Construction of the quasiparticle action
We wish to construct a lowest-order action
$$I_{\mathrm{C0}}(\psi )=I_\mathrm{C}^{(2)}(\psi )+d^4x\overline{\psi }_a_{ab}\psi _bd^4x\overline{\psi }_a𝒟_{ab}\psi _b$$
(38)
that will serve as a starting point for our partially resummed perturbation theory. As before, $`I_\mathrm{C}^{(2)}`$ is the quadratic part of the original closed-time-path action, while the term involving $`_{ab}`$ is a counterterm which will be subtracted from the interaction part, so as to leave the whole theory unchanged. Specifying the form of $`_{ab}`$ is equivalent to specifying the resulting differential operator $`𝒟_{ab}`$. To begin, we construct the real-time components of $`𝒟_{ab}`$, with $`a,b=1,2`$. The unperturbed propagator matrix $`S^{(ab)}(t,t^{})`$ is a solution of the equations
$$𝒟_{ac}(t,_t)S^{(cb)}(t,t^{})=S^{(ac)}(t,t^{})𝒟_{cb}(t^{},\stackrel{}{_t^{}})=i\delta _{ab}\delta (tt^{}).$$
(39)
(For economy of notation, we shall usually not indicate explicitly the dependence of these quantities on $`𝒌`$ and $`\mu `$.) The form that might usefully be chosen for $`𝒟_{ab}`$ is constrained to a considerable extent by the requirement that this equation have solutions for $`S^{(ab)}`$ which have the structure exhibited in (30) and (32) for the full 2-point functions and inherit the various properties that we discussed in Section III. Observe first that the $`\delta (tt^{})`$ on the right of (39) arises from differentiating $`\theta (tt^{})`$ and $`\theta (t^{}t)`$. We can ensure that only these $`\delta `$ functions will appear by restricting $`𝒟`$ to the form
$$𝒟=\left(\begin{array}{cc}i\gamma ^0_t& 0\\ 0& i\gamma ^0_t\end{array}\right)+\mathrm{},$$
(40)
where the ellipsis indicates terms without time derivatives. The coefficients $`\pm i\gamma _0`$ are determined by the boundary condition (37). In principle, an ansatz using more time derivatives (along with further boundary conditions to eliminate unwanted $`\delta `$ functions and their derivatives) might be possible, but we have not found such a generalization tractable. Next, if the solution of (39) is to have the symmetry expressed by (33), then $`𝒟`$ must satisfy
$$\stackrel{~}{𝒟}_{ab}(t,_t)=𝒟_{ba}(t,_t),$$
(41)
which ensures that $`I_{\mathrm{C0}}`$ is CP invariant. Similarly, the symmetry expressed by (34) implies
$$\left(\begin{array}{cc}\overline{𝒟}_{11}(t,_t)& \overline{𝒟}_{12}(t,_t)\\ \overline{𝒟}_{21}(t,_t)& \overline{𝒟}_{21}(t,_t)\end{array}\right)=\left(\begin{array}{cc}𝒟_{22}(t,_t)& 𝒟_{12}(t,_t)\\ 𝒟_{21}(t,_t)& 𝒟_{11}(t,_t)\end{array}\right).$$
(42)
Finally, as for a scalar field, causality requires
$$𝒟_{11}(t,_t)+𝒟_{12}(t,_t)+𝒟_{21}(t,_t)+𝒟_{22}(t,_t)=0.$$
(43)
The net effect of these considerations is that $`𝒟`$ can be written as
$$𝒟(t,_t)=\left(\begin{array}{cc}i\gamma ^0_t+𝒟_1(t)& i𝒟_2(t)\\ i\stackrel{~}{𝒟}_2(t)& i\gamma ^0_t\overline{𝒟}_1(t)\end{array}\right),$$
(44)
where $`𝒟_1(t)`$ and $`𝒟_2(t)`$ are subject to the constraints
$`\stackrel{~}{𝒟}_1(t)`$ $`=`$ $`𝒟_1(t)`$ (45)
$`\overline{𝒟}_2(t)`$ $`=`$ $`𝒟_2(t)`$ (46)
$`𝒟_1(t)+i𝒟_2(t)`$ $`=`$ $`\overline{𝒟}_1(t)i\stackrel{~}{𝒟}_2(t).`$ (47)
When $`𝒟`$ has this structure, and the propagator $`S(t,t^{})`$ is written in the form (30) in terms of a function $`H(t,t^{})`$ that approximates $`(t,t^{})`$, then equations (39) for the propagator reduce to
$`\left[i\gamma ^0_t+𝒟_1(t)+i𝒟_2(t)\right]H(t,t^{})`$ $`=`$ $`0`$ (48)
$`\left[i\gamma ^0_t+𝒟_1(t)\right]\stackrel{~}{H}(t^{},t)+i𝒟_2(t)\overline{H}(t^{},t)`$ $`=`$ $`0.`$ (49)
At this point, the most general procedure would be to expand $`𝒟_i(t)`$ in terms of a complete basis of Dirac matrices
$$𝒟_i(t)=\underset{p=1}{\overset{16}{}}d_i^p(t,𝒌)\mathrm{\Gamma }_p,$$
(50)
the thirty-two undetermined functions $`d_i^p(t,𝒌)`$ being the analogues of the functions $`\alpha _k(t)`$, $`\beta _k(t)`$ and $`\gamma _k(t)`$ that appeared in the scalar theory. This general problem is one that we have not found tractable. To simplify matters, we use instead the smallest subset of the Dirac algebra that closes under multiplication and under $`\overline{}`$ and $`\stackrel{~}{}`$ conjugation, and that includes the matrices $`\gamma ^0`$ and $`𝜸𝒌`$ that appear in the free theory. For the operator $`𝒟`$, a convenient basis is $`\{1,\gamma ^0,\mathrm{\Gamma }_+,\mathrm{\Gamma }_{}\}`$, where
$$\mathrm{\Gamma }_\pm =\frac{1}{2|𝒌|}\left(1\pm \gamma ^0\right)𝜸𝒌$$
(51)
and $`1`$ denotes the unit matrix. The conjugates of these matrices are $`\overline{\gamma ^0}=\stackrel{~}{\gamma ^0}=\gamma ^0`$, $`\overline{\mathrm{\Gamma }}_\pm =\mathrm{\Gamma }_{}`$ and $`\stackrel{~}{\mathrm{\Gamma }}_\pm =\mathrm{\Gamma }_\pm `$. Our ansatz for $`𝒟`$ is then
$`𝒟_1(t)`$ $`=`$ $`i\left[\lambda (t)\nu (t)\right]\gamma ^0+\left[\sigma (t)ϵ(t)\right]\mathrm{\Gamma }_++\left[\sigma ^{\mathrm{}}(t)+ϵ^{}(t)\right]\mathrm{\Gamma }_{}\left[\tau (t)+i\eta (t)\right]`$ (52)
$`𝒟_2(t)`$ $`=`$ $`i\nu (t)\gamma ^0+ϵ(t)\mathrm{\Gamma }_+ϵ^{}(t)\mathrm{\Gamma }_{}+i\eta (t).`$ (53)
Although we have not indicated it explicitly, the coefficients depend on $`𝒌`$ and $`\mu `$ as well as on $`t`$. For complex functions of $`\mu `$, we define $`f^{\mathrm{}}(\mu )=f^{}(\mu )`$, so that $`\stackrel{~}{f}(\mu )=f(\mu )=f^{\mathrm{}}(\mu )`$. This is the most general ansatz that satisfies the restrictions (46) on $`𝒟_i(t)`$, provided that
$`\lambda ^{\mathrm{}}=\lambda ,\tau ^{\mathrm{}}=\tau ,`$ $`\nu ^{}=\nu ,\eta ^{}=\eta `$ (54)
$`\lambda +\lambda ^{}`$ $`=`$ $`\nu +\nu ^{\mathrm{}}`$ (55)
$`\sigma \stackrel{~}{\sigma }`$ $`=`$ $`ϵ\stackrel{~}{ϵ}`$ (56)
$`\tau \tau ^{}`$ $`=`$ $`i(\eta \eta ^{\mathrm{}}).`$ (57)
## V Solution for the real-time quasiparticle propagators
Having constructed the unperturbed action (38) in terms of the differential operator $`𝒟(t,_t)`$ given by (44) together with the ansätze (52) and (53), we require a formal solution to equations (48) and (49) for the matrix-valued function $`H(t,t^{};𝒌)`$ from which the quasiparticle propagator $`S^{(ab)}(t,t^{};𝒌)`$ is to be constructed via (32). It proves convenient to reorganize our basis of Dirac matrices into the set $`\{\gamma _\pm \frac{1}{2}(1\pm \gamma ^0),\mathrm{\Gamma }_\pm \}`$, expanding $`H(t,t^{})`$ as
$$H(t,t^{})=\mathrm{exp}\left(_t^{}^t\lambda (t^{\prime \prime })𝑑t^{\prime \prime }\right)\left[A(t,t^{})\gamma _++B(t,t^{})\mathrm{\Gamma }_{}+C(t,t^{})\mathrm{\Gamma }_+D(t,t^{})\gamma _{}\right].$$
(58)
Noting that (48) governs the dependence of $`H(t,t^{})`$ on its first time argument, while (49) refers to the second time argument, we introduce the notation $`\dot{A}(t,t^{})_tA(t,t^{})`$ and $`\stackrel{̊}{A}(t,t^{})_t^{}A(t,t^{})`$. With this notation, (48) can be written as
$`i\left(\begin{array}{c}\dot{A}\\ \dot{B}\end{array}\right)`$ $`=`$ $`𝑻\left(\begin{array}{c}A\\ B\end{array}\right)`$ (59)
$`i\left(\begin{array}{c}\dot{C}\\ \dot{D}\end{array}\right)`$ $`=`$ $`𝑻\left(\begin{array}{c}C\\ D\end{array}\right),`$ (60)
where
$$𝑻=\left(\begin{array}{cc}\tau & \sigma \\ \sigma ^{\mathrm{}}& \tau \end{array}\right).$$
(61)
For orientation, we note that, in the absence of the counterterm $`_{ab}`$, we would have $`\tau =m`$ and $`\sigma =|𝒌|`$. (Equations somewhat analogous to these have been obtained, for example, by Sahni in the course of solving the Dirac equation in certain curved spacetimes.) The matrix $`𝑻`$ has the generalized Hermiticity property $`𝑻^{}=𝑻`$, where the operation is defined by taking the transpose, and replacing each element by its conjugate. (This reduces to the usual Hermitian conjugate when the chemical potential vanishes.) If the two-component vectors $`\phi `$ and $`\chi `$ are solutions to $`i\dot{\phi }=𝑻\phi `$, then it is simple to see that the inner product $`(\phi ,\chi )=\phi ^{}\chi `$ is preserved by the time evolution. The eigenvalues of $`𝑻(t)`$ are $`\pm \mathrm{\Omega }(t)`$, where the time-dependent frequency $`\mathrm{\Omega }=\sqrt{\tau ^2+\sigma ^{\mathrm{}}\sigma }`$ satisfies $`\mathrm{\Omega }^{\mathrm{}}=\mathrm{\Omega }`$, and the corresponding normalized and mutually orthogonal eigenvectors are
$$u^{(+)}(t)=\frac{1}{\sqrt{2\mathrm{\Omega }(\mathrm{\Omega }\tau )}}\left(\begin{array}{c}\sigma \\ \mathrm{\Omega }\tau \end{array}\right)u^{()}(t)=\frac{1}{\sqrt{2\mathrm{\Omega }(\mathrm{\Omega }\tau )}}\left(\begin{array}{c}(\mathrm{\Omega }\tau )\\ \sigma ^{\mathrm{}}\end{array}\right)$$
(62)
Given some fixed time $`t_0`$, one can formally write exact solutions
$$\varphi _{t_0}^{(\pm )}(t)=𝓣\mathrm{exp}\left[i_{t_0}^t𝑻(t^{})𝑑t^{}\right]u^{(\pm )}(t_0),$$
(63)
which are positive- and negative-frequency solutions
$$\varphi _{t_0}^{(\pm )}(t)\mathrm{exp}\left[i_{t_0}^t\mathrm{\Omega }(t^{})𝑑t^{}\right]u^{(\pm )}(t_0),$$
(64)
at times near $`t_0`$. More generally, we may choose an orthonormal basis
$$\varphi _1(t)=\left(\begin{array}{c}f(t)\\ g(t)\end{array}\right),\varphi _2(t)=\left(\begin{array}{c}g^{\mathrm{}}(t)\\ f^{\mathrm{}}(t)\end{array}\right),$$
(65)
with $`f^{\mathrm{}}(t)f(t)+g^{\mathrm{}}(t)g(t)=1`$. It is readily verified that $`\varphi _2(t)`$ is a solution if $`\varphi _1(t)`$ is, and that if $`\varphi _1(t)=\varphi _{t_0}^{(+)}(t)`$ for some $`t_0`$, then $`\varphi _2(t)=\varphi _{t_0}^{()}(t)`$. Thus, the solution to (59) and (60) can be written as
$`\left(\begin{array}{c}A(t,t^{})\\ B(t,t^{})\end{array}\right)`$ $`=`$ $`P_1(t^{})\left(\begin{array}{c}f(t)\\ g(t)\end{array}\right)+P_2(t^{})\left(\begin{array}{c}g^{\mathrm{}}(t)\\ f^{\mathrm{}}(t)\end{array}\right)`$ (66)
$`\left(\begin{array}{c}C(t,t^{})\\ D(t,t^{})\end{array}\right)`$ $`=`$ $`Q_1(t^{})\left(\begin{array}{c}f(t)\\ g(t)\end{array}\right)+Q_2(t^{})\left(\begin{array}{c}g^{\mathrm{}}(t)\\ f^{\mathrm{}}(t)\end{array}\right),`$ (67)
where $`P_i(t^{})`$ and $`Q_i(t^{})`$ are to be determined by solving (49) and applying suitable boundary conditions.
To solve (49), it is helpful to define
$`W(t,t^{})`$ $`=`$ $`B^{\mathrm{}}(t,t^{})C(t,t^{})`$ (68)
$`X(t,t^{})`$ $`=`$ $`D^{\mathrm{}}(t,t^{})+A(t,t^{})`$ (69)
$`Y(t,t^{})`$ $`=`$ $`B^{\mathrm{}}(t,t^{})+C(t,t^{})`$ (70)
$`Z(t,t^{})`$ $`=`$ $`D^{\mathrm{}}(t,t^{})A(t,t^{}).`$ (71)
In terms of these functions, (49) is
$`i\left(\begin{array}{c}\stackrel{̊}{W}\\ \stackrel{̊}{X}\end{array}\right)`$ $`=`$ $`𝑻\left(\begin{array}{c}W\\ X\end{array}\right)`$ (72)
$`i\left(\begin{array}{c}\stackrel{̊}{Y}\\ \stackrel{̊}{Z}\end{array}\right)`$ $`=`$ $`i(\lambda +\lambda ^{})\left(\begin{array}{c}Y\\ Z\end{array}\right)+𝑻^{\mathbf{}}\left(\begin{array}{c}Y\\ Z\end{array}\right)+𝑬\left(\begin{array}{c}W\\ X\end{array}\right).`$ (73)
The matrix $`𝑻`$ is the same as the one defined in (61), while
$`𝑻^{\mathbf{}}`$ $`=`$ $`\left(\begin{array}{cc}\tau ^{}& \sigma +ϵ^{\mathrm{}}ϵ\\ \sigma ^{\mathrm{}}+ϵ^{}ϵ^{\mathrm{}}& \tau ^{}\end{array}\right)`$ (74)
$`𝑬`$ $`=`$ $`\left(\begin{array}{cc}i(\nu \nu ^{\mathrm{}}\eta \eta ^{\mathrm{}})& ϵ^{\mathrm{}}+ϵ\\ (ϵ^{}+ϵ^{\mathrm{}})& i(\nu \nu ^{\mathrm{}}+\eta +\eta ^{\mathrm{}})\end{array}\right).`$ (75)
The appearance of a new matrix $`𝑻^{\mathbf{}}`$ (which also satisfies $`𝑻_{}^{\mathbf{}}{}_{}{}^{}=𝑻^{\mathbf{}}`$) and a new damping constant $`\lambda +\lambda ^{}`$ in (73) reflects the fact (previously noted in ) that particles and antiparticles acquire different thermal masses and widths in the presence of a nonzero chemical potential. The appearance of the inhomogeneous term proportional to $`𝑬`$ in (73) implies that the time dependence of the propagators is not exhausted by that of the single-particle mode functions, and we hope to interpret the additional time dependence in terms of time-dependent occupation numbers for the single-particle modes.
Evidently, the solution of (73) will require the introduction of a new set of antiparticle mode functions, which make matters rather complicated. In the remainder of this paper, we shall simplify our calculations by specializing to the case of zero chemical potential. In that case, conjugation reduces to complex conjugation, while the restrictions on the counterterm parameters noted in (56) imply that $`\tau `$, $`\lambda `$ and $`\eta `$ are real, and that $`\nu =\lambda `$. We now have $`𝑻^{\mathbf{}}=𝑻=𝑻^{}`$ and
$$𝑬=\left(\begin{array}{cc}2i\eta & 2ϵ\\ 2ϵ^{}& 2i\eta \end{array}\right)=𝑬^{}.$$
(76)
We look for a solution to (72) and (73) in the form
$`\left(\begin{array}{c}W(t,t^{})\\ X(t,t^{})\end{array}\right)`$ $`=`$ $`R_1(t)\left(\begin{array}{c}f(t^{})\\ g(t^{})\end{array}\right)+R_2(t)\left(\begin{array}{c}g^{}(t^{})\\ f^{}(t^{})\end{array}\right)`$ (77)
$`\left(\begin{array}{c}Y(t,t^{})\\ Z(t,t^{})\end{array}\right)`$ $`=`$ $`S_1(t,t^{})\left(\begin{array}{c}f(t^{})\\ g(t^{})\end{array}\right)+S_2(t,t^{})\left(\begin{array}{c}g^{}(t^{})\\ f^{}(t^{})\end{array}\right),`$ (78)
and find that they are satisfied if
$`S_1(t,t^{})`$ $`=`$ $`S_1(t)L(t^{})+R_1(t)K_1(t^{})+R_2(t)K_2(t^{})`$ (79)
$`S_2(t,t^{})`$ $`=`$ $`S_2(t)L(t^{})+R_1(t)K_2^{}(t^{})R_2(t)K_1(t^{}),`$ (80)
where
$`L(t^{})`$ $`=`$ $`\mathrm{exp}\left[2{\displaystyle _0^t^{}}\lambda (t^{\prime \prime })𝑑t^{\prime \prime }\right],`$ (81)
$`K_i(t^{})`$ $`=`$ $`L(t^{}){\displaystyle _0^t^{}}L^1(t^{\prime \prime })e_i(t^{\prime \prime })𝑑t^{\prime \prime }`$ (82)
and the $`e_i(t)`$ are related to the matrix elements $`e_{ij}(t)=\varphi _i^{}(t)𝑬(t)\varphi _j(t)`$ by
$`e_{11}`$ $`=`$ $`e_{22}^{}=2\left[ϵf^{}gϵ^{}fg^{}i\eta (t^{})(f^{}fg^{}g)\right]ie_1`$ (83)
$`e_{12}`$ $`=`$ $`e_{21}^{}=2\left[ϵf^{}f^{}+ϵ^{}g^{}g^{}+i\eta f^{}g^{}\right]ie_2.`$ (84)
At this point, the solution of (48) (expressed by (59) and (60)) has left us with four undetermined functions $`P_i(t^{})`$ and $`Q_i(t^{})`$, while the solution of (49) (expressed by (72) and (73)) produced another four undetermined functions $`R_i(t)`$ and $`S_i(t)`$. However, the two solutions are related by (70) and by comparing them, we determine all eight functions up to constants of integration. Moreover, the values that these constants can take are constrained by the equal-time conditions (35) and (37), which express very general properties of the 2-point functions. In fact, it turns out that that the propagator is determined up to two constants of integration, that we shall denote by $`a_1`$ (which is real) and $`a_2`$ (which is complex). Our final result for $`H(t,t^{})`$ is expressed by the original ansatz (58) with the time-dependent coefficients given by
$`A(t,t^{})`$ $`=`$ $`\left[1N(t^{})\right]f(t)f^{}(t^{})+N(t^{})g^{}(t)g(t^{})\mathrm{\Delta }(t^{})f(t)g(t^{})\mathrm{\Delta }^{}(t^{})g^{}(t)f^{}(t^{})`$ (85)
$`B(t,t^{})`$ $`=`$ $`\left[1N(t^{})\right]g(t)f^{}(t^{})N(t^{})f^{}(t)g(t^{})\mathrm{\Delta }(t^{})g(t)g(t^{})+\mathrm{\Delta }^{}(t^{})f^{}(t)f^{}(t^{})`$ (86)
$`C(t,t^{})`$ $`=`$ $`\left[1N(t^{})\right]f(t)g^{}(t^{})N(t^{})g^{}(t)f(t^{})+\mathrm{\Delta }(t^{})f(t)f(t^{})\mathrm{\Delta }^{}(t^{})g^{}(t)g^{}(t^{})`$ (87)
$`D(t,t^{})`$ $`=`$ $`\left[1N(t^{})\right]g(t)g^{}(t^{})+N(t^{})f^{}(t)f(t^{})+\mathrm{\Delta }(t^{})g(t)f(t^{})+\mathrm{\Delta }^{}(t^{})f^{}(t)g^{}(t^{}).`$ (88)
The real function $`N(t^{})`$ and the complex function $`\mathrm{\Delta }(t^{})`$ are given, in terms of the quantities introduced above, by
$`N(t^{})`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[1K_1(t^{})a_1L(t^{})\right]`$ (89)
$`\mathrm{\Delta }(t^{})`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[K_2(t^{})+a_2L(t^{})\right],`$ (90)
but it is more illuminating to observe that they obey the differential equations
$`_t^{}N_k(t^{})`$ $`=`$ $`2\lambda _k(t^{})N_k(t^{})+\left[\lambda _k(t^{}){\displaystyle \frac{1}{2}}e_{1k}(t^{})\right]`$ (91)
$`_t^{}\mathrm{\Delta }_k(t^{})`$ $`=`$ $`2\lambda _k(t^{})\mathrm{\Delta }_k(t^{})+{\displaystyle \frac{1}{2}}e_{2k}(t^{}),`$ (92)
with initial conditions $`N_k(0)=\frac{1}{2}(1a_{1k})`$ and $`\mathrm{\Delta }_k(0)=\frac{1}{2}a_{2k}`$, in which we have reinstated the dependence on spatial momentum $`k`$ that has been implicit throughout. These are the equations that we might hope to interpret as kinetic equations for the occupation numbers of single-particle modes. Let us, indeed, specialize to the case of thermal equilibrium, and temporarily delete the counterterm $`_{ab}`$. The mode functions can be written as
$$\left(\begin{array}{c}f_k(t)\\ g_k(t)\end{array}\right)=\frac{e^{i\mathrm{\Omega }_kt}}{\sqrt{2\mathrm{\Omega }_k(\mathrm{\Omega }_km)}}\left(\begin{array}{c}|𝒌|\\ \mathrm{\Omega }_km\end{array}\right),$$
(93)
with $`\mathrm{\Omega }_k=\sqrt{k^2+m^2}`$ and we have $`K_{1k}=K_{2k}=0`$ and $`L_k=1`$. We then find that that the time-ordered function $`S^{(11)}(t,t^{},𝒌)=H(t,t^{};𝒌)\theta (tt^{})+\stackrel{~}{H}(t^{},t;𝒌)\theta (t^{}t)`$, with $`H(t,t^{};𝒌)`$ given by (58) agrees with the corresponding function obtained in , provided that we can identify $`N_k=\frac{1}{2}(1a_1)=N_k^{\mathrm{eq}}`$ and $`\mathrm{\Delta }_k=\frac{1}{2}a_2=0`$, where $`N_k^{\mathrm{eq}}=\left[\mathrm{exp}(\beta \mathrm{\Omega }_k)+1\right]^1`$ is the usual Fermi-Dirac distribution. The other 2-point functions do not agree with those of , because these authors made use of a different time path (which is legitimate in thermal equilibrium, but not in the nonequilibrium situation considered here). At this stage, these values of the constants of integration $`a_1`$ and $`a_2`$ are merely guesses that yield this agreement with . The actual values that are required by our formalism are determined by computing imaginary-time correlators and applying appropriate boundary conditions. This computation is the subject of the following section, where we shall find our guesses confirmed.
## VI Imaginary and mixed-time propagators
In terms of the imaginary-time field operator $`\psi (𝒙,\tau )=e^{\widehat{H}\tau }\psi (𝒙,0)e^{\widehat{H}\tau }`$, the imaginary- and mixed-time 2-point functions are (for $`a=1,2`$)
$`𝒮_{\alpha \beta }^{(33)}(𝒙,\tau ;𝒙^{},\tau ^{})`$ $`=`$ $`\psi _\alpha (𝒙,\tau )\overline{\psi }_\beta (𝒙^{},\tau ^{})\theta (\tau \tau ^{})\overline{\psi }_\beta (𝒙^{},\tau ^{})\psi _\alpha (𝒙,\tau )\theta (\tau ^{}\tau )`$ (94)
$`𝒮_{\alpha \beta }^{(a3)}(𝒙,t;𝒙^{},\tau ^{})`$ $`=`$ $`\overline{\psi }_\beta (𝒙^{},\tau ^{})\psi _\alpha (𝒙,t)`$ (95)
$`𝒮^{(3a)}(𝒙,\tau ;𝒙^{},t^{})`$ $`=`$ $`\stackrel{~}{𝒮}^{(a3)}(𝒙^{},t^{};𝒙,\tau ).`$ (96)
Because the time-path ordering makes imaginary times later than real times, the functions $`𝒮^{(13)}`$ and $`𝒮^{(23)}`$ are identical. With $`\mu =0`$, antiperiodicity of the path integration variables, $`\psi _3(𝒙,\beta )=\psi _1(𝒙,0)`$, and the fact that the field operator $`\psi (𝒙,0)`$ is unique supply the two boundary conditions
$`𝒮^{(a3)}(t,\beta ;𝒌)`$ $`=`$ $`𝒮^{(a1)}(t,0;𝒌)`$ (97)
$`𝒮^{(a3)}(t,0;𝒌)`$ $`=`$ $`𝒮^{(a2)}(t,0;𝒌)`$ (98)
that we shall use to determine the constants of integration $`a_1`$ and $`a_2`$. If the sources for real-time fields are set to zero, then the time path reduces to just its imaginary-time segment, and antiperiodicity yields
$$𝒮^{(33)}(0,\tau ^{};𝒌)=𝒮^{(33)}(\beta ,\tau ^{};𝒌).$$
(99)
Finally, uniqueness of $`\psi (𝒙,0)`$ also implies
$$𝒮^{(a3)}(0,\tau ^{};𝒌)=𝒮^{(33)}(0,\tau ^{};𝒌)$$
(100)
and the two latter boundary conditions serve to fix constants of integration that arise in the calculation of $`𝒮^{(a3)}`$ and $`𝒮^{(33)}`$.
In order to construct a tractable perturbation theory, we have insisted that the unperturbed action (38) be local in time. With this restriction, there are no terms of the form $`\overline{\psi }_a(t)𝒟_{a3}\psi _3(\tau )`$, ($`a=1,2`$) so we have
$$I_{\mathrm{c0}}(\psi )=d^3x\left[\underset{a,b=1}{\overset{2}{}}𝑑t\overline{\psi }_a𝒟_{ab}\psi _b+_0^\beta 𝑑\tau \overline{\psi }_3𝒟_{33}\psi _3\right].$$
(101)
In the second term, which approximates the path integral representation of the initial density operator, $`\overline{\psi }_3𝒟_{33}\psi _3`$ is $`i`$ times the Euclidean version of the free part of (18), with $`m=m(0)`$, supplemented by a counterterm $`_{33}`$. The form of $`𝒟_{33}`$ is determined by the assumed CP invariance, $`\stackrel{~}{𝒟}_{33}(_\tau )=𝒟_{33}(_\tau )`$, together with the requirement that the quasiparticle energy $`\mathrm{\Omega }_0`$ be equal to the $`t0`$ limit of the energy that appears in the real-time mode functions. This yields
$$𝒟_{33}(_\tau )=i\left[\gamma ^0_\tau \sigma _0\mathrm{\Gamma }_+\sigma _0^{}\mathrm{\Gamma }_{}+\tau _0\right],$$
(102)
where $`\sigma _0`$ and $`\tau _0`$ are the $`t0`$ limits of the parameters appearing in (52). The new propagators satisfy
$`𝒟_{33}(_\tau )S^{(33)}(\tau ,\tau ^{})`$ $`=`$ $`S^{(33)}(\tau ,\tau ^{})𝒟_{33}(\stackrel{}{}_\tau ^{})=i\delta (\tau \tau ^{})`$ (103)
$`\left[𝒟_{11}(_t)+𝒟_{12}\right]S^{(13)}(t,\tau ^{})`$ $`=`$ $`S^{(13)}(t,\tau ^{})𝒟_{33}(\stackrel{}{}_\tau ^{})=0,`$ (104)
and these equations can be solved by the method explained in the previous section. The imaginary-time propagator can be expressed as $`S^{(33)}(\tau ,\tau ^{};𝒌)=H_3(\tau ,\tau ^{};𝒌)\theta (\tau \tau ^{})+\stackrel{~}{H}_3(\tau ^{},\tau ;𝒌)\theta (\tau ^{}\tau )`$, and we look for solutions of the form
$`H_3(\tau ,\tau ^{})`$ $`=`$ $`A_3(\tau ,\tau ^{})\gamma _++B_3(\tau ,\tau ^{})\mathrm{\Gamma }_{}+C_3(\tau ,\tau ^{})\mathrm{\Gamma }_+D_3(\tau ,\tau ^{})\gamma _{}`$ (105)
$`S^{(13)}(t,\tau ^{})`$ $`=`$ $`\mathrm{exp}\left({\displaystyle _0^t}\lambda (t^{\prime \prime })𝑑t^{\prime \prime }\right)\left[A_{13}(t,\tau ^{})\gamma _++B_{13}(t,\tau ^{})\mathrm{\Gamma }_{}+C_{13}(t,\tau ^{})\mathrm{\Gamma }_+D_{13}(t,\tau ^{})\gamma _{}\right].`$ (106)
Defining positive- and negative-frequency mode functions in imaginary time by
$$\left(\begin{array}{c}f_\mathrm{I}(\tau )\\ g_\mathrm{I}(\tau )\end{array}\right)=\frac{e^{\mathrm{\Omega }_0\tau }}{\sqrt{2\mathrm{\Omega }_0(\mathrm{\Omega }_0\tau _0)}}\left(\begin{array}{c}\sigma _0\\ \mathrm{\Omega }_0\tau _0\end{array}\right)\left(\begin{array}{c}\overline{g}_\mathrm{I}(\tau )\\ \overline{f}_\mathrm{I}(\tau )\end{array}\right)=\frac{e^{\mathrm{\Omega }_0\tau }}{\sqrt{2\mathrm{\Omega }_0(\mathrm{\Omega }_0\tau _0)}}\left(\begin{array}{c}(\mathrm{\Omega }_0\tau _0)\\ \sigma _0^{}\end{array}\right),$$
(107)
the solutions subject to the boundary conditions (99) and (100) are given by
$`A_3(\tau ,\tau ^{})`$ $`=`$ $`(1N^{\mathrm{eq}})f_\mathrm{I}(\tau )\overline{f}_\mathrm{I}(\tau ^{})+N^{\mathrm{eq}}\overline{g}_\mathrm{I}(\tau )g_\mathrm{I}(\tau ^{})`$ (108)
$`B_3(\tau ,\tau ^{})`$ $`=`$ $`(1N^{\mathrm{eq}})g_\mathrm{I}(\tau )\overline{f}_\mathrm{I}(\tau ^{})N^{\mathrm{eq}}\overline{f}_\mathrm{I}(\tau )g_\mathrm{I}(\tau ^{})`$ (109)
$`C_3(\tau ,\tau ^{})`$ $`=`$ $`(1N^{\mathrm{eq}})f_\mathrm{I}(\tau )\overline{g}_\mathrm{I}(\tau ^{})N^{\mathrm{eq}}\overline{g}_\mathrm{I}(\tau )f_\mathrm{I}(\tau ^{})`$ (110)
$`D_3(\tau ,\tau ^{})`$ $`=`$ $`(1N^{\mathrm{eq}})g_\mathrm{I}(\tau )\overline{g}_\mathrm{I}(\tau ^{})+N^{\mathrm{eq}}\overline{f}_\mathrm{I}(\tau )f_\mathrm{I}(\tau ^{})`$ (111)
and
$`A_{13}(t,\tau ^{})`$ $`=`$ $`N^{\mathrm{eq}}f(t)\overline{f}_\mathrm{I}(\tau ^{})(1N^{\mathrm{eq}})g^{}(t)g_\mathrm{I}(\tau ^{})`$ (112)
$`B_{13}(t,\tau ^{})`$ $`=`$ $`N^{\mathrm{eq}}g(t)\overline{f}_\mathrm{I}(\tau ^{})+(1N^{\mathrm{eq}})f^{}(t)g_\mathrm{I}(\tau ^{})`$ (113)
$`C_{13}(t,\tau ^{})`$ $`=`$ $`N^{\mathrm{eq}}f(t)\overline{g}_\mathrm{I}(\tau ^{})+(1N^{\mathrm{eq}})g^{}(t)f_\mathrm{I}(\tau ^{})`$ (114)
$`D_{13}(t,\tau ^{})`$ $`=`$ $`N^{\mathrm{eq}}g(t)\overline{g}_\mathrm{I}(\tau ^{})(1N^{\mathrm{eq}})f^{}(t)f_\mathrm{I}(\tau ^{}),`$ (115)
with $`N^{\mathrm{eq}}=\left[e^{\beta \mathrm{\Omega }_0}+1\right]^1`$. Finally, the boundary conditions (97) and (98) are satisfied provided, as promised, that $`a_1=12N^{\mathrm{eq}}`$ and $`a_2=0`$.
## VII Concrete realization: A simple model
In the preceding sections, we have constructed a quasiparticle action, and the propagators that correspond to it, in terms of several time- and momentum-dependent coefficients that so far are undetermined. In general terms, these coefficients are to be determined self-consistently by asking the counterterm $`_{ab}`$ to cancel some part of the higher-order contributions to the self energy, which arise within an interacting theory that we also left unspecified. Thus, the full 2-point functions can be expressed through the Schwinger-Dyson equation
$$𝒮^{(ab)}(t,t^{};𝒌)=S^{(ab)}(t,t^{},𝒌)i𝑑t^{\prime \prime }𝑑t^{\prime \prime \prime }S^{(ac)}(t,t^{\prime \prime };𝒌)\mathrm{\Sigma }_{cd}(t^{\prime \prime },t^{\prime \prime \prime };𝒌)𝒮^{(db)}(t^{\prime \prime \prime },t^{};𝒌),$$
(116)
in terms of a self energy that has contributions both from the counterterm and from loop corrections
$$\mathrm{\Sigma }_{ab}(t,t^{})=_{ab}(t,_t)\delta (tt^{})+\mathrm{\Sigma }_{ab}^{\mathrm{loop}}(t,t^{}).$$
(117)
If the counterterm could be chosen so that $`\mathrm{\Sigma }_{ab}=0`$, then the propagator $`S^{(ab)}`$ would be the same as the full 2-point function $`𝒮^{(ab)}`$ and the perturbation series would be completely resummed. In practice, of course, we can achieve at best a partial resummation by cancelling the first few terms in the expansion of $`\mathrm{\Sigma }^{\mathrm{loop}}`$. Moreover, since $`\mathrm{\Sigma }^{\mathrm{loop}}`$ is non-local in time, and may well have a more complicated spinor structure than the counterterm we have constructed, it will not be possible to cancel even these terms exactly. We can effect a selective resummation by cancelling only part of $`\mathrm{\Sigma }^{\mathrm{loop}}`$, but the choice of which part to cancel will depend on details of a specific application and of the supplementary approximations that will inevitably be required. Here, we illustrate how the process can be made to yield sensible results by studying the simplest possible model, in which our fermion interacts with a real scalar field, the interaction being specified by
$$_{\mathrm{int}}=g\overline{\psi }\psi \varphi .$$
(118)
We suppose that the $`\varphi `$ particles have a mass $`M`$ that is greater than $`2m`$, so that the decay and annihilation processes $`\varphi \psi +\overline{\psi }`$ are kinematically allowed. We anticipate that these on-shell processes will give absorptive parts to the fermion self energies, yielding a nonzero thermal width $`\lambda `$, and that the equation (91), which gives the rate of change of the occupation numbers $`N_k(t)`$, will, within a suitable approximation, be recognisable as a kinetic equation of the Boltzmann type.
We study explicitly the one-loop, real-time self energy, corresponding to the emission and reabsorption of a $`\varphi `$ particle. By setting $`+\mathrm{\Sigma }^{1\mathrm{loop}}0`$, we obtain a complicated set of constraints, which implicitly specify the functions $`\lambda _k(t)`$, $`ϵ_k(t)`$, etc. In fact, these functions enter $`\mathrm{\Sigma }^{1\mathrm{loop}}`$ through the mode functions $`f_k(t)`$ and $`g_k(t)`$ and the auxiliary functions $`N_k(t)`$ and $`\mathrm{\Delta }_k(t)`$, for which we have no concrete expressions in hand. We know only that they are solutions of (59), (91) and (92). In principle, we have a closed set of equations that we might attempt to solve numerically. To see more clearly what these equations imply, however, we introduce some further approximations. First, we will suppose that time evolution is sufficiently slow for an adiabatic approximation to be reasonable. Then the mode functions will be written as
$$\left(\begin{array}{c}f_k(t)\\ g_k(t)\end{array}\right)\frac{1}{\sqrt{2\mathrm{\Omega }_k(t)(\mathrm{\Omega }_k(t)\tau _k(t))}}\left(\begin{array}{c}\sigma _k(t)\\ \mathrm{\Omega }_k(t)\tau _k(t)\end{array}\right)\mathrm{exp}\left[i_0^t\mathrm{\Omega }_k(t^{})𝑑t^{}\right].$$
(119)
Further, when the fermions have a nonzero thermal width, the propagator will decay at large time separations, very roughly as $`e^{\lambda |tt^{}|}`$. Assuming that $`\tau _k(t)`$ and $`\sigma _k(t)`$ do not change too much over a thermal lifetime $`\lambda _k(t)^1`$, it will be reasonable to approximate the product $`f(t)f^{}(t^{})`$ that appears in (85) as
$$f(t)f^{}(t^{})\frac{\sigma _k(\overline{t})\sigma _k^{}(\overline{t})\mathrm{exp}\left[i\mathrm{\Omega }_k(\overline{t})(tt^{})\right]}{2\mathrm{\Omega }_k(\overline{t})\left[\mathrm{\Omega }_k(\overline{t})\tau _k(\overline{t})\right]},$$
(120)
where $`\overline{t}=(t+t^{})/2`$, with corresponding approximations for other products of mode functions. Finally, since $`\mathrm{\Sigma }^{1\mathrm{loop}}`$ is proportional to $`g^2`$, so are the functions $`\lambda _k(t)`$, etc that appear in $``$. At the lowest order of perturbation theory, it is therefore reasonably consistent to set these functions to zero in the propagators that we use in evaluating $`\mathrm{\Sigma }^{1\mathrm{loop}}`$, and this is what we do. In particular, we then have $`\sigma _k(\overline{t})=|𝒌|`$ and $`\tau _k(\overline{t})=m(\overline{t})`$. Clearly, these approximations are valid, at best, only for a weakly interacting system in a state close to thermal equilibrium. It is therefore important to emphasise that our purpose in introducing them is to obtain simple analytical results that illustrate essential features of our formalism. The formalism itself is by no means restricted to situations where these approximations are valid. For different reasons, discussed below, we will set $`\mathrm{\Delta }(t^{})=0`$.
At this level of approximation, the propagators used in evaluating $`\mathrm{\Sigma }^{1\mathrm{loop}}`$ (and those that multiply $`\mathrm{\Sigma }_{ab}`$ in (116)) are essentially those of the equilibrium theory, except that we allow for time-dependent masses and occupation numbers. It is useful to introduce the projection operators
$`\mathrm{\Lambda }(𝒌)`$ $`=`$ $`\mathrm{\Omega }_k(\overline{t})\gamma ^0𝜸𝒌+m(\overline{t})`$ (121)
$`\stackrel{~}{\mathrm{\Lambda }}(𝒌)`$ $`=`$ $`\mathrm{\Omega }_k(\overline{t})\gamma ^0𝜸𝒌+m(\overline{t})`$ (122)
$`\mathrm{\Lambda }_\pm (𝒌)`$ $`=`$ $`\pm \left({\displaystyle \frac{m^2(\overline{t})k^2}{\mathrm{\Omega }_k(\overline{t})}}\right)\gamma ^0𝜸𝒌+m(\overline{t}),`$ (123)
which have the properties
$`\mathrm{\Lambda }(𝒌)\stackrel{~}{\mathrm{\Lambda }}(𝒌)\mathrm{\Lambda }(𝒌)`$ $`=`$ $`\mathrm{\Lambda }(𝒌)\mathrm{\Lambda }_{}(𝒌)\mathrm{\Lambda }(𝒌)=0`$ (124)
$`\stackrel{~}{\mathrm{\Lambda }}(𝒌)\mathrm{\Lambda }(𝒌)\stackrel{~}{\mathrm{\Lambda }}(𝒌)`$ $`=`$ $`\stackrel{~}{\mathrm{\Lambda }}(𝒌)\mathrm{\Lambda }_+(𝒌)\stackrel{~}{\mathrm{\Lambda }}(𝒌)=0.`$ (125)
After a Fourier transform on the time difference $`tt^{}`$, each of the propagators appearing in the second term of (116) can be written in the form
$$S^{(ab)}(𝒌,\omega ;\overline{t})S_1^{(ab)}(𝒌,\omega ;\overline{t})\mathrm{\Lambda }(𝒌)+S_2^{(ab)}(𝒌,\omega ;\overline{t})\stackrel{~}{\mathrm{\Lambda }}(𝒌),$$
(126)
where $`S_1^{(ab)}(𝒌,\omega ;\overline{t})`$ has poles at $`\omega =\mathrm{\Omega }_k(\overline{t})\pm i\lambda _k(\overline{t})`$, while $`S_2^{(ab)}(𝒌,\omega ;\overline{t})`$ has poles at $`\omega =\mathrm{\Omega }_k(\overline{t})\pm i\lambda _k(\overline{t})`$ and, within the approximations described above, $`\lambda _k(\overline{t})`$ is to be regarded as infinitesimal. In particular, we shall make explicit use of $`S^{(12)}(𝒌,\omega ;\overline{t})`$, in which the poles combine to yield
$$S^{(12)}(𝒌,\omega ;\overline{t})\frac{\pi }{\omega }\left[\omega \gamma ^0𝜸𝒌+m(\overline{t})\right]\left[N_k(\overline{t})\delta \left(\omega \mathrm{\Omega }_k(\overline{t})\right)+\left(1N_k(\overline{t})\right)\delta \left(\omega +\mathrm{\Omega }_k(\overline{t})\right)\right]$$
(127)
and of the corresponding scalar propagator
$$g^{(12)}(𝒌,\omega ;\overline{t})\frac{\pi }{\omega _k(\overline{t})}\left[n_k(\overline{t})\delta \left(\omega \omega _k(\overline{t})\right)+\left(1+n_k(\overline{t})\right)\delta \left(\omega +\omega _k(\overline{t})\right)\right],$$
(128)
with $`\omega _k(\overline{t})=\sqrt{k^2+M^2(\overline{t})}`$.
We now consider how our vaguely-stated criterion $`+\mathrm{\Sigma }^{\mathrm{loop}}0`$ can be implemented in practice, to yield a selective resummation of the perturbation series. To be clear, let us first reiterate the rôle of the simplifying approximations introduced above. In principle, the functions $`\tau _k(t)`$, $`\sigma _k(t)`$, $`\lambda _k(t)\mathrm{}`$ are to be determined self-consistently by solving a set of equations of the form $`(\tau ,\sigma ,\lambda ,\mathrm{})\mathrm{\Sigma }^{\mathrm{loop}}(\tau ,\sigma ,\lambda ,\mathrm{})`$. However, for the purposes of discovering how $``$ might sensibly be interpreted and eventually of making contact with kinetic theory, we plan instead to study the simplified set of equations $`(\tau ,\sigma ,\lambda ,\mathrm{})\mathrm{\Sigma }^{\mathrm{loop}}(m,|𝒌|,0,\mathrm{})`$, for which we can find closed-form expressions. Of the functions that we wish to determine, $`\tau `$ and $`\sigma `$ clearly encode the dispersion relation for the quasiparticles whose mode functions are given by (63). These functions appear only in $`𝒟_1`$ (for which our ansatz was given in (52)) and thus in $`_{11}`$ and $`_{22}`$. By arranging for them to cancel appropriate parts of $`\mathrm{\Sigma }_{11}^{\mathrm{loop}}`$ and $`\mathrm{\Sigma }_{22}^{\mathrm{loop}}`$, we can endow our quasiparticles with dispersion relations that approximate those of the true elementary excitations of the nonequilibrium state. We do not do this explicitly, however, preferring to focus on dissipative aspects of the problem, and especially on the evolution of occupation numbers. In fact, for simplicity, we shall set $`\tau _k(t)=m(t)`$ and $`\sigma _k(t)=|𝒌|`$ in $``$ as well as in $`\mathrm{\Sigma }^{\mathrm{loop}}`$. The terms proportional to $`\mathrm{\Delta }_k(t)`$ in our propagators (85) - (88) are awkward and (as discussed in the next section) difficult to interpret. Within our present approximations, it is possible to eliminate these terms in the following way. We saw at the end of section VI that the constant $`a_{2k}=\mathrm{\Delta }_k(0)`$ vanishes when the initial state is one of thermal equilibrium. It is therefore consistent to set $`\mathrm{\Delta }_k(t)=0`$ at all times, provided that the quantity $`e_{2k}(t)`$ vanishes in (92). The expression given in (84) does not vanish in general, but with the approximations made here it can be made to do so by choosing $`ϵ_k(t)`$ to be purely imaginary, say $`ϵ_k(t)=i\widehat{ϵ}_k(t)`$, and by choosing $`\eta _k(t)=\tau _k(t)\widehat{ϵ}_k(t)/\sigma _k(t)=m(t)\widehat{ϵ}_k(t)/|𝒌|`$.
With these simplifications, the functions that remain to be determined are $`\lambda _k(t)`$ and $`\widehat{ϵ}_k(t)`$, which appear in the counterterm
$$_{12}(𝒌,t)=\frac{i}{2\mathrm{\Omega }_k(t)}\left[\left(\lambda _k(t)+\frac{\mathrm{\Omega }_k(t)}{|𝒌|}\widehat{ϵ}_k(t)\right)\mathrm{\Lambda }(𝒌)\left(\lambda _k(t)\frac{\mathrm{\Omega }_k(t)}{|𝒌|}\widehat{ϵ}_k(t)\right)\stackrel{~}{\mathrm{\Lambda }}(𝒌)\right].$$
(129)
Ideally, we would like this to cancel the one-loop contribution
$$\mathrm{\Sigma }_{12}^{1\mathrm{loop}}(𝒌,\omega ;\overline{t})ig^2\frac{d\omega ^{}d\omega ^{\prime \prime }}{2\pi }\frac{d^3k^{}d^3k^{\prime \prime }}{(2\pi )^3}\delta (\omega \omega ^{}\omega ^{\prime \prime })\delta (𝒌𝒌^{}𝒌^{\mathbf{\prime \prime }})S^{(12)}(𝒌^{},\omega ^{};\overline{t})g^{(12)}(𝒌^{\prime \prime },\omega ^{\prime \prime };\overline{t}),$$
(130)
but $`_{12}`$, being derived from a counterterm that is local in time, has no dependence on $`\omega `$, while $`\mathrm{\Sigma }_{12}^{1\mathrm{loop}}`$ cannot be expressed in the form of (129) as a linear combination of $`\mathrm{\Lambda }(𝒌)`$ and $`\stackrel{~}{\mathrm{\Lambda }}(𝒌)`$. Clearly, we must be a little less ambitious. First, we shall attempt to effect the desired cancellation on shell: that is, to cancel only $`\mathrm{\Sigma }_{12}^{1\mathrm{loop}}(𝒌,\pm \mathrm{\Omega }_k(\overline{t}))`$. For $`\omega =\pm \mathrm{\Omega }`$, we find
$`\mathrm{\Sigma }_{12}^{1\mathrm{loop}}(𝒌,\omega ;\overline{t}){\displaystyle \frac{ig^2}{16\pi ^2}}{\displaystyle d^3k^{}\frac{\delta (\omega _p\mathrm{\Omega }\mathrm{\Omega }^{})}{\mathrm{\Omega }^{}\omega _p}}`$ $`[(\mathrm{\Omega }^{}\gamma ^0+{\displaystyle \frac{𝒌𝒌^{}}{|𝒌|^2}}𝜸𝒌+m)n(1N^{})\theta (\omega )`$ (132)
$`(\mathrm{\Omega }^{}\gamma ^0+{\displaystyle \frac{𝒌𝒌^{}}{|𝒌|^2}}𝜸𝒌+m)(1+n)N^{}\theta (\omega )].`$
Here, the kinematics is that of on-shell decay or pair annihilation, involving two fermions of momenta $`𝒌`$ and $`𝒌^{}`$, with energies $`\mathrm{\Omega }=\mathrm{\Omega }_k(\overline{t})`$ and $`\mathrm{\Omega }^{}=\mathrm{\Omega }_k^{}(\overline{t})`$, and a scalar with momentum $`𝒑=𝒌+𝒌^{}`$ and energy $`\omega _p=\sqrt{|𝒑|^2+M^2(\overline{t})}`$. The fermion occupation numbers $`N=N_k(\overline{t})`$, and $`N^{}=N_k^{}(\overline{t})`$ are those that we originally introduced in (85) - (88), while $`n=n_p(\overline{t})`$ is the corresponding quantity for the scalar.
When $`\omega `$ is close to $`+\mathrm{\Omega }_k`$, the propagators given approximately by (126) can be further reduced by retaining only the term containing a pole, namely $`S^{(ab)}S_1^{(ab)}\mathrm{\Lambda }(𝒌)`$. The one-loop correction term in the Schwinger-Dyson equation (116) then involves the projection $`\mathrm{\Lambda }(𝒌)\left[_{12}+\mathrm{\Sigma }_{12}^{1\mathrm{loop}}\right]\mathrm{\Lambda }(𝒌)`$. Using the properties (124), we see that this can be made to vanish by expressing $`\mathrm{\Sigma }_{12}^{1\mathrm{loop}}`$ as a linear combination of $`\mathrm{\Lambda }(𝒌)`$, $`\stackrel{~}{\mathrm{\Lambda }}(𝒌)`$ and $`\mathrm{\Lambda }_{}(𝒌)`$ and requiring the coefficients of $`\mathrm{\Lambda }(𝒌)`$ in $`\mathrm{\Sigma }_{12}^{1\mathrm{loop}}`$ and $`_{12}`$ to cancel. Similarly, when $`\omega `$ is close to $`\mathrm{\Omega }_k`$, we express $`\mathrm{\Sigma }_{12}^{1\mathrm{loop}}`$ as a linear combination of $`\stackrel{~}{\mathrm{\Lambda }}(𝒌)`$, $`\mathrm{\Lambda }(𝒌)`$ and $`\mathrm{\Lambda }_+(𝒌)`$, and require the coefficients of $`\stackrel{~}{\mathrm{\Lambda }}(𝒌)`$ to cancel. In this way, we obtain
$`\lambda _k(\overline{t})`$ $`=`$ $`{\displaystyle \frac{g^2}{64\pi ^2}}(M^24m^2){\displaystyle d^3k^{}\frac{\delta (\omega _p\mathrm{\Omega }\mathrm{\Omega }^{})}{\mathrm{\Omega }\mathrm{\Omega }^{}\omega _p}\left[n+N^{}\right]}`$ (133)
$`{\displaystyle \frac{\mathrm{\Omega }_k(\overline{t})}{|𝒌|}}\widehat{ϵ}_k(\overline{t})`$ $`=`$ $`{\displaystyle \frac{g^2}{64\pi ^2}}(M^24m^2){\displaystyle d^3k^{}\frac{\delta (\omega _p\mathrm{\Omega }\mathrm{\Omega }^{})}{\mathrm{\Omega }\mathrm{\Omega }^{}\omega _p}\left[n(1N^{})(1+n)N^{}\right]}.`$ (134)
Here, we have made use of the kinematic identity $`\mathrm{\Omega }\mathrm{\Omega }^{}𝒌𝒌^{}m^2=\frac{1}{2}(M^24m^2)`$. Reassuringly, we have arrived at a thermal quasiparticle width that is positive-definite when $`M>2m`$, so that the on-shell decay and annihilation processes that give rise to it are kinematically allowed. When $`M<2m`$, the thermal width vanishes, because (130) contains products of $`\delta `$-functions that cannot be satisfied simultaneously. The expression (133) agrees with the equilibrium damping rate calculated in (see also )for the same model.
Finally, we can evaluate the right hand side of equation (91), which we hoped to interpret as governing the evolution of time-dependent occupation numbers $`N_k(t)`$. With the approximations used in this section, we have $`e_{1k}(t)=2\mathrm{\Omega }_k(t)\widehat{ϵ}_k(t)/|𝒌|`$, and the equation becomes
$$\frac{dN_k(t)}{dt}=\frac{g^2}{32\pi ^2}(M^24m^2)d^3k^{}\frac{\delta (\omega _p\mathrm{\Omega }\mathrm{\Omega }^{})}{\mathrm{\Omega }\mathrm{\Omega }^{}\omega _p}\left[n(1N)(1N^{})(1+n)NN^{}\right].$$
(135)
We recognise the standard form of a relativistic Boltzmann equation (restricted to a spatially homogeneous system) in which the gain and loss terms have the correct statistical factors, $`n(1N)(1N^{})`$ and $`(1+n)NN^{}`$ respectively, to represent a fermion of momentum $`𝒌`$ being produced by the decay of a $`\varphi `$ in the thermal bath or annihilating with an antifermion in the bath. Again, the ‘collision’ integral would be replaced by 0 if $`M<2m`$. Of course, our results for $`\lambda _k(t)`$ and $`dN_k(t)/dt`$ were obtained only at the lowest non-trivial order of perturbation theory. At higher orders, we would expect non-zero answers in both cases, arising from scattering processes that are allowed even for $`M<2m`$.
## VIII Discussion
We have described a selective resummation of the perturbation series for the nonequilibrium 2-point functions of spin-$`\frac{1}{2}`$ fermions. The general philosophy of this resummation is to describe the nonequilibrium state as nearly as possible in terms of its own quasiparticle excitations. The propagators for these excitations, unlike the free-particle propagators used in standard perturbation theory ought, roughly speaking, to incorporate nonzero widths and occupation numbers that evolve with time, reflecting the evolution of the nonequilibrium state. The resummation is achieved through the use of a counterterm that transfers some contributions of higher-order self energies into the lowest-order theory about which we perturb. In a nonequilibrium situation, this is tractable only if the unperturbed action is local in time, and this places strong constraints on what can be resummed in practice. For example, resummations somewhat similar in spirit to ours, but restricted to scalar theories in thermal equilibrium, are described in and applied to the ‘warm inflation’ scenario in . In thermal equilibrium, the 2-point functions depend only on $`tt^{}`$, and after a Fourier transformation, one can (in principle) construct a counterterm analogous to $``$ that subtracts the whole frequency-dependent self energy at whatever order of perturbation theory one has the energy to compute. Generalizing this to a nonequilibrium state would mean replacing (48) and (49) with integro-differential equations containing arbitrary non-local kernels in place of the functions $`\tau _k(t)`$, $`\lambda _k(t)`$, etc. (or perhaps with infinite-order differential equations having infinitely many time-dependent coefficients, all to be determined self-consistently) and this does not seem to be a practical proposition. For fermions, this may be particularly unfortunate, because the high-temperature plasma has, at least in some important cases, ‘hole’ or ‘plasmino’ excitations in addition to the particle and antiparticle poles that are present at $`T=0`$. These have been known for some time from one-loop calculations and their properties have recently been explored in terms independent of perturbation theory by Weldon . The counterterm we have constructed is linear in $`_t`$ and cannot accommodate a dispersion relation with these multiple branches (though a generalization that mimics them might be possible).
Despite this inevitable deficiency, the resummation appears to make reasonable sense. Of the initially arbitrary parameters that we introduced in (52) and (53) to describe the quasiparticle excitations, $`\tau _k(t)`$ and $`\sigma _k(t)`$ have clear interpretations in terms of the quasiparticle dispersion relation, while $`\lambda _k(t)`$ is a thermal width which, after enough weak-coupling and adiabatic approximations, turns out to agree with the fermion damping rate calculated in equilibrium. In the same approximation, the function $`N_k(t)`$ that appears in our resummed propagators does indeed correspond to a quasiparticle occupation number, evolving according to a kinetic equation of the Boltzmann type. It is worth emphasizing that, while this kinetic equation provides reassurance that our formalism has a sensible interpretation, its derivation is by no means the purpose of the formalism presented here. There are indeed, many routes to equations of this kind. One may, for example, investigate directly the time evolution of the expectation value of a time-dependent number operator (see, e.g. ). Another route that has been pursued extensively in connection with the calculation of transport coefficients of high-temperature plasmas is to extract a transport equation for a Wigner density through truncation and gradient expansion of the Schwinger-Dyson equations . Attempts to calculate transport coefficients directly from the Kubo formula of linear response theory reveal, on the other hand, infrared singularities that require the resummation of large classes of diagrams and this resummation turns out to be equivalent to solving a Boltzmann equation . (The relationship between these approaches is discussed in ). The transport equations that arise in these calculations go well beyond the simple one-loop approximation exhibited in (135), but they apply to systems very close to equilibrium. Our own goal of constructing a resummed perturbation theory to describe the evolution of highly-excited states that may be far from equilibrium is rather different. The functions $`N_k(t)`$ that arise in the course of solving for the resummed propagators are not necessarily equivalent to the Wigner distribution, and the equation (91) that defines them reduces to a Boltzmann equation only after approximations that one may in general hope to avoid.
Finally, the investigation reported here indicates that the structure of nonequilibrium fermion propagators is more complicated than might be expected from the equilibrium theory. Close to equilibrium, we found that their spinor structure can be expressed in terms of the two projection operators (121) and (122), as can be deduced on general grounds within the equilibrium theory (see, e.g. ). Away from equilibrium, this is no longer true. In general, there are at least terms proportional to $`\gamma ^0𝜸𝒌`$ and possibly other terms that we have not succeeded in resumming. To arrive at propagators that contain only $`\mathrm{\Lambda }(𝒌)`$ and $`\stackrel{~}{\mathrm{\Lambda }}(𝒌)`$, we need the coefficients $`B(t,t^{})`$ and $`C(t,t^{})`$ in (58) to be equal. This will be true of the solutions presented in (86) and (87) if the function $`\mathrm{\Delta }(t)`$ vanishes, and if the products of mode functions $`g(t)f^{}(t^{})`$ and $`f(t)g^{}(t^{})`$ are equal. Both of these conditions would be automatic if we were to assume time-translation invariance, so that the propagators depend only on $`tt^{}`$, which is, of course, true in thermal equlibrium. The equation (92) satisfied by $`\mathrm{\Delta }(t)`$ is superficially similar to (91), but does not bear the same interpretation as a kinetic equation. In fact, the terms in the propagator that involve this function are products (loosely speaking) of two positive-frequency or two negative-frequency mode functions that have no counterpart in equilibrium. In general, it seems that such terms should be present in nonequilibrium self energies, but we have found no simple interpretation for them.
###### Acknowledgements.
The final version of this work has benefitted from discussions with participants at the program on Non-equilibrium Dynamics in Quantum Field Theory (1999) at the Institute for Nuclear Theory, University of Washington. IDL would like to thank, in particular, Larry Yaffe and Francoise Guerin for helpful comments and the INT for its hospitality. DBM thanks the University of Leeds for financial support through the award of a William Wright Smith scholarship.
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# Abstract
## Abstract
The construction of anyonic operators and algebra is generalized by using quons operators. Therefore, the particular version of fractional supersymmetry is constructed on the two-dimensional lattice by associating two generalized anyons of different kinds. The fractional supersymmetry Hamiltonian operator is obtained on the two-dimensional lattice and the quantum algebra $`U_q(sl_2)`$ is realized.
## 1 Introduction
A grown interest has been devoted, recently, to fractional statistics$`^{\text{[1, 2]}}`$ describing particles with fractional spin$`^{\text{[3, 4, 5]}}`$. These particles are called anyons and interpolate between bosons and fermions, it has been shown also that anyons are non-local particles defined on the two-dimensional space. Mathematically, the group symmetry associated to the anyonic systems involves some special Lie algebras appearing as a quantum deformation of the usual Lie algebras and Lie groups$`^{\text{[6, 7, 8, 9, 10]}}`$. Indeed, one can prove that quantum groups$`^{\text{[11]}}`$ are the mathematical objects allowing the description of these particular systems$`^{\text{[12, 13]}}`$. Many works have been devoted, in the last few years, to the construction and the realization of quantum groups. We note for example the anyonic algebra$`^{\text{[14]}}`$ leading to the obtaining of these quantum algebras. The anyonic operators are introduced on the two-dimensional lattice as a non-local operators seen as a generalization of generators of the Jordan-Wigner transformation$`^{\text{[15]}}`$.
The present work concerns the study on the two-dimensional lattice of N=2 fractional supersymmetry (FSUSY)$`^{\text{[16, 17, 18, 19]}}`$. We obtain exactly the same algebraic structure of this latter starting from the introduction of the generalized anyonic algebra. Then we introduce some special operators which can be seen as a generalized anyonic operators. We notice that this construction is different from the one used to reproduce the N=2 FSUSY algebra basing on the quonic algebra$`^{\text{[19]}}`$. This algebraic study allows us to find the Hamiltonian operator describing one N=2 FSUSY system.
This paper is organized as follows : In section $`2`$, as a generalization, we construct new exotic operators in using quonic operators as basic ones. Therefore we realize the generalized anyonic operators and the corresponding algebra. In the third section, the last new operators and algebra will be used to construct N=2 FSUSY on the two-dimensional lattice corresponding to anyonic systems as a certain coupling of two generalized anyonic oscillators of different kinds ($`\gamma `$ and $`\delta `$)$`^{\text{[14]}}`$. We proceed as in the work where we have constructed the N=2 FSUSY through two different quons. Another result consists on the realization of the quantum algebra $`U_q(sl_2)`$ using the supercharges already constructed. Finally, section $`4`$ presents some concluding remarks.
## 2 Generalized anyonic algebra
Let us recall at first the anyonic operators which are seen as non-local two-dimensional operators interpolating between bosonic and fermionic ones. The generalization of these operators allows us to introduce the generalized anyons. We will use the famous angle function $`\mathrm{\Theta }(x,y)`$ appearing in the construction and the description of anyons in the work .
We start by giving a brief review on this angle function. One designs by $`\gamma _x`$ the cut associated to each point $`x`$ that we denote by $`x_\gamma `$ on the two-dimensional lattice $`\mathrm{\Omega }`$. Denoting by $`\mathrm{\Omega }^{}`$ the dual lattice of $`\mathrm{\Omega }`$; it is a set of points $`x^{}=x+0^{}`$ with $`0^{}=(\frac{1}{2}ϵ,\frac{1}{2}ϵ)`$ the origin of $`\mathrm{\Omega }^{}`$ and $`ϵ`$ its spacing which eventually will be sent to zero. In this case $`\gamma _x`$ will be on $`\mathrm{\Omega }^{}`$ from minus infinity to $`x^{}`$ along the $`x`$-axis. One considers another type of cuts. We choose the set of cuts $`\delta _x`$ coming from plus infinity to $`{}_{}{}^{}x=x0^{}`$ along the $`x`$-axis. Consequently, the two types of cuts $`\gamma _x`$ and $`\delta _x`$ involve an ordering and opposite ordering respectively of points $`x`$ on $`\mathrm{\Omega }`$. This is described by the following proposition
$$x_\delta <y_\delta x_\gamma >y_\gamma x>y\{\begin{array}{cc}x_2>y_2,& \\ x_1>y_1,x_2=y_2& \end{array}$$
(1)
Owing to the equation $`(1)`$ the angle functions satisfy
$$\begin{array}{ccc}\hfill \mathrm{\Theta }_{\gamma _x}(x,y)\mathrm{\Theta }_{\gamma _y}(y,x)& =& \{\begin{array}{ccc}\hfill \pi \text{}x>y& & \\ \hfill \pi \text{}x<y& & \end{array}\hfill \\ \hfill \mathrm{\Theta }_{\delta _x}(x,y)\mathrm{\Theta }_{\delta _y}(y,x)& =& \{\begin{array}{ccc}\hfill \pi \text{}x>y& & \\ \hfill \pi \text{}x<y& & \end{array}\hfill \\ \hfill \mathrm{\Theta }_{\delta _x}(x,y)\mathrm{\Theta }_{\gamma _x}(x,y)& =& \{\begin{array}{ccc}\hfill \pi \text{}x>y& & \\ \hfill \pi \text{}x<y& & \end{array}\hfill \\ \hfill \mathrm{\Theta }_{\delta _x}(x,y)\mathrm{\Theta }_{\gamma _y}(y,x)& =& 0,x,y\mathrm{\Omega }.\hfill \end{array}$$
(2)
Now let us introduce the operators $`K_i(x_\alpha )`$ those called the disorder ones. They are expressed by
$$K_i(x_\alpha )=e^{i\nu \underset{yx}{}\mathrm{\Theta }_{\alpha _x}(x,y)[N_i(y)\frac{1}{2}]}$$
(3)
with $`\alpha _x=\gamma _x`$ or $`\delta _x`$ and $`i=1,2,\mathrm{},n,n𝐍`$. In the equality $`(3)`$ $`N_i(y)`$ is nothing but the number operator of quons on the two-dimensional lattice, defined by the generators $`a_i^{}(x)`$ and $`a_i(x)`$ as follows
$$\begin{array}{ccc}\hfill a_i^{}(x)a_i(x)& =& [N_i(x)]_{q_i}\hfill \\ \hfill a_i(x)a_i^{}(x)& =& [N_i(x)+\mathrm{𝟏}]_{\mathrm{q}_\mathrm{i}}\hfill \end{array}$$
(4)
where $`[x]_q=\frac{q^x1}{q1}`$. The operators $`a_i^{}(x)`$ and $`a_i(x)`$ are respectively the creation and annihilation quonic operators on the two-dimensional lattice $`\mathrm{\Omega }`$, satisfying the relations
$$\begin{array}{ccc}\hfill [a_i(x),a_j^{}(y)]_{q_i^{\delta _{ij}}}& =& \delta _{ij}\delta (x,y)\hfill \\ \hfill [a_i(x),a_j(y)]_{q_i^{\delta _{ij}}}& =& \begin{array}{cc}0& x,y\text{}i,j\end{array}\hfill \\ \hfill [a_i^{}(x),a_j^{}(y)]_{q_i^{\delta _{ij}}}& =& \begin{array}{cc}0& x,y\text{}i,j\end{array}\hfill \\ \hfill [N_i(x),a_j(y)]& =& \delta _{ij}\delta (x,y)a_i(x)\hfill \\ \hfill [N_i(x),a_j^{}(y)]& =& \delta _{ij}\delta (x,y)a_i^{}(x)\hfill \end{array}$$
(5)
The Dirac function is defined by
$$\delta (x,y)=\{\begin{array}{cc}1& \text{if }x=y\\ 0& \text{ if }x0\end{array}$$
(6)
In the relations $`(5)`$, one can consider the operators $`a_i^{}(x)`$, $`a_i(x)`$ and $`N_i(x)`$; those are generators of the oscillator algebra describing a system of quons. They are seen as a non-local particles and we define them to satisfy the equality
$$(a_i(x))^{d_i}=(a_i^{}(x))^{d_i}=0$$
(7)
where we take the deformation parameter $`q_i`$ to be a root of unity, $`q_i^{d_i}=1`$ then $`q_i=e^{i\frac{2\pi }{d_i}}`$. One can show that the irreducible representation space (Fock space) of the algebra (relations $`(5)`$) is given by the set
$$F_{i_x}=\{|n_{i_x}>,n_{i_x}=0,1,\mathrm{},d_i1\}$$
(8)
where the notation $`i_x`$ means that this Fock space is introduced in each site of the $`\mathrm{\Omega }`$. the actions of the generators $`a_i(x)`$, $`a_i^{}(x)`$ and $`N_i(x)`$ are expressed by the following equalities
$$\begin{array}{cc}\hfill a_i^{}(x)|n_{i_x}>=|n_{i_x}+1>,& a_i^{}(x)|d_i1>=0\hfill \\ \hfill a_i(x)|n_{i_x}>=[n_{i_x}]_{q_i}|n_{i_x}1>,& a_i(x)|0>=0.\hfill \end{array}$$
(9)
Now we introduce an operators allowing the description of the N=2 FSUSY on the two-dimensional lattice. This realization is seen as an anyonic one, in the sense that the generators will depend on a function which seems to be the angle function used by the authors in the work when they describe the system of anyons. This leads to the definition of the operators
$$A_i(x_\alpha )=K_i(x_\alpha )a_i(x),$$
(10)
where $`K_i(x_\alpha )`$ is the disorder operator introduced in the equation $`(3)`$.
One proves that these operators obey the following commutation relations
$$\begin{array}{ccc}\hfill [A_i(x_\gamma ),A_i(y_\gamma )]_{q_ip^1}& =& 0\text{}x>y\hfill \\ \hfill [A_i^{}(x_\gamma ),A_i^{}(y_\gamma )]_{q_ip^1}& =& 0\text{}x>y\hfill \\ \hfill [A_i(x_\gamma ),A_i^{}(y_\gamma )]_{q_ip}& =& 0\text{}x>y\hfill \\ \hfill [A_i^{}(x_\gamma ),A_i(y_\gamma )]_{q_ip}& =& 0\text{}x>y\hfill \\ \hfill [A_i(x_\gamma ),A_i^{}(x_\gamma )]_{q_i}& =& 1\hfill \\ \hfill [A_i(x_\gamma ),A_j(y_\gamma )]& =& 0\text{}ij\text{}x,y\mathrm{\Omega }\hfill \\ \hfill [A_i^{}(x_\gamma ),A_j^{}(y_\gamma )]& =& 0\text{}ij\text{}x,y\mathrm{\Omega }\hfill \\ \hfill [A_i^{}(x_\gamma ),A_j(y_\gamma )]& =& 0\text{}ij\text{}x,y\mathrm{\Omega }\hfill \\ \hfill [A_i(x_\gamma ),A_j^{}(y_\gamma )]& =& 0\text{}ij\text{}x,y\mathrm{\Omega }.\hfill \end{array}$$
(11)
for the anyonic operators of type $`\gamma `$. We point out that the same results can be obtained for the kind $`\delta `$ in replacing $`p`$ by $`p^1`$ in $`(11)`$. It is obvious thus to get the commutation relations between the different kinds of generalized anyonic operators, we have so
$$\begin{array}{ccc}\hfill [A_i(x_\delta ),A_j(y_\gamma )]_{q_i^{\delta _{ij}}}& =& 0,\text{ }x,y\mathrm{\Omega }\hfill \\ \hfill [A_i(x_\delta ),A_j^{}(y_\gamma )]_{q_i^{\delta _{ij}}}& =& \delta _{ij}\delta (x,y)p^{[\underset{z<x}{}\underset{z>x}{}][N_i(z)\frac{1}{2}]}\hfill \end{array}$$
(12)
with $`p=e^{i\nu \pi }`$, $`\nu `$ is seen as the statistical parameter$`^{\text{[1, 2]}}`$.
By considering the generalized anyonic oscillators $`A_i(x_\alpha )`$ and $`A_i^{}(x_\alpha )`$, we can demonstrate the following relations
$$(A_i(x_\alpha ))^{d_i}=(A_i^{}(x_\alpha ))^{d_i}=0$$
(13)
with $`\alpha =\gamma ,\delta `$ and $`i=1,2,\mathrm{},N`$. The equation $`(13)`$ can be seen as a generalization of the hard-core condition found when we have study the generalized statistics in the works . Indeed, for $`d_i=2`$ one recover the result allowing us, in this work, to think about some connexion between the generalized statistics and the anyonic ones. Returning to the present paper, the relation $`(13)`$ can be seen a nilpotency condition leading to the study of quons on the two-dimensional lattice. This equality is obtained starting from the function $`K_i(x_\alpha )`$ discussed in the work , which is the subject of the following equations
$$\begin{array}{ccc}\hfill K_i^{}(x_\alpha )a_i(y)& =& e^{i\nu \mathrm{\Theta }_{\alpha _x}(x,y)}a_i(y)K_i^{}(x_\alpha )\hfill \\ \hfill K_i^{}(x_\alpha )a_i^{}(y)& =& e^{i\nu \mathrm{\Theta }_{\alpha _x}(x,y)}a_i^{}(y)K_i^{}(x_\alpha )\hfill \\ \hfill K_i(x_\alpha )a_i(y)& =& e^{i\nu \mathrm{\Theta }_{\alpha _x}(x,y)}a_i(y)K_i(x_\alpha )\hfill \\ \hfill K_i(x_\alpha )a_i^{}(y)& =& e^{i\nu \mathrm{\Theta }_{\alpha _x}(x,y)}a_i^{}(y)K_i(x_\alpha )\hfill \\ \hfill K_i^{}(x_\alpha )K_i(y_\alpha )& =& K_i(y_\alpha )K_i^{}(x_\alpha )\hfill \end{array}$$
(14)
It is proved that, on the above Fock space, these algebraic relations are coherent with the following equalities
$$\begin{array}{ccc}\hfill A_i^{}(x_\alpha )|n_{i_x}>& =& e^{i\frac{\nu }{2}\underset{yx}{}\mathrm{\Theta }_{\alpha _x}(x,y)}|n_{i_x}+1>\hfill \\ \hfill A_i(x_\alpha )|n_{i_x}>& =& [n_{i_x}]_{q_i}e^{i\frac{\nu }{2}\underset{yx}{}\mathrm{\Theta }_{\alpha _x}(x,y)}|n_{i_x}1>\hfill \\ \hfill A_i^{}(x_\alpha )|d_i1>& =& 0\hfill \\ \hfill A_i(x_\alpha )|0>& =& 0.\hfill \end{array}$$
(15)
We note that the operators $`A_i^{}(x_\alpha )`$ and $`A_i(x_\alpha )`$ are seen, owing to the relations $`(14)`$ and $`(15)`$, as generalized creation and annihilation anyonic operators respectively.
## 3 N=2 FSUSY on the two-dimensional lattice
The construction of generalized anyonic operators allows us to realize the N=2 FSUSY. This realization involves two different generalized of anyons and is analogous to the one obtained via two different quons in the work . We introduce then the fractional supercharges on $`\mathrm{\Omega }`$ as
$$\begin{array}{ccc}\hfill Q_+(x)& =& A_1^{}(x_\gamma )A_2(x_\delta )\hfill \\ \hfill Q_{}(x)& =& A_1(x_\delta )A_2^{}(x_\gamma ).\hfill \end{array}$$
(16)
basing on the equality $`(13)`$; one can get
$$(Q_\pm (x))^{d_2}=0,$$
(17)
where $`d_2`$ is required to be $`<d_1`$.
In using the construction given in the relation $`(16)`$, we can get the action of the above fractional supercharges on the Fock space defined by the tensor product
$$F_x=F_{1_x}F_{2_x}$$
(18)
with $`F_{1_x}`$ and $`F_{2_x}`$ correspond respectively to the different generalized anyons used to introduce the fractional supercharges.
Now we can get easily the relations
$$\begin{array}{ccc}\hfill Q_+(x)|n_{1_x}>|n_{2_x}>=p^{\frac{1}{2}}[n_{2_x}]_{q_2}|n_{1_x}+1>|n_{2_x}1>& & \\ \hfill Q_{}(x)|n_{1_x}>|n_{2_x}>=p^{\frac{1}{2}}[n_{1_x}]_{q_1}|n_{1_x}1>|n_{2_x}+1>& & \end{array}$$
(19)
Through this realization of the fractional supercharges $`Q_\pm (x)`$, we can show that these new operators satisfy the following commutation relation in the case of $`x>y`$ on $`\mathrm{\Omega }`$
$$q_1Q_+(x)Q_{}(y)q_2Q_{}(y)Q_+(x)=\delta (x,y)P[q_1[N_1(x)]_{q_1}q_2[N_2(x)]_{q_2}]$$
(20)
where the operator $`P`$ is written as
$$P=p^{[\underset{z<x}{}\underset{z>x}{}][N_1(z)+N_2(z)1]}$$
(21)
One can remark that the equality $`(20)`$ is not invariant under the hermitian conjugate. This is related to the fact that these generators involve a complex numbers $`q_{1,2}`$ and $`p`$ which are different from $`\pm 1`$. To avoid this difficulty we introduce the hermitian conjugate of the generators $`Q_\pm (x)`$ and after calculation it is easy to obtain the conjugate equation of $`(20)`$ as
$$q_1^1Q_{}^{}(y)Q_+^{}(x)q_2^1Q_+^{}(x)Q_{}^{}(y)=\delta (x,y)P^1[q_1^1[N_1(x)]_{q_1^1}q_2^1[N_2(x)]_{q_2^1}],$$
(22)
We can verify also that
$$(Q_\pm ^{}(x))^{d_2}=0.$$
(23)
for $`d_2<d_1`$.
Another reason to introduce this operation is to construct the Hamiltonian operator corresponding to this system. We can thus express the FSUSY Hamiltonian operator as
$$\begin{array}{ccc}\hfill q_1Q_+(x)Q_{}(y)+q_1^1Q_{}^{}(y)Q_+^{}(x)q_2Q_{}(y)Q_+(x)q_2^1Q_+^{}(x)Q_{}^{}(y)=& & \\ \hfill \delta (x,y)[Pq_1[N_1(x)]_{q_1}+P^1q_1^1[N_1(x)]_{q_1^1}Pq_2[N_2(x)]_{q_2}P^1q_2^1[N_2(x)]_{q_2^1}].& & \end{array}$$
(24)
Using the relation
$$[N_i(x)]_{q_i^1}=q_i^{1N_i(x)}[N_i(x)]_{q_i}.$$
(25)
The RHS of $`(24)`$ will be rewritten as
$$RHS=\delta (x,y)[P^1q_1^{N_1(x)}+Pq_1][N_1(x)]_{q_1}[P^1q_2^{N_2(x)}+Pq_2][N_2(x)]_{q_2},$$
(26)
The last equation can be interpreted as a Hamiltonian of the generalized anyonic system investigated. So, we can rewrite this Hamiltonian operator as
$$H(x)=\underset{i,j=1,2}{}ϵ_{ij}\frac{\mathrm{sin}[\nu \pi \underset{z\mathrm{\Omega }}{}\mathrm{}(z)+\frac{\pi }{d_i}(2N_i(x)+1)]\mathrm{sin}[\nu \pi \underset{z\mathrm{\Omega }}{}\mathrm{}(z)+\frac{\pi }{d_i}]}{\mathrm{sin}\frac{\pi }{d_i}},$$
(27)
where $`(ϵ_{ij})=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)`$ and $`\underset{z\mathrm{\Omega }}{}\mathrm{}(z)=(\underset{z<x}{}\underset{z>x}{})(N_1(z)+N_2(z)1)`$. Up to now, we can recapulate our result which consists on the complet description of the N=2 FSUSY basing on a given anyonic realization.
Furthermore a global version of this realization can be readily constructed as follows
$$H=\underset{x\mathrm{\Omega }}{}H(x),$$
(28)
where the global supercharges are defined as
$$Q_\pm =\underset{x\mathrm{\Omega }}{}Q_\pm (x),$$
(29)
In addition, we can link the N=2 FSUSY obtained on the two-dimensional lattice to the ”local” algebra $`U_q(sl_2)`$ in considering $`q_1=q_2=q`$. In this particular case we define the three local generators as
$$\begin{array}{ccc}\hfill J_\pm (x)=P^{\frac{1}{2}}q^{\frac{N_2(x)}{2}}Q_\pm (x)& & \\ \hfill J_3(x)=\frac{1}{2}(N_1(x)N_2(x)).& & \end{array}$$
(30)
We can easily check that these local densities of quantum group generators satisfy the following commutation relations
$$\begin{array}{ccc}\hfill [J_+(x),J_{}(y)]=\delta (x,y)[2J_3(x)]_q& & \\ \hfill [J_3(x),J_\pm (y)]=\pm \delta (x,y)J_\pm (x).& & \end{array}$$
(31)
Thus to define the global generators it is sufficient to write
$$\begin{array}{ccc}\hfill J_\pm =\underset{x\mathrm{\Omega }}{}J_\pm (x)& & \\ \hfill J_3=\underset{x\mathrm{\Omega }}{}J_3(x)& & \end{array}$$
(32)
and close the $`U_q(sl_2)`$ algebra as
$$\begin{array}{ccc}\hfill [J_+,J_{}]=[2J_3]_q& & \\ \hfill [J_3,J_\pm ]=\pm J_\pm .& & \end{array}$$
(33)
Consequently, we have close the algebra of $`U_q(sl_2)`$ generated by $`J_\pm `$ and $`J_3`$ which are built out of generalized anyonic oscillators.
## 4 Conclusion
To conclude, we can summarize the lines of this paper in saying that we have constructed a generalized anyonic operators on the two-dimensional lattice in using the q-bosonic operators or called quonic operators, as a generalization of ones defined in the paper . Moreover, we have realized the N=2 FSUSY on the two-dimensional lattice. Where the supercharges are constructed by coupling two different generalized anyonic operators, and the FSUSY Hamiltonian operator of the corresponding system is given. Thus from the N=2 FSUSY realized the well known algebra $`U_q(sl_2)`$ is derived.
It will be interesting also to rewrite the obtained FSUSY Hamiltonian, in the connexion with the gauge theory, as a form of which we held the term coincidering with the Chern-Simons term. This point is the matter that we are preparing for the next paper.
## Acknowledgments
The author J. Douari would like to thank the Max-Planck-Institut für Physik Komplexer Systeme for link hospitality during the stage in which one part of this paper was done and the International Atomic Energy Agency and UNESCO for hospitality at the Abdus Salam International Centre for Theoretical Physics, Trieste. This work was achieved within the framework of the Associateship Scheme of the Abdus Salam International Centre for Theoretical Physics. And special thanks to Professors Manfred Scheunert, Ruibin Zhang and G. Thompson for an inspiring comment and many useful discussions.
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# Homogeneous Lorentz manifolds with simple isometry group
## 1. Introduction
###### 1.1 Definition.
* A *Minkowski form* on a real vector space $`V`$ is a nondegenerate quadratic form that is isometric to the form $`x_1^2+x_2^2+\mathrm{}+x_{n+1}^2`$ on $`^{n+1}`$, where $`dimV=n+12`$.
* A *Lorentz metric* on a smooth manifold $`M`$ is a choice of Minkowski metric on the tangent space $`T_pM`$, for each $`pM`$, such that the form varies smoothly as $`p`$ varies.
A. Zeghib \[Ze1\] classified the compact homogeneous spaces that admit an invariant Lorentz metric. In this note, we remove the assumption of compactness, but add the restriction that the transitive group $`G`$ is almost simple. Our starting point is a special case of a theorem of N. Kowalsky.
###### 1.2 Theorem (N. Kowalsky, cf. \[Ko3, Thm. 5.1\]).
Let $`G/H`$ be a nontrivial homogeneous space of a connected, almost simple Lie group $`G`$ with finite center. If there is a $`G`$-invariant Lorentz metric on $`G/H`$, then either
1. there is also a $`G`$-invariant Riemannian metric on $`G/H`$; or
2. $`G`$ is locally isomorphic to either $`\mathrm{SO}(1,n)`$ or $`\mathrm{SO}(2,n)`$, for some $`n`$.
As explained in the following elementary proposition, it is easy to characterize the homogeneous spaces that arise in Conclusion (1) of Theorem 1.2, although it is probably not reasonable to expect a complete classification.
###### 1.3 Notation.
We use $`𝔤`$ to denote the Lie algebra of a Lie group $`G`$, and $`𝔥𝔤`$ to denote the Lie algebra of a Lie subgroup $`H`$ of $`G`$.
###### 1.4 Proposition (cf. \[Ko3, Thm. 1.1\]).
Let $`G/H`$ be a homogeneous space of a Lie group $`G`$, such that $`𝔤`$ is simple and $`dimG/H2`$. The following are equivalent.
1. The homogeneous space $`G/H`$ admits both a $`G`$-invariant Riemannian metric and a $`G`$-invariant Lorentz metric.
2. The closure of $`\mathrm{Ad}_GH`$ is compact, and leaves invariant a one-dimensional subspace of $`𝔤`$ that is not contained in $`𝔥`$.
The two main results of this note examine the cases that arise in Conclusion (2) of Theorem 1.2. It is well known \[Ko2, Egs. 2 and 3\] that $`\mathrm{SO}(1,n)^{}/\mathrm{SO}(1,n1)^{}`$ and $`\mathrm{SO}(2,n)^{}/\mathrm{SO}(1,n)^{}`$ have invariant Lorentz metrics. Also, for any discrete subgroup $`\mathrm{\Gamma }`$ of $`\mathrm{SO}(1,2)`$, the Killing form provides an invariant Lorentz metric on $`\mathrm{SO}(1,2)^{}/\mathrm{\Gamma }`$. We show that these are essentially the only examples.
Note that $`\mathrm{SO}(1,1)`$ and $`\mathrm{SO}(2,2)`$ fail to be almost simple. Thus, in 1.2(2), we may assume
* $`G`$ is locally isomorphic to $`\mathrm{SO}(1,n)`$, and $`n2`$; or
* $`G`$ is locally isomorphic to $`\mathrm{SO}(2,n)`$, and $`n3`$.
###### 2.3 Proposition.
Let $`G`$ be a Lie group that is locally isomorphic to $`\mathrm{SO}(1,n)`$, with $`n2`$. If $`H`$ is a closed subgroup of $`G`$, such that
* the closure of $`\mathrm{Ad}_GH`$ is not compact, and
* there is a $`G`$-invariant Lorentz metric on $`G/H`$,
then either
1. after any identification of $`𝔤`$ with $`𝔰𝔬(1,n)`$, the subalgebra $`𝔥`$ is conjugate to a standard copy of $`𝔰𝔬(1,n1)`$ in $`𝔰𝔬(1,n)`$, or
2. $`n=2`$ and $`H`$ is discrete.
###### 3.5 Theorem.
Let $`G`$ be a Lie group that is locally isomorphic to $`\mathrm{SO}(2,n)`$, with $`n3`$. If $`H`$ is a closed subgroup of $`G`$, such that
* the closure of $`\mathrm{Ad}_GH`$ is not compact, and
* there is a $`G`$-invariant Lorentz metric on $`G/H`$,
then, after any identification of $`𝔤`$ with $`𝔰𝔬(2,n)`$, the subalgebra $`𝔥`$ is conjugate to a standard copy of $`𝔰𝔬(1,n)`$ in $`𝔰𝔬(2,n)`$.
N. Kowalsky announced a much more general result than Theorem 3.5 in \[Ko2, Thm. 4\], but it seems that she did not publish a proof before her premature death. She announced a version of Proposition 2.3 (with much more general hypotheses and a somewhat weaker conclusion) in \[Ko2, Thm. 3\], and a proof appears in her Ph.D. thesis \[Ko1, Cor. 6.2\].
###### 1.5 Remark.
It is easy to see that there is a $`G`$-invariant Lorentz metric on $`G/H`$ if and only if there is an $`(\mathrm{Ad}_GH)`$-invariant Minkowski form on $`𝔤/𝔥`$. Thus, although Proposition 2.3 and Theorem 3.5 are geometric in nature, they can be restated in more algebraic terms. It is in such a form that they are proved in §2 and §3.
Proposition 2.3 and Theorem 3.5 are used in work of S. Adams \[Ad3\] on nontame actions on Lorentz manifolds. See \[Zi, Ko3, AS, Ze2, Ad1, Ad2\] for some other research concerning actions of Lie groups on Lorentz manifolds.
###### 1.6 Acknowledgments.
The author would like to thank the Isaac Newton Institute for Mathematical Sciences for providing the stimulating environment where this work was carried out. It is also a pleasure to thank Scot Adams for suggesting this problem and providing historical background. The research was partially supported by a grant from the National Science Foundation (DMS-9801136).
## 2. Homogeneous spaces of $`\mathrm{SO}(1,n)`$
The following lemma is elementary.
###### 2.1 Lemma.
Let $`\pi `$ be the standard representation of $`𝔤=𝔰𝔬(1,k)`$ on $`^{k+1}`$, and let $`𝔤=𝔨+𝔞+𝔫`$ be an Iwasawa decomposition of $`𝔤`$.
1. The representation $`\pi `$ has only one positive weight (with respect to $`𝔞`$), and the corresponding weight space is 1-dimensional.
2. There are subspaces $`V`$ and $`W`$ of $`^{k+1}`$, such that
1. $`dim(^{k+1}/V)=1`$;
2. $`dimW=1`$;
3. $`\pi (𝔫)VW`$;
4. for all nonzero $`u𝔫`$, we have $`\pi (u)^2^{k+1}=W`$; and
5. for all nonzero $`u𝔫`$ and $`v^{k+1}`$, we have $`\pi (u)^2v=0`$ if and only if $`vV`$.
###### 2.2 Corollary.
Let $`𝔥`$ be a subalgebra of a real Lie algebra $`𝔤`$, let $`Q`$ be a Minkowski form on $`𝔤/𝔥`$, and define $`\pi :N_G(𝔥)\mathrm{GL}(𝔤/𝔥)`$ by $`\pi (g)(v+𝔥)=(\mathrm{Ad}_Gg)v+𝔥`$.
1. Suppose $`T`$ is a connected Lie subgroup of $`G`$ that normalizes $`H`$, such that $`\pi (T)\mathrm{SO}(Q)`$ and $`\mathrm{Ad}_GT`$ is diagonalizable over $``$. Then, for any ordering of the $`T`$-weights on $`𝔤`$, the subalgebra $`𝔥`$ contains codimension-one subspaces of both $`𝔤^+`$ and $`𝔤^{}`$, where $`𝔤^+`$ is the sum of all the positive weight spaces of $`T`$, and $`𝔤^{}`$ is the sum of all the negative weight spaces of $`T`$.
2. If $`U`$ is a connected Lie subgroup of $`G`$ that normalizes $`H`$, such that $`\pi (U)\mathrm{SO}(Q)`$ and $`\mathrm{Ad}_GU`$ is unipotent, then there are subspaces $`V/𝔥`$ and $`W/𝔥`$ of $`𝔤/𝔥`$, such that
1. $`dim(𝔤/V)=1`$;
2. $`dim(W/𝔥)=1`$;
3. $`[V,𝔲]W`$;
4. for each $`u𝔲`$, either $`W=𝔥+(\mathrm{ad}_𝔤u)^2𝔤`$, or $`[𝔤,u]𝔥`$; and
5. for all $`u𝔲`$, we have $`(\mathrm{ad}_𝔤u)^2V𝔥`$.
###### 2.3 Proposition.
Let $`H`$ be a Lie subgroup of $`G=\mathrm{SO}(1,n)`$, with $`n2`$, such that
* the closure of $`H`$ is not compact; and
* there is an $`(\mathrm{Ad}_GH)`$-invariant Minkowski form on $`𝔤/𝔥`$.
Then either
1. $`H^{}`$ is conjugate to a standard copy of $`\mathrm{SO}(1,n1)^{}`$ in $`\mathrm{SO}(1,n)`$, or
2. $`n=2`$ and $`H^{}`$ is trivial.
###### Proof.
Let $`\overline{H}`$ be the Zariski closure of $`H`$, and note that the Minkowski form is also invariant under $`\mathrm{Ad}_G\overline{H}`$. Replacing $`H`$ by a finite-index subgroup, we may assume $`\overline{H}`$ is Zariski connected.
Let $`G=KAN`$ be an Iwasawa decomposition of $`G`$.
Case 1 . Assume $`n3`$ and $`A\overline{H}`$. From Corollary 2.2(1), we see that $`𝔥`$ contains codimension-one subspaces of both $`𝔫`$ and $`𝔫^{}`$. (Note that this implies $`H^{}`$ is nontrivial.) This implies that $`\overline{H}`$ is reductive. (Because $`(HN)^{}\mathrm{unip}\overline{H}`$ is a unipotent subgroup that intersects $`N`$ nontrivially (and $`\text{-rank}G=1`$), it must be contained in $`N`$, so $`\mathrm{unip}\overline{H}N`$. Similarly, $`\mathrm{unip}\overline{H}N^{}`$. Therefore $`\mathrm{unip}\overline{H}NN^{}=e`$.) Then, since $`\overline{H}`$ contains a codimension-one subgroup of $`N`$, and since $`A\overline{H}`$, it follows that $`\overline{H}`$ is conjugate to either $`\mathrm{SO}(1,n1)`$ or $`\mathrm{SO}(1,n)`$. Because $`H^{}`$ is a nontrivial, connected, normal subgroup of $`\overline{H}`$, we conclude that $`H^{}`$ is conjugate to either $`\mathrm{SO}(1,n1)^{}`$ or $`\mathrm{SO}(1,n)^{}`$. Because $`𝔤/𝔥0`$ (else $`dim𝔤/𝔥=0<2`$, which contradicts the fact that there is a Minkowski form on $`𝔤/𝔥`$), we see that $`H^{}`$ is conjugate to $`\mathrm{SO}(1,n1)^{}`$.
Case 2 . Assume $`n3`$ and $`\overline{H}`$ does not contain any nontrivial hyperbolic elements. The Levi subgroup of $`\overline{H}`$ must be compact, and the radical of $`\overline{H}`$ must be unipotent, so choose a compact $`M`$ and a nontrivial unipotent subgroup $`U`$ such that $`\overline{H}=MU`$. Replacing $`H`$ by a conjugate, we may assume, without loss of generality, that $`UN`$.
Let us show, for every nonzero $`u𝔲`$, that $`[𝔤,u]𝔥`$. From the Morosov Lemma \[Ja, Thm. 17(1), p. 100\], we know there exists $`v𝔤`$, such that $`[v,u]`$ is hyperbolic (and nonzero). If $`[v,u]𝔥`$, this contradicts the fact that $`\overline{H}`$ does not contain nontrivial hyperbolic elements.
Let $`V/𝔥`$ and $`W/𝔥`$ be subspaces of $`𝔤/𝔥`$ as in Corollary 2.2(2). Because $`(\mathrm{ad}_𝔤u)^2𝔤=𝔫`$ for every nonzero $`u𝔫`$, we have $`W=𝔫+𝔥`$ (see 2.2(2d)), so $`dim𝔫/(𝔥𝔫)=1`$ (see 2.2(2b)) and
(2.4)
$$[𝔲,V]W=𝔫+𝔥𝔫+\overline{𝔥}=𝔫+𝔪$$
(see 2.2(2c)).
Assume, for the moment, that $`n4`$. Then
$`dim𝔲+dim(V𝔫^{})`$ $``$ $`dim(𝔥𝔫)+dim(V𝔫^{})(dim𝔫1)+(dim𝔫^{}1)`$
$`=`$ $`(n2)+(n2)n>dim𝔫.`$
This implies that there exist $`u𝔲`$ and $`vV𝔫^{}`$, such that $`u,v𝔰𝔩(2,)`$, with $`[u,v]`$ hyperbolic (and nonzero). This contradicts the fact that $`𝔪+𝔫`$ has no nontrivial hyperbolic elements.
We may now assume that $`n=3`$. For any nonzero $`u𝔫`$, we have
$$dim[u,V]dim[u,𝔤]1=dim𝔫+1>dim𝔫,$$
so $`[𝔲,V]𝔫`$. Then, from (2.4), we conclude that $`𝔪0`$, so $`𝔪`$ acts irreducibly on $`𝔫`$. This contradicts the fact that $`𝔥𝔫`$ is a codimension-one subspace of $`𝔫`$ that is normalized by $`𝔪`$.
Case 3 . Assume $`n=2`$. We may assume $`H^{}`$ is nontrivial (otherwise Conclusion (2) holds). We must have $`dim𝔤/𝔥2`$, so we conclude that $`dimH^{}=1`$ and $`dim𝔤/𝔥=2`$. Because $`\mathrm{SO}(1,1)`$ consists of hyperbolic elements, this implies that $`\mathrm{Ad}_Gh`$ acts diagonalizably on $`𝔤/𝔥`$, for every $`hH`$. Therefore $`H^{}`$ is conjugate to $`A`$, and, hence, to $`\mathrm{SO}(1,1)^{}`$. ∎
## 3. Homogeneous spaces of $`\mathrm{SO}(2,n)`$
###### 3.1 Theorem (Borel-Tits \[BT2, Prop. 3.1\]).
Let $`H`$ be an $`F`$-subgroup of a reductive algebraic group $`G`$ over a field $`F`$ of characteristic zero. Then there is a parabolic $`F`$-subgroup $`P`$ of $`G`$, such that $`\mathrm{unip}H\mathrm{unip}P`$ and $`HN_G(\mathrm{unip}H)P`$.
###### 3.2 Notation.
Let $`k=n/2`$. Identifying $`^{k+1}`$ with $`^{2k+2}`$ yields an embedding of $`\mathrm{SU}(1,k)`$ in $`\mathrm{SO}(2,2k)`$. Then the inclusion $`^{2k+2}^{2k+3}`$ yields an embedding of $`\mathrm{SU}(1,k)`$ in $`\mathrm{SO}(2,2k+1)`$. Thus, we may identify $`\mathrm{SU}(1,n/2)`$ with a subgroup of $`\mathrm{SO}(2,n)`$.
We use the following well-known result to shorten one case of the proof of Theorem 3.5.
###### 3.3 Lemma (\[OW, Lem. 6.8\]).
If $`L`$ is a connected, almost-simple subgroup of $`\mathrm{SO}(2,n)`$, such that $`\text{-rank}L=1`$ and $`dimL>3`$, then $`L`$ is conjugate under $`\mathrm{O}(2,n)`$ to a subgroup of either $`\mathrm{SO}(1,n)`$ or $`\mathrm{SU}(1,n/2)`$.
###### 3.4 Corollary.
Let $`L`$ be a connected, reductive subgroup of $`G=\mathrm{SO}(2,n)`$, such that $`\text{-rank}L=1`$. Then $`dimUn1`$, for every connected, unipotent subgroup $`U`$ of $`L`$.
Furthermore, if $`dimU=n1`$, then either
1. $`L`$ is conjugate to $`\mathrm{SO}(1,n)^{}`$; or
2. $`n`$ is even, and $`L`$ is conjugate under $`\mathrm{O}(2,n)`$ to $`\mathrm{SU}(1,n/2)`$.
###### 3.5 Theorem.
Let $`H`$ be a Lie subgroup of $`G=\mathrm{SO}(2,n)`$, with $`n3`$, such that
* the closure of $`H`$ is not compact, and
* there is an $`(\mathrm{Ad}_GH)`$-invariant Minkowski form on $`𝔤/𝔥`$.
Then $`H^{}`$ is conjugate to a standard copy of $`\mathrm{SO}(1,n)^{}`$ in $`\mathrm{SO}(2,n)`$.
###### Proof.
Let $`\overline{H}`$ be the Zariski closure of $`H`$, and note that the Minkowski form is also invariant under $`\mathrm{Ad}_G\overline{H}`$. Replacing $`H`$ by a finite-index subgroup, we may assume $`\overline{H}`$ is Zariski connected.
Let $`G=KAN`$ be an Iwasawa decomposition of $`G`$. For each real root $`\varphi `$ of $`𝔤`$ (with respect to the Cartan subalgebra $`𝔞`$), let $`𝔤_\varphi `$ be the corresponding root space, and let $`\mathrm{proj}_\varphi :𝔤𝔤_\varphi `$ and $`\mathrm{proj}_{\varphi \varphi }:𝔤𝔤_\varphi +𝔤_\varphi `$ be the natural projections. Fix a choice of simple real roots $`\alpha `$ and $`\beta `$ of $`𝔤`$, such that $`dim𝔤_\alpha =1`$ and $`dim𝔤_\beta =n2`$ (so the positive real roots are $`\alpha `$, $`\beta `$, $`\alpha +\beta `$, and $`\alpha +2\beta `$). Replacing $`N`$ by a conjugate under the Weyl group, we may assume $`𝔫=𝔤_\alpha +𝔤_\beta +𝔤_{\alpha +\beta }+𝔤_{\alpha +2\beta }`$. From the classification of parabolic subgroups \[BT1, Prop. 5.14, p. 99\], we know that the only proper parabolic subalgebras of $`𝔤`$ that contain $`𝔫_𝔤(𝔫)`$ are
(3.6) $`𝔫_𝔤(𝔫)`$, $`𝔭_\alpha =𝔫_𝔤(𝔫)+𝔤_\alpha `$, and $`𝔭_\beta =𝔫_𝔤(𝔫)+𝔤_\beta `$.
Case 1 . Assume $`\overline{𝔥}`$ contains nontrivial hyperbolic elements. Let $`𝔱=\overline{𝔥}𝔞`$. Replacing $`H`$ by a conjugate, we may assume $`𝔱0`$.
Subcase 1.1 . Assume $`𝔱\{\mathrm{ker}(\alpha +\beta ),\mathrm{ker}\beta \}`$.
Subsubcase 1.1.1 . Assume $`\overline{H}`$ is reductive. We may assume $`𝔱=\mathrm{ker}(\alpha +\beta )`$ (if necessary, replace $`H`$ with its conjugate under the Weyl reflection corresponding to the root $`\alpha `$). Then, from Corollary 2.2(1), we see that $`𝔥`$ contains a codimension-one subspace of $`𝔤_{\alpha +2\beta }+𝔤_\beta +𝔤_\alpha `$. (Note that this implies $`H^{}`$ is nontrivial.)
Let $`𝔫^{}=𝔤_{\alpha +\beta }+𝔤_{\alpha +2\beta }+𝔤_\beta +𝔤_\alpha `$, so $`𝔫^{}`$ is the Lie algebra of a maximal unipotent subgroup of $`G`$. (In fact, $`𝔫^{}`$ is the image of $`𝔫`$ under the Weyl reflection corresponding to the root $`\alpha `$.) From the preceding paragraph, we know that
$$dim(\overline{𝔥}𝔫^{})dim(𝔤_{\alpha +2\beta }+𝔤_\beta +𝔤_\alpha )1=n1.$$
Therefore, Corollary 3.4 implies that $`\overline{H}`$ is conjugate (under $`\mathrm{O}(2,n)`$) to either $`\mathrm{SO}(1,n)`$ or $`\mathrm{SU}(1,n/2)`$. It is easy to see that $`\overline{H}`$ is not conjugate to $`\mathrm{SU}(1,n/2)`$. (See \[OW, proof of Thm. 1.5\] for an explicit description of $`𝔰𝔲(1,n/2)𝔫`$. If $`n`$ is even, then $`n>3`$, so $`𝔰𝔲(1,n/2)`$ does not contain a codimension-one subspace of any $`(n2)`$-dimensional root space, but $`\overline{𝔥}`$ does contain a codimension-one subspace of $`𝔤_\beta `$.) Therefore, we conclude that $`\overline{H}`$ is conjugate to $`\mathrm{SO}(1,n)`$. Then, because $`H^{}`$ is a nontrivial, connected, normal subgroup of $`\overline{H}`$, we conclude that $`H^{}=(\overline{H})^{}`$ is conjugate to $`\mathrm{SO}(1,n)^{}`$.
Subsubcase 1.1.2 . Assume $`\overline{H}`$ is not reductive. Let $`P`$ be a maximal parabolic subgroup of $`G`$ that contains $`\overline{H}`$ (see Theorem 3.1). By replacing $`P`$ and $`H`$ with conjugate subgroups, we may assume that $`P`$ contains the minimal parabolic subgroup $`N_G(N)`$. Therefore, the classification of parabolic subalgebras (3.6) implies that $`P`$ is either $`P_\alpha `$ or $`P_\beta `$.
Subsubsubcase 1.1.2.1 . Assume $`𝔱=\mathrm{ker}(\alpha +\beta )`$. From Corollary 2.2(1), we see that $`𝔥`$ (and hence also $`𝔭`$) contains codimension-one subspaces of $`𝔤_{\alpha +2\beta }+𝔤_\beta +𝔤_\alpha `$ and $`𝔤_{\alpha 2\beta }+𝔤_\beta +𝔤_\alpha `$. Because $`𝔭_\alpha `$ does not contain such a subspace of $`𝔤_{\alpha 2\beta }+𝔤_\beta +𝔤_\alpha `$, we conclude that $`P=P_\beta `$. Furthermore, because the intersection of $`𝔭_\beta `$ with each of these subspaces does have codimension one, we conclude that $`𝔥`$ has precisely the same intersection; therefore $`(𝔤_{\alpha +2\beta }+𝔤_\beta )+(𝔤_\beta +𝔤_\alpha )𝔥`$. Hence $`𝔥[𝔤_\alpha ,𝔤_\beta ]=𝔤_{\alpha +\beta }`$. We now have
$$(\mathrm{ad}_𝔤𝔤_{\alpha +\beta })^2𝔤=𝔤_\alpha +𝔤_{\alpha +\beta }+𝔤_{\alpha +2\beta }0(mod𝔥),$$
so Corollary 2.2(2d) implies
$$𝔥[𝔤,𝔤_{\alpha +\beta }][𝔤_{\alpha \beta },𝔤_{\alpha +\beta }]\mathrm{ker}\beta .$$
This contradicts the fact that $`\overline{𝔥}𝔞=𝔱=\mathrm{ker}(\alpha +\beta )`$.
Subsubsubcase 1.1.2.2 . Assume $`𝔱=\mathrm{ker}\beta `$. From Corollary 2.2(1), we see that $`𝔥`$ (and hence also $`𝔭`$) contains a codimension-one subspace of $`𝔤_\alpha +𝔤_{\alpha \beta }+𝔤_{\alpha 2\beta }`$. Because neither $`𝔭_\alpha `$ nor $`𝔭_\beta `$ contains such a subspace, this is a contradiction.
Subcase 1.2 . Assume $`𝔱\{\mathrm{ker}\alpha ,\mathrm{ker}(\alpha +2\beta )\}`$. We may assume $`𝔱=\mathrm{ker}\alpha `$ (if necessary, replace $`H`$ with its conjugate under the Weyl reflection corresponding to the root $`\beta `$). From Corollary 2.2(1), we see that $`𝔥`$ contains a codimension-one subspace of $`𝔤_\beta +𝔤_{\alpha +\beta }+𝔤_{\alpha +2\beta }`$. Because any codimension-one subalgebra of a nilpotent Lie algebra must contain the commutator subalgebra, we conclude that $`𝔥`$ contains $`𝔤_{\alpha +2\beta }`$. Then we have
$$(\mathrm{ad}_𝔤𝔤_{\alpha +2\beta })^2𝔤=𝔤_{\alpha +2\beta }0(mod𝔥),$$
so Corollary 2.2(2d) implies
$$𝔥[𝔤,𝔤_{\alpha +2\beta }]𝔤_\beta +𝔤_{\alpha +\beta }+𝔤_{\alpha +2\beta }.$$
Similarly, we also have $`𝔥𝔤_\beta +𝔤_{\alpha \beta }+𝔤_{\alpha 2\beta }`$. It is now easy to show that $`𝔥𝔤_\varphi `$ for every real root $`\varphi `$, so $`𝔥=𝔤`$. This contradicts the fact that $`𝔤/𝔥0`$.
Subcase 1.3 . Assume $`𝔱`$ contains a regular element of $`𝔞`$. Replacing $`H`$ by a conjugate under the Weyl group, we may assume that $`𝔫`$ is the sum of the positive root spaces, with respect to $`𝔱`$. Then, from Corollary 2.2(1), we see that $`𝔥`$ contains codimension-one subspaces of both $`𝔫`$ and $`𝔫^{}`$. Therefore, $`𝔥`$ contains codimension-one subspaces of $`𝔤_\beta +𝔤_{\alpha +\beta }+𝔤_{\alpha +2\beta }`$ and $`𝔤_\beta +𝔤_{\alpha \beta }+𝔤_{\alpha 2\beta }`$, so the argument of Subcase 3 applies.
Case 2 . Assume $`\overline{𝔥}`$ does not contain nontrivial hyperbolic elements. The Levi subgroup of $`\overline{H}`$ must be compact, and the radical of $`\overline{H}`$ must be unipotent, so choose a compact $`M`$ and a nontrivial unipotent subgroup $`U`$ such that $`\overline{H}=MU`$. Choose subspaces $`V/𝔥`$ and $`W/𝔥`$ of $`𝔤/𝔥`$ as in Corollary 2.2(2).
Let $`P`$ be a proper parabolic subgroup of $`G`$, such that $`U\mathrm{unip}P`$ and $`HP`$ (see Theorem 3.1). Replacing $`H`$ and $`P`$ by conjugates, we may assume, without loss of generality, that $`P`$ contains the minimal parabolic subgroup $`N_G(N)`$ (so $`\mathrm{unip}PN`$). From the classification of parabolic subalgebras (3.6), we know that there are only three possibilities for $`P`$. We consider each of these possibilities separately.
First, though, let us show that
(3.7) for every nonzero $`u𝔲`$, we have $`[𝔤,u]𝔥`$.
From the Morosov Lemma \[Ja, Thm. 17(1), p. 100\], we know there exists $`v𝔤`$, such that $`[v,u]`$ is hyperbolic (and nonzero). If $`[v,u]𝔥`$, this contradicts the fact that $`\overline{𝔥}`$ does not contain nontrivial hyperbolic elements.
Subcase 2.1 . Assume $`P=N_G(N)`$ is a minimal parabolic subgroup of $`G`$.
Subsubcase 2.1.1 . Assume $`\mathrm{proj}_\beta 𝔲0`$. Choose $`u𝔲`$, such that $`\mathrm{proj}_\beta u0`$, and let $`Z=(\mathrm{ad}_𝔤u)^2𝔤_{\alpha 2\beta }`$. (So $`dimZ=1`$, $`\mathrm{proj}_\alpha Z0`$, and $`\mathrm{proj}_{\alpha \beta }Z=0`$.) From Corollary 2.2(2d), we know that $`ZW`$. Then, because $`\mathrm{proj}_\alpha 𝔥\mathrm{proj}_\alpha 𝔭=0`$, we conclude, from Corollary 2.2(2b), that $`W=𝔥+Z`$.
Because $`W=𝔥+Z𝔭+Z`$, we have $`\mathrm{proj}_{\alpha \beta }W=0`$. Therefore, because $`\mathrm{proj}_\beta u0`$, we conclude, from Corollary 2.2(2c), that $`\mathrm{proj}_{\alpha 2\beta }V=0`$, so Corollary 2.2(2a) implies that $`V=\mathrm{ker}(\mathrm{proj}_{\alpha 2\beta })`$. In particular, we have $`𝔤_\beta V`$, so Corollary 2.2(2c) implies $`[𝔤_\beta ,u]W`$. Therefore, we have
$`[𝔤_\beta ,\mathrm{proj}_\beta u]`$ $``$ $`[𝔤_\beta ,u+(𝔤_\alpha +𝔤_{\alpha +\beta }+𝔤_{\alpha +2\beta })]=[𝔤_\beta ,u]+[𝔤_\beta ,𝔤_\alpha +𝔤_{\alpha +\beta }+𝔤_{\alpha +2\beta }]`$
$``$ $`W+(𝔤_\alpha +𝔤_{\alpha +\beta })=𝔥+Z+(𝔤_\alpha +𝔤_{\alpha +\beta })𝔪+𝔫+Z.`$
Because $`\mathrm{proj}_\alpha [𝔤_\beta ,\mathrm{proj}_\beta u]=0`$, we conclude that $`[𝔤_\beta ,\mathrm{proj}_\beta u]𝔪+𝔫`$. This contradicts the fact that $`𝔪+𝔫`$ does not contain nontrivial hyperbolic elements.
Subsubcase 2.1.2 . Assume $`\mathrm{proj}_\beta 𝔲=0`$. Replacing $`H`$ by a conjugate under $`N`$, we may assume $`𝔪𝔤_0`$, so $`\mathrm{proj}_\beta \overline{𝔥}=0`$.
We have $`𝔲𝔤_\alpha +𝔤_{\alpha +\beta }+𝔤_{\alpha +2\beta }`$, so $`(\mathrm{ad}_𝔤u)^2𝔤𝔤_\alpha +𝔤_{\alpha +\beta }+𝔤_{\alpha +2\beta }`$ for every $`u𝔲`$. Thus, Corollary 2.2(2d) implies $`W(𝔤_\alpha +𝔤_{\alpha +\beta }+𝔤_{\alpha +2\beta })+𝔥`$.
We have
$$\mathrm{proj}_{\beta \beta }W\mathrm{proj}_{\beta \beta }(𝔤_\alpha +𝔤_{\alpha +\beta }+𝔤_{\alpha +2\beta })+\mathrm{proj}_{\beta \beta }𝔥=0,$$
so Corollary 2.2(2c) implies that $`\mathrm{proj}_{\beta \beta }\left((\mathrm{ad}_𝔤𝔲)V\right)=0`$.
Subsubsubcase 2.1.2.1 . Assume $`\mathrm{proj}_\alpha u0`$, for some $`u𝔲`$. From the conclusion of the preceding paragraph, we know that $`\mathrm{proj}_\beta \left((\mathrm{ad}_𝔤u)V\right)=0`$. Because $`\mathrm{proj}_\beta u=0`$ and $`\mathrm{proj}_\alpha 0`$, this implies $`\mathrm{proj}_{\alpha \beta }V=0`$, so $`V=\mathrm{ker}(\mathrm{proj}_{\alpha \beta })`$ (see 2.2(2a)). In particular, $`𝔤_\alpha V`$, so Corollary 2.2(2c) implies
$`[𝔤_\alpha ,𝔤_\alpha ]`$ $``$ $`[𝔲+(𝔤_{\alpha +\beta }+𝔤_{\alpha +2\beta }),𝔤_\alpha ][𝔲,V]+[𝔤_{\alpha +\beta }+𝔤_{\alpha +2\beta },𝔤_\alpha ]`$
$``$ $`W+𝔤_\beta 𝔥+𝔫𝔪+𝔫.`$
This contradicts the fact that $`𝔪+𝔫`$ does not contain nontrivial hyperbolic elements.
Subsubsubcase 2.1.2.2 . Assume $`\mathrm{proj}_{\alpha +\beta }u0`$, for some $`u𝔲`$. From Subsubsubcase 3, we may assume $`\mathrm{proj}_\alpha u=0`$. Because $`0=\mathrm{proj}_{\beta \beta }\left((\mathrm{ad}_𝔤u)V\right)`$ has codimension $`1`$ in $`\mathrm{proj}_{\beta \beta }\left((\mathrm{ad}_𝔤u)𝔤\right)`$ (see 2.2(2a)), which contains the 2-dimensional subspace $`\mathrm{proj}_{\beta \beta }\left([u,𝔤_{\alpha 2\beta }+𝔤_\alpha ]\right)`$, we have a contradiction.
Subsubsubcase 2.1.2.3 . Assume $`𝔲=𝔤_{\alpha +2\beta }`$. (This argument is similar to Subsubsubcase 3.) Because $`\mathrm{proj}_\beta \left((\mathrm{ad}_𝔤𝔲)V\right)=0`$, we know that $`\mathrm{proj}_{\alpha \beta }V=0`$, so $`V=\mathrm{ker}(\mathrm{proj}_{\alpha \beta })`$ (see 2.2(2a)). In particular, $`𝔤_{\alpha 2\beta }V`$, so Corollary 2.2(2c) implies
$$[𝔤_{\alpha +2\beta },𝔤_{\alpha 2\beta }][𝔲,V]W𝔥+𝔫𝔪+𝔫.$$
This contradicts the fact that $`𝔪+𝔫`$ does not contain nontrivial hyperbolic elements.
Subcase 2.2 . Assume $`P=P_\alpha `$. We may assume there exists $`x𝔥`$, such that $`\mathrm{proj}_\alpha x0`$ (otherwise, $`HN_G(N)`$, so Subcase 3 applies). Note that, because $`U\mathrm{unip}P`$, we have $`\mathrm{proj}_\alpha 𝔲=0`$.
Subsubcase 2.2.1 . Assume $`\mathrm{proj}_{\alpha +\beta }𝔲0`$. Choose $`u𝔲`$, such that $`\mathrm{proj}_{\alpha +\beta }u0`$. Then $`[x,u][𝔥,𝔲]𝔲`$, and $`[[x,u],u]`$ is a nonzero element of $`𝔤_{\alpha +2\beta }`$, so we see that $`𝔤_{\alpha +2\beta }[𝔲,𝔲]`$. Because every unipotent subgroup of $`\mathrm{SO}(1,k)`$ is abelian, we conclude that $`\mathrm{ad}_𝔤𝔤_{\alpha +2\beta }`$ acts trivially on $`𝔤/𝔥`$, which means $`𝔥[𝔤,𝔤_{\alpha +2\beta }]`$. This contradicts (3.7).
Subsubcase 2.2.2 . Assume $`\mathrm{proj}_{\alpha +\beta }𝔲=0`$. We may assume, furthermore, that $`\mathrm{proj}_\alpha 𝔥0`$ (otherwise, by replacing $`H`$ with its conjugate under the Weyl reflection corresponding to the root $`\alpha `$, we could revert to Subcase 3). Then, because $`[𝔥,𝔲]𝔲`$, we must have $`\mathrm{proj}_\beta 𝔲=0`$. Thus, $`𝔲=𝔤_{\alpha +2\beta }`$. From Corollary 2.2(2d), we have
$$W=[𝔤,𝔤_{\alpha +2\beta },𝔤_{\alpha +2\beta }]+𝔥=𝔤_{\alpha +2\beta }+𝔥𝔲+\overline{𝔥}=\overline{𝔥},$$
so
$`W(𝔤_\beta +𝔤_{\alpha +\beta })`$ $``$ $`\overline{𝔥}(𝔤_\beta +𝔤_{\alpha +\beta })=(\overline{𝔥}𝔫)(𝔤_\beta +𝔤_{\alpha +\beta })`$
$`=`$ $`𝔲(𝔤_\beta +𝔤_{\alpha +\beta })=𝔤_{\alpha +2\beta }(𝔤_\beta +𝔤_{\alpha +\beta })=0.`$
On the other hand, from Corollary 2.2(2c), we know that $`W`$ contains a codimension-one subspace of $`[𝔤,𝔤_{\alpha +2\beta }]`$, so $`W`$ contains a codimension-one subspace of $`𝔤_\beta +𝔤_{\alpha +\beta }`$. This is a contradiction.
Subcase 2.3 . Assume $`P=P_\beta `$. Note that, because $`U\mathrm{unip}P`$, we have $`\mathrm{proj}_\beta 𝔲=0`$.
From Corollary 2.2(2d), we have
$`W`$ $`=`$ $`𝔥+(\mathrm{ad}_𝔤u)^2𝔤𝔥+(𝔤_\alpha +𝔤_{\alpha +\beta }+𝔤_{\alpha +2\beta })`$
$`=`$ $`𝔥+\mathrm{unip}𝔭_\beta (𝔪+𝔲)+\mathrm{unip}𝔭_\beta =𝔪+\mathrm{unip}𝔭_\beta .`$
Subsubcase 2.3.1 . Assume there is some nonzero $`u𝔲`$, such that $`\mathrm{proj}_\alpha u=0`$. Replacing $`H`$ by a conjugate (under $`G_\beta `$), we may assume $`\mathrm{proj}_{\alpha +\beta }u0`$.
Let $`V^{}=V(𝔤_\alpha +𝔤_{\alpha \beta })`$. Because $`V^{}`$ contains a codimension-one subspace of $`𝔤_\alpha +𝔤_{\alpha \beta }`$ (see Corollary 2.2(2a)), one of the following two subsubsubcases must apply.
Subsubsubcase 2.3.1.1 . Assume there exists $`vV^{}`$, such that $`\mathrm{proj}_{\alpha \beta }v=0`$. From Corollary 2.2(2c), we have $`[u,v]W`$. Then, because $`[u,v]`$ is a nonzero element of $`𝔤_\beta `$, we conclude that
$$0W𝔤_\beta (𝔪+\mathrm{unip}𝔭_\beta )𝔤_\beta =0.$$
This contradicts the fact that $`M`$, being compact, has no nontrivial unipotent elements.
Subsubsubcase 2.3.1.2 . Assume $`\mathrm{proj}_{\alpha \beta }V^{}=𝔤_{\alpha \beta }`$. For $`vV^{}`$, we have $`\mathrm{proj}_0[u,v]=[\mathrm{proj}_{\alpha +\beta }u,\mathrm{proj}_{\alpha \beta }v]`$. Thus, there is some $`vV^{}`$, such that $`\mathrm{proj}_0[u,v]`$ is hyperbolic (and nonzero). On the other hand, from Corollary 2.2(2c), we have $`[u,v]W=𝔪+\mathrm{unip}𝔭_\beta `$. This contradicts the fact that $`𝔪\overline{𝔥}`$ does not contain nonzero hyperbolic elements.
Subsubcase 2.3.2 . Assume $`\mathrm{proj}_\alpha u0`$, for every nonzero $`u𝔲`$. Fix some nonzero $`u𝔲`$. Because $`dim𝔲_\alpha =1`$, we must have $`dim𝔲=1`$ (so $`𝔲=u`$). Replacing $`H`$ by a conjugate (under $`G_\beta `$), we may assume $`\mathrm{proj}_{\alpha +\beta }u=0`$. Also, we may assume $`\mathrm{proj}_{\alpha +2\beta }u0`$ (otherwise, we could revert to Subsubcase 3 by replacing $`H`$ with its conjugate under the Weyl reflection corresponding to the root $`\beta `$).
Let $`𝔱=[u,𝔤_\alpha +𝔤_{\alpha 2\beta }]`$. Because $`𝔤_\alpha ,𝔤_\alpha `$ and $`𝔤_{\alpha +2\beta },𝔤_{\alpha 2\beta }`$ centralize each other, we see that $`𝔱=[𝔤_\alpha ,𝔤_\alpha ]+[𝔤_{\alpha +2\beta },𝔤_{\alpha 2\beta }]`$ is a two-dimensional subspace of $`𝔤`$ consisting entirely of hyperbolic elements. Because $`V`$ contains a codimension-one subspace of $`𝔤_\alpha +𝔤_{\alpha 2\beta }`$ (see Corollary 2.2(2a)), and $`[u,V]W`$ (see Corollary 2.2(2c)), we see that $`W`$ contains a codimension-one subspace of $`𝔱`$, so $`W`$ contains nontrivial hyperbolic elements. This contradicts the fact that $`W𝔪+\mathrm{unip}𝔭_\beta `$ does not contain nontrivial hyperbolic elements. ∎
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# Equilibrium Properties of Temporally Asymmetric Hebbian Plasticity
## Abstract
A theory of temporally asymmetric Hebb (TAH) rules which depress or potentiate synapses depending upon whether the postsynaptic cell fires before or after the presynaptic one is presented. Using the Fokker-Planck formalism, we show that the equilibrium synaptic distribution induced by such rules is highly sensitive to the manner in which bounds on the allowed range of synaptic values are imposed. In a biologically plausible multiplicative model, we find that the synapses in asynchronous networks reach a distribution that is invariant to the firing rates of either the pre- or post-synaptic cells. When these cells are temporally correlated, the synaptic strength varies smoothly with the degree and phase of synchrony between the cells.
Recent experimental evidence indicates that synaptic modification in cortical neurons depends on the precise temporal relation between pre- and postsynaptic firing . Presynaptic spikes that precede postsynaptic firing lead to synaptic potentiation, while those that follow postsynaptic firing elicit synaptic depression. The temporal window for inducing these changes is on the order of 10 msec. Several recent theoretical studies addressed the potential implications of this temporally asymmetric Hebbian (TAH) synaptic plasticity on learning . The present study is motivated by recent work by Abbott and coworkers who applied TAH learning in a large population of excitatory presynaptic cells asynchronously driving a single postsynaptic cell . Their simulations showed that the distribution of synapses converged to a bimodal distribution. The synapses were either almost zero or had values close to their upper limit. Moreover, when the firing rate of the presynaptic cells was increased, the number of strong synapses decreased so that there was very little change in the output rate. Thus, the TAH rule seems to provide a mechanism for keeping the mean output rate invariant. Since Hebb rules are presumed to underlie many developmental and learning processes in neuronal systems, it is important to understand the equilibrium properties of networks with TAH plasticity and how they depend upon the particular implementation of these rules.
In this Letter, we study the TAH rule using Fokker-Planck theory . Surprisingly, we find that the behavior of the system depends crucially on how the boundaries on the allowed range of synaptic efficacies are incorporated. In particular, the salient features found in are unique to an additive learning rule in which the magnitude of the update does not explicitly depend on the current value of the synapse. A very different behavior is found with a multiplicative rule where the magnitude of the update decreases as either the upper or lower bounds are approached.
TAH plasticity is described as a change to the synaptic efficacy $`w`$ between two cells. A pair of spikes in the input cell and the output cell, at times $`t_i`$ and $`t_o`$, respectively, induces a change in $`w`$:
$$\mathrm{\Delta }^\pm w=\pm \lambda f_\pm (w)K(|t_ot_i|).$$
(1)
The weight $`w`$ is increased by $`\mathrm{\Delta }^+w`$ when $`t_o>t_i`$ and decreased by $`\mathrm{\Delta }^{}w`$ when $`t_i>t_o`$. The temporal dependence of the update is defined by the filter $`K`$ which for simplicity is taken to be $`K(t)\mathrm{exp}(t/\tau )`$. The coefficient $`\lambda `$ sets the scale of the synaptic change at each update. The factors $`f_\pm (w)`$ determine the relative magnitude of the changes in the positive and negative directions.
We consider two particular examples of these update rules. The first is an additive update rule where the magnitude of the changes is independent of $`w`$, so that:
$$f_+=1;f_{}=\alpha .$$
(2)
The parameter $`\alpha >0`$ denotes a possible asymmetry between increasing and decreasing the synaptic efficacy. If the update results in a synaptic weight outside the bounds $`0<w<1`$, the weight is clipped to the boundary values. For the second example, which we will call the multiplicative rule:
$$f_+(w)=1w;f_{}(w)=\alpha w.$$
(3)
This results in a synaptic increase (decrease) whose magnitude scales linearly with the distance to the upper (lower) boundary, similar to the model in .
To evaluate the equilibrium properties of these rules, the firing activity in the two cells needs to be specified. We consider the case when the input and output activity are stationary stochastic processes. The firing of the input cell is characterized by an instantaneous rate function $`\nu _i(t)=_{t_i}\delta (tt_i)`$, where $`t_i`$ are the spike times of the input cell with mean rate $`\nu _i=r_{in}`$. Similarly, the activity of the output cell is given by $`\nu _o(t)=_{t_o}\delta (tt_o)`$ with mean rate $`r_{out}`$. The correlation between these two spike trains is described by the normalized time delayed crosscorrelation function $`C(t)=\nu _i(t^{})\nu _o(t^{}+t)/r_{in}r_{out}1`$. Note that for uncorrelated spike trains, $`C(t)=0`$.
In the limit of small step sizes ($`\lambda <<1`$), Eq. (1) can be averaged in order to describe the behavior of $`w`$ on times of order $`1/\lambda `$ as a continuous random walk, similar to the approach in (see also ). This random walk has a mean drift:
$$v=\frac{dw}{dt}=r_{in}r_{out}\left[(f_+f_{})\tau +f_+T_+f_{}T_{}\right]$$
(4)
where the weighted correlation times, $`T_\pm `$, are
$$T_\pm =_0^{\mathrm{}}𝑑tK(t)C(\pm t).$$
(5)
The first term in Eq. (4) is the contribution from uncorrelated firing activity in the cells and is proportional to $`\tau =_0^{\mathrm{}}𝑑tK(t)`$. The other terms represent the contribution from the synchrony between the two spike trains and are proportional to $`T_\pm `$. The expression for the diffusion constant $`D(w)`$ of this random walk is more complex and will be presented elsewhere. Here we note that $`D`$ is small since it is proportional to the $`\lambda `$. According to Fokker-Planck theory, the equilibrium density $`P(w)`$ can be described as a Gibbs distribution with a plasticity potential $`U(w)`$ where in the limit of small $`\lambda `$:
$$U(w)\lambda \mathrm{log}[P(w)]2_0^w𝑑w^{}v(w^{})/D(w^{})$$
(6)
Thus, $`P(w)`$ will be concentrated near the global minima of $`U(w)`$. Depending upon the implementation of the model, the minimum can be located at an interior point where the drift $`v(w)`$ vanishes, or at the boundaries $`w=0`$ and $`w=1`$. To evaluate whether the distribution of $`w`$ contains a peak at $`0<w<1`$ or at the boundaries, the specific form of the correlation function $`C(\pm t)`$ needs to be considered.
We first consider the simple example where the spike train $`\{t_i\}`$ is a homogeneous Poisson process and the output spike train is a shifted version of the input train, i.e., $`t_o=t_i+\mathrm{\Delta }t`$ where $`\mathrm{\Delta }t`$ is the temporal shift between the two spike trains. In this case, $`r=r_{in}=r_{out}`$, and $`C(t)=r^1\delta (t\mathrm{\Delta }t)`$, and $`T_\pm =r^1\mathrm{exp}(|\mathrm{\Delta }t|/\tau )\theta (\pm \mathrm{\Delta }t)`$, where $`\theta (x)=1`$ if $`x>0`$ and $`0`$ otherwise. For the additive model in Eq. (2), this leads to a net drift:
$$v=\{\begin{array}{cc}(1\alpha )\tau r^2+re^{\mathrm{\Delta }t/\tau },& \mathrm{\Delta }t>0\\ (1\alpha )\tau r^2\alpha re^{\mathrm{\Delta }t/\tau },& \mathrm{\Delta }t<0.\end{array}$$
(7)
Here $`v`$ (as well as $`D`$) is independent of $`w`$ and the potential $`U(w)`$ is $`U(w)2vw/D`$. In the limit of small $`\lambda `$, the equilibrium distribution will be a $`\delta `$-function centered at $`0`$ when $`v<0`$ and at $`1`$ if $`v>0`$.
These results are confirmed by the simulations shown in Figure 1(a) and 2(a), where we have taken $`\alpha =1.05`$ and $`0.95`$. For $`\alpha =1.05`$, the magnitude of the negative change is slightly larger than the positive one. The mean synapse is zero except when $`0<\mathrm{\Delta }t<\mathrm{\Delta }t_0`$ with $`\mathrm{\Delta }t_050`$ msec. In this range, the positive correlation between the input and output cells overcome the negative bias in the update rule to generate a positive drift so that $`w1`$. The transition at $`\mathrm{\Delta }t_0`$ is precisely the point where $`T_+=(\alpha 1)\tau `$, see Eq. (7). This behavior is highly sensitive to whether $`\alpha `$ is larger or smaller than 1. For $`\alpha =0.95`$, the mean synapse is at zero only in the range $`\mathrm{\Delta }t_0<\mathrm{\Delta }t<0`$ where the negative correlation is larger than the positive bias. Otherwise $`w1`$.
In contrast, for the multiplicative model the drift velocity is given by:
$$v=\{\begin{array}{cc}\left[1(1+\alpha )w\right]\tau r^2+r(1w)e^{\mathrm{\Delta }t/\tau },\hfill & \mathrm{\Delta }t>0\hfill \\ \left[1(1+\alpha )w\right]\tau r^2\alpha wre^{\mathrm{\Delta }t/\tau },\hfill & \mathrm{\Delta }t<0\hfill \end{array}$$
(8)
Here the drift depends on $`w`$. It is positive for small $`w`$ and becomes negative for large values of $`w`$. In this case, $`U`$ has an approximately parabolic shape with a minimum located at $`w=w_0`$ where the drift velocity vanishes:
$$w_0=\{\begin{array}{cc}1\alpha \left[1+\alpha +(\tau r)^1e^{\mathrm{\Delta }t/\tau }\right]^1,\hfill & \mathrm{\Delta }t>0\hfill \\ \left[1+\alpha (1+(\tau r)^1e^{\mathrm{\Delta }t/\tau })\right]^1,\hfill & \mathrm{\Delta }t<0.\hfill \end{array}$$
(9)
This leads to a distribution $`P(w)`$ with a narrow peak at $`w_0`$, as shown in Fig. 2(b). For large values of $`\mathrm{\Delta }t`$, the input and output cells are essentially uncorrelated, in which case $`w_0=1/(1+\alpha )0.5`$ for $`\alpha 1`$. For positive $`\mathrm{\Delta }t50`$ msec, the positive correlation between the two cells gives rise to a mean $`w>0.5`$ as shown in Fig. 1(b). Conversely, for small negative $`\mathrm{\Delta }t`$, the reverse corelation leads to a mean $`w<0.5`$. Thus, through this learning rule, the synapse smoothly encodes the temporal phase relationship between the presynaptic and postsynaptic cells. A similar dependence is found when one varies the degree of synchrony between the two cells rather than its phase.
Let us now consider the situation where a large population of $`N`$ cells with modifiable excitatory synapses $`w_i`$ drive a single postsynaptic cell. In the numerical simulations below, the output neuron obeys dynamics commonly known as “Integrate and Fire”, where the potential of the cell is described by the equation: $`\tau _m\dot{V}=VI_s`$. $`\tau _m`$ is the passive time constant of the cell, and when the potential $`V`$ reaches the threshold $`V=1`$ it is reset to zero. $`I_s(t)`$ is the synaptic current generated by the $`N`$ excitatory cells. Each spike in the presynaptic cells triggers a contribution to the output cell’s conductance that decays with synaptic time constant $`\tau _s`$, yielding:
$$I_s(t)=g_s\underset{i=1}{\overset{N}{}}w_i(t)\underset{t_i<t}{}e^{(t_it)/\tau _s)}(VV_s).$$
(10)
$`V_s`$ is the reversal potential for the excitatory synapses. The synaptic efficacy $`w_i(t)`$ describes the increase in the output cell’s conductance in units of $`g_s`$, immediately after a spike in the $`i`$-th cell. The peak conductances, $`w_i`$, are in turn modified by the TAH dynamics described above.
Here we present a general theoretical analysis of the system which is independent of the details of the output cell dynamics. In the limit of small $`\lambda `$ Eqs. (4)-(6) holds for each of the $`i`$-th synapses, with its correlation times $`T_\pm ^i`$ and $`T_\pm ^{}_{}{}^{}i`$ defined using the correlation $`C_i(t)`$ between the $`i`$-th cell and the output. As before, we will assume that the inputs are described by independent homogeneous Poisson processes, all with mean rate $`r_{in}`$. However, the statistics of firing in the output cell as described by its mean rate $`r_{out}`$ and $`\{C_i(t)\}_{i=1}^N`$, is determined by its response to the incoming spikes rather than determined externally as in the previous example. We first describe the behavior of the additive model. Due to causality and the Poisson nature of the input spikes, $`T_{}^i=0`$. For $`T_+^i`$, we make the following plausible assumptions (i) since the output cell is driven by a large number of asynchronous inputs, $`T_+^i`$ is positive but small; and (ii) its value increases roughly linearly with the synaptic efficacy of the presynaptic $`i`$-th cell, namely,
$$T_+^i\tau \chi w_i,\chi >0$$
(11)
where (iii) $`\chi =\chi (r_{out})`$ is expected to be a monotonically decreasing function of the firing rate of the output cell. $`\chi `$ is larger at smaller $`r_{out}`$ when the output cell spends more time near threshold and thus is more sensitive to the timing of the incoming spikes. Eq. (11) implies that the drift for each synapse, $`v_i1\alpha +\chi w_i`$. Thus, if $`\alpha 1`$ is positive and of order $`1`$, the system will converge to a state where all the $`w_i`$ are zero, or when $`\alpha <1`$, $`w_i1`$. An interesting situation occurs when $`0<\alpha 11`$ so that the weak negative bias can balance the weak positive correlations. In this case, the diffusion constant $`D`$ is approximately constant, and the potential is given by:
$$U(w)\frac{4}{1+\alpha ^2}\left[(\alpha 1)w\frac{1}{2}\chi w^2\right]$$
(12)
which has local minima at the boundaries, $`w=0,1`$. This equation has to be solved self-consistently since $`r_{out}`$ which determines $`\chi `$ is itself dependent on $`U`$. Over a wide range of input rates the self-consistent solution is an unsaturated state, in which $`P(w)`$ has significant weight both near $`0`$ and near $`1`$, which in turn implies that $`U(0)=U(1)`$ up to order $`\lambda `$. Hence by Eq. (12), this state is characterized by an output rate $`r_{out}^{}`$ such that
$$\chi (r_{out}^{})=2(\alpha 1).$$
(13)
More precisely, as $`r_{in}`$ increases, $`r_{out}`$ increases slightly by an amount of order $`\lambda `$ inducing a decrease in $`\chi `$ of that order. This leads to a small relative increase in $`U(1)`$, which in turn reduces $`P(1)`$ by an amount which roughly compensates for the increase in $`r_{in}`$, maintaining Eq. (13). This behavior is confirmed by simulations whose results are shown in Fig. 2(a) and 3(a). As $`r_{in}`$ increases from 10 to 40 Hz the output rate remains approximately constant at $`22Hz`$, and $`\chi 0.1`$ in agreement with eq. (13). The mean efficacy $`w`$ decreases to compensate for the increase in $`r_{in}`$, as found in .
The multiplicative model results in a very different state. In fact, since the correlations are weak for large $`N`$, the equilibrium behavior is similar to the previous example with only a single input and output cell with large $`\mathrm{\Delta }t`$. In particular, like Eq. (9), the potential $`U(w)`$ has a single minimum at the point of zero drift:
$$w_0\frac{1}{1+\alpha }$$
(14)
Thus, the distribution is highly concentrated near $`w_0`$, as seen in Fig. 2(b), and is largely independent of the mean rates of both the input and output cells. As $`r_{in}`$ increases, the output rate also increases and is similar in behavior to a cell with fixed synapses, $`w_iw_0`$. Note that as $`r_{out}`$ increases, $`\chi `$ decreases but this results in only a small decrease in the mean synaptic efficacy. Finally, we note that in contrast to the additive model where the boundaries are local minima of $`U`$ and the point of zero drift is a maximum of $`U`$, in the multiplicative model $`U(w)`$ has a single minimum at the point of zero drift. Hence, the equilibration time of the additive model will in general be much slower than that of the multiplicative model since synapses have to overcome the potential barrier of the synaptic potential. Indeed, this difference in equilibration times is seen in our numerical simulations.
We have shown here that the multiplicative TAH rule leads to a very different equilibrium distribution of synapses compared with an additive rule. Most importantly, the multipicative model is not sensitive to moderate changes in the parameters of the plasticity rule and does not suffer from slow convergence. Furthermore, experimental results reveal a dependence of the magnitudes of the synaptic changes on the amplitude of the initial synaptic efficacy, which supports a multiplicative TAH rule . The observed mean fractional negative change remained constant over a wide range of synaptic efficacies, and is consistent with the assumption that the negative change is proportional to $`w`$ as described by a multiplicative $`f_{}(w)`$ in Eq. (3). The fractional positive changes monotonically decreased with synaptic efficacy and vanish smoothly at some maximum value, again in qualitaive agreement with the multiplicative model. Although the observed shape of the $`w`$ dependence deviates from the simple linear form of $`f_+(w)`$ assumed here, the qualitative properties of the rule are unaffected by this difference. In conclusion, networks with the multiplicative TAH rule in the asynchronous state should display an equilibrium synaptic distribution that is largely insensitive to the firing rates of the pre- and post-synaptic cells. On the other hand, a coherent temporal modulation of the firing of the inputs to a target cell leads to a synaptic distribution which encodes the degree and phase of the synchrony between the cells in a smooth manner. The functional implications of this behavior in the development and learning properties of neuronal systems remains to be explored.
JR is supported in part by the NSF grant DMS-9804447. HS is supported in part by a grant of the Israeli Science Foundation and the Israel-USA Binational Science Foundation. HS and DDL also acknowledge the support of Bell Laboratories, Lucent Technologies.
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# 1 Intoduction
## 1 Intoduction
The first observation of $`B_c`$ meson at FNAL motivates new theoretical study in the field of heavy quarkonium physics. The unique properties of $`B_c`$ meson, containing heavy quarks of different flavours, may be studied, first of all, via decay modes with $`J/\psi `$ meson in the final state. We consider here two-particle hadronic decays of $`B_c`$ meson into S-wave
and P-wave charmonium states
$$B_cX_{c\overline{c}}+\pi (\rho ),$$
(1)
where $`X_{c\overline{c}}=\eta _c`$ or $`J/\psi `$ for S-wave states, and $`X_{c\overline{c}}=h_c`$, $`\chi _{c0}`$, $`\chi _{c1}`$, $`\chi _{c2}`$ for P-wave states. These decays are of considerable interest as a clean signal of $`B_c`$ meson production in high energy collisions .
Because of the large momentum transfer to the spectator quark ($`|k^2|m_c^2\mathrm{\Lambda }_{QCD}`$) in the decays (1.1), the hard scattering formalism is more appropriate than the spectator model, which based on the transition form-factor calculation under the overlapping of the quarkonium wave functions. In contrast to the spectator model, the hard scattering formalism results in the approximate double enhancement of the decay amplitude for the decays $`B_cJ/\psi (\eta _c)\pi `$, as it was found in .
## 2 The model
In the model under consideration it is assumed that heavy quarkonium ($`B_c`$ or $`X_{c\overline{c}}`$) is a non-relativistic quark-antiquark system with small binding energy. Such a way, the $`B_c`$ meson decay amplitude is factorized in a hard part which describes the process $`\overline{b}+c\overline{c}+c+\pi (\rho )`$ and an amplitude describing the binding of the initial and the final heavy quark pairs into $`B_c`$ meson and $`X_{c\overline{c}}`$ state.
The general covariant formalism for calculating the production and decay rates of S-wave and P-wave heavy quarkonium in the non-relativistic expansion was developed some times ago . In the leading non-relativistic approximation, the mass of $`B_c`$ meson $`m_1`$ is simply the sum of b-quark and c-quark masses $`m_b+m_c`$, and the mass of $`X_{c\overline{c}}`$ is equal $`2m_c`$. The amplitude for decay the bound-state ($`\overline{b}c`$) in a state with momentum $`p_1`$, total angular momentum $`J_1`$, total orbital momentum $`L_1`$ and total spin $`S_1`$ into bound-state ($`\overline{c}c`$) in a state with momentum $`p_2`$, total angular momentum $`J_2`$, total orbital momentum $`L_2`$ and total spin $`S_2`$ is given by
$`A(p_1,p_2)={\displaystyle \frac{d\stackrel{}{q}_1}{(2\pi )^3}\underset{L_{1z}S_{1z}}{}\mathrm{\Psi }_{L_{1z}S_{1z}}(\stackrel{}{q}_1)}<L_1L_{1z};S_1S_{1z}|J_1J_{1z}>\times `$ (2)
$`\times {\displaystyle }{\displaystyle \frac{d\stackrel{}{q}_2}{(2\pi )^3}}{\displaystyle \underset{L_{2z}S_{2z}}{}}\mathrm{\Psi }_{L_{2z}S_{2z}}(\stackrel{}{q}_2)<L_2L_{2z};S_2S_{2z}|J_2J_{2z}>M(p_1,p_2,q_1,q_2),`$
where $`M(p_1,p_2,q_1,q_2)`$ is the hard amplitude which is described by the diagrams in Fig. 1,2.
To introduce the operators $`\mathrm{\Gamma }_{SS_z}(p,q)`$ which projects the quark-antiquark pairs onto the bound states with fixed quantum number up to second order in $`q_1`$ and $`q_2`$:
$`\mathrm{\Gamma }_{S_1S_{1z}}(p_1,q_1)={\displaystyle \frac{\sqrt{m_1}}{4m_cm_b}}({\displaystyle \frac{m_c}{m_1}}\widehat{p}_1\widehat{q}_1+m_c)\widehat{A}_1({\displaystyle \frac{m_b}{m_1}}\widehat{p}_1+\widehat{q}_1m_b),`$ (3)
where $`\widehat{A}_1=\gamma _5`$ for $`S_1=0`$ and $`\widehat{A}_1=\widehat{\epsilon }(S_{1z})`$ for $`S_1=1`$;
$`\mathrm{\Gamma }_{S_2S_{2z}}^{}(p_2,q_2)={\displaystyle \frac{\sqrt{m_2}}{4m_c^2}}({\displaystyle \frac{m_c}{m_2}}\widehat{p}_2+\widehat{q}_2m_c)\widehat{A}_2({\displaystyle \frac{m_c}{m_2}}\widehat{p}_2\widehat{q}_2+m_c),`$ (4)
where $`\widehat{A}_2=\gamma _5`$ for $`S_2=0`$, $`\widehat{A}_2=\widehat{\epsilon }(S_{2z})`$ for $`S_2=1`$ and $`\epsilon (S_{1z,2z})`$ are spin-one polarization four-vectors.
Using projection operators (2.2) and (2.3) the hard amplitude $`M(p_1,p_2,q_1,q_2)`$ may be presented as follows:
$$M(p_1,p_2,q_1,q_2)=\text{Tr}\left[\mathrm{\Gamma }^{}(p_2,q_2)\gamma ^\beta \mathrm{\Gamma }(p_1,q_1)𝒪_\beta \right],$$
(5)
where for decay with $`\pi `$ meson in final state
$`𝒪_\beta `$ $`=𝒪_\beta ^1+𝒪_\beta ^2,`$ (6)
$`𝒪_\beta ^1`$ $`={\displaystyle \frac{G_F}{\sqrt{2}}}{\displaystyle \frac{16\pi \alpha _s}{3}}V_{bc}f_\pi a_1\widehat{p}_3(1\gamma _5)\left({\displaystyle \frac{\widehat{x}_1+m_c}{x_1^2m_c^2}}\right){\displaystyle \frac{\gamma _\beta }{k^2}},`$ (7)
$`𝒪_\beta ^2`$ $`={\displaystyle \frac{G_F}{\sqrt{2}}}{\displaystyle \frac{16\pi \alpha _s}{3}}V_{bc}f_\pi a_1{\displaystyle \frac{\gamma _\beta }{k^2}}\left({\displaystyle \frac{\widehat{x}_2+m_b}{x_2^2m_b^2}}\right)\widehat{p}_3(1\gamma _5),`$ (8)
and
$$\widehat{x}_1=\frac{m_b}{m_1}\widehat{p}_1+\widehat{q}_1\widehat{p}_3,\widehat{x}_2=\frac{m_c}{m_2}\widehat{p}_2+\widehat{q}_2+\widehat{p}_3$$
$$\widehat{k}=\frac{m_c}{m_2}\widehat{p}_2\frac{m_c}{m_1}\widehat{p}_1\widehat{q}_1+\widehat{q}_2.$$
The factor $`a_1`$ comes from hard gluon corrections to the four-fermion effective vertex. Since $`q_1/m_1`$ and $`q_2/m_2`$ are small quantities, we can expand $`M(p_1,p_2,q_1,q_2)`$ around $`q_1=q_2=0`$ in a Taylor expansion:
$`M(p_1,p_2,q_1,q_2)=M(p_1,p_2,0,0)+q_{1\alpha }{\displaystyle \frac{M}{q_{1\alpha }}}|_{q_{1,2}=0}+q_{2\alpha }{\displaystyle \frac{M}{q_{2\alpha }}}|_{q_{1,2}=0}+`$
$`+{\displaystyle \frac{1}{2}}q_{1\alpha }q_{1\beta }{\displaystyle \frac{^2M}{q_{1\alpha }q_{2\beta }}}|_{q_{1,2}=0}+\mathrm{}`$ (9)
Here the each term correspond to quantum numbers $`L_1=L_2=0`$, $`L_1=1`$ and $`L_2=0`$, $`L_1=0`$ and $`L_2=1`$, and so forth. Thus for the S-wave ($`L=0`$) and P-wave ($`L=1`$) states, the amplitude $`A(p_1,p_2)`$ will depend on the quarkonium radial wave-functions through the following relations:
$`{\displaystyle \frac{d\stackrel{}{q}}{(2\pi )^3}\mathrm{\Psi }_{00}(\stackrel{}{q})}={\displaystyle \frac{R_s(0)}{\sqrt{4\pi }}},`$ (10)
$`{\displaystyle \frac{d\stackrel{}{q}}{(2\pi )^3}\mathrm{\Psi }_{1L_z}(\stackrel{}{q})q_\alpha }=i\sqrt{{\displaystyle \frac{3}{4\pi }}}R_p^{}(0)\epsilon _\alpha (p,L_z),`$ (11)
where $`\epsilon _\alpha (p,L_z)`$ is the polarization vector for the spin-one particle.
In the case of production charmonium state $`{}_{}{}^{1}P_{1}^{}`$ one has
$$\underset{L_{2z}}{}\epsilon ^\alpha (p_2,L_{2z})<1L_{2z},00|1,J_{2z}>=\epsilon ^\alpha (p_2,J_{2z}).$$
(12)
The summation over polarization may be done using following expression
$$\underset{J_{2z}=1}{\overset{1}{}}\epsilon ^\alpha (p2,J_{2z})\epsilon ^\beta (p2,J_{2z})=𝒫^{\alpha \beta }(p_2),$$
(13)
where
$$𝒫^{\alpha \beta }(p_2)=g^{\alpha \beta }+\frac{p_2^\alpha p_2^\beta }{m_2^2}.$$
In the case of production $`{}_{}{}^{3}P_{J}^{}(J=0,1,2)`$ states one has
$`{\displaystyle \underset{S_{2z},L_{2z}}{}}\epsilon ^\alpha (p_2,L_{2z})<1L_{2z},1S_{2z}|J_2,J_{2z}>\epsilon ^\beta (S_{2z})=`$
$`=\{\begin{array}{c}\frac{1}{\sqrt{3}}(g^{\alpha \beta }\frac{p_2^\alpha p_2^\beta }{m_2^2})\text{ for }J_2=0,\hfill \\ \frac{i}{\sqrt{2}m_2}\epsilon ^{\alpha \beta \mu \nu }p_{2\mu }\epsilon _\nu (p_2,J_{2z})\text{ for }J_2=1,\hfill \\ \epsilon ^{\beta \alpha }(p_2,J_{2z})\text{ for }J_2=2.\hfill \end{array}`$ (17)
The polarization sums for $`J_2=2`$ is given by the following expression :
$`{\displaystyle \underset{J_{2z}=2}{\overset{2}{}}}\epsilon _{\alpha \beta }(p_2,J_{2z})\epsilon _{\mu \nu }^{}(p_2,J_{2z})`$ $`=`$ $`{\displaystyle \frac{1}{2}}(𝒫_{\alpha \mu }(p_2)𝒫_{\beta \nu }(p_2)+𝒫_{\alpha \nu }(p_2)𝒫_{\beta \mu }(p_2))`$ (18)
$`{\displaystyle \frac{1}{3}}𝒫_{\alpha \beta }(p_2)𝒫_{\mu \nu }(p_2))`$
## 3 The results
Omitting the details of calculations we presented below the results for decay widths of ground-state $`(\overline{b}c)`$ system into different states of charmonium plus $`\pi `$ or $`\rho `$ meson. In the limit of vanishing $`\pi `$ meson mass we obtained the simple analytical formulae:
$`\mathrm{\Gamma }(B_c\psi \pi )={\displaystyle \frac{128}{9\pi }}F{\displaystyle \frac{|R_2(0)|^2}{m_2^3}}{\displaystyle \frac{(1+x)^3}{(1x)^5}},`$ (19)
$`\mathrm{\Gamma }(B_c\eta _c\pi )={\displaystyle \frac{32}{9\pi }}F{\displaystyle \frac{|R_2(0)|^2}{m_2^3}}{\displaystyle \frac{(1+x)^3}{(1x)^5}}(x^22x+3)^2,`$ (20)
$`\mathrm{\Gamma }(B_ch_c\pi )={\displaystyle \frac{128}{3\pi }}F{\displaystyle \frac{|R_2^{}(0)|^2}{m_2^5}}{\displaystyle \frac{(1+x)^3}{(1x)^7}}(x^2x+2)^2,`$ (21)
$`\mathrm{\Gamma }(B_c\chi _{c0}\pi )={\displaystyle \frac{128}{9\pi }}F{\displaystyle \frac{|R_2^{}(0)|^2}{m_2^5}}{\displaystyle \frac{(1+x)^3}{(1x)^7}}(3x^312x^2+14x7)^2,`$ (22)
$`\mathrm{\Gamma }(B_c\chi _{c1}\pi )={\displaystyle \frac{256}{3\pi }}F{\displaystyle \frac{|R_2^{}(0)|^2}{m_2^5}}{\displaystyle \frac{(1+x)^3}{(1x)^5}}(x^2x1)^2,`$ (23)
$`\mathrm{\Gamma }(B_c\chi _{c2}\pi )={\displaystyle \frac{256}{9\pi }}F{\displaystyle \frac{|R_2^{}(0)|^2}{m_2^5}}{\displaystyle \frac{(1+x)^5}{(1x)^7}},`$ (24)
where
$$x=\frac{m_2}{m_1},\text{ and }F=\alpha _s^2G_F^2V_{bc}^2f_\pi ^2|R_1(0)|^2a_1^2.$$
To perform the numerical calculations we use following set of parameters: $`G_F=1.166\times 10^5`$ GeV<sup>-2</sup>, $`\alpha _s=0.33`$, $`V_{bc}=0.04`$, $`f_\pi =0.13`$ GeV, $`m_\pi =0.14`$ GeV, $`m_{B_c}=6.3`$ Gev, $`m_\psi =3.1`$ Gev, $`m_{\eta _c}=2.98`$ Gev, $`m_{h_c}=3.5`$ Gev, $`m_{\chi _{c0}}=3.4`$ Gev, $`m_{\chi _{c1}}=3.5`$ Gev, $`m_{\chi _{c2}}=3.55`$ Gev, $`|R_{s1}(0)|^2=1.27`$ GeV<sup>3</sup>, $`|R_{s2}(0)|^2=0.94`$ GeV<sup>3</sup>, $`|R_{p2}^{}(0)|^2=0.08`$ GeV<sup>5</sup>.
With above mentioned set of parameters one gets the following result
$$\mathrm{\Gamma }(B_cJ/\psi +\pi )=7.5\times 10^{15}a_1^2\text{ GeV}.$$
(25)
The decay widths into different charmonium states plus $`\pi `$ meson may be presented through decay width for $`B_cJ/\psi \pi `$ as it is shown in Table 1.
| $`X_{c\overline{c}}`$ | $`{}_{}{}^{2S+1}X_{J}^{}`$ | $`\frac{\mathrm{\Gamma }(B_cX_{c\overline{c}}\pi )}{\mathrm{\Gamma }(B_cJ/\psi \pi )}`$ | $`\frac{\mathrm{\Gamma }(B_cX_{c\overline{c}}\rho )}{\mathrm{\Gamma }(B_cX_{c\overline{c}}\pi )}`$ |
| --- | --- | --- | --- |
| $`J/\psi `$ | $`{}_{}{}^{3}S_{1}^{}`$ | 1.00 | 4.0 |
| $`\eta _c`$ | $`{}_{}{}^{1}S_{0}^{}`$ | 1.17 | 3.2 |
| $`h_c`$ | $`{}_{}{}^{1}P_{1}^{}`$ | 0.50 | 3.7 |
| $`\chi _{c0}`$ | $`{}_{}{}^{3}P_{0}^{}`$ | 0.29 | 3.6 |
| $`\chi _{c1}`$ | $`{}_{}{}^{3}P_{1}^{}`$ | 0.10 | 5.6 |
| $`\chi _{c2}`$ | $`{}_{}{}^{3}P_{2}^{}`$ | 0.28 | 4.3 |
Table 1.
The another source of $`J/\psi `$ mesons is two-particle decay of $`B_c`$ meson with $`\rho `$ meson in the final state: $`B_cX_{c\overline{c}}\rho `$. The calculating for the decay widths $`\mathrm{\Gamma }(B_cX_{c\overline{c}}\rho )`$ may be done the same way as for widths $`\mathrm{\Gamma }(B_cX_{c\overline{c}}\pi )`$ using substitution $`f_\pi \widehat{p}_3m_\rho f_\rho \widehat{\epsilon }_3`$ in (2.6) and (2.7), where $`\epsilon _3^\mu `$ is $`\rho `$ meson polarization four-vector. Taking into account that $`f_\rho =0.22`$ GeV and $`m_\rho =0.77`$ GeV, we have obtained the decay widths $`\mathrm{\Gamma }(B_cX_{c\overline{c}}\rho )`$ which are presented in the Table 1 too as the ratio $`\mathrm{\Gamma }(B_cX_{c\overline{c}}\rho )/\mathrm{\Gamma }(B_cX_{c\overline{c}}\pi )`$.
We found surprisingly large value for the decay width of $`B_c`$ meson into P-wave charmonia, which is 50 % of the decay width into S-wave states. Because of $`J/\psi `$ meson production in the decays of $`B_c`$ meson is very suitable process from viewpoint of an experimental study ,
it is interesting to compare the direct $`J/\psi `$ production ($`B_cJ/\psi \pi (\rho )`$) and the cascade $`J/\psi `$ production rates. The second one comes from radiative decays of the P-wave states $`\chi _{c0},\chi _{c1}`$ and $`\chi _{c2}`$, which have following branching ratios into $`J/\psi `$ plus $`\gamma `$: $`\text{Br}(\chi _{c0}J/\psi +\gamma )=0.007,`$ $`\text{Br}(\chi _{c1}J/\psi +\gamma )=0.27`$ and $`\text{Br}(\chi _{c2}J/\psi +\gamma )=0.14`$ . Thus we have obtained
$$\frac{\mathrm{\Gamma }(B_c\chi _{c0,c1,c2}\pi J/\psi \gamma )}{\mathrm{\Gamma }(B_cJ/\psi \pi )}=0.068$$
(26)
and
$$\frac{\mathrm{\Gamma }(B_c\chi _{c0,c1,c2}\rho J/\psi \gamma )}{\mathrm{\Gamma }(B_cJ/\psi \rho )}=0.082$$
(27)
The ratio for the sum of the $`B_c`$ meson decay widths into $`J/\psi `$ meson plus $`\pi `$ or $`\rho `$ meson is equal to
$$\frac{\mathrm{\Gamma }(B_cJ/\psi \pi (\rho ),\text{cascade})}{\mathrm{\Gamma }(B_cJ/\psi \pi (\rho ),\text{direct})}=0.08$$
(28)
In compare with the previous estimations we should expect the enhancement by 8 % the branching ratio for the $`B_c`$ decays into $`J/\psi `$ in the two-particle hadronic decays via the cascade processes. This fact makes the probability of $`B_c`$ meson observation through the decays $`B_cX_{c\overline{c}}\pi (\rho )`$
in the current experiments more real.
The author thanks V.V. Kiselev and A.K. Likhoded for the valuable discussion. This work is supported by the Program ”Universities of Russia”, Grant 02.01.03 and the Russian Ministry of Education, Grant 98-0-6.2-53.
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# Nucleon-nucleon interaction in the Skyrme model
## 1 Introduction
Nucleon-nucleon $`(NN)`$ interactions are relatively simple at large distances and become rapidly more complex as one moves inward. In the best phenomenological models existing at present, that reproduce low-energy observables accurately, they are described by the consensual one-pion exchange potential (OPEP), supplemented by theoretical two-pion exchange potentials (TPEP) and parametrized at short distances . The OPEP is responsible for a strong tensor component, which is mostly important in few-body systems, such as the deuteron. Two-pion exchange, on the other hand, gives rise to the central potential, that survives to all averages and is responsible for most properties of large systems and nuclear matter. Quantum Chromodynamics (QCD) is the basic framework for the study of strong processes and should have, in principle, an important role in the description of nuclear forces. However, at present, the non-Abelian character of this theory prevents low energy calculations and one has to resort to effective theories, which must reflect the main features of QCD. Thus, in Nuclear Physics applications, besides the usual space-time invariances, one requires these theories to have approximate chiral symmetry. The latter is usually restricted to the $`SU(2)\times SU(2)`$ sector, for most processes involve only the quarks $`u`$ and $`d`$. This symmetry is explicitly broken by the small quark masses and, at the effective level, by the pion mass.
Chiral symmetry has no influence over the OPEP, but is crucial to the TPEP, which depends on an intermediate pion-nucleon ($`\pi N`$) amplitude . In the case of $`NN`$ interactions, the importance of this symmetry was stressed already in the early seventies, by Brown and Durso and by Chemtob, Durso and Riska , who used it to constrain the form of the TPEP. In that decade it also became popular to describe nuclear processes by means of the linear $`\sigma `$ model , containing a fictitious particle called $`\sigma `$ that, to some extent, simulates the TPE. The elimination of this unobserved degree of freedom gave rise to non-linear theories, which underlie modern descriptions of the interaction. The first theoretical framework to incorporate non-linear chiral dynamics into the $`NN`$ problem was proposed by Skyrme. This remarkable model for the nucleon, developed in the sixties and revived in the eighties , describes baryons as topological solitons, objects extended in space that rotate according to the laws of Quantum Mechanics. The quark condensate appears as an intrinsic feature, corresponding to a non-vanishing classical content of the vacuum, whose intensity is given by the pion decay constant $`f_\pi `$. Skyrmions correspond to distortions of this condensate that carry topological charges. One then works with pion fields which are unusually strong, in the sense that their amplitudes may be comparable to $`f_\pi `$. Thus, in spite of its well know limitations , the Skyrme model remains a unique laboratory for studying chiral symmetry in the non-perturbative regime.
In the early nineties Weinberg restated the role of perturbative chiral symmetry in nuclear interactions and motivated interest in the TPEP. Initially, several authors explored the pion-nucleon sector of non-linear Lagrangians , but the corresponding potentials could not reproduce even the medium range attraction in the scalar channel. This happened because the TPEP is based on an intermediate $`\pi N`$ amplitude, that can only be well described with the help of other degrees of freedom . Accordingly, in a later stage, agreement with empirical $`\pi N`$ information was enforced and descriptions could reproduce $`NN`$ scattering data .
In the case of perturbative calculations, the delta is by far the most important non-nucleonic degree of freedom and is largely responsible for the intermediate range scalar attraction. As the Skyrme model incorporates the delta from the very beginning, one expects that it should yield a good qualitative $`NN`$ potential. However, it fails to do so.
Skyrme himself considered $`NN`$ interactions, already in the sixties, using the so called product ansatz (PA) . The basic idea underlying the PA is that solutions corresponding to baryon number $`B=1`$ can be used as building blocks to construct approximate solutions with an arbitrary value of $`B`$. The great advantage of this approach is that the baryon number of the composite system is automatically equal to the number of individual $`B=1`$ skyrmions, irrespectively of their relative positions. In the PA, the skyrmions that constitute a larger system are assumed to retain their shape all the time, what is known as sudden approximation. In this framework, the construction of the $`NN`$ potential is rather simple and the fact that each nucleon has a profile function which falls off rapidly with the distance allows one to assume that, for medium and large distances, the $`B=2`$ system is not considerably different from the superposition of two with $`B=1`$.
In the eighties, the product ansatz was used by Jackson et al. and Vinh Mau et al. to calculate the $`NN`$ potential, who found out a fully repulsive central component, in disagreement with very well established phenomenology. This puzzle motivated several attempts to construct improved versions of those early works. Among them, one notes the symmetrized product ansatz by Nyman and Riska , which could produce an intermediate range scalar attraction. However, as pointed out by Sternheim and Kälbermann , there is a violation of baryon number conservation in this ansatz. Exact numerical calculations were also used, which allowed one to evaluate the reliability of the sudden approximation at short distances . Lattice calculations, using a method developed by Manton and collaborators, gave rise to a torus-like baryon density, believed to correspond to the true $`B=2`$ ground state and having almost twice the nucleon mass . The scalar potential associated with this configuration does show some medium to long range attraction . However, it is worth recalling that lattice results depend on the definitions adopted for collective coordinates and, as the full treatment is rather difficult, one usually resorts to approximations .
In this work we consider the scalar interaction between skyrmions, in order to explore the possibility of obtaining the central attraction at large distances by relaxing some of the constraints present in usual calculations. We employ the sudden approximation because it gives rise to a constructive interaction, in which undeformed nucleons are the main building blocks, as in perturbative calculations. Our presentation is divided as follows. In section 2 we study the asymptotic behaviour of the scalar potential in the standard product ansatz approximation, in order to understand why it does not yield attraction. In section 3 we discuss the construction of alternative solutions, which must be constrained to have the correct baryon number. Finally, in section 4 we analyze a possible solution to the problem, and present concluding remarks.
## 2 Central potential
The structure of the central potential has been studied recently, in the framework of chiral perturbation theory . In momentum space, the leading contribution has the generic form
$$V_C(t)=\frac{2}{f_\pi ^2m_\pi ^2}\left[f_\pi ^2\left(a_{00}^++a_{01}^+t\right)\right]\sigma (m;t),$$
(1)
where $`f_\pi `$ is the pion decay constant, the $`a_{0i}^+`$ are subthreshold coefficients and $`\sigma (m;t)`$ is the scalar form factor, that depends on both the momentum transferred $`t`$ and the baryon mass $`m`$. This form factor is defined generically in terms of the symmetry breaking Lagrangian $`_{sb}`$ as
$$N(𝒑^{})|_{sb}|N(𝒑)=\sigma (t)\overline{u}(𝒑^{})u(𝒑).$$
(2)
In configuration space, eq. (1) becomes
$$V_C(d)=\frac{2}{f_\pi ^2m_\pi ^2}\left[f_\pi ^2\left(a_{00}^++a_{01}^+\mathbf{}^2\right)\right]\sigma (m;d),$$
(3)
where $`d`$ is the internucleon distance and $`\sigma (m;d)`$ is the Fourier transform of $`\sigma (m;t)`$. In order to allow this result to be compared with the corresponding one in the Skyrme model, we note that, in the large $`N_c`$ limit, the nucleon and the delta are degenerate and very heavy. In ref. , the heavy baryon limit of $`\sigma (m;d)`$ was considered and found to be
$$\sigma (m\mathrm{};d)\sigma ^{HB}(d)=\frac{9m_\pi ^6}{128\pi ^2}\left(\frac{g_A}{f_\pi }\right)^2\left[\frac{d}{dx}\frac{e^x}{x}\right]^2,$$
(4)
with $`x=m_\pi d`$. The central potential is
$`V_C^{HB}(d)`$ $`=`$ $`{\displaystyle \frac{2}{f_\pi ^2m_\pi ^2}}\left[{\displaystyle \frac{9m_\pi ^6g_A^2}{128\pi ^2}}\right][a_{00}^+(1+{\displaystyle \frac{2}{x}}+{\displaystyle \frac{1}{x^2}})`$ (5)
$`+`$ $`4m_\pi ^2a_{01}^+(1+{\displaystyle \frac{3}{x}}+{\displaystyle \frac{11}{2x^2}}+{\displaystyle \frac{6}{x^3}}+{\displaystyle \frac{3}{x^4}})]{\displaystyle \frac{e^{2x}}{x^2}}`$
and, at very large distances, it behaves as
$$V_C^{HB}(d)K\left(\frac{e^x}{x}\right)^2.$$
(6)
The sign of the constant $`K`$ is determined by the values of the subthreshold coefficients in the combination $`(a_{00}^++4m_\pi ^2a_{01}^+)`$. In table 1 we display empirical values for the $`a_{0i}^+`$ and it is possible to note that the correct sign of $`V_C`$ comes mainly from $`a_{01}^+`$, since $`a_{00}^+`$ in isolation would give rise to a repulsive interaction.
In order to study the central potential in the Skyrme model, we recall that the standard soliton Lagrangian density is written as
$$=_\sigma +_4,$$
(7)
where
$$_\sigma =\frac{f_\pi ^2}{4}\text{Tr}(_\mu U^\mu U^{})+m_\pi ^2\frac{f_\pi ^2}{4}\text{Tr}(U+U^{}2)$$
(8)
corresponds to the non-linear $`\sigma `$ model and
$$_4=\frac{1}{32e^2}\text{Tr}[_\mu UU^{},_\nu UU^{}]^2$$
(9)
is the stabilizing term. In these expressions, $`e`$ is a free parameter, called Skyrme constant, whereas the dynamical variable $`U`$ is a $`2\times 2`$ unitary matrix, given by
$$U=e^{i𝝉\widehat{𝝅}F}=\mathrm{cos}F+i𝝉\widehat{𝝅}\mathrm{sin}F,$$
(10)
where $`𝝉`$ are the isospin Pauli matrices and $`F`$ is the chiral angle, whose boundary conditions determine the baryon number of a particular configuration. The function $`F`$ and the isospin direction $`\widehat{𝝅}`$ are related to the pion field $`𝝅`$ of the non-linear $`\sigma `$ model by $`𝝅=f_\pi \mathrm{sin}F\widehat{𝝅}.`$
In the $`B=1`$ case, a static solution is obtained using the condition $`\widehat{𝝅}=\widehat{𝒓}`$, the so called hedgehog ansatz, with boundary conditions $`F(r=0)=\pi `$ and $`F(r\mathrm{})=0`$ . The quantization of this baryon is achieved by rotating the static solution with the help of collective coordinates, as a rigid body. This procedure endows the skyrmion with spin and isospin and corresponds to multiplying the pion field by the rigid body rotation matrix $`D`$,
$$\pi _i\pi _\alpha ^q=D_{\alpha i}\pi _i.$$
(11)
The matrices $`D`$ satisfy the completeness relations $`D_{\alpha i}D_{\alpha j}=\delta _{ij},D_{\alpha i}D_{\beta i}=\delta _{\alpha \beta }`$ and, in the case of nucleons, the correspondence with the ordinary formalism is achieved by using
$$N|D_{\alpha i}|N=\frac{1}{3}N|\tau _\alpha \sigma _i|N,$$
(12)
$`\sigma _i`$ being the spin Pauli matrices.
The scalar form factor in the Skyrme model can be obtained directly from eqs. (8) and (10) and reads
$$\sigma ^{Sk}(d)=N|_{sb}(d)|N=m_\pi ^2f_\pi ^2\left[\mathrm{cos}F(d)1\right].$$
(13)
On the other hand, the asymptotic form of the chiral angle is determined by $`_\sigma `$ as
$$F_{\mathrm{}}(d)=\left(\frac{3g_Am_\pi ^2}{8\pi f_\pi ^2}\right)\left(\frac{d}{dx}\frac{e^x}{x}\right)$$
(14)
and hence, for large distances, the leading term in eq. (13) yields $`\sigma _{\mathrm{}}^{Sk}(d)=\sigma ^{HB}(d)`$. In order to test this relationship further, we write
$`\sigma _{\mathrm{}}^{Sk}(d)`$ $`=`$ $`m_\pi ^2f_\pi N|\sqrt{f_\pi ^2𝝅^q𝝅^q}f_\pi |N{\displaystyle \frac{m_\pi ^2}{2}}𝝅^2`$ (15)
$`=`$ $`{\displaystyle \frac{m_\pi ^2}{2}}\left[N|\pi _\alpha ^q|NN|\pi _\alpha ^q|N+N|\pi _\alpha ^q|\mathrm{\Delta }\mathrm{\Delta }|\pi _\alpha ^q|N\right].`$
Using eqs. (11) and (12) in the last expression, one concludes that $`N`$ and $`\mathrm{\Delta }`$ intermediate states determine respectively 1/3 and 2/3 of the total value of $`\sigma _{\mathrm{}}^{Sk}(d)`$. This relative proportion is identical to that found recently in the framework of chiral perturbation theory , indicating that the Skyrme model does provide a rather reliable description of the scalar form factor in the heavy baryon limit.
For systems with $`B=2`$, the standard point of departure for constructing approximate solutions is the product ansatz (PA). It uses two undistorted $`B=1`$ hedgehog solutions, whose centers are located at two fixed points equidistant from origin along the $`z`$ axis, so that the hedgehog space coordinates are given by $`𝒚=𝒓+d\widehat{𝒛}/2`$ and $`𝒘=𝒓d\widehat{𝒛}/2`$. Denoting the composite field by $`U(𝒚,𝒘)`$, one writes
$$U(𝒚,𝒘)=U(𝒚)U(𝒘).$$
(16)
In this configuration, the $`B=2`$ condition is automatically fulfilled, for any distance between their centers . As the PA keeps the identities of constituent skyrmions, it allows the direct incorporation of spin and isospin, through collective rotations of individual hedgehogs.
The potential is a function of the distance $`d`$ and given by
$$V(d)=d^3r_{int}(𝒓,d\widehat{𝒛}),$$
(17)
where $`_{int}`$ is obtained by using the field $`U(𝒚,𝒘)`$ in the Skyrme Lagrangian, eqs. (7-9), and subtracting the self energies $`[U(𝒚)]`$ and $`[U(𝒘)]`$. This potential works well in the isospin-dependent channels, since the OPEP is reproduced for distances larger than $`2`$ fm and it is also possible to identify the roles of $`\rho `$ and $`A1`$ mesons . On the other hand, problems occur in the scalar-isoscalar channel, where the interaction is repulsive at all distances, in sharp contradiction with phenomenology.
Using the definitions $`F_r^{}=dF(r)/dr`$, $`s_r=\mathrm{sin}F(r)`$ and $`c_r=\mathrm{cos}F(r)`$, the central potential is given by
$`V_C^{pa}(d)`$ $`=`$ $`{\displaystyle \frac{2f_\pi }{e}}{\displaystyle \frac{4\pi }{3}}{\displaystyle _0^{\mathrm{}}}dz{\displaystyle _0^{\mathrm{}}}\rho d\rho \{{\displaystyle \frac{3m_\pi ^2}{16e^2f_\pi ^2}}(1c_y)(1c_w)`$
$`+`$ $`[(F_y^2+{\displaystyle \frac{s_y^2}{y^2}})(F_w^2+{\displaystyle \frac{s_w^2}{w^2}})+{\displaystyle \frac{2s_y^2s_w^2}{y^2w^2}}(\widehat{𝒚}\widehat{𝒘})^2(F_y^2{\displaystyle \frac{s_y^2}{y^2}})(F_w^2{\displaystyle \frac{s_w^2}{w^2}})]\}.`$
In order to study its asymptotic structure, we note that the pion fields exist effectively only in the neighbourhood of the hedgehog centers. When the distance $`d`$ is large the skyrmion located at $`(0,0,d/2)`$ is in the presence of the asymptotic region of $`U(𝒚)`$, we expand $`F_y`$, $`F_y^{}`$ and $`\widehat{𝒚}\widehat{𝒘}`$ around the point $`𝒘=0`$ and write
$`F_y`$ $``$ $`\alpha e^{m_\pi w_z}\left(1+{\displaystyle \frac{f_1}{x}}+{\displaystyle \frac{f_2}{x^2}}\right){\displaystyle \frac{e^x}{x}},`$ (19)
$`F_y^{}`$ $``$ $`\alpha e^{m_\pi w_z}\left(1+{\displaystyle \frac{g_1}{x}}+{\displaystyle \frac{g_2}{x^2}}\right){\displaystyle \frac{e^x}{x}},`$ (20)
$`\widehat{𝒚}\widehat{𝒘}`$ $``$ $`{\displaystyle \frac{w_z}{\sqrt{\rho ^2+w_z^2}}}\left(1+{\displaystyle \frac{m_\pi \rho ^2}{w_zx}}{\displaystyle \frac{3m_\pi ^2\rho ^2}{2x^2}}\right),`$ (21)
where $`f_i`$, $`g_i`$ are dimensionless polynomials of $`w_z(zd/2)`$ and $`\rho `$, which are not displayed here. These expressions were tested order by order, by using them in eq. (2) and checking that the potential did had the asymptotic structure, as in eq. (6). We found out that it was necessary to expand $`F(𝒚)`$ up to order $`d^2`$, in order to have accurate results.
Replacing eqs. (19-21) into (2), we obtain an asymptotic contribution of the form
$$V_C^{pa}(d)K\left[1+\frac{\alpha _1}{x}+\frac{\alpha _2}{x^2}\right]\frac{e^{2x}}{x^2},$$
(22)
for both $`_\sigma `$ and $`_4`$, separately. The values of the parameters $`K`$ and $`\alpha _i`$ are displayed in table 2, based on the numerical constants $`m_\pi =139`$ MeV, $`f_\pi =93`$ MeV and $`e=4.0`$. For the sake of comparison, we also present the values of those parameters in the case of the phenomenological Argonne potential .
Inspecting this table, one notes that the part of the potential due to $`_\sigma `$ is attractive, but is superseeded by a repulsive contribution coming from the stabilizing term. The net sign of the potential is, then, the outcome of a large cancellation. On the other hand, the dependence of eq. (22) on $`d`$ is similar to those of both the perturbative chiral calculation, eq. (6), and of the phenomenological Argonne potential. We stress that this correct geometry is a general feature of the model, because it depends only on the form of the $`B=1`$ solution, and not on the specific ansatz used to obtain the $`B=2`$ result. In fig. 1 we display the ratios between the full and asymptotic PA potentials, as given by eqs. (2) and (22), with the purpose of illustrating their convergence. The interplay between the attractive contribution from $`_\sigma `$ and the repulsive one from $`_4`$ can also be seen in fig. 2, where we present the function $`dV_C^{pa}(d)/dz`$, corresponding to the integrand in $`z`$ of expression (2), for several values of $`d`$. One notes that the contributions in the neighbourhood of the skyrmion centers are large and positive but, on the other hand, a negative region develops as the distance $`d`$ increases.
These features of the central potential allow us to identify clearly the stabilizing term as the responsible for its repulsive character. Therefore, mechanisms which can reduce the importance of $`_4`$ may help in producing an attractive interaction. In the next section we discuss a class of such mechanisms, associated with deformations of the QCD vacuum.
## 3 Constructing $`𝐁=\mathrm{𝟐}`$ solutions
We consider here $`B=2`$ solutions, constructed by using the hedgehog $`B=1`$ skyrmions as building blocks, in the framework of the sudden approximation. In general, an ansatz is a prescription of the form
$$U(𝒚,𝒘)=f[U(𝒚),U(𝒘)],$$
(23)
where $`f`$ is a function, chosen according to physical criteria. The construction of such a function should follow some guidelines:
1. the baryon number of the composite configuration must be two for all distances $`d`$;
2. the composite pion field must have the correct quantum numbers, being pseudoscalar, isovector and odd under G-parity;
3. the composite Lagrangian must be chiral symmetric, even under G-parity and invariant under the exchange of the two constituent skyrmions.
The standard constructive approximate solution to the $`B=2`$ system is based on the PA, as discussed in the previous section. In this approach, the composite pion field, obtained from eqs. (10) and (16), is given by
$$𝑷_{pa}=\frac{1}{f_\pi }\left(\sigma _y𝝅_w+\sigma _w𝝅_y𝝅_y\times 𝝅_w\right),$$
(24)
where $`𝝅_r`$ is the pion field of the hedgehog with coordinate $`𝒓`$ and $`\sigma _rf_\pi \mathrm{cos}F(𝒓)`$. The function
$$S_{pa}=\frac{1}{f_\pi }\left(\sigma _y\sigma _w𝝅_y𝝅_w\right).$$
(25)
is the composite analogous of $`\sigma `$ and satisfies $`S_{pa}^2+𝑷_{pa}^2=f_\pi ^2`$.
The field $`𝑷_{pa}`$ has a rather serious drawback as a candidate for the pion field, namely that it contains an azimuthal term which is both even under G-parity and antisymmetric under hedgehog exchange. Hence it does not have good pion quantum numbers, violating requirements 2 and 3 stated above.
This motivated us to try to understand whether this problem could be responsible for the absence of attraction found in the central potential. We considered several alternative possibilities, inspired in the PA. The basic idea is to propose a composite field $`𝑷`$, use it to define a function $`S`$ by
$$S^2=f_\pi ^2𝑷^2,$$
(26)
construct the unitary field as
$$U=\left[S+i𝝉𝑷\right]/f_\pi ,$$
(27)
and feed it into the Skyrme Lagrangian. We begin by describing briefly some unsuccessful attempts, in order to prevent readers from repeating them.
The simplest exchange-symmetric ansatz would be the average $`𝑷=\left(𝝅_y+𝝅_w\right)/2`$. However, when $`d=0`$, one has $`F_y=F_w=F_r`$ and hence $`𝑷=f_\pi \mathrm{sin}F_r\widehat{𝒓}`$ corresponds to a $`B=1`$ field, which must be disregarded.
This suggests that, in order to obtain $`B=2`$, it is mandatory to mix $`𝝅`$ and $`\sigma `$. In the case of the PA, which yields $`B=2`$ at all distances, we note that the chiral constraint between $`𝑷_{pa}`$ and $`S_{pa}`$ allows one to write
$$U_{pa}=\left[S_{pa}+i𝝉𝑷_{pa}\right]/f_\pi e^{i𝝉𝒖F_{pa}},$$
(28)
where $`𝒖`$ is a unit vector, taken as pointing always away from the origin of the coordinate system, and $`F_{pa}`$ is a profile function. In fig. 3 we display the behaviour of this angle along the axes $`z`$ and $`\rho `$, for various values of the internucleon distance $`d`$. The solid lines correspond to the case $`d=0`$, which is spherically symmetric and it is possible to see that, along both directions, the chiral angle varies smoothly from $`2\pi `$ at the origin to $`0`$ at infinity. In the case $`d=0.5`$ fm, shown in dotted lines, one notes that a discontinuity has appeared along the $`\rho `$ axis. This discontinuity increases with distance and, at $`d_{crit}=0.86`$ fm, the chiral angle is such that $`F_{pa}(\rho =0,z0)=2\pi `$ and $`F_{pa}(\rho 0,z=0)=0`$. Therefore, at this critical point, it is more natural to set $`F_{pa}(0,0)=0`$ and to work with two separate solutions, such as illustrated by the dashed and dot-dashed lines, corresponding to $`0.9`$ and $`2.0`$ fm, respectively. This suggests that, from $`d_{crit}`$ onwards, each of the interacting skyrmions acquires a considerable individuality.
The combination
$$𝑷=\frac{1}{f_\pi }\left(\sigma _y𝝅_w+\sigma _w𝝅_y\right)$$
(29)
is interesting, for it has an explicit physical meaning. As the function $`\sigma (𝒓)`$ is associated with the quark condensate that surrounds the baryon labeled by $`𝒓`$, this field $`𝑷`$ represents each skyrmion immersed in the distorted vacuum of the other one. The condition (26) allows one to determine $`S`$ up to a sign. In the case $`B=1`$ the field $`\sigma `$ changes sign when one goes from infinity to the origin and the same happens when $`B=2`$. The sign of $`S`$ is also important and, in order to fix it, we note that the behaviours of eqs. (24) and (29) along the $`z`$ axis are identical, since the azimuthal component vanishes. We then forced the condition $`S=S_{pa}`$ along this axis. However, this ansatz, based on eq. (29), gives rise to a baryon number which varies with $`d`$, as shown in fig. 4, and had to be abandoned.
This discussion illustrates the fact that it is not trivial to build an ansatz with a good topology. We thus decided to adopt simultaneously the pion field as given by eq. (29) and the function $`S_{ap}`$ of the PA, eq. (25), for its topology is automatically correct. With this option, the unitarity constraint reads
$$S_{ap}^2+𝑷^2=f_\pi ^2\eta ^2,$$
(30)
where
$$\eta =\sqrt{1+\left[\left(𝝅_y𝝅_w\right)^2𝝅_y^2𝝅_w^2\right]/f_\pi ^4}$$
(31)
and the pion field becomes in fact $`𝑷/\eta `$. This form for the dynamical variable is the same as that proposed by Nyman and Riska, in their symmetrized product ansatz . This ansatz has a topology similar to the PA, as illustrated in fig. 3. The corresponding baryon number density is given in appendix A and, in the classical case, yields $`B=2`$ for all distances when integrated over space, as shown in fig. 5.
## 4 Results and conclusions
In order to derive the potential, we use the quantized fields of eq. (11), obtained by rotating the constituent skyrmions. This idea of rotating individual hedgehogs corresponds to an approximation and deserves some attention. The quantization of a hedgehog, as discussed in sect. 2, amounts to multiplying the classical field by the matrix $`D`$, which depends on three free parameters. In the case $`B=1`$, this procedure does not change the baryon current. This can be seen by writing the baryon density for quantized fields as
$$B^0=\frac{1}{12\pi ^2}ϵ_{abc}ϵ_{\alpha \beta \gamma }\frac{1}{\sigma }_aD_{\alpha i}\pi _i_bD_{\beta j}\pi _j_cD_{\gamma k}\pi _k$$
(32)
and using the result
$$D_{\alpha a}D_{\beta b}=\frac{1}{3}\delta _{\alpha \beta }\delta _{ab}+\frac{1}{2}ϵ_{\alpha \beta \gamma }D_{\gamma c}ϵ_{cab}+\text{isotensors}$$
(33)
in order to obtain
$$B^0=\frac{1}{12\pi ^2}ϵ_{ijk}ϵ_{abc}\frac{1}{\sigma }_a\pi _i_b\pi _j_c\pi _k.$$
(34)
This shows that the baryon density is the same for both quantized and classical fields.
Analogously, in the case $`B=2`$, quantization would require a matrix $`\overline{D}`$, depending on six collective coordinates. However, the determination of this general matrix may prove to be very difficult and, in the spirit of the the sudden approximation, one normally uses $`\overline{D}D^{(y)}I^{(w)}+I^{(y)}D^{(w)}`$, where $`I`$ is an identity matrix and $`D^{(y)},D^{(w)}`$ are operators over the skyrmions labeled by $`𝒚`$ and $`𝒘`$ respectively. The price one pays for this approximation is that it leads to a quantized baryon number which depends on $`D^{(y)}`$ and $`D^{(w)}`$. This happens because the relation equivalent to eq. (33) does not hold for the approximate matrix $`\overline{D}`$ and hence does not represent a major shortcoming for the symmetrized ansatz (SA). Indeed, as pointed out by Sternheim and Kälbermann , this poses problems for short distances only.
A collective rotation of the pion field $`𝑷`$, as in the $`B=1`$ case, would leave $`S/\eta `$ unmodified, as a classical function. However, this would also mean to treat the scalar product $`𝝅_y𝝅_w`$ as a classical quantity and would lead to serious contradictions, for the OPEP content of the isospin dependent channels relies on the quantum character of such a scalar product in eq. (25).
Therefore, the individual rotation of each pion field is more consistent with a constructive approach, although not free of problems. In principle, every pion field $`𝝅`$ in the composite Skyrme Lagrangian should be quantized. When applying this prescription to the dynamical variable of the SA, one has to deal with the functions $`\eta ^2`$ and $`\eta ^4`$, which depend on pion fields, coupled to operators $`D`$. The meaning of the quantized $`\eta ^2`$ is that of a power series in $`D`$, which involves products of arbitrary numbers of these matrices and hence can only be handled by resorting to truncation. With this limitation in mind, we treat $`\eta ^2`$ as a polynomial in $`𝝅`$ and thus its expectation value between two-nucleon states can be evaluated without ambiguities.
In order to test the implications of this assumption, in the sequence we present results with two versions of $`\eta `$, namely, a classical one,
$$\eta _c^2=\left\{1\left[1(\widehat{𝒚}\widehat{𝒘})^2\right]s_y^2s_w^2\right\}^1,$$
(35)
and a quantized one, truncated at the first order in the $`D`$ expansion, given by
$`NN|\eta _q^2|NN`$ $`=`$ $`NN|\left\{1s_y^2s_w^2+D_{\alpha i}^{(y)}D_{\alpha j}^{(w)}D_{\beta k}^{(y)}D_\beta \mathrm{}^{(w)}\left(\pi _y\right)_i\left(\pi _w\right)_j\left(\pi _y\right)_k\left(\pi _w\right)_{\mathrm{}}/f_\pi ^4\right\}^1|NN`$ (36)
$``$ $`\left\{1{\displaystyle \frac{2}{3}}s_y^2s_w^2\right\}^1.`$
Replacing the pion field of the symmetrized ansatz into the interaction Lagrangian used to calculate the potential, one has
$$_{int}=^S+D_{\alpha m}^{(y)}D_{\alpha n}^{(w)}_{mn}^V,$$
(37)
where the labels $`S`$ and $`V`$ stand respectively for isoscalar and isospin dependent parts of $`_{int}`$. Using this result in eq. (17), one gets
$$V(d)=V_C+𝝉^{(y)}𝝉^{(w)}\left[𝝈^{(y)}𝝈^{(w)}V_{SS}+S_{12}V_T\right],$$
(38)
where $`V_C`$, $`V_{SS}`$ and $`V_T`$ are the usual central, spin-spin and tensor components. All terms receive both G-parity odd and even contributions.
The G-parity odd components of spin-spin and tensor terms of the potential are shown in fig. 6, for the two possible choices of $`\eta `$, compared to the PA and pure OPEP results. One sees that all curves coincide for distances larger than 2 fm, indicating that all ansätze reproduce asymptotically the OPEP. The results for the G-parity even terms are of minor importance here, as we are interested in the long range behaviour of the potential, but they are included for completeness in fig. 7. One should note that in the case of $`\eta _c`$, the potentials present a singularity at $`d0.6`$ fm, due to a root of $`\eta _c`$.
Results for the scalar component $`V_C`$ are presented in figs. 8 and 9. In the former we display the behaviour of the PA (left) and the predictions from the SA with $`\eta =1`$ (right), which is non-unitary and considered just for pedagogical purposes. Inspecting it one learns that the SA includes a contribution from $`_2`$, that was not present in the PA. Moreover the contribution from $`_4`$ is negative, and so is the net result for $`V_C`$.
The two valid options for the SA considered here, based on $`\eta _c`$ and $`\eta _q`$, are given in fig. 9. In both cases we observe that the unitarity constraint restores the repulsion due to $`_4`$, but in such a way that the net result is asymptotically attractive. On the other hand, the amount of overall attraction found in the SA depends on the specific quantization prescription adopted. At very large distances, the curves corresponding to $`\eta _c`$ and $`\eta _q`$ have the same geometry and yield respectively the following approximate values for the intensity of the potential: $`K_c=14`$ MeV and $`K_q=57`$ MeV. Comparing them with the empirical values in table 1, it is possible to see that predictions from the SA are qualitatively reasonable.
In summary, we have shown that the SA provides the correct baryon number for the two-nucleon system in the Skyrme model, as well as attractive central $`NN`$ potentials. The correct quantitative feature is somewhere between the values obtained for the two versions of the normalization function $`\eta `$.
We conclude that it is indeed relevant to the central potential to eliminate the term with a wrong G-parity from $`𝑷_{pa}`$. We expect that a deeper and more careful study of the quantization procedure will lead to a more accurate evaluation of the amount of attraction coming from the SA.
## Acknowledgements
The work by I.P.C. was supported by Fundação de Amparo à Pesquisa do Estado de São Paulo - FAPESP (Brazilian agency), grant 94/1801-4. We would like to thank the hospitality of the Nuclear Theory Group, in the University of Washington, and of the Division de Physique Théorique, in the Institute de Physique Nucléaire/Orsay, where part of this work was performed.
## Appendix A Appendix: The baryon number
The explicit calculation of the zero-component of the baryon current yields
$$B^0=B_1^0+B_2^0,$$
(39)
where
$`B_1^0(𝒓;d)={\displaystyle \frac{1}{2\pi ^2\eta ^2}}\left(c_y{\displaystyle \frac{s_w}{w}}+c_w{\displaystyle \frac{s_y}{y}}\right)^2\left(F_y^{}+F_w^{}\right),`$
$`B_2^0(𝒓;d)={\displaystyle \frac{1}{2\pi ^2}}{\displaystyle \frac{1}{\left(c_yc_w\left(\widehat{𝒚}\widehat{𝒘}\right)s_ys_w\right)\eta ^2}}\left(c_y{\displaystyle \frac{s_w}{w}}+c_w{\displaystyle \frac{s_y}{y}}\right)`$
$`\{(1(\widehat{𝒚}\widehat{𝒘})^2)[c_yc_w(c_yF_y^{}{\displaystyle \frac{s_y}{y}})(c_wF_w^{}{\displaystyle \frac{s_w}{w}})s_y^2s_w^2F_y^{}F_w^{}]`$
$`\left[{\displaystyle \frac{s_ys_w}{\eta ^2}}J_y((c_y{\displaystyle \frac{s_w}{w}}+c_w{\displaystyle \frac{s_y}{y}})(c_ws_y+(\widehat{𝒚}\widehat{𝒘})c_ys_w)+(1(\widehat{𝒚}\widehat{𝒘})^2)s_yc_y(F_w^{}c_w{\displaystyle \frac{s_w}{w}}))\right][yw]\},`$
with $`J_y=\left(\left(\widehat{𝒚}\widehat{𝒘}\right)^21\right)c_yF_y^{}s_w+\left(\widehat{𝒚}\widehat{𝒘}\right)\left(s_y\frac{s_w}{w}\left(\widehat{𝒚}\widehat{𝒘}\right)s_w\frac{s_y}{y}\right)`$.
The numerical integration of $`B_2^0`$ is tricky due to the presence of the function $`S`$ in the denominator. Results are shown in fig. 5, as functions of separation distance $`d`$. It shows that the SA presents the correct topology for the $`NN`$ system.
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# Exotic Mesons, Theory and Experiment Presented at MESON2000 (Cracow, 19-23 May 2000)
## 1 Theoretical Expectations
### 1.1 Exotics Defined
An “exotic meson” has J<sup>PC</sup> or flavor quantum numbers forbidden to the $`|q\overline{q}`$ states of the nonrelativistic quark model.
The current experimental candidates are “spin-parity exotics”, which have J<sup>PC</sup> forbidden to $`q\overline{q}`$ mesons. In principle one might also find flavor exotics in a multiquark sector, for example in I=2, but no such (widely accepted) experimental candidates are known at present .
As a caveat we emphasise that every meson is a linear superposition of all allowed basis states, spanning $`|q^2\overline{q}^2`$, $`|q\overline{q}g`$, $`|gg,\mathrm{}`$ (where not strictly forbidden), with amplitudes that are determined by QCD interactions. For convenience we usually classify resonances as “quarkonia”, “hybrids”, “glueballs” and so forth, and are implicitly assuming that one type of basis state dominates the state expansion of each resonance. Of course this may not be the case in general, and the amount of “configuration mixing” is an open and rather controversial topic in hadron physics. Exotics are special in that the $`|q\overline{q}`$ component must be zero, due to the quantum numbers of the state.
### 1.2 What became of multiquarks?
Multiquark systems such as $`q^2\overline{q}^2`$ were once expected to contribute a rich spectrum of resonances to the meson spectrum, and in the 1970s there were many detailed calculations of the spectrum of multiquark resonances in various models. Now one hears little about this subject. What became of multiquarks?
The answer is that they “fell apart”. Even in the early work on $`q^2\overline{q}^2`$ multiquarks it was realized that their decay couplings would be very different from conventional $`q\overline{q}`$ mesons; the latter decay mainly through the production of a second $`q\overline{q}`$ pair, whereas the $`q^2\overline{q}^2`$ system can simply be rearranged into a state of two $`(q\overline{q})+(q\overline{q})`$ mesons. If the expected energy of a continuously deformed $`q^2\overline{q}^2(q\overline{q})(q\overline{q})`$ system is monotonically decreasing, one would not expect to find a $`q^2\overline{q}^2`$ resonance. This was the situation found variationally in the scalar sector by Weinstein and Isgur for most light quark masses .
Life can be more complicated, and Weinstein and Isgur also found that weakly bound deuteronlike $`\mathrm{K}\overline{\mathrm{K}}`$ states existed in their model. Presumably many more such weakly bound quasinuclear states exist, both in meson-meson and meson-baryon sectors. This subject of multihadron systems is at least as rich as the table of nuclear levels.
One often hears that the $`q^6`$ system may have a bound state in the $`u^2d^2s^2`$ I=0, J=0 flavor sector, known as the “H dibaryon” . Caution is appropriate here. Some experiments that are nominally searching for the H dibaryon have “widened their net” to include states very close to $`\mathrm{\Lambda }\mathrm{\Lambda }`$ threshold; if one is found, it would more likely be a weakly bound $`\mathrm{\Lambda }\mathrm{\Lambda }`$ hypernucleus. It is important not to equate these two ideas. A $`\mathrm{\Lambda }\mathrm{\Lambda }`$ hypernucleus would certainly be a very interesting discovery, especially in its implications for models of the intermediate ranged baryon-baryon attraction, but it is not the H dibaryon envisaged in bag model calculations. The H dibaryon calculations assumed an SU(3) flavor-singlet $`u^2d^2s^2`$ system, and $`\mathrm{\Lambda }\mathrm{\Lambda }`$ is a quite different flavor state. Quark model calculations actually find the $`\mathrm{\Lambda }\mathrm{\Lambda }`$ interaction to be repulsive, rather like the NN core.
It does appear likely that real multiquark $`Q^2\overline{q}^2`$ clusters will exist given sufficiently heavy quarks ($`Q=c`$, or perhaps only $`b`$) , but these are unfortunately not easily accessible to experiment.
### 1.3 Hybrid mesons
A hybrid meson is usually “defined” as a resonance whose dominant valence component is $`|q+\overline{q}+excitedglue`$. This deliberately vague definition covers our present ignorance over how one can most accurately describe gluonic excitation. Possibilities include models with explicit transverse gluon quanta, such as the bag model, as well as an excitation of the flux tube that one sees in LGT simulations. Fortunately for experimenters, many of the model calculations reach rather similar predictions for the properties of these states.
One general conclusion is that, unlike $`q\overline{q}`$, all J<sup>PC</sup> can be constructed from hybrid basis states. This conclusion is seen most rigorously in the list of gauge-invariant local operators that one may construct from a product of $`\overline{\psi },\psi `$ and $`F_{\mu \nu }^a`$, since one may couple to physical states by operating on the vacuum $`|0`$ with such an “interpolating field”. This list of $`\overline{\psi }\psi F`$ operators covers all J<sup>PC</sup>, including the so-called exotic combinations
$$\mathrm{J}^{PC}|_{exotic}=0^{},0^+,1^+,2^+,3^+,\mathrm{}$$
that one cannot construct from a $`\overline{\psi }\psi `$ quarkonium operator.
At present the experimental J<sup>PC</sup> exotics are usually considered to be hybrid meson candidates, simply because theorists know of no other general class of J<sup>PC</sup> exotic resonance, excepting multiquark systems that purportedly “fall apart” into light $`q\overline{q}`$ mesons. (Possible exceptions which merit future investigation are weakly bound quasinuclear states, which might exist near threshold in S-wave in attractive meson-meson channels.)
In any case, if a J<sup>PC</sup> exotic meson is found, we can be certain that we have discovered something beyond the naive quark model. This is an extremely important possibility experimentally, and assuming that such states are clearly identified we may hope that the pattern of their spectroscopy will eventually make it clear just what has been discovered!
## 2 Specific models of hybrids
### 2.1 Introduction
Much of the work on hybrids has made use of very specific models of “excited glue”. These models are the bag model, the flux tube model, and the rather underexplored constituent gluon model. Finally, masses and other properties of J<sup>PC</sup> exotic hybrids may be predicted by QCD sum rules and LGT using interpolating fields, and these approaches do not make model assumptions about the nature of gluonic excitation. We will discuss some of the more fundamental results of these models, especially as regards masses, quantum numbers and decay properties.
### 2.2 Bag model hybrids
Many early hybrids studies used the bag model, which assumed that quarks and gluons could be treated as spherical cavity modes of Dirac and Maxwell quanta, confined within the cavity by the choice of color boundary conditions. The “zeroth-order” bag model states were color singlet product basis states of quark, antiquark and gluon modes, for example
$$|q\overline{q},|q\overline{q}g,|gg,|q^2\overline{q}^2,\mathrm{}.$$
The quark-gluon and gluon self interactions mixed these basis states, so that the physical levels were linear combinations of these “bare” basis states, just as we anticipated in our introduction. The distinction between “conventional $`q\overline{q}`$ meson” and “nonexotic hybrid” in the bag model was thus rather vague, and was clearest as a theoretical identification as the strength of the QCD coupling constant was made small. The bag model gave a rather good description of the light “conventional $`q\overline{q}`$ meson” spectrum as $`|q\overline{q}+O(\sqrt{\alpha _s})|q\overline{q}g`$ states, and the hybrids appeared as an extra set of $`|q\overline{q}g+O(\sqrt{\alpha _s})(|q\overline{q}+|q\overline{q}g^2+\mathrm{})`$ states, which should appear as an “overpopulation” of the experimental meson spectrum relative to the naive $`q\overline{q}`$ quark model.
In the bag model the lowest quark mode is a conventional J$`{}_{}{}^{P}=1/2^+`$, but the lowest gluon mode is a (perhaps surprising) J$`{}_{}{}^{P}=1^+`$ TE gluon. Combining these lowest lying $`q,\overline{q}`$ and $`g`$ modes, one finds hybrid basis states with
$$\mathrm{J}^{PC}|_{bagmodelhybrids}=(0^{},1^{})1^+=1^{},0^+,1^+,2^+.$$
Thus the bag model predicts that the lowest lying hybrid multiplet should consist of these 4 J<sup>PC</sup>, of which the $`1^+`$ combination is exotic. Hybrid mass estimates required detailed calculations in which configuration mixing with the other quark+gluon basis states was included to $`O(\sqrt{\alpha _s})`$, and the resulting truncated Hamiltonian was diagonalized. The results depended somewhat on the bag model parameters assumed, with masses of $`1.5`$ GeV being typical . Spin dependent splittings ordered the levels as $`0^+<1^+<1^{}<2^+`$, with a total multiplet splitting of ca. 500 MeV with the usual bag model parameters. Since each of these J<sup>PC</sup> levels is a flavor nonet in the $`u,d,s`$ system, many hybrid states are predicted that might be experimentally accessible.
One may also form baryon hybrids, since the basis states $`|qqqg`$ contain color singlets. The corresponding bag model calculations of the spectrum of baryon hybrids predict a lowest multiplet of $`u,d`$ “hybrid baryons” with a mean mass of about 2 GeV, and a J<sup>P</sup>, flavor content of $`(1/2^+\mathrm{N})^2`$, $`(3/2^+\mathrm{N})^2`$, $`(5/2^+\mathrm{N})`$, $`(1/2^+\mathrm{\Delta })`$, $`(3/2^+\mathrm{\Delta })`$. The spin-splittings due to quark-gluon and gluon-gluon forces predict rather large overall multiplet splitting of ca. 500 MeV, resulting in a $`(1/2^+\mathrm{N})`$ near 1.5 GeV as the lightest hybrid baryon. This result led to the speculation that the N(1440) Roper might be the lightest hybrid baryon. Of course there are no J<sup>P</sup> exotics in the baryons, since all J<sup>P</sup> can be made from $`qqq`$. Since $`|qqq|qqqg`$ configuration mixing is large in the bag model, the distinction between hybrid and conventional baryons is problematic. One must simply conclude that, in the context of this model, there should be an overpopulation of baryons relative to the simple $`|qqq`$ quark model, due to the presence of the extra $`|qqqg`$ basis states.
### 2.3 Flux-tube model hybrids
In LGT simulations a roughly cylindrical region of chaotic glue fields can be observed between widely separated static color sources. This “flux tube” leads to the confining linear potential between color-singlet $`q`$ and $`\overline{q}`$ that is familiar from quark potential models. The “flux tube model” is an approximate description of this state of glue, which is treated as a string of point masses “beads” connected by a linear potential. This system is treated quantum mechanically, and has normal modes of excitation which are transverse to the axis between the (fixed) endpoints of the string. The orbital angular momentum of a transverse string excitation may be combined with the $`q\overline{q}`$ spin and orbital angular momentum using rigid body wavefunctions, which leads to predictions for the quantum numbers of these flux-tube hybrids. Since the assumption about the nature of excited glue is quite different from the $`1^+`$ TE gluon mode of the bag model, one finds a different spectrum of hybrid states. The lowest flux-tube hybrids are predicted to span 8 J<sup>PC</sup> levels, all degenerate in the simplest version of the model, with
$$\mathrm{J}^{PC}|_{fluxtubehybrids}=0^\pm ,1^\pm ,2^\pm ;1^{\pm \pm }.$$
The first 6 of these levels have $`S_{q\overline{q}}=1`$ and the last 2 have $`S_{q\overline{q}}=0`$.
In the earliest mass estimates in the flux tube model various approximations were made, such as a small oscillation approximation and an adiabatic quark motion approximation. After several studies of this system, Isgur, Kokoski and Paton reached their well known estimate of 1.9(1) GeV for the mass of this lightest hybrid multiplet. This work has since been improved upon by Barnes, Close and Swanson using a Hamiltonian Monte Carlo algorithm that does not make the small oscillation and adiabatic approximations. It appears that these approximations gave opposite and comparable mass shifts, so their final result was a very similar 1.8-1.9 GeV for this lightest hybrid multiplet.
Since each of these 8 J<sup>PC</sup> levels has a flavor nonet of associated states, the flux tube model predicts a very rich spectrum, with an additional 72 meson resonances expected in the vicinity of 2.0 GeV, in addition to the conventional $`q\overline{q}`$ quark model states!
Finally, the very interesting question of the masses and quantum numbers of hybrid baryons in the flux tube model has only recently been considered, by Capstick and Page .They find that the lightest hybrid baryon multiplet contains degenerate $`(1/2^+\mathrm{N})^2`$ and $`(3/2^+\mathrm{N})^2`$ states at a mass of 1870(100) MeV, with $`(1/2^+\mathrm{\Delta })`$, $`(3/2^+\mathrm{\Delta })`$ and $`(5/2^+\mathrm{\Delta })`$ partners slightly higher in mass. These conclusions are not so different from the bag model, which predicted a similar hybrid baryon content, with a lightest $`(1/2^+\mathrm{N})`$ hybrid near 1.6 GeV. The most obvious distinction (other than the 300 MeV difference in the lightest hybrid baryon’s mass) is the high mass $`(5/2^+\mathrm{N})`$ (bag model) versus $`(5/2^+\mathrm{\Delta })`$ (flux-tube model).
### 2.4 LGT and QCD Sum Rules
These approaches both estimate exotic masses by evaluating correlation functions of the form $`0|𝒪(\stackrel{}{x},\tau )𝒪^{}(0,0)|0`$, where $`𝒪^{}`$ is an operator that can couple to the state of interest from the vacuum. When summed over $`\stackrel{}{x}`$, at large $`\tau `$ this quantity approaches $`\kappa \mathrm{exp}(M_𝒪\tau )`$, where $`M_𝒪`$ is the mass of the lightest state created from the vacuum by the operator $`𝒪^{}`$. Thus by choosing various operators with exotic quantum numbers one may extract mass estimates for the lightest states with those quantum numbers.
Both methods are subject to systematic errors due to approximations. The QCD sum rules relate these operators to calculable pQCD contributions and to VEVs of other operators that are not calculated, but are inferred from experiment. Different choices for these parameters, algebra errors and uncertainties in higher-mass contributions have led to a moderately wide scatter of results. For example, for the $`1^+`$ exotic, which is of greatest phenomenological interest, the earliest work of Balitsky et al. in 1982 estimated a mass of $`1`$ GeV. Subsequently in 1984 Govaerts et al. estimated 1.3 GeV, Latorre et al. estimated 1.7(1) GeV and 2.1 GeV in 1987 . The most recent work of Chetyrkin and Narison finds $`1.6`$-$`1.7`$ GeV, with the radial hybrid only about 0.2 GeV higher. This reference also considers decay couplings; the partial width to $`\pi \rho `$ is found to be about 300 MeV, but to $`\pi \eta ^{}`$ is only about 3 MeV. As we shall see, this is not what has been reported for either experimental exotic $`\pi _1`$ candidate.
Other exotic quantum numbers have been considered in QCD sum rules. For example, the $`0^{}`$ has been considered by several of these references, and is found to have a rather high mass of ca. 3 GeV.
Recently LGT groups have presented results for the masses of exotic mesons. The MILC collaboration gave results for light $`1^+`$ and $`0^+`$ exotics, and and UKQCD considered these and $`2^+`$ as well. (These are the three exotics predicted to be lightest, and degenerate, in the zeroth-order flux tube model.) At present the LGT results appear consistent with the expectations of the flux-tube model; signals in all these channels are observed, with the mass of the $`1^+`$ (the best determined) being about 2.0(1) GeV. The $`0^+`$ and $`2^+`$ may lie somewhat higher, but this is unclear with present statistics.
The application of LGT to nonrelativistic heavy quark systems has been the topic of much recent research, and considerably smaller statistical errors follow from the use of a QCD action derived from a heavy quark expansion. This NRQCD has been applied to $`1^+`$ heavy-quark exotic hybrids, with very interesting results; the $`1^+`$ $`b\overline{b}`$-H is predicted to lie at $`10.99(1)`$ GeV, and the $`1^+`$ $`c\overline{c}`$-H charmonium hybrid is predicted to lie at $`4.39(1)`$ GeV. With such small statistical errors in these heavy hybrid mass estimates, there is strong motivation for a careful, high statistics scan of $`R`$ near these masses, since models of hybrids anticipate that the multiplet containing the $`1^+`$ will also possess a $`1^{}`$ state nearby in mass.
## 3 Hybrid decays
There appears to be universal agreement that hybrids should exist, and that the lightest of these states with $`u,d`$ quarks should include a $`1^+`$ resonance with a mass in the $`1.5`$-$`2`$ GeV region, with the higher mass preferred by LGT and the flux tube model. For the experimental detection of these states we are faced with the crucial question of what their strong decay properties are. In the worst case they might be so broad as to be difficult to identify, a problem familiar from the $`f_0`$ sector.
Several models of strong decays have been applied to hybrids, and their results have motivated and directed experimental studies. The best known is the flux-tube decay model, which was applied to exotic hybrids by Isgur, Kokoski and Paton and subsequently to nonexotic hybrids by Close and Page . This model assumes that decays take place by <sup>3</sup>P<sub>0</sub> $`q\overline{q}`$ pair production along the length of the flux tube. For the unexcited flux tubes of conventional mesons the predictions are quite similar to the conventional <sup>3</sup>P<sub>0</sub> model; for hybrids in the flux tube model this decay assumption allows the calculation of hybrid meson decay amplitudes.
The orbital angular momentum gives the $`q\overline{q}`$ source produced during a decay a phase dependence around the original $`q\overline{q}`$ axis, and the hadronic final states produced are those which have similar angular dependence. Naively favorable modes such as $`\pi \pi `$, $`\rho \pi `$ and so forth are predicted to be produced quite weakly due to poor spatial overlap with this $`\mathrm{exp}(i\varphi )`$-dependent $`q\overline{q}`$ source. The favored modes are those that have a large L$`{}_{z}{}^{}=1`$ axial projection, such as an S+P meson pair. This is the origin of the flux-tube S+P selection rule, which in the I=1 $`1^+`$ case favors the unusual modes $`\pi f_1`$ and $`\pi b_1`$ over $`\eta \pi `$, $`\eta ^{}\pi `$ and $`\rho \pi `$, despite their greater phase space.
In addition to the flux tube decay calculations, there are also QCD sum rule results (cited above), a decay model that assumes a specific relation between the flux tube excitation and the color vector potential , and constituent gluon decay amplitude calculations . There is general agreement (with some variation between models) that in most cases the flux-tube result of S+P mode dominance in hybrid strong decays is correct.
## 4 Experimental exotic meson candidates
### 4.1 Introduction
Since there are only two experimental candidiate exotic meson resonances, the $`\pi _1(1400)`$ and the $`\pi _1(1600)`$, this section is relatively brief. I will first review the better established $`\pi _1(1600)`$, and then discuss the $`\pi _1(1400)`$. Both resonances are reported to have rather different properties than theorists expected for exotic hybrid mesons. Although I will compare these experimental exotic candidiates to theoretical predictions for exotic hybrids, they might of course be another kind of non-$`q\overline{q}`$ state or even a misinterpreted nonresonant scattering effect.
### 4.2 $`\pi _1(1600)`$
The best established exotic candidate is the $`\pi _1(1600)`$. Evidence for this state has been reported in three channels, $`b_1\pi `$ (VES ), $`\eta ^{}\pi `$ (VES ) and $`\rho \pi `$ (VES , E852 at BNL ). The $`\rho \pi `$ channel is the least controversial, since there are two independent experiments involved, and clear resonant phase motion against several well-established $`q\overline{q}`$ states is evident. The mass and width reported by VES and BNL are consistent,
$$M_{\pi _1}=\{\begin{array}{cc}1.61(2)\mathrm{GeV}\hfill & \text{VES, all three modes}\hfill \\ 1.593\pm 0.008\genfrac{}{}{0pt}{}{+0.029}{0.047}\mathrm{GeV}\hfill & \text{BNL E852, }\rho \pi \text{,}\hfill \end{array}$$
(1)
$$\mathrm{\Gamma }_{\pi _1}=\{\begin{array}{cc}0.29(3)\mathrm{GeV}\hfill & \text{VES, all three modes}\hfill \\ 0.168\pm 0.020\genfrac{}{}{0pt}{}{+0.150}{0.012}\mathrm{GeV}\hfill & \text{BNL E852, }\rho \pi \text{.}\hfill \end{array}$$
(2)
One difficulty with interpreting this state as a hybrid is the $`300`$-$`400`$ MeV mass difference between flux-tube and LGT estimates of $`M1.9`$-$`2.0`$ GeV and the $`\pi _1(1600)`$ mass.
The evidence for the $`\pi _1(1600)`$ in VES data in the three channels $`b_1\pi `$, $`\eta ^{}\pi `$ and $`\rho \pi `$, is shown in Fig.2. The relative branching fractions reported by VES for these final states are
$$\mathrm{\Gamma }(\pi _1(1600)f)=\{\begin{array}{cc}1\hfill & b_1\pi \hfill \\ 1.0\pm 0.3\hfill & \eta ^{}\pi \hfill \\ 1.6\pm 0.4\hfill & \rho \pi \text{ .}\hfill \end{array}$$
(3)
We can see immediately that there is a serious problem here, as the reported relative branching fractions are inconsistent with the predictions of the flux-tube model for hybrid decays (Table 1). The theoretical expectation is that $`b_1\pi `$ should be dominant, with $`\rho \pi `$ weak and $`\eta \pi `$ and $`\eta ^{}\pi `$ very small . Some $`\rho \pi `$ coupling is expected in the flux tube model due to different $`\rho `$ and $`\pi `$ spatial wavefunctions , but this is expected to be a much smaller effect in the $`\eta \pi `$ and $`\eta ^{}\pi `$ modes. Indeed, there is a generalized G-parity argument that says these would be zero except differences in spatial wavefunctions. Either these three modes are not all due to a hybrid exotic, or our understanding of hybrid decays is inaccurate. Future experimental studies of $`\rho \pi `$ will be especially interesting here, since this channel is easily accessible for example in photoproduction at the planned HallD facility at Jefferson Lab .
### 4.3 $`\pi _1(1400)`$
This state is reported in $`\eta \pi `$, which is a channel with a long and complicated history. Prima facie $`\eta \pi `$ appears to be a very attractive channel in which to search for exotics, because there is no spin degree of freedom, and all the odd-L $`\eta \pi `$ channels are J<sup>PC</sup>-exotic.
The $`\eta \pi ^o`$ channel was studied by GAMS in 1980 before the idea of J<sup>PC</sup> exotic resonances was widely accepted, and this collaboration was rather surprised to find a significant (exotic) P-wave. Of course the question was whether this exotic partial wave showed resonant phase motion or was simply a nonresonant background; GAMS had insufficient statistics to decide, but speculated that it was probably nonresonant. In a subsequent 1988 study of $`\pi ^{}p\eta \pi ^on`$ they concluded that $`\eta \pi ^o`$ did indeed support a 1.4 GeV $`1^+`$ exotic P-wave resonance, although their analysis does not agree with subsequent studies of the same $`\pi ^{}p\eta \pi ^on`$ reaction. This was followed by a KEK experiment that concluded that $`\eta \pi ^{}`$ did show evidence for a resonant P-wave, albeit with a mass and width consistent with the $`a_2(1320)`$. Since the D-wave $`a_2(1320)`$ dominated this reaction, there were concerns that the reported exotic P-wave was actually due to feedthrough of the large $`a_2(1320)`$ signal in the partial wave analysis. This is now believed to be the case, perhaps due to the angular asymmetry of the detector. A subsequent VES experiment also found a resonant signal in the exotic $`\eta \pi ^{}`$ P-wave , but at a rather higher mass. Their $`\pi _1(1400)`$ was confirmed in 1997 by BNL experiment E852 . Finally, the Crystal Barrel Collaboration also find that their $`\eta \pi \pi `$ Dalitz plots in both charged ($`\eta \pi ^o\pi ^{}`$) and neutral ($`\eta \pi ^o\pi ^o`$) final states show evidence for a broad resonant P-wave exotic, and fits give a $`\pi _1(1400)`$ mass and width consistent with VES and BNL. The BNL and Crystal Barrel masses and widths are
$$M_{\pi _1}=\{\begin{array}{cc}1.370\pm 0.016\genfrac{}{}{0pt}{}{+0.050}{0.030}\mathrm{GeV}\hfill & \text{BNL E852, }\eta \pi ^{}\hfill \\ 1.400\pm 0.02\pm 0.020\mathrm{GeV}\hfill & \text{C.Bar, neutral and charged }\eta \pi \hfill \\ 1.360\pm 0.025\mathrm{GeV}\hfill & \text{C.Bar, neutral }\eta \pi \text{,}\hfill \end{array}$$
(4)
$$\mathrm{\Gamma }_{\pi _1}=\{\begin{array}{cc}0.385\pm 0.040\genfrac{}{}{0pt}{}{+0.065}{0.105}\mathrm{GeV}\hfill & \text{BNL E852, }\eta \pi ^{}\hfill \\ 0.310\pm 0.050\genfrac{}{}{0pt}{}{+0.050}{0.030}\mathrm{GeV}\hfill & \text{C.Bar, neutral and charged }\eta \pi \hfill \\ 0.220\pm 0.090\mathrm{GeV}\hfill & \text{C.Bar, neutral }\eta \pi \text{.}\hfill \end{array}$$
(5)
There has been much concern expressed regarding various possible experimental problems with this rather light and broad $`\pi _1(1400)`$, for example the size and energy dependence of backgrounds, and possible nonresonant inelastic scattering mechanisms that might mimic a resonance . One should certainly be extremely careful in establishing the lightest exotic meson. Nonetheless it is evident that VES, BNL and Crystal Barrel have all found evidence for a $`\pi _1(1400)`$ with comparable mass and width in $`\eta \pi `$, despite theoretical expectations that a light, broad exotic hybrid should not exist.
## 5 Acknowledgements
It is a great pleasure to thank the organisers for their kind invitation to discuss exotics at Meson 2000. This research was supported in part by the DOE Division of Nuclear Physics, at ORNL, managed by UT-Battelle, LLC, for the US Department of Energy under Contract No. DE-AC05-00OR22725, and by the Deutsche Forschungsgemeinschaft DFG at the University of Bonn and the Forschungszentrum Jülich under contract Bo 56/153-1.
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# 1 Introduction
## 1 Introduction
Non-commutative field theories, open string theories and open membrane theories have been extensively studied recently. These theories are realized on the world-volume of Dp branes with a nonzero NS $`B`$ field and M5 branes with nonzero $`C`$ field.
While space non-commutativity can be accommodated within field theory , space-time non-commutativity seems to require string theory for consistency .<sup>1</sup><sup>1</sup>1 Dual supergravity descriptions of these theories have been constructed in . Further interesting progress in the study of theories with space-time non-commutativity has been made in . However, it was argued in that light-like non-commutativity (that is, $`[x^\mu ,x^\nu ]=i\theta ^{\mu \nu }`$, where $`\theta ^{\mu \nu }`$ is light-like, e.g. $`\theta ^{+2}0`$) can be realized within field theories. One of the aims of this paper is to construct dual supergravity descriptions of such field theories in various dimensions. The field theories are realized on the world-volume of Dp branes with light-like NS $`B`$ field in a particular decoupling limit.
Of particular interest is a six dimensional field theory with light-like noncommutativity realized on the world-volume of M5 branes with light-like $`C`$ field. This theory is conjectured to have a DLCQ matrix description as a quantum mechanics on the resolved moduli space of instantons with the light-like $`C`$ field corresponding to the resolution parameter. We will construct a dual supergravity description of this theory. We will see that the supergravity descriptions of field theories with light-like non-commutativity are closely related to those without non-commutativity, and we will discuss the implications.
Another class of interesting theories are the non-commutative little string theories. These theories are realized on the world-volume of NS5 branes with light-like RR $`A`$ fields. We will construct dual supergravity descriptions of these theories.
A second aim of this paper is to construct dual supergravity descriptions of theories realized on the worldvolume of NS5 branes, whose excitations include light-open Dp branes (ODp) . Such backgrounds are obtained from NS5 branes with near critical RR fields.
The paper is organized as follows. In section 2 we construct Dp brane, M5 brane and NS5 brane solutions corresponding to theories in various dimensions with light-like non-commutativity. We also discuss some salient aspects of these solutions and what we learn from them about the corresponding field theories. In section 3 we construct supergravity solutions of NS5 branes with RR backgrounds and obtain dual supergravity descriptions of ODp theories. As an application, we compute the absorption cross section of a graviton polarized along the world-volume.
## 2 Branes with light-like background fields
### 2.1 Construction of D-brane solutions
Our aim is to construct supergravity backgrounds with Dp brane charge and a $`B`$-field with components $`B_2`$. We will start with the D3 brane case. The construction can be done by performing an infinite Lorentz boost in the $`x^1`$ direction on the known D3 brane background in the presence of a $`B_{12}`$ field. A finite Lorentz boost gives the following background
$$ds^2=f^{\frac{1}{2}}\left[d\stackrel{~}{x}_0^2+dx_3^2+\frac{f}{H}(d\stackrel{~}{x}_1^2+dx_2^2)\right]+f^{\frac{1}{2}}\left(dr^2+r^2d\mathrm{\Omega }_5^2\right),$$
(1)
$$f=1+\frac{\alpha _{}^{}{}_{}{}^{2}R^4}{r^4},H=1+\frac{\alpha _{}^{}{}_{}{}^{2}R^4}{r^4}\mathrm{cos}^2\alpha ,$$
$$e^{2\varphi }=g_s^2\frac{f}{H},F_{\stackrel{~}{0}\stackrel{~}{1}23r}=\frac{1}{g_s}\mathrm{cos}\alpha \frac{f}{H}_rf^1,$$
$$B=\mathrm{tan}\alpha H^1d\stackrel{~}{x}_1dx_2,A=\frac{1}{g_s}\mathrm{sin}\alpha f^1d\stackrel{~}{x}_0dx_3,$$
(2)
where
$$\stackrel{~}{x}_0=\mathrm{cosh}\gamma x_0\mathrm{sinh}\gamma x_1,\stackrel{~}{x}_1=\mathrm{sinh}\gamma x_0+\mathrm{cosh}\gamma x_1,$$
(3)
or $`\stackrel{~}{x}_+=e^\gamma x_+,\stackrel{~}{x}_{}=e^\gamma x_{}`$, with $`x_\pm =\pm x_0+x_1`$, and $`A`$ is the RR 2-form. Note also that $`F_{\stackrel{~}{0}\stackrel{~}{1}23r}=F_{0123r}`$.
We obtain
$$ds^2=f^{\frac{1}{2}}\left[dx_0^2+dx_1^2+dx_3^2+\frac{f}{H}dx_2^2+\frac{\alpha _{}^{}{}_{}{}^{2}R^4}{Hr^4}\mathrm{sin}^2\alpha (\mathrm{cosh}\gamma dx_1\mathrm{sinh}\gamma dx_0)^2\right]$$
$$+f^{\frac{1}{2}}\left(dr^2+r^2d\mathrm{\Omega }_5^2\right),$$
(4)
$$e^{2\varphi }=g_s^2\frac{f}{H},F_{0123r}=\frac{1}{g_s}\mathrm{cos}\alpha \frac{f}{H}_rf^1,$$
$$B_{02}=\mathrm{tan}\alpha \mathrm{sinh}\gamma H^1,B_{12}=\mathrm{tan}\alpha \mathrm{cosh}\gamma H^1,$$
(5)
$$A_{03}=\frac{1}{g_s}\mathrm{sin}\alpha \mathrm{cosh}\gamma f^1,A_{13}=\frac{1}{g_s}\mathrm{sin}\alpha \mathrm{sinh}\gamma f^1.$$
(6)
The fields take the following asymptotic values
$$B_{02}^{\mathrm{}}=\mathrm{tan}\alpha \mathrm{sinh}\gamma E,B_{12}^{\mathrm{}}=\mathrm{tan}\alpha \mathrm{cosh}\gamma B,$$
(7)
$$A_{03}^{\mathrm{}}=\frac{1}{g_s}\mathrm{sin}\alpha \mathrm{cosh}\gamma =\frac{\sqrt{B^2E^2}\mathrm{cosh}\gamma }{g_s\sqrt{1+B^2E^2}},$$
$$A_{13}^{\mathrm{}}=\frac{1}{g_s}\mathrm{sin}\alpha \mathrm{sinh}\gamma =\frac{\sqrt{B^2E^2}\mathrm{sinh}\gamma }{g_s\sqrt{1+B^2E^2}},$$
where we have used $`B^2E^2=\mathrm{tan}^2\alpha `$. To have only light-like B-field components, we now take the infinite boost limit, $`\gamma \mathrm{}`$. At the same time, we must take the limit $`\alpha 0`$ with
$$e^\gamma \mathrm{tan}\alpha =\mathrm{finite}b.$$
In this limit the asymptotic values of the gauge fields simply become
$$B_{02}^{\mathrm{}}=E=b,B_{12}^{\mathrm{}}=B=b,A_{03}^{\mathrm{}}=\frac{b}{g_s}=A_{13}^{\mathrm{}}.$$
(8)
We obtain the background
$$ds^2=f^{\frac{1}{2}}\left[dx_+dx_{}+dx_2^2+dx_3^2+\frac{\alpha _{}^{}{}_{}{}^{2}R^4b^2}{r^4f}dx_{}^2\right]+f^{\frac{1}{2}}\left(dr^2+r^2d\mathrm{\Omega }_5^2\right),$$
(9)
$$e^{2\varphi }=g_s^2,F_{0123r}=\frac{1}{g_s}_rf^1,$$
$$B=bf^1dx_{}dx_2,A=\frac{1}{g_s}bf^1dx_{}dx_3.$$
(10)
Note that this is a constant dilaton solution representing a wave travelling on the D3 brane, but with $`B`$ and $`A`$ fields. The same background is obtained by a similar procedure by starting with the purely electric field configuration $`B_{03}0`$ and performing a Lorentz boost. An alternative derivation by dualities is described in subsection 2.3.
Let us consider the decoupling limit $`\alpha ^{}0`$. Before taking the limit, it is convenient to redefine the coordinate $`x_+x_+b^2x_{}`$. The metric (9) takes the form
$$ds^2=f^{\frac{1}{2}}\left[dx_+dx_{}+dx_2^2+dx_3^2\frac{b^2}{f}dx_{}^2\right]+f^{\frac{1}{2}}\left(dr^2+r^2d\mathrm{\Omega }_5^2\right),$$
(11)
We set as usual $`r=\alpha ^{}u`$, with $`u`$ fixed. In addition, in order to have non-vanishing gauge fields $`B`$ and $`A`$ after $`\alpha ^{}0`$, we must rescale $`b`$ by introducing $`\stackrel{~}{b}=\alpha ^{}b`$= fixed. So $`u,R,\stackrel{~}{b},g_s`$ remain fixed. We get the following background
$$ds^2=\alpha ^{}\left(\frac{u^2}{R^2}\left[dx_{}dx_++dx_2^2+dx_3^2\frac{\stackrel{~}{b}^2}{R^4}u^4dx_{}^2\right]+R^2\frac{du^2}{u^2}+R^2d\mathrm{\Omega }_5^2\right),$$
$$e^{2\varphi }=g_s^2,F_{0123r}=\alpha _{}^{}{}_{}{}^{2}\frac{4}{g_sR^4}u^3,R^4=4\pi g_sN=2g_{\mathrm{YM}}^2N,$$
(12)
$$B=\alpha ^{}\stackrel{~}{b}\frac{u^4}{R^4}dx_{}dx_2,A=\alpha ^{}\frac{\stackrel{~}{b}}{g_s}\frac{u^4}{R^4}dx_{}dx_3,x_\pm =x_1\pm x_0.$$
Note that $`g_s`$ was maintained fixed and $`b\mathrm{}`$ in the decoupling limit, in accordance with the field theory analysis of .
The above solutions can be easily generalized to the case of the Dp branes, with $`p=2,\mathrm{},5`$. A similar procedure gives the following solution representing a wave on the Dp brane:
$`ds^2`$ $`=`$ $`f^{\frac{1}{2}}\left[dx_+dx_{}+dx_2^2+\mathrm{}+dx_p^2{\displaystyle \frac{b^2}{f}}dx_{}^2\right]+f^{\frac{1}{2}}\left(dr^2+r^2d\mathrm{\Omega }_{8p}^2\right),`$ (13)
$`e^{2\varphi }`$ $`=`$ $`g_s^2f^{\frac{3p}{2}},f=1+{\displaystyle \frac{c_p\alpha _{}^{}{}_{}{}^{5p}g_{\mathrm{YM}}^2N}{r^{7p}}},c_p=2^{72p}\pi ^{\frac{93p}{2}}\mathrm{\Gamma }\left({\displaystyle \frac{7p}{2}}\right),`$ (14)
$`B_2`$ $`=`$ $`bf^1,A_3\mathrm{}^{(p1)}={\displaystyle \frac{b}{g_s}}f^1`$ (15)
where $`A^{(p1)}`$ is an RR $`(p1)`$-form. Note that for D2-brane we have a one-form $`A_{}^1`$. There is also the usual RR form which gives the Dp-brane charge. The gauge coupling is
$$g_{\mathrm{YM}}^2=(2\pi )^{p2}g_s\alpha _{}^{}{}_{}{}^{\frac{p3}{2}}.$$
(16)
The decoupling limit is obtained by rescaling variables as follows:
$$r=\alpha ^{}u,\stackrel{~}{b}=\alpha ^{}b,$$
and taking the limit $`\alpha ^{}0`$ with $`u,\stackrel{~}{b},g_{\mathrm{YM}}^2`$ fixed.
We get
$$ds^2=\alpha ^{}\frac{u^{\frac{p3}{2}}}{\sqrt{\lambda }}\left(u^{5p}\left[dx_{}dx_++dx_2^2+\mathrm{}+dx_p^2\stackrel{~}{b}^2\frac{u^{7p}}{\lambda }dx_{}^2\right]+\lambda \frac{du^2}{u^2}+\lambda d\mathrm{\Omega }_{8p}^2\right),$$
$$e^\varphi =g_{\mathrm{YM}}^2\frac{u^{\frac{1}{4}(p3)(7p)}}{(2\pi )^{p2}\lambda ^{\frac{1}{4}(p3)}}=\frac{1}{N}\frac{(\lambda u^{p3})^{\frac{1}{4}(7p)}}{c_p(2\pi )^{p2}},\lambda c_pg_{\mathrm{YM}}^2N,$$
$$B=\alpha ^{}\stackrel{~}{b}\frac{u^{7p}}{\lambda }dx_{}dx_2.$$
(17)
### 2.2 Light-Like Non-commutative SYM
The Dp brane background (17) should provide a dual supergravity description of the $`p+1`$ dimensional light-like non-commutative super Yang-Mills theory, where $`\theta ^{\mu \nu }`$ has only non-vanishing $`\theta ^{+2}=\theta ^{2+}`$ components. Since the parameter $`\stackrel{~}{b}`$ can be scaled away from the metric by $`x_{}x_{}/\stackrel{~}{b}`$, $`x_+x_+\stackrel{~}{b}`$, the curvature invariants are independent of $`\stackrel{~}{b}`$. They are functions only of the coordinate $`u`$, and therefore they must be the same as in the (commutative) $`\stackrel{~}{b}=0`$ case. Thus the regime of validity of supergravity approximation is as in the commutative case.
The curvature of the metric in eq. (17) has the behavior
$$l_s^2\frac{1}{g_{\mathrm{eff}}},$$
(18)
where we have defined a dimensionless effective gauge coupling $`g_{\mathrm{eff}}`$
$$g_{\mathrm{eff}}^2=g_{\mathrm{YM}}^2Nu^{p3}.$$
(19)
As usual, when $`g_{\mathrm{eff}}1`$ the perturbative field theory description is valid, while when $`g_{\mathrm{eff}}1`$ the dual supergravity description is valid. The expressions for the dilaton and curvature thus indicate that the phase structure of the non-commutative light-like theory is similar to that of the ordinary Dp-branes.
Let us now consider the case $`p=3`$ in detail. The D3 brane background (12) provides a dual description of $`3+1`$ dimensional super Yang-Mills theory with light-like non-commutativity. The solution (12) is the first example of a constant dilaton solution describing a non-commutative field theory. Remarkably, this geometry has also constant curvature invariants, e.g. $``$, $`_{\mu \nu }^{\mu \nu }`$, and $`_{\mu \nu \rho \sigma }^2`$, which have just the same values as in the $`AdS_5\times S^5`$ case. The fact that they are constants can be understood from the invariance of the metric under the scaling $`ucu`$, with an appropriate rescaling of $`x_+,x_{},x_2,x_3`$. In particular, although the Weyl tensor has some non-vanishing components, the square Weyl tensor $`C_{\mu \nu \rho \lambda }^2`$ is identically zero. Along with the fact that the dilaton is constant, this suggests a mild energy dependence of the gauge coupling. In the low energy region $`u0`$, the metric approaches $`AdS_5\times S^5`$, which is consistent with the expectation that the low-energy theory should be described by the usual $`𝒩=4`$ SYM theory.
For $`g_s1`$, the dilaton is large and we have to use the S-dual picture. The S-dual background is simply obtained by $`g_s1/g_s`$ and exchanging the gauge fields $`B`$ and $`A`$, which gives $`B_3`$ and $`A_2`$ non-zero components. Therefore the strong coupling limit of SYM theory with light-like non-commutativity is another NCSYM theory with non-commutativity in light-like directions, but with $`\theta ^{+3}`$ instead of $`\theta ^{+2}`$. The exchange of 2-3 directions can be undone by an $`SO(2)`$ rotation in the plane $`x_2`$-$`x_3`$, since light-like NCSYM theory is invariant under such $`SO(2)`$ transformations (see also ). More general $`SL(2,R)`$ transformations only mix the $`B_2`$ and $`A_3`$ components and introduce a constant RR scalar field $`\chi `$ (which in the low-energy Yang-Mills theory gives rise to a $`\mathrm{\Theta }`$ term $`\mathrm{\Theta }\stackrel{~}{F}F`$).
The supergravity background (12) is notably simple, in fact, given by a simple perturbation of the $`AdS_5\times S^5`$ background:
$$g_{\mu \nu }g_{\mu \nu }+\delta g_{\mu \nu },\delta g_{}=\stackrel{~}{b}^2\frac{u^6}{R^6},$$
(20)
$$\delta B_2=\stackrel{~}{b}\frac{u^4}{R^4},\delta A_3=\frac{\stackrel{~}{b}}{g_s}\frac{u^4}{R^4}.$$
(21)
The form of the perturbations (20), (21) implies that as an expansion in the non-commutativity parameter $`\stackrel{~}{b}`$ the first-order perturbation around the $`𝒩=4`$ background is exact. Note in comparison that for spatial non-commutativity, say in coordinates $`x_2,x_3`$, there are infinite number of terms in the expansion in the non-commutativity parameter around the $`𝒩=4`$ background. This suggests that there could exist a simple modification of the $`𝒩=4`$ SYM theory Lagrangian which turns it into a light-like non-commutative SYM theory. The pure gauge theory part of the dimension six SCFT operator corresponding to $`\delta B_2`$ is
$$𝒪_6=\left[F^{2m}F_{mk}F^k+\frac{1}{4}F_{lm}F^{ml}F^2\right].$$
(22)
A question of interest is whether the resulting theory, after adding a perturbation of the form $`𝒪_6\delta B_2`$ to the SYM action, is exactly equivalent to the non-commutative field theory action. This is not the case and one can check that with light-like non-commutativity there are still infinite number of terms. In particular, in the abelian case one can compute the NCSYM action using Seiberg-Witten map , and show that the Lagrangian contains arbitrary powers in $`\theta ^{+2}`$.<sup>2</sup><sup>2</sup>2However, it is interesting to note that in the particular case that $`F_2=0`$, the series truncates and the full Lagrangian contains only the usual SYM action plus a linear term in $`\theta ^{+2}`$.
Thus, while in the perturbative SYM description there are an infinite number of terms, there seems to be a simplification at strong t’ Hooft coupling $`g_sN`$. The S-duality symmetry of the theory is not useful in this limit, since the simplification seems to take place at strong ’t Hooft coupling $`g_sN`$ and large $`N`$, but $`g_s1`$. A possible explanation for the simplification at strong coupling is an existence of a simple resummation of the perturbation series.
The strong coupling expansion is different from that of the ordinary $`𝒩=4`$ SYM theory, since the geometries of the corresponding supergravity backgrounds are different. From the field theory point of view, now the Lagrangian contains a dimension 6 operator, and the resulting theory is expected to be non-renormalizable. In the spatial NCSYM theory it is important to have an infinite series of terms. If there are only a finite number of terms it is hard to see why renormalization works. Renormalizability in the light-like non-commutativity seems to work in the same way as in the spatial NCSYM theory, in the sense that the divergences in the planar diagrams are taken care of as in the SYM case. For the non-planar case, as long as $`(p_0p_1)^20`$ there are no divergences, whereas when $`(p_0p_1)^2=0`$ we get the divergences that are interpreted as IR divergences <sup>3</sup><sup>3</sup>3We would like to thank J. Gomis for a discussion on this point.. It is worth noting that in a gauge $`A_{}=0`$ the new interaction (22) will always involve multiplicative factors $`p_{}`$, which may lead to an improvement of the IR behavior of non-planar diagrams. Clearly, a more detailed analysis is needed in order to understand the structure of perturbation theory for the $`𝒩=4`$ SYM theory with light-like non-commutativity.
### 2.3 M5 brane with light-like C field
To find a solution representing an M5 brane in the presence of a light-like C-field, we can proceed as above and boost the M5 brane solution with $`C_{345}`$ component , in the direction $`x_3`$. This is
$$ds_{11}^2=(kf)^{\frac{1}{3}}\left[\frac{1}{f}\left(d\stackrel{~}{x}_0^2+dx_1^2+dx_2^2\right)+\frac{1}{k}\left(d\stackrel{~}{x}_3^2+dx_4^2+dx_5^2\right)+dr^2+r^2d\mathrm{\Omega }_4^2\right],$$
(23)
$$f=1+\frac{l_p^3R^3}{r^3},k=1+\mathrm{cos}^2\alpha \frac{l_p^3R^3}{r^3},R^3=\frac{\pi N}{\mathrm{cos}\alpha },$$
$$dC_3=\mathrm{sin}\alpha df^1d\stackrel{~}{x}_0dx_1dx_2+\mathrm{cos}\alpha 3R^3l_p^3ϵ_46\mathrm{tan}\alpha dk^1d\stackrel{~}{x}_3dx_4dx_5,$$
(24)
where by $`ϵ_4`$ we denote the volume form of the 4-sphere, and we have made the Lorentz boost
$$\stackrel{~}{x}_0=\mathrm{cosh}\gamma x_0\mathrm{sinh}\gamma x_3,\stackrel{~}{x}_3=\mathrm{sinh}\gamma x_0+\mathrm{cosh}\gamma x_3.$$
(25)
Now we take the limit $`\gamma \mathrm{}`$, $`\alpha 0`$ with $`e^\gamma \mathrm{tan}\alpha =b`$=fixed. We get
$$ds_{11}^2=f^{\frac{1}{3}}[dx_+dx_{}+dx_1^2+dx_2^2+dx_4^2+dx_5^2$$
$$\frac{b^2}{f}dx_{}^2]+f^{\frac{2}{3}}(dr^2+r^2d\mathrm{\Omega }_4^2),$$
(26)
$$dC_3=bdf^1dx_{}dx_1dx_2+3R^3l_p^3ϵ_46bdf^1dx_{}dx_4dx_5.$$
(27)
where we have redefined $`x^+x^+b^2x^{}`$. This represents a gravitational wave moving parallel to the M5 brane.
The decoupling limit is taken by rescaling variables as follows:
$$r=l_p^3u^2,b=l_p^3\stackrel{~}{b}.$$
(28)
and then taking $`l_p0`$ with fixed $`u,R,\stackrel{~}{b}`$. We obtain
$$ds_{11}^2=l_p^2(\frac{u^2}{(\pi N)^{\frac{1}{3}}}[dx_+dx_{}+dx_1^2+dx_2^2+dx_4^2+dx_5^2\frac{\stackrel{~}{b}^2}{\pi N}u^6dx_{}^2]$$
$$+(\pi N)^{\frac{1}{3}}[\frac{4du^2}{u^2}+d\mathrm{\Omega }_4^2]),$$
(29)
$$dC_3=l_p^3\left(\frac{6\stackrel{~}{b}u^5}{\pi N}dudx_{}dx_1dx_2+3\pi Nϵ_4\frac{36\stackrel{~}{b}u^5}{\pi N}dudx_{}dx_4dx_5\right).$$
(30)
At low energy (small $`u`$) the background (29), (30) reduces to $`AdS_7\times S^4`$ as expected. The curvature invariants are the same as in the $`AdS_7\times S^4`$ case (again by virtue of the symmetry under rescalings of $`\stackrel{~}{b}`$, combined with a rescaling of coordinates). The background (29), (30) provides a dual description of the $`(0,2)`$ theory perturbed by a dimension nine operator. This theory is conjectured to have a matrix-like description as the quantum mechanics on the resolved moduli space of instantons. The light-like $`C`$ field is interpreted as the resolution parameters (the Fayet-Iliopoulos parameters of the $`0+1`$ Yang-Mills theory) (see also ).
An alternative derivation of the solution (26) is by using as starting point the supergravity solution of M5-branes in the presence of C field with rank 4
$`ds^2`$ $`=`$ $`h^{\frac{2}{3}}\left[f^{\frac{1}{3}}\left(dx_0^2+hdx_{1,2,3,4}^2+h^2(dx_5Cdx_0)^2\right)+f^{\frac{2}{3}}(dr^2+r^2d\mathrm{\Omega }_4^2)\right],`$ (31)
$`f`$ $`=`$ $`1+{\displaystyle \frac{\pi Nl_p^3}{\mathrm{cos}^2\theta r^3}},h^1=\mathrm{sin}^2\theta f^1+\mathrm{cos}^2\theta ,C=\mathrm{sin}^2\theta f^1,`$ (33)
$`C_{012}`$ $`=`$ $`\mathrm{cos}\theta \mathrm{sin}\theta f^1h,C_{345}=\mathrm{tan}\theta f^1h,`$ (35)
$`C_{034}`$ $`=`$ $`\mathrm{cos}\theta \mathrm{sin}\theta f^1h,C_{125}=\mathrm{tan}\theta f^1h.`$ (37)
Dimensional reduction along the direction $`x_5`$ gives the D4 brane supergravity background in presence of a B-field with magnetic components. The infinite boost limit can be taken by introducing
$$\stackrel{~}{x}_0=x_0\mathrm{cos}\theta ,\stackrel{~}{x}_5=\frac{x_5}{\mathrm{cos}\theta },$$
and taking the limit $`\theta \pi /2`$ with fixed $`\stackrel{~}{x}_0,\stackrel{~}{x}_5,r`$ and fixed $`l_P`$. In this way one reproduces the background (26). The Dp brane backgrounds can then be obtained by dimensional reduction along either of the coordinates $`(1,2,3,4)`$ and T-dualities.
### 2.4 NS5-branes with light-like RR fields
By an S-duality transformation on the D5-brane solution, given by eq. (15) with $`p=5`$, one finds the solution representing type IIB NS5-branes in the presence of a light-like RR $`A`$ field:
$`ds^2`$ $`=`$ $`dx_{}dx_++{\displaystyle \underset{i=2}{\overset{5}{}}}dx_i^2{\displaystyle \frac{b^2}{f}}dx_{}^2+f(dr^2+r^2d\mathrm{\Omega }_3^2),`$ (38)
$`f`$ $`=`$ $`1+{\displaystyle \frac{Nl_s^2}{r^2}},e^{2\varphi }=g_s^2f,l_s\sqrt{\alpha ^{}},`$ (40)
$`A_2`$ $`=`$ $`{\displaystyle \frac{b}{g_s}}f^1,A_{345}={\displaystyle \frac{b}{g_s}}f^1.`$ (42)
Using T-duality one can also find the type IIA NS5-branes with light-like RR fields. T-duality in the direction $`x_2`$ gives an NS5-brane in the presence of light-like RR one- and 4-forms,
$$A_{}=\frac{b}{g_s}f^1,A_{2345}=\frac{b}{g_s}f^1,$$
(43)
with the same metric and dilaton fields. T-duality in the direction $`x_3`$ gives an NS5-brane in the presence of light-like RR three-form, with components
$$A_{23}=\frac{b}{g_s}f^1,A_{45}=\frac{b}{g_s}f^1.$$
(44)
The decoupling limit for these NS5-branes with light-like gauge fields is taken in the same way as that for the usual NS5-branes , namely, $`g_s0`$ and $`l_s`$=fixed, but in addition we have to rescale $`\stackrel{~}{b}=g_sb,r=g_sl_su`$ with fixed $`\stackrel{~}{b}`$ and $`u`$. Setting $`u=N^{\frac{1}{2}}e^{z/r_0}`$, with $`r_0=l_s\sqrt{N}`$, we get for the type IIB NS5-branes (42)
$`ds^2`$ $`=`$ $`dx_{}dx_++{\displaystyle \underset{i=2}{\overset{5}{}}}dx_i^2\stackrel{~}{b}^2e^{2z/r_0}dx_{}^2+dz^2+r_0^2d\mathrm{\Omega }_3^2,`$ (45)
$`A_2`$ $`=`$ $`\stackrel{~}{b}e^{2z/r_0},A_{345}=\stackrel{~}{b}e^{2z/r_0},\varphi ={\displaystyle \frac{z}{r_0}}.`$ (47)
For the Type IIA NS5-branes, the metric and dilaton are the same, and the gauge field components in eqs. (43), (44) become $`A_{}=\stackrel{~}{b}e^{2z/r_0}`$, etc. These backgrounds provide dual descriptions of non-commutative little string theories. The deformation parameters are the light-like RR fields. The phase structure of the theories is the same as the ordinary little string theories. They are characterized by the linear dilaton behavior.
The Type IIB NS5-brane has a DLCQ description as a low energy SCFT of the Coulomb branch of a $`1+1`$ dimensional gauge theory . In this case, the deformation parameters are identified as mass parameters. For the Type IIA NS5-branes, the DLCQ deformation corresponds to turning on a Fayet-Iliopoulos term in the corresponding $`1+1`$ dimensional gauge theory .
## 3 NS5-branes in the presence of RR fields
In this section we will study the theory on the Type II NS5-branes in the presence of different RR field strengths, which can be either electric or magnetic. The supergravity equations of motion require that for NS5-branes in the presence of an RR magnetic (electric) $`(p+1)`$-form there is also a RR electric (magnetic) $`(5p)`$-form with $`p=0,\mathrm{},5`$. Therefore the theory of NS5-branes in the presence of an electric (magnetic) $`(p+1)`$-form is the same as the theory on the NS5-branes in the presence of a magnetic (electric) RR $`(5p)`$-form. For NS5-branes with a RR 3-form there is no difference between electric and magnetic, so there is only one case to be studied.
The supergravity solution for NS5-branes in the presence of an RR $`(p+1)`$\- form is given by
$`ds^2`$ $`=`$ $`h^{1/2}\left[dx_0^2+{\displaystyle \underset{i=1}{\overset{p}{}}}dx_i^2+h{\displaystyle \underset{j=p+1}{\overset{5}{}}}dx_j^2+f(dr^2+r^2d\mathrm{\Omega }_3^2)\right],`$ (48)
$`f`$ $`=`$ $`1+{\displaystyle \frac{Nl_s^2}{\mathrm{cos}\theta r^2}},h^1=\mathrm{sin}^2\theta f^1+\mathrm{cos}^2\theta ,`$ (50)
$`A_{0\mathrm{}p}`$ $`=`$ $`{\displaystyle \frac{\mathrm{sin}\theta }{g_s}}f^1,A_{(p+1)\mathrm{}5}={\displaystyle \frac{\mathrm{tan}\theta }{g_s}}f^1h,e^{2\varphi }=g_s^2fh^{(1p)/2}.`$ (52)
For $`p=5`$, $`A_{(p+1)\mathrm{}5}`$ denotes the RR scalar field. A way to find these solutions is to start with M5-branes in the presence of a $`C`$ field (23), smeared in some transverse direction. By reducing on this transverse direction, one finds the type IIA NS5-branes with an electric RR 3-form. Other solutions are generated by T-duality.
The decoupling limit of the above supergravity solution can be defined as the limit $`l_s0`$, keeping the following quantities fixed:
$$\alpha _{\mathrm{eff}}^{}=\frac{l_s^2}{\mathrm{cos}\theta },u=\frac{r}{l_s^2},\stackrel{~}{g}l_{\mathrm{eff}}^{p3}=g_sl_s^{p3},$$
(53)
$$\stackrel{~}{x}_{0,\mathrm{},p}=\frac{1}{l_{\mathrm{eff}}}x_{0,\mathrm{},p},\stackrel{~}{x}_{(p+1),\mathrm{},5}=\frac{l_{\mathrm{eff}}}{l_s^2}x_{(p+1),\mathrm{},5},l_{\mathrm{eff}}\sqrt{\alpha _{\mathrm{eff}}^{}}.$$
(54)
In this limit the supergravity solution becomes
$`l_s^2ds^2`$ $`=`$ $`(1+a^2u^2)^{1/2}\left[d\stackrel{~}{x}_0^2+{\displaystyle \underset{i=1}{\overset{p}{}}}d\stackrel{~}{x}_i^2+{\displaystyle \frac{_{j=p+1}^5d\stackrel{~}{x}_j^2}{1+a^2u^2}}+{\displaystyle \frac{N}{u^2}}(du^2+u^2d\mathrm{\Omega }_3^2)\right],`$ (55)
$`A_{0\mathrm{}p}`$ $`=`$ $`{\displaystyle \frac{l_s^{(p+1)}}{\stackrel{~}{g}}}a^2u^2,A_{(p+1)\mathrm{}5}={\displaystyle \frac{l_s^{(5p)}}{\stackrel{~}{g}}}{\displaystyle \frac{a^2u^2}{1+a^2u^2}},`$ (57)
$`e^{2\varphi }`$ $`=`$ $`\stackrel{~}{g}^2{\displaystyle \frac{(1+a^2u^2)^{(p1)/2}}{a^2u^2}},a^2={\displaystyle \frac{\alpha _{\mathrm{eff}}^{}}{N}}.`$ (59)
These backgrounds provide a supergravity dual description for the ODp theories investigated in <sup>4</sup><sup>4</sup>4The supergravity description of ODp with $`p=1,2`$ and $`p=2,3`$ has also been considered in and , respectively. (with coupling $`g_{YM}^2=\stackrel{~}{g}l_{\mathrm{eff}}^{p3}`$). The scalar curvature of the metric is given by
$$l_s^2=\frac{1}{N}\frac{c_1+c_2a^2u^2+c_3a^4u^4}{(1+a^2u^2)^{\frac{5}{2}}},$$
(60)
where $`c_1,c_2,c_3`$ are numerical constants depending only on $`p`$. Therefore for large $`N`$ the curvature is small and one can trust the supergravity description.
As an application, let us now consider the absorption cross section of polarized gravitons. This calculation has already been done for the type IIA NS5-branes in the presence of an RR 3-form and type IIB NS5-branes in the presence of a magnetic RR 2-form in , which correspond to the supergravity dual of OD2 and OD3 theories, respectively. In general one can show that in these backgrounds the scattering potential for a graviton polarized along the brane directions is
$$V(\rho )=1+(\frac{3}{4}\omega ^2R^2)\frac{1}{\rho ^2},R^2=\frac{Nl_s^2}{\mathrm{cos}\theta },$$
(61)
where $`\rho =\omega r`$ and $`\omega `$ is the energy of incoming waves. Therefore we see that after the decoupling limit the absorption cross section can be nonzero only for waves with energy $`\omega ^2`$ larger than $`\frac{1}{N\alpha _{\mathrm{eff}}^{}}`$. Essentially the same effect appears in the little string theories. Following , one can see that these theories have a mass gap of order $`M_{\mathrm{gap}}^2\frac{1}{N\alpha _{\mathrm{eff}}^{}}`$ . To compare the decoupling limit (64) with that of ordinary little string theories, it is convenient to describe the decoupling limit in terms of $`g_s`$. For $`p2`$, one can take the decoupling limit of the NS5-branes in the presence of electric $`(p+1)`$-form as follows
$$g_s0,\stackrel{~}{g}^{\frac{2}{3p}}=\frac{g_s^{\frac{2}{3p}}}{\mathrm{cos}\theta },r=g_s^{\frac{1}{3p}}l_su.$$
(62)
with fixed $`l_s,u,\stackrel{~}{g}`$. In this limit the supergravity background (52) reduces to the same expression (59). This description is equivalent to a rescaling of the coordinates. In this way one has $`g_s0`$ and $`l_s`$ fixed, as in the little string theory.
The ODp theories have all the same physics at low energies: for odd $`p`$ they flow to SYM theory in (5+1)-dimension in the IR; for even $`p`$, the theories flow to a fixed point in the IR described by the (0,2) conformal theory.
In the ultraviolet regime, where the effects of nonzero RR fields become important, the different ODp theories exhibit different behaviors, according to the value of $`p`$.
For $`p2`$ case, the string coupling $`e^\varphi `$ in eq. (59) is small in the ultraviolet regime and one can trust the supergravity solution. In this region $`ua^1`$ the NS5 brane supergravity reduces to a metric describing ordinary Dp-branes smeared in $`5p`$ directions. In the particular case of the OD0 theory, the supergravity solution (59) provides a supergravity description of a DLCQ compactification of M-theory with $`N`$ units of DLCQ momentum, in the presence of a transverse M5-brane. The relation between M-theory and type IIA parameters is as follows:
$$\alpha _{\mathrm{eff}}^{}\stackrel{~}{g}^{2/3}=M_{\mathrm{eff}}^2,\alpha _{\mathrm{eff}}^{}\stackrel{~}{g}^2=R_{11}^2,$$
(63)
where $`M_{\mathrm{eff}}`$ if the effective eleven-dimensional Planck mass.
For $`p=3`$ the dilaton in (59) is constant at large $`u`$, i.e. $`e^\varphi =\stackrel{~}{g}`$. For $`\stackrel{~}{g}1`$ the theory can be described by smeared D3-branes, while for $`\stackrel{~}{g}1`$ we have to use the S-dual picture describing D5-branes in the presence of a magnetic $`B`$ field with rank two. Therefore, in the UV regime, strongly coupled OD3 theory and large $`N`$ 5+1 dimensional NCSYM theory exhibit a similar behavior. Note that under S-duality the parameters of the theory change as
$$\alpha _{\mathrm{eff}}^{}\stackrel{~}{g}\alpha _{\mathrm{eff}}^{},\stackrel{~}{g}\frac{1}{\stackrel{~}{g}},$$
(64)
For $`p=4`$ the dilaton is large at $`ua^1`$. In this case it means that the proper supergravity description is in terms of eleven-dimensional supergravity. From M-theory point of view, the supergravity solution is the bound state of two M5-branes in the directions (0,1,2,3,4,5) and (0,1,2,3,4,6) in the decoupling limit.
For $`p=5`$, the dilaton is also large at $`ua^1`$. The S-dual picture is not useful, since the transformed dilaton field $`\varphi ^{}`$ is also large in this regime. Indeed, due to the nonzero RR 0-form (of order one), under S-duality we find $`e^\varphi ^{}\stackrel{~}{g}au`$. This is in agreement with the discussion of . It can be understood from the fact that, under S-duality the system maps to a similar configuration of NS5-branes in the presence of electric RR 6-form. If $`N`$ is the number of NS5-branes and $`M`$ the charge induced by RR 6-form one has the relation $`\frac{1}{\mathrm{cos}\theta }=g_s\frac{M}{N}`$. In the decoupling limit (53), one obtains
$$\stackrel{~}{g}=\frac{N}{M}.$$
(65)
A similar relation is found for D1-branes in the presence of electric B field .
A T-duality transformation on the background (59) implies the following relation between the parameters of ODp theory and OD(p-1) theory:
$$R\frac{\alpha _{\mathrm{eff}}^{}}{R},\stackrel{~}{g}^2\frac{\alpha _{\mathrm{eff}}^{}}{R^2}\stackrel{~}{g}^2.$$
(66)
Thus, from eqs. (63), (64) and (66), we see that the parameters of OM, NCOS and ODp theories are related in the same way as the corresponding parameters of type IIA, type IIB and M-theory, as expected.
Acknowledgement: We wish to thank O. Aharony, J. Gomis and A. Tseytlin for valuable discussions. Y. O. would like to thank the USC/Caltech center for theoretical physics for hospitality during the course of this work. J. R. would like to thank CERN for hospitality.
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# Spatiotemporally localized solitons in resonantly absorbing Bragg reflectors
## Abstract
We predict the existence of spatiotemporal solitons (“light bullets”) in two-dimensional self-induced transparency media embedded in a Bragg grating. The “bullets” are found in an approximate analytical form, their stability being confirmed by direct simulations. These findings suggest new possibilities for signal transmission control and self-trapping of light.
Light propagation in periodic dielectric structures exhibits a variety of unique regimes that are technologically promising: nonlinear filtering, switching and distributed feedback amplification. Of particular interest are gap solitons (GSs), i.e., moving or standing self-localized field structures centered in a band gap of the grating. These self-localized field structures arise due to the interplay between the medium nonlinearity and resonant Bragg reflections. Their spectrum is tuned away from the Bragg resonance by the nonlinearity at sufficiently high field intensities. Theoretical studies of gap solitons in Bragg gratings with Kerr nonlinearity have been followed up by their experimental observation in a nonlinear optical fiber with the grating written on it. Gap solitons have also been theoretically studied in gratings with second harmonic generation.
A principally different mechanism giving rise to gap solitons has been found in studies of models consisting of a periodic array of thin layers of resonant two-level systems (TLS’s) separated by half-wavelength nonabsorbing dielectric layers, i.e., a resonantly absorbing Bragg reflector (RABR) (Fig. 1). Such a RABR has been shown, for any Bragg reflectivity, to produce a vast family of stable solitons, both standing and moving ones. As opposed to the 2$`\pi `$ solitons arising in self-induced transparency (SIT), i.e., near-resonant field-TLS interaction in a uniform medium, gap solitons in a RABR can have an arbitrary pulse area. The existence of GS solutions can only be consistently demonstrated in a RABR with thin active TLS layers. By contrast, a recent attempt to obtain such solutions in a periodic structure uniformly filled with active TLS’s is physically doubtful, and fails for many parameter values.
The potential applications of GSs are based on the system’s ability to filter out (by Bragg reflection) all pulses except for those satisfying the GS dispersion condition, as well as the control of pulse shape and velocity. It would be clearly desirable to supplement these advantages by immunity to transverse diffraction of the pulse, i.e. to achieve simultaneous transverse as well as longitudinal self-localization in the structure. This calls for the consideration of “light bullets” (LBs), multi-dimensional solitons that are localized in both space and time. In the last decade they have been theoretically investigated in various nonlinear optical media , and the first experimental observation of a quasi two-dimensional (2D) bullet was recently reported. In a recent work , we have predicted that uniformly 2D and 3D self-induced transparency (SIT) media can support stable light bullets.
In the present work, we aim to extend this investigation to RABR’s of the kind shown in Fig. 1, so as to combine LB and GS properties in resonantly absorbing media. We find that a RABR with any Bragg reflectivity and any absorption cross-section can support the propagation of attenuated stable LBs, which are closely related to the bullets we have found in uniform SIT media . It should be noted that 2D LBs supported by a combination of the Bragg reflector with a different (second-harmonic-generation) nonlinearity were theoretically investigated in Ref. .
We start by considering a 2D SIT medium with a spatially-varying refractive index $`n(z,x)`$, which is described by the lossless Maxwell-Bloch equations ,
$`i_{xx}+n^2_\tau +_z+i(1n^2)𝒫`$ $`=`$ $`0,`$ (2)
$`𝒫_\tau W`$ $`=`$ $`0,`$ (3)
$`W_\tau +{\displaystyle \frac{1}{2}}(^{}𝒫+𝒫^{})`$ $`=`$ $`0.`$ (4)
Here $``$ and $`𝒫`$ are the slowly varying amplitudes of the electric field and polarization, $`W`$ is the inversion, $`z`$ and $`x`$ are the longitudinal and transverse coordinates (measured in units of the resonant-absorption length), and $`\tau `$ is time (measured in units of the input pulse duration). The Fresnel number, which governs the transverse diffraction in the 2D and 3D propagation, has been incorporated in $`x`$ and the detuning of the carrier frequency $`\omega _0`$ from the central atomic-resonance frequency was absorbed in $``$ and $`𝒫`$. The Fresnel number $`F`$, detuning $`\mathrm{\Delta }\mathrm{\Omega }`$, and wave vector $`k_0\omega _0/c`$ can be brought back into Eqs. (Spatiotemporally localized solitons in resonantly absorbing Bragg reflectors) by the transformation $`(\tau ,z,x)\left(2/k_0\right)(\tau ,z,x)\text{exp}(i\mathrm{\Delta }\mathrm{\Omega }\tau )`$, $`𝒫(\tau ,z,x)\left(4/k_0^2\right)𝒫(\tau ,z,x)\text{exp}(i\mathrm{\Delta }\mathrm{\Omega }\tau )`$, $`W(\tau ,z,x)\left(4/k_0^2\right)W`$, $`\tau \left(2/k_0\right)\tau `$, $`z\left(2/k_0\right)z`$, and $`x\sqrt{2/k_0F}x`$. To neglect the polarization dephasing and inversion decay we assume pulse durations that are short on the time scale of the relaxation processes. Equations (Spatiotemporally localized solitons in resonantly absorbing Bragg reflectors) are then compatible with the local constraint $`|𝒫|^2+W^2=1`$, which corresponds to the Bloch-vector conservation . In the absence of the $`x`$-dependence and for $`n(z,x)=1`$, Eq. (2) reduces to the sine-Gordon (SG) equation which has the soliton solution $`(\tau ,z)=2\alpha \text{sech}(\alpha \tau z/\alpha +\mathrm{\Theta }_0)`$, where $`\alpha `$ and $`\mathrm{\Theta }_0`$ are real parameters.
We proceed to search for LBs in a 2D medium subject to a resonant periodic longitudinal modulation, i.e., the above-mentioned RABR (Fig. 1). Accordingly, we assume the following periodic modulation of the refractive index
$$n^2(z)=n_0^2\left[1+a_1\mathrm{cos}(2k_cz)\right].$$
(5)
Here, $`n_0`$ and $`a_1`$ are constants and $`k_c=\omega _c/c`$, with $`\omega _c`$ being the central frequency of the band gap. A RABR is then constructed by placing very thin layers (much thinner than $`1/k_c`$) of two-level atoms, whose resonance frequency is close to $`\omega _c`$, at the maxima of this modulated refractive index. We aim to consider the propagation of an electromagnetic wave with a frequency close to $`\omega _c`$ through a 2D RABR. Due to the Bragg reflections, the electric field $``$ is decomposed into forward- and backward-propagating components $`_F`$ and $`_B`$, which satisfy equations that are a straightforward generalization of the 1D version in Refs. ,
$`i\mathrm{\Sigma }_{\tau xx}^++i\mathrm{\Sigma }_{zxx}^{}+\mathrm{\Sigma }_{\tau \tau }^+\mathrm{\Sigma }_{zz}^+`$ (7)
$`+\eta \mathrm{\Sigma }_{xx}^++\eta ^2\mathrm{\Sigma }^+2𝒫_\tau 2i\eta 𝒫`$ $`=`$ $`0,`$ (8)
$`i\mathrm{\Sigma }_{\tau xx}^{}+i\mathrm{\Sigma }_{zxx}^++\mathrm{\Sigma }_{\tau \tau }^{}\mathrm{\Sigma }_{zz}^{}`$ (9)
$`\eta \mathrm{\Sigma }_{xx}^{}+\eta ^2\mathrm{\Sigma }^{}+2𝒫_z`$ $`=`$ $`0,`$ (10)
$`𝒫_\tau +i\mathrm{\Delta }\mathrm{\Omega }𝒫\mathrm{\Sigma }^+W`$ $`=`$ $`0,`$ (11)
$`W_\tau +{\displaystyle \frac{1}{2}}(\mathrm{\Sigma }^+𝒫+\mathrm{\Sigma }^+𝒫^{})`$ $`=`$ $`0.`$ (12)
Here $`\mathrm{\Sigma }^\pm 2\tau _0\mu n_0(_F\pm _B)/\mathrm{}`$, $`\tau _0n_0\mu ^1\sqrt{\mathrm{}/2\pi \omega _c\rho _0}`$, $`\mu `$ is the transition dipole moment, and $`\rho _0`$ is the density of the two-level atoms. The parameter $`\eta `$ is a ratio of resonant-absorption length in the two-level medium to the Bragg reflection length, and can be expressed as $`\eta =a_1\omega _c\tau _0/4`$.
In the 1D case, a family of exact soliton solutions to Eqs. (5) was found in Ref. : $`\mathrm{\Sigma }^{(\pm )}=2(\alpha ;\sqrt{\alpha ^2+2})\text{sech}\mathrm{\Theta }(\tau ,z)\mathrm{exp}\left(i\mathrm{\Phi }\right)`$, $`𝒫=2\text{sech}\mathrm{\Theta }(\tau ,z)\mathrm{tanh}\mathrm{\Theta }(\tau ,z)\mathrm{exp}\left(i\mathrm{\Phi }\right)`$, $`W=2\text{sech}^2\mathrm{\Theta }(\tau ,z)1`$, where $`\mathrm{\Theta }(\tau ,z)\alpha \tau +\sqrt{\alpha ^2+2}z+\mathrm{\Theta }_0`$, $`\mathrm{\Phi }\eta M\tau +\eta Nz+\varphi `$, $`M(\alpha ^2+1)`$, $`N\alpha \sqrt{\alpha ^2+2}`$, and $`\mathrm{\Delta }\mathrm{\Omega }=\eta (\alpha ^2+1)`$. The shape of the fields $`\mathrm{\Sigma }^+`$ and $`\mathrm{\Sigma }^{}`$ in this solution is similar to the SG soliton in the uniform 1D SIT medium. Inspired by this analogy and the fact that there exist LBs in the uniform 2D SIT medium which reduces to the SG soliton in 1D , we search for a LB solution to the 2D equations (5), which reduces to the exact soliton in 1D. To this end, we try the following approximation,
$`\mathrm{\Sigma }^+`$ $`=`$ $`2\alpha \sqrt{\text{sech}\mathrm{\Theta }_1\text{sech}\mathrm{\Theta }_2}e^{i\eta M\tau +i\eta Nz+i\pi /4},`$ (14)
$`\mathrm{\Sigma }^{}`$ $`=`$ $`2\sqrt{\alpha ^2+2}\sqrt{\text{sech}\mathrm{\Theta }_1\text{sech}\mathrm{\Theta }_2}e^{i\eta M\tau +i\eta Nz+i\pi /4},`$ (15)
$`𝒫`$ $`=`$ $`\sqrt{\text{sech}\mathrm{\Theta }_1\text{sech}\mathrm{\Theta }_2}\{(\mathrm{tanh}\mathrm{\Theta }_1+\mathrm{tanh}\mathrm{\Theta }_2)^2+`$ (18)
$`{\displaystyle \frac{1}{4}}\alpha ^2C^4[(\mathrm{tanh}\mathrm{\Theta }_1\mathrm{tanh}\mathrm{\Theta }_2)^22(\text{sech}^2\mathrm{\Theta }_1+`$
$`\text{sech}^2\mathrm{\Theta }_2)]^2\}^{1/2}e^{i\eta M\tau +i\eta Nz+i\nu },W=[1|𝒫|^2]^{1/2},`$
with $`\mathrm{\Theta }_1(\tau ,z)\alpha \tau +\sqrt{\alpha ^2+2}z+\mathrm{\Theta }_0+Cx`$, $`\mathrm{\Theta }_2(\tau ,z)\alpha \tau +\sqrt{\alpha ^2+2}z+\mathrm{\Theta }_0Cx`$, the phase $`\nu `$ and coefficient $`C`$ being real constants.
The ansatz (1) satisfies Eqs. (8) and (10) exactly, while Eqs. (12) are satisfied to order $`|\alpha |C^2`$, which requires that $`|\alpha |C^21`$. The ansatz is relevant for arbitrary $`\eta `$, admitting both weak ($`\eta 1`$) and strong ($`\eta >1`$) reflectivities of the Bragg grating, provided that the detuning remains small with respect to the gap frequency, or $`\eta \omega _c/(\alpha ^2+1)`$. Comparison with numerical simulations of Eqs. (5), using (1) as an initial configuration, tests this analytical approximation and shows that it is indeed fairly close to a numerically exact solution; in particular, the shape of the bullet remains within 98% of its originally presumed shape after having propagated a large distance, as is shown in Fig. 2.
We stress that 2D or 3D LB solutions of the variable-separated form $`\mathrm{\Sigma }^+\mathrm{\Sigma }^{}f(\tau ,z)g(x)`$ do not exist in RABR. Indeed, substituting this into Eqs. (8) and (10) yields only a trivial solution of the form $`\mathrm{\Sigma }^\pm \mathrm{exp}\left(iA\tau +iBx\right)`$, with constant $`A`$ and $`B`$.
We briefly discuss experimental conditions under which LBs can be observed in RABRs. The incident pulse should have uniform transverse intensity within its diameter $`d`$. For the transverse diffraction to be strong enough, one needs $`\alpha _{\mathrm{eff}}d^2/\lambda _0<1`$, where $`\alpha _{\mathrm{eff}}`$ and $`\lambda _0`$ are the inverse resonant-absorption length and carrier wavelength, respectively. For $`\alpha _{\mathrm{eff}}10^3`$ m<sup>-1</sup> and $`\lambda _010^4`$ m, one thus requires a diameter $`d<10^4`$ m, which implies that the transverse medium size $`L_x`$a few$`\mu `$m.
In order to realize a RABR, thin layers of rare-earth ions embedded in a semiconductor structure with a spatially-periodic RI may be used. The two-level atoms in the layers should be rare-earth-ions with density of $`10^{15}10^{16}`$ cm<sup>-3</sup>, whose resonant-absorption time and inverse length are respectively $`\tau _010^{13}10^{12}`$ s and $`\alpha _{\mathrm{eff}}10^410^5`$ m<sup>-1</sup>. The parameter $`\eta `$ can vary from 0 to 100 and the detuning is $`10^{12}10^{13}`$ s<sup>-1</sup>. In a RABR with transverse size $`10\mu `$m, LBs depicted in Fig. 2 are localized on the time and transverse-length scales $`10^{13}`$ s and $`1\mu `$m. Cryogenic conditions strongly extend the dephasing time $`T_2`$ and thus the LB lifetime, well into the $`\mu `$sec range . The construction of suitable structures constitutes a feasible experimental challenge.
To conclude, we have studied light bullets in SIT media embedded in a resonantly-absorbing Bragg reflector. Light bullets in a multi-dimensional Bragg reflector have the potential of serving as a novel type of optical filters, which stably transmit selected signal frequencies through their spectral gap and block others. They can also be used to both spatially and temporarily localize light in certain frequency bands.
M.B. acknowledges support through a fellowship from the Israeli Council for Higher Education. Support from ISF, Minerva and BSF is acknowledged by G.K.
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# Relativistic Scalar Aharonov-Bohm Scattering
## I INTRODUCTION
The Aharonov-Bohm (AB) effect , the scattering of a charged particle by an infinitely long and arbitrarily thin solenoid, presents a very peculiar situation of nonrelativistic (NR) quantum dynamics, with charged particles feeling the vector potential in regions where the electromagnetic field is null. It is an exactly solvable quantum mechanical problem which, due to the singular nature of the potential, requires the use of renormalization procedures to make its perturbative treatment meaningful. In fact, as noticed by Corinaldesi and Rafeli , the bare perturbation theory leads to an incomplete result in the Born approximation and to a divergent one in second order. For spinless particles in quantum mechanics, the necessary renormalization is accomplished by adding a delta function potential , while in the second quantized version it is implemented by introducing a $`\varphi ^4`$ self-interaction, with an appropriated strength, in a scalar Chern-Simons Lagrangian .
In this paper, we discuss the relativistic scalar AB scattering, that is the scattering of a relativistic charged spin zero particle by a thin fixed solenoid in the viewpoint of the first quantization comparing with the analysis in the framework of the field theory. In section II, it is shown that the problem can be solved exactly through a mimic of either the original AB solution or the Berry’s magnetization scheme. We then develop (section III) a perturbative analysis of the Klein-Gordon equation in the presence of the solenoid using the Feshbach-Villars two-component formalism and show that the bare perturbation treatment presents similar problems as those of the nonrelativistic case and possesses an analogous renormalization. In section IV, the results of the field theoretical perturbative approach, corresponding to the two-body sector of the Chern-Simons theory, are presented and compared with the scattering amplitude obtained in the relativistic quantum mechanics. Finally, some conclusions are outlined.
## II EXACT SOLUTION IN THE FRAMEWORK OF THE FIRST QUANTIZATION
The Klein-Gordon equation in the presence of an external electromagnetic field, fixing the Coulomb gauge ($`𝐀=0`$), can be written (in natural units, $`\mathrm{}=c=1`$) as
$$\left(_t^2^2+m^2+𝒰\right)\varphi =0$$
(1)
where
$$𝒰=eA^0i2e𝐀+e^2𝐀^2.$$
(2)
For an ideal AB solenoid (a line carrying magnetic flux $`\mathrm{\Phi }`$) at origin, the magnetic field
$$𝐁=\mathrm{\Phi }\delta (𝐫)\widehat{𝐳},$$
(3)
with $`𝐫=(x^1,x^2,0)`$, and one may choose
$$A^0=A^3=0,A^i=\frac{\mathrm{\Phi }}{2\pi }\frac{ϵ^{ij}x^j}{𝐫^2},i=1,2,$$
(4)
where $`ϵ^{ij}`$ is the anti-symmetric symbol ($`ϵ^{\mathrm{\hspace{0.17em}1\hspace{0.17em}2}}=1`$). The potential (2) then becomes
$$𝒰_{Sol}=i\left(\frac{e\mathrm{\Phi }}{\pi }\right)\frac{𝐫\times }{r^2}+\left(\frac{e\mathrm{\Phi }}{\pi }\right)^2\frac{1}{r^2}$$
(5)
showing that the problem, owing to the symmetry, is actually a two-dimensional one.
Considering a particle with (positive) energy $`w_𝐩`$ ($`𝐩`$ in the $`x^1`$ negative direction), the wave function can be separated as $`\varphi (𝐫,t)=\mathrm{exp}(iw_𝐩t)\varphi (𝐫)`$ and (1) reduces (in cylindrical coordinates) to
$$\left[\frac{1}{r}\frac{}{r}\left(r\frac{}{r}\right)+\frac{1}{r^2}\left(\frac{}{\theta }+i\alpha \right)^2+p^2\right]\varphi (r,\theta )=0$$
(6)
where $`\alpha =e\mathrm{\Phi }/2\pi `$ is the magnetic flux parameter and $`p^2=w_𝐩^2m^2`$. This equation coincides with the one solved by Aharonov and Bohm for the nonrelativistic particle if one replaces the dispersion relation by $`p^2=2mE`$, $`E`$ being the nonrelativistic energy. Thus, the exact solution of ( 6), vanishing when $`r0`$ as required by the “impenetrability” condition, is given by
$$\varphi (r,\theta )=\underset{l=\mathrm{}}{\overset{+\mathrm{}}{}}(i)^{|l\alpha |}J_{|l\alpha |}(pr)\mathrm{exp}(il\theta )$$
(7)
where $`J_{|l\alpha |}`$ denotes a Bessel function of first kind.
This exact solution of the Klein-Gordon equation in the presence of the solenoid can also be obtained, in a rather distinct way, by applying the Berry’s magnetization scheme to its solution in the absence of the flux line, corresponding to the incident plane wave. Using the Fourier expansion of a plane wave,
$$\varphi _0(r,\theta )=\mathrm{exp}\left[ipr\mathrm{cos}(\theta )\right]=\underset{l=\mathrm{}}{\overset{+\mathrm{}}{}}(i)^{|l|}J_{|l|}(pr)\mathrm{exp}(il\theta ),$$
(8)
and the Poisson summation formula
$$\underset{l=\mathrm{}}{\overset{+\mathrm{}}{}}f(l)=\underset{m=\mathrm{}}{\overset{+\mathrm{}}{}}_{\mathrm{}}^+\mathrm{}𝑑\eta f(\eta )\mathrm{exp}(2\pi im\eta ),$$
(9)
one obtains the whirling-wave expansion of the free solution
$`\varphi _0(r,\theta )`$ $`=`$ $`{\displaystyle \underset{m=\mathrm{}}{\overset{+\mathrm{}}{}}}w_m(r,\theta ),`$ (10)
$`w_m(r,\theta )`$ $`=`$ $`{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑\eta \mathrm{exp}({\displaystyle \frac{1}{2}}i\pi |\eta |)J_{|\eta |}(pr)\mathrm{exp}\left[i\eta (\theta +2\pi m)\right];`$ (11)
notice that $`w_m`$ is not single-valued but satisfies $`w_m(r,\theta +2\pi )=w_{m+1}(\theta )`$, therefore guaranteeing that $`\varphi _0(r,\theta )`$ has unique value.
The ingenious idea of Berry to rescue the Dirac’s magnetization prescription, by which incorporating the phase factor $`\mathrm{exp}\left(iq_{𝐫_0}^𝐫𝐀(𝐫^{})𝑑𝐫^{}\right)`$ to the free wave function would lead to the solution in the presence of the magnetic field, and to produce a single-valued wave function was to apply the Dirac’s procedure to each whirling-wave separately and to resum the “magnetized” expansion. Doing so, one finds
$`\varphi _0^D(r,\theta )`$ $`=`$ $`{\displaystyle \underset{m=\mathrm{}}{\overset{+\mathrm{}}{}}}w_m^D(r,\theta ),`$ (12)
$`w_m^D(r,\theta )`$ $`=`$ $`{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑\eta \mathrm{exp}({\displaystyle \frac{1}{2}}i\pi |\eta |)J_{|\eta |}(pr)\mathrm{exp}\left[i(\eta +\alpha )(\theta +2\pi m)\right],`$ (13)
and making the change of variables $`\eta ^{}=\eta +\alpha `$, the inverse Poisson transform leads to the exact solution (7). Notice that the Berry’s procedure relies on the fact that the free solution is a plane wave an so it can be applied equally well to both relativistic and nonrelativistic Aharonov-Bohm scattering. An extension of the whirling wave formulation of Berry to the case of the Dirac equation has been presented recently .
The AB scattering amplitude can be obtained by analyzing the asymptotic behavior of the exact solution (7). It can be shown that
$$\varphi (r,\theta )\stackrel{r1}{}\mathrm{e}^{ipr\mathrm{cos}\theta }+\mathrm{e}^{i\pi /4}𝒜_{\mathrm{AB}}(|𝐩|,\theta )\frac{\mathrm{e}^{ipr}}{\sqrt{r}}$$
(14)
where the scattering amplitude is given by
$$𝒜_{\mathrm{AB}}(|𝐩|,\theta )=\frac{i}{\sqrt{2\pi p}}\mathrm{sin}(\pi \alpha )\left[\mathrm{tan}\left(\frac{\theta }{2}\right)i\mathrm{sgn}(\alpha )\right],$$
(15)
with $`\mathrm{sgn}(\alpha )=|\alpha |/\alpha `$. It should be remarked that the exact solution, and consequently the AB scattering amplitude, can be also obtained within the two-component formalism, which is employed to construct the perturbative expansion of the $`S`$ matrix in the next section.
## III PERTURBATIVE ANALYSIS OF THE RELATIVISTIC AB SCATTERING
In order to be able to use the standard perturbation theory one has to cast the Klein-Gordon equation as a differential equation of first order in time, that is in a form similar to the Schrödinger equation. This can be done following the prescription below.
### A Two-component formalism of the Klein-Gordon equation
In the Feshbach-Villars representation , the Klein-Gordon wave function in the presence of an external field $`A^\mu `$ is written in the form
$$\mathrm{\Psi }=\left(\begin{array}{c}\chi \\ \zeta \end{array}\right)\{\begin{array}{c}\chi =\frac{1}{\sqrt{2m}}\left[i_t\varphi \left(eA^0m\right)\varphi \right]\\ \zeta =\frac{1}{\sqrt{2m}}\left[i_t\varphi +\left(eA^0+m\right)\varphi \right]\end{array}.$$
(16)
This two-component wave function satisfies a Schrödinger like equation,
$$i_t\mathrm{\Psi }=H\mathrm{\Psi },$$
(17)
with the “Hamiltonian” given by
$$H=\left(m\frac{(ie𝐀)^2}{2m}\right)\tau _3\frac{(ie𝐀)^2}{2m}i\tau _2+eA^0$$
(18)
where $`\tau _i`$ are the Pauli matrices. Although $`H`$ is not Hermitian, $`\tau _3H^{}\tau _3=H`$ and the norm $`𝑑𝐱\mathrm{\Psi }^{}\tau _3\mathrm{\Psi }`$ is conserved, i. e., $`\tau _3`$ plays the role of a metric tensor.
This formalism allows the use of perturbation theory in the standard way and we make the partition $`H=H_0+H_{\mathrm{int}}=H_0+V+H_1\left[\tau _3+i\tau _2\right]`$ where $`V=eA^0`$ and
$`H_0`$ $`=`$ $`\left(m{\displaystyle \frac{^2}{2m}}\right)\tau _3{\displaystyle \frac{^2}{2m}}i\tau _2,`$ (19)
$`H_1`$ $`=`$ $`\left({\displaystyle \frac{ie}{2m}}\left[𝐀+𝐀\right]+{\displaystyle \frac{e^2}{2m}}𝐀^2\right).`$ (20)
The free positive and negative energy solutions are given by
$$\mathrm{\Psi }_𝐩^{(+)}(𝐱,t)=\mathrm{e}^{iw_𝐩t}\mathrm{\Psi }_𝐩^{(+)}(𝐱)=\frac{1}{2\pi }\frac{w_𝐩+m}{\sqrt{4mw_𝐩}}\left(\begin{array}{c}1\\ \frac{mw_𝐩}{m+w_𝐩}\end{array}\right)\mathrm{e}^{i𝐩𝐱iw_𝐩t},$$
(21)
$$\mathrm{\Psi }_𝐩^{()}(𝐱,t)=\mathrm{e}^{+iw_𝐩t}\mathrm{\Psi }_𝐩^{()}(𝐱)=\frac{1}{2\pi }\frac{w_𝐩+m}{\sqrt{4mw_𝐩}}\left(\begin{array}{c}\frac{mw_𝐩}{m+w_𝐩}\\ 1\end{array}\right)\mathrm{e}^{i𝐩𝐱+iw_𝐩t},$$
(22)
which satisfy the normalization conditions
$$𝑑𝐱\mathrm{\Psi }_𝐩^{}^{(\pm )}\tau _3\mathrm{\Psi }_𝐩^{(\pm )}=\pm \delta (𝐩^{}𝐩).$$
(23)
Notice that the Berry’s magnetization procedure can be applied to the positive energy solution (21) leading to the exact solution described in the last section and thus to the AB scattering amplitude (15).
In fact, the two-component formalism is the appropriate scenario to implement the perturbative treatment of the scalar AB scattering in the framework of the relativistic quantum mechanics. The $`S`$ matrix elements can be calculated through the formula
$$S_{fi}=<f|\mathrm{T}\mathrm{exp}\left[i𝑑tH_I(t)\right]|i>,$$
(24)
where T stands for the time ordering operator and $`H_I(t)`$ is the interaction picture of the interaction Hamiltonian $`H_{\mathrm{int}}`$, from which we can define the scattering amplitude $`𝒜_{fi}`$ as
$$S_{fi}=2\pi i\delta (w_fw_i)𝒜_{fi}.$$
(25)
### B The bare scattering amplitude
The S matrix, expressing the scattering of a positive energy particle, in the first Born approximation, becomes
$$S_{fi}^{(1)}=i𝑑t𝑑𝐱\mathrm{\Psi }_𝐩^{^{}}^{}(𝐱)\tau _3H_{\mathrm{int}}(𝐱,𝐭)\mathrm{\Psi }_𝐩(𝐱)=i\mathrm{\hspace{0.17em}2}\pi \delta (w_𝐩^{^{}}w_𝐩)𝒜_{fi}^{(1)}$$
(26)
with the reduced amplitude given by
$$𝒜_{fi}^{(1)}=𝑑𝐱\mathrm{e}^{i𝐩^{^{}}𝐱}\left[\frac{m}{w_𝐩}H_1(𝐱,)+V(𝐱)\right]\mathrm{e}^{i𝐩𝐱}.$$
(27)
Considering the Coulomb gauge, $`𝐀=0`$, and for the moment neglecting the higher order term $`e^2𝐀^2`$, one obtains
$$𝒜_{fi}=𝑑𝐱\left[\frac{e}{w_𝐩}𝐀(𝐱)𝐩+V(𝐱)\right]\mathrm{e}^{i(𝐩^{}𝐩)𝐱}=\frac{e}{w_𝐩}\stackrel{~}{𝐀}(𝐪)𝐩+\stackrel{~}{V}(𝐪)$$
(28)
where $`𝐪=𝐩^{}𝐩`$.
For the AB potential, $`V=0`$ and
$$\stackrel{~}{A}^i(𝐪)=\frac{i\mathrm{\Phi }}{2\pi }ϵ^{ij}\underset{\lambda 0}{lim}_{q^j}𝑑𝐫\frac{\mathrm{e}^{i𝐪𝐫}}{𝐫^2+\lambda ^2}=i\mathrm{\Phi }ϵ^{ij}\underset{\lambda 0}{lim}_{q^j}K_0(\lambda |𝐪|)=i\mathrm{\Phi }\frac{ϵ^{ij}q^j}{𝐪^2},$$
(29)
$`K_0`$ denoting the modified Bessel function, and therefore the reduced amplitude is given by
$$𝒜_{fi}^{(1)}=\frac{ie\mathrm{\Phi }}{w_𝐩}\frac{𝐩\times 𝐪}{𝐪^2}=\frac{ie\mathrm{\Phi }}{2w_𝐩}\mathrm{cot}\left(\frac{\theta _S}{2}\right)$$
(30)
where $`\theta _S`$ is the scattering angle. Notice, since $`\theta _S`$ is the angle between the incoming ($`𝐩`$) and the outgoing ($`𝐩^{}`$ ) momenta, it is related to the AB polar angle by $`\theta _S=\pi \theta `$. One immediately notices the absence of the nonanalytic term occurring in the first order of the expansion of the exact result (15).
To calculate higher-order contributions one should be aware of the fact that, since the interaction with the external field contains terms of order $`e`$ and $`e^2`$, the perturbative expansion in the coupling constant does not correspond “order by order” to the Born expansion. Thus, to find the second order contribution to the amplitude, besides the contribution coming from the second term of the expansion (24), one has to account for the $`𝐀^2`$ term present in the first Born approximation; one obtains
$$𝒜_{fi(𝐀^2)}^{(2)}=𝑑𝐫\mathrm{e}^{i𝐩^{}𝐫}\left[\frac{m}{w_𝐩}\frac{e^2}{2m}𝐀^2(𝐫)\right]\mathrm{e}^{i𝐩𝐫}=\frac{m}{w_𝐩}\frac{e^2}{2m}\stackrel{~}{𝒞}(𝐪)$$
(31)
where, the Fourier transform of $`𝐀^2`$,
$$\stackrel{~}{𝒞}(𝐪)=\left(\frac{\mathrm{\Phi }}{2\pi }\right)^2𝑑𝐫\frac{\mathrm{e}^{i𝐪𝐫}}{r^2}=\left(\frac{\mathrm{\Phi }}{2\pi }\right)^2\underset{\lambda 0}{lim}𝑑𝐫\frac{\mathrm{e}^{i𝐪𝐫}}{r^2+\lambda ^2}=\left(\frac{\mathrm{\Phi }}{2\pi }\right)^2\underset{\lambda 0}{lim}K_0(\lambda |𝐪|),$$
(32)
leads to a divergent contribution. Notice that, it is precisely due to the absence of the $`𝐀^2`$ term in the interaction of the spin half particles with the field that no divergence occurs in that case; furthermore the Pauli magnetic interaction makes naturally the first Born approximation correct. Wishing to be able to suppress the divergence of the scalar case, we keep the regularized contribution
$$𝒜_{fi(𝐀^2)}^{(2)Reg}=\frac{e^2\mathrm{\Phi }^2}{4(2\pi )^3}\frac{1}{w_p}\left\{\mathrm{ln}(\lambda ^2)\mathrm{ln}(p^2)\mathrm{ln}[2(1\mathrm{cos}\theta _S)]+2(\mathrm{ln}2\gamma )\right\}.$$
(33)
where $`\gamma `$ is the Euler’s constant.
The other second order contribution comes from the second Born approximation,
$$S_{fi}^{(2)}=𝑑t_1𝑑t_2\theta (t_2t_1)𝑑𝐫_\mathrm{𝟐}\mathrm{\Psi }_𝐩^{}^{(+)}(𝐫_\mathrm{𝟐})\tau _3H_{int}(𝐫_\mathrm{𝟐},t_2)H_{int}(𝐫_\mathrm{𝟐},t_1)\mathrm{\Psi }_𝐩^{(+)}(𝐫_\mathrm{𝟐}),$$
(34)
considering only the part of the interaction Hamiltonian which is first order in $`e`$. Using the completeness relation
$$𝑑𝐤\left[\mathrm{\Psi }_𝐤^{(+)}(𝐫_\mathrm{𝟏})\mathrm{\Psi }_𝐤^{(+)}(𝐫_\mathrm{𝟐})\tau _3\mathrm{\Psi }_𝐤^{()}(𝐫_\mathrm{𝟏})\mathrm{\Psi }_𝐤^{()}(𝐫_\mathrm{𝟐})\tau _3\right]=𝐈\delta (𝐫_\mathrm{𝟏}𝐫_\mathrm{𝟐})$$
(35)
and the identity
$$\theta (t)E^n\mathrm{e}^{iEt}=\frac{1}{2\pi i}𝑑\omega \frac{\omega ^n\mathrm{e}^{i\omega t}}{E\omega iϵ},$$
(36)
one gets
$$S_{fi(𝐀𝐩)}^{(2)}=2\pi i\delta (w_fw_i)𝒜_{fi(𝐀𝐩)}^{(2)}$$
(37)
with
$`𝒜_{fi(𝐀𝐩)}^{(2)}`$ $`=`$ $`{\displaystyle \frac{e^2}{(2\pi )^4}}{\displaystyle \frac{1}{w_p}}{\displaystyle }{\displaystyle \frac{d𝐤}{\omega _k}}\{{\displaystyle \frac{[𝐤\stackrel{~}{𝐀}(𝐩^{}𝐤)][𝐩\stackrel{~}{𝐀}(𝐤𝐩)]}{\omega _p\omega _k+iϵ}}`$ (38)
$`{\displaystyle \frac{[𝐤\stackrel{~}{𝐀}(𝐩^{}+𝐤)][𝐩\stackrel{~}{𝐀}(𝐤+𝐩)]}{\omega _p+\omega _k+iϵ}}\}.`$ (39)
where $`w_k=\sqrt{𝐤^2+m^2}`$. From (29), it follows that $`𝐤\stackrel{~}{𝐀}(𝐩)=i\mathrm{\Phi }(𝐤\times 𝐩)/p^2`$, so the contribution (39) to the second order amplitude is given by
$`𝒜_{fi(𝐀𝐩)}^{(2)}`$ $`=`$ $`{\displaystyle \frac{e^2\mathrm{\Phi }^2}{(2\pi )^4}}{\displaystyle \frac{1}{w_p}}{\displaystyle 𝑑𝐤\frac{(𝐤\times 𝐩^{})(𝐤\times 𝐩)}{(𝐤𝐩^{})^2(𝐤𝐩)^2}\frac{2}{p^2k^2+iϵ}}`$ (40)
$`=`$ $`{\displaystyle \frac{e^2\mathrm{\Phi }^2}{4(2\pi )^3}}{\displaystyle \frac{1}{w_p}}\left\{\mathrm{ln}[2(1\mathrm{cos}\theta _S)]+i\pi \right\},`$ (41)
where the integral was performed as presented in . The $`\theta _S`$ dependent part of this finite contribution is canceled out with one of the terms of $`𝒜_{fi(𝐀^2)}^{(2)}`$. One thus obtains, adding up both contributions, a divergent second order term for the bare scattering amplitude. To recover a finite result, coinciding with the expansion of the exact amplitude, one has to implement some renormalization procedure.
### C Perturbative renormalization
To regain the correct perturbative expansion we have to search for an appropriate counterterm, an additional interaction, which should suppress the divergence of the second order and contributes in the lowest order recovering the exact result. At first sight, one can try to follow the case of spin half particles and look for a magnetic type of interaction to do the job. However, in the scalar case the magnetic interaction, given by $`H_{mag}=gϵ^{\mu \nu \rho }F_{\nu \rho }j_\mu `$, where $`F_{\nu \rho }`$ is the field strength and $`j_\mu `$ the particle’s current, can only furnish the correct result in the leading nonrelativistic order, and thus can not describe the full relativistic case. One of the reasons for that is the fact that such interaction does not have the same matrix structure in the two-component formalism as the $`𝐀^2`$ term, which is responsible for the divergence one wishes to eliminate.
The simplest additional interaction having the same matrix structure of the magnetic potential term one can consider, leading to the same logarithmic divergence in second order, is a pure delta function external potential. In fact, by adding to (2) the term
$$𝒰_{(delta)}=g\delta (𝐫)$$
(42)
one obtains, in first order, the contribution
$$𝒜_{fi(delta)}^{(1)}=\frac{mg}{2}\frac{1}{w_p}$$
(43)
while the cutoff regularized “delta-delta” contribution to the second order becomes
$`𝒜_{fi(delta)}^{(2)Reg}`$ $`=`$ $`{\displaystyle \frac{m^2g^2}{(2\pi )^4}}{\displaystyle \frac{1}{w_p}}{\displaystyle ^{\mathrm{\Lambda }^2}}𝑑𝐤{\displaystyle \frac{1}{w_pw_k+iϵ}}`$ (44)
$`=`$ $`{\displaystyle \frac{m^2g^2}{4(2\pi )^3}}{\displaystyle \frac{1}{w_p}}\left\{\mathrm{ln}(\mathrm{\Lambda }^2)+\mathrm{ln}(p^2)i\pi \right\}`$ (45)
The crossing terms involving the delta and $`𝐀𝐩`$ interactions need not to be considered; the sum of their contributions vanishes, a fact which can be inferred from symmetry arguments.
Including these contributions from the delta potential, the first and the regularized second order parts of the scattering amplitude become
$$𝒜_{fi}^{(1)}=i\frac{e\mathrm{\Phi }}{2w_p}\mathrm{cot}\left(\frac{\theta _S}{2}\right)+\frac{mg}{2w_p}$$
(46)
$`𝒜_{fi}^{(2)}`$ $`=`$ $`{\displaystyle \frac{1}{4(2\pi )^3}}{\displaystyle \frac{1}{w}}_p(m^2g^2e^2\mathrm{\Phi }^2)[\mathrm{ln}(p^2)i\pi ]`$ (47)
$`+{\displaystyle \frac{1}{4(2\pi )^3}}{\displaystyle \frac{1}{w}}_p\left[m^2g^2\mathrm{ln}(\mathrm{\Lambda }^2)e^2\mathrm{\Phi }^2(\mathrm{ln}(\lambda ^2)2\mathrm{ln}2+2\gamma )\right]`$ (48)
One sees then that the agreement with the expansion of the exact result can be reached if the strength of the delta interaction is fixed satisfying
$$m^2g^2=e^2\mathrm{\Phi }^2=4\pi ^2\alpha ^2$$
(49)
and the cutoffs are adjusted so that
$$\mathrm{\Lambda }\lambda =2\mathrm{exp}(\gamma );$$
(50)
in doing so, the first order term, multiplied by the appropriate kinematical factor, reproduces the correct result and the second order, proportional to $`\alpha ^2`$, vanishes as it should.
## IV FIELD THEORETICAL APPROACH
The scalar nonrelativistic AB scattering corresponds, in the field theoretical approach, to the two-body sector of the theory of a Chern-Simons field coupled with a self-interacting scalar field for which the Lagrangian density is given by
$$_{NR}=\psi ^{}\left(iD_t+\frac{𝐃^2}{2m}\right)\psi \frac{v_0}{4}(\psi ^{}\psi )^2+\frac{\mathrm{\Theta }}{2}_t𝐀\times 𝐀\mathrm{\Theta }A_0\times 𝐀,$$
(51)
where $`D_t=_t+ieA_0`$ and $`𝐃=ie𝐀`$ are covariant derivatives, $`\mathrm{\Theta }`$ is the Chern-Simons parameter and $`v_0`$ is the bare self-coupling. The renormalized nonrelativistic two particle scattering amplitude, in the center of mass (CM) frame, is given, up to one loop, by
$$𝒜_{NR}=vi\frac{2e^2}{m\mathrm{\Theta }}\mathrm{cot}\theta +\frac{m}{8\pi }\left(v^2\frac{4e^4}{m^2\mathrm{\Theta }^2}\right)\left[\mathrm{ln}\left(\frac{\mu ^2}{𝐩^2}\right)+i\pi \right],$$
(52)
where $`\mu `$ is an arbitrary mass scale that breaks the scale invariance of the amplitude. By choosing the critical value $`v_c^+=+2e^2/m|\mathrm{\Theta }|`$, which corresponds to a repulsive quartic interaction, this amplitude reduces to the first order ($`e^2`$) Aharonov-Bohm amplitude for identical particles which is given by
$$_{\mathrm{AB}}(|𝐩|,\theta )=i\frac{4\pi \beta }{m}\left[\mathrm{cot}\theta _Si\mathrm{sgn}(\beta )\right],$$
(53)
where $`\beta =e^2/2\pi \mathrm{\Theta }`$ coincides with the flux parameter $`\alpha `$ if the identification $`\mathrm{\Theta }=e/\mathrm{\Phi }`$ is made; this symmetrized amplitude is the same, adjusting the kinematical factor, as the one obtained from the first order expansion of the exact NR result (15).
It has been shown that the relativistic Chern-Simons theory,
$$=(D_\mu \varphi )^{}(D^\mu \varphi )m^2\varphi ^{}\varphi \frac{\lambda }{4}(\varphi ^{}\varphi )^2+\frac{\mathrm{\Theta }}{2}ϵ_{\sigma \mu \nu }A^\sigma ^\mu A^\nu ,$$
(54)
reduces to the nonrelativistic case in the leading approximation . We have calculated in this theory the $`|𝐩|/m`$ expansion of the two particle amplitude, up to one loop order, using an intermediate cutoff procedure introduced in ref. . The renormalized CM amplitude, including the factor $`1/4w_𝐩^2`$ (where $`w_𝐩=\sqrt{m^2+𝐩^2}`$) which makes the states to have the same normalization as in the NR case, can be written as $`𝒜=𝒜^{(0)}+𝒜^{(1)}`$ , where
$$𝒜^{(0)}=\frac{\lambda }{4w_p^2}\frac{ie^2}{\mathrm{\Theta }w_p}\mathrm{cot}\left(\frac{\theta _S}{2}\right)+\left[\theta _S\theta _S\pi \right]$$
(55)
is the exact tree level contribution and the one loop term, up to order $`𝐩^2/m^2`$, is given by
$`𝒜^{(1)}`$ $``$ $`{\displaystyle \frac{m}{8\pi }}\left({\displaystyle \frac{\lambda ^2}{16m^4}}{\displaystyle \frac{4e^4}{m^2\mathrm{\Theta }^2}}\right)\left[\mathrm{ln}\left({\displaystyle \frac{4m^2}{𝐩^2}}\right)+i\pi \right]`$ (58)
$`{\displaystyle \frac{m}{8\pi }}\left({\displaystyle \frac{3\lambda ^2}{32m^4}}{\displaystyle \frac{2e^4}{m^2\mathrm{\Theta }^2}}\right){\displaystyle \frac{𝐩^2}{m^2}}\left[\mathrm{ln}\left({\displaystyle \frac{4m^2}{𝐩^2}}\right)+i\pi \right]`$
$`+{\displaystyle \frac{m}{8\pi }}\left({\displaystyle \frac{\lambda ^2}{4m^4}}{\displaystyle \frac{14e^4}{3m^2\mathrm{\Theta }^2}}\right){\displaystyle \frac{m}{8\pi }}\left({\displaystyle \frac{25\lambda ^2}{96m^4}}+{\displaystyle \frac{74e^4}{15m^2\mathrm{\Theta }^2}}\right){\displaystyle \frac{𝐩^2}{m^2}}.`$
One sees that the leading term of the $`|𝐩|/m`$ expansion of $`𝒜`$ coincides with $`𝒜_{NR}`$ if one identifies $`v=\lambda /4m^2`$ and choose $`\mu ^2=4m^2`$. Independently of the fixing of $`\mu `$, by taking $`\lambda _c^+=4m^2v_c^+`$ the one loop contribution for the leading order (in $`|𝐩|/m`$) vanishes and the tree level one reproduces the AB scattering. However, the subdominant terms do not vanish for $`\lambda =\lambda _c^+`$ and constitute additional relativistic corrections to the AB effect which are originated from field theoretical effects as vacuum polarization and vertex radiative corrections. Notice, in this respect, that the vectorial interaction vertex expressing the bare coupling between the matter and the Chern-Simon fields possesses, in the relativistic case, an energy factor in the zeroth component which is not present in the NR Lagrangian.
## V CONCLUSIONS
Using a two-component formalism, in this work we have studied perturbatively the first quantized AB scattering of relativistic scalar particles. We proved that, to eliminate divergences due to the $`A^2`$ coupling, it is necessary to add a contact delta interaction which, in a field theoretical language, corresponds to a $`(\varphi ^{}\varphi )^2`$ self-interaction. Now, it is known that in the fermionic case Pauli’s magnetic interaction $`𝐁\psi ^{}𝐬\psi `$, with $`𝐬`$ standing for the spin operator, provides the necessary ingredient to make the final result well defined. This immediately suggest that in the scalar case the linearized form $`\sigma (\varphi ^{}\varphi )`$, where in lowest order the external field $`\sigma =g\delta (x)`$, should be added to the original Lagrangian, as we did. In the nonrelativistic case the purely magnetic coupling $`gϵ^{\mu \nu \rho }F_{\nu \rho }j_\mu `$, where $`F_{\nu \rho }`$ is the field strength and $`j_\mu `$ the particle’s current, equally provides the cancellation of the divergence, but in the relativistic domain it has to be disregarded since it does not have the appropriated momentum dependence and, in the field theory context, it is nonrenormalizable. Comparison between the first-quantized and the field theoretical perturbative expansions shows that the latter has additional contributions coming from vacuum polarization and vertex radiative corrections. These terms, absent in a direct nonrelativistic approach, show that the original Aharonov-Bohm problem is an idealized situation since vacuum polarization makes the magnetic field necessarily nonvanishing outside the solenoid.
Acknowledgments
This work was partially supported by CAPES, CNPq and FAPESP, Brazilian agencies.
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# The critical behavior of frustrated spin models with noncollinear order
\[
## Abstract
We study the critical behavior of frustrated spin models with noncollinear order, including stacked triangular antiferromagnets and helimagnets. For this purpose we compute the field-theoretic expansions at fixed dimension to six loops and determine their large-order behavior. For the physically relevant cases of two and three components, we show the existence of a new stable fixed point that corresponds to the conjectured chiral universality class. This contradicts previous three-loop field-theoretical results but is in agreement with experiments.
\]
The critical behavior of frustrated spin systems with noncollinear or canted order has been the object of intensive theoretical and experimental studies (see, e.g., Ref. ). In spite of these efforts, the critical behavior of these systems is still unclear, field-theoretic (FT) renormalization-group (RG) methods, Monte Carlo simulations and experiments obtaining different results.
In physical magnets noncollinear order is due to frustration that may arise either because of the special geometry of the lattice, or from the competition of different kinds of interactions. Typical examples of systems of the first type are three-dimensional stacked triangular antiferromagnets (STA), where magnetic ions are located at each site of a three-dimensional stacked triangular lattice. Examples are $`\mathrm{ABX}_3`$-type compounds, where A denotes elements such as Cs and Rb, B stands for magnetic ions such as Mn, Cu, Ni, and Co, and X for halogens as Cl, Br, and I. They may be modeled by using short-ranged Hamiltonians for $`N`$-component spin variables defined on a stacked triangular lattice as
$$_{\mathrm{STA}}=J\underset{ij_{xy}}{}\stackrel{}{s}_i\stackrel{}{s}_jJ^{}\underset{ij_z}{}\stackrel{}{s}_i\stackrel{}{s}_j,$$
(1)
where $`J<0`$, the first sum is over nearest-neighbor pairs within triangular layers ($`xy`$ planes), and the second one is over orthogonal interlayer nearest neighbors. In these spin systems the Hamiltonian is minimized by noncollinear configurations, showing a 120<sup>o</sup> spin structure. Frustration is partially released by mutual spin canting, and the degeneracy of the ground-state is limited to global O($`N`$) spin rotations and reflections. As a consequence, at criticality there is a breakdown of the symmetry from O($`N`$) in the high-temperature phase to O($`N2`$) in the low-temperature phase, implying a matrix-like order parameter. Frustration due to the competition of interactions may be realized in helimagnets where a magnetic spiral is formed along a certain direction of the lattice (see, e.g., Ref. ). The rare-earth metals Ho, Dy and Tb provide examples of such systems.
The critical behavior of two- and three-component frustrated spin models with noncollinear order is controversial. Many experiments (see, e.g., Ref. ) are consistent with a second-order phase transition belonging to a new (chiral) universality class. This is partially supported by Monte Carlo simulations (see, e.g., Ref. and references therein). On the other hand, three-loop perturbative calculations at fixed dimension $`d=3`$ and within the framework of the $`ϵ`$-expansion indicate a first-order transition, since no stable chiral fixed points are found for $`N=2`$ and $`N=3`$. These three-loop analyses show the presence of a stable chiral fixed point only for $`N>N_c`$ with $`N_c>3`$: $`N_c=3.91`$ and $`N_c=3.39`$ .
To explain these contradictory results it has been suggested that these systems undergo weak first-order transitions, that effectively appear as second-order ones in numerical and experimental works. This hypothesis has been supported by studies based on approximate solutions of the Wilson RG equations , and by Monte Carlo investigations of modified lattice spin systems which, according to general universality ideas, should belong to the same universality class of the Hamiltonian (1), and which show a first-order transition.
For larger values of $`N`$, all theoretical approaches predict a second-order phase transition, but there are still substantial discrepancies between Monte Carlo and three-loop FT calculations (see the discussion of Ref. for $`N=6`$).
All these considerations show that a satisfactory theoretical understanding has not yet been reached. It is not clear whether experiments are observing first-order transitions in disguise or field theory is unable to describe these rather complex systems. Of course, one may think that the observed disagreement is due to the shortness of the available series, thereby calling for an extension of the perturbative expansions to clarify the issue.
FT studies of systems with noncollinear order are based on the O($`N`$)$`\times `$O($`M`$) symmetric Hamiltonian
$`={\displaystyle d^dx}`$ $`\{{\displaystyle \frac{1}{2}}{\displaystyle \underset{a}{}}[(_\mu \varphi _a)^2+r\varphi _a^2]+{\displaystyle \frac{1}{4!}}u_0\left({\displaystyle \underset{a}{}}\varphi _a^2\right)^2`$ (3)
$`+{\displaystyle \frac{1}{4!}}v_0{\displaystyle \underset{a,b}{}}[(\varphi _a\varphi _b)^2\varphi _a^2\varphi _b^2]\},`$
where $`\varphi _a`$ ($`1aM`$) are $`M`$ sets of $`N`$-component vectors. We will consider the case $`M=2`$, that, for $`v_0>0`$, describes frustrated systems with noncollinear ordering such as STA’s. Negative values of $`v_0`$ correspond to simple ferromagnetic or antiferromagnetic ordering, and to magnets with sinusoidal spin structures .
For $`N=2`$, which is the case relevant for frustrated two-component spin models, an $`ϵ`$-expansion analysis indicates the presence of four fixed points: the Gaussian one, an $`XY`$ fixed point, an O(4)-symmetric and a mixed fixed point. Using nonperturbative arguments , one can show that the $`XY`$ fixed point is the only stable one among them. However, the region relevant for frustrated models, $`v_0>0`$, is outside the domain of attraction of the $`XY`$ fixed point, which would imply a first-order transition. However, it is still possible that other fixed points are present in the region $`v_0>0`$, although they are not predicted by the $`ϵ`$-expansion. For $`N=3`$, one may easily show the existence of an O(6) fixed point for $`v_0=0`$, which is expected to be unstable . According to the three-loop analyses of Refs. no other fixed points are found for $`N=3`$, which would imply that the transition is of first order as well.
In order to investigate the existence of new fixed points, we have considered the fixed-dimension perturbative approach, extending the three-loop series of Ref. to six loops. As we shall see, the results of our six-loop analysis are somehow surprising, contradicting most of the earlier FT works. Indeed, the analysis of the longer series provides a rather robust evidence for the existence of a new stable fixed point in the $`XY`$ and Heisenberg cases, with critical exponents that are in agreement with the experimental results.
In the fixed-dimension FT approach one expands in powers of the quartic couplings and renormalizes the theory by introducing a set of zero-momentum conditions for the two-point and four-point correlation functions. All perturbative series are finally expressed in terms of the zero-momentum four-point renormalized couplings $`u`$ and $`v`$ normalized so that, at tree level, $`uu_0`$ and $`vv_0`$. The fixed points of the theory are given by the common zeros of the $`\beta `$-functions $`\beta _u(u,v)`$ and $`\beta _v(u,v)`$. In the case of a continuous transition, when $`m0`$, the couplings $`u,v`$ are driven toward an infrared-stable zero $`u^{},v^{}`$ of the $`\beta `$-functions. On the other hand, the absence of stable fixed points is usually considered as an indication of a (weak) first-order transition.
In Tables I and II we present the six-loop expansion of the $`\beta `$-functions for $`M=2`$, associated respectively with the rescaled couplings $`\overline{u}=3u/(16\pi R_{2N})`$ and $`\overline{v}=3v/(16\pi )`$, where $`R_K9/(8+K)`$. Since FT perturbative expansions are asymptotic, the resummation of the series is essential to obtain accurate estimates of the physical quantities. For this purpose we studied the large-order behavior of the expansion in $`\overline{u}`$ and $`\overline{v}`$ at fixed $`z\overline{v}/\overline{u}`$. For $`z\overline{v}/\overline{u}`$ fixed and $`M=2`$, the singularity of the Borel transform closest to the origin, $`\overline{u}_b`$, is given by
$`{\displaystyle \frac{1}{\overline{u}_b}}`$ $`=`$ $`aR_{2N}\mathrm{for}\mathrm{\hspace{0.33em}\hspace{0.33em}4}R_{2N}>z>0,`$ (4)
$`{\displaystyle \frac{1}{\overline{u}_b}}`$ $`=`$ $`a\left(R_{2N}{\displaystyle \frac{1}{2}}z\right)\mathrm{for}z<0,z>4R_{2N},`$ (5)
where $`a=0.14777422\mathrm{}`$ and $`R_K=9/(8+K)`$. Moreover, we find that for $`z>2R_{2N}`$ the Borel transform has a singularity on the positive real axis, which however is not the closest one for $`z<4R_{2N}`$. Thus, for $`z>2R_{2N}`$ the series is not Borel summable.
In order to determine the fixed points we use the same method applied in Ref. to the analysis of the RG functions of the cubic model. We resum the perturbative series by means of an appropriate conformal mapping that takes into account the large-order behavior of the perturbative series at fixed $`z`$ and turns the original series into a convergent sequence of approximations. To understand the systematic errors we vary two different parameters, $`b`$ and $`\alpha `$, in the analysis. We apply this method also for those values of $`z`$ for which the series is not Borel summable. Although in this case the sequence of approximations is only asymptotic, it should provide reasonable estimates as long as $`z<4R_{2N}`$, since we are taking into account the leading large-order behavior.
In Figs. 1 and 2 we report our results for the zeros of the $`\beta `$-functions, obtained from the analysis of the $`l`$-loop series, $`l=3,4,5,6`$. For each $`\beta `$-function we consider 18 different approximants with $`b=3,6,\mathrm{},18`$ and $`\alpha =0,2,4`$ and we determine the lines in the $`(\overline{u},\overline{v})`$ plane on which they vanish. Then, we divide the domain $`0\overline{u}4`$ and $`0\overline{v}5`$ into $`40^2`$ rectangles, marking those in which at least two approximants of each $`\beta `$ function vanish. No fixed point is observed at three loops, consistently with Ref. . As the number of loops increases, a new fixed point—quite stable with respect to the number of loops—clearly appears for
$`\overline{u}^{}=1.9(1),\overline{v}^{}=4.10(15),`$ $`\mathrm{for}N=2,`$ (6)
$`\overline{u}^{}=1.8(1),\overline{v}^{}=3.00(15),`$ $`\mathrm{for}N=3.`$ (7)
We have been conservative in setting the error bars: all zeros of the approximants with $`3b18`$ and $`0\alpha 4`$ lie within the reported confidence interval. Notice that the fixed points belong to the region in which the series are not Borel summable, but still satisfy $`\overline{v}^{}/\overline{u}^{}<4R_{2N}`$. Therefore, we expect our resummations to be reliable, and the stability of the results with respect to the order of the series confirms it.
We then compute the eigenvalues of the stability matrix. They vary significantly with the two parameters $`\alpha `$ and $`b`$ and turn out to be complex in most of the cases. Nonetheless, the sign of the real part of the eigenvalues is always positive showing the stability of the new fixed points. A reasonable estimate of the exponent $`\omega `$ is however impossible.
We also compute the critical exponents, by estimating the corresponding six-loops series at the fixed point, following Ref. . The results are in substantial agreement with the experimental and Monte Carlo (MC) estimates, see Table III.
Finally, we compare the six-loop results with the critical exponents computed to $`O(1/N^2)`$ in the framework of the large-$`N`$ expansion. For example,
$`\nu =1{\displaystyle \frac{16}{\pi ^2}}{\displaystyle \frac{1}{N}}\left({\displaystyle \frac{56}{\pi ^2}}{\displaystyle \frac{640}{3\pi ^4}}\right){\displaystyle \frac{1}{N^2}}+O\left({\displaystyle \frac{1}{N^3}}\right).`$ (8)
We find $`\nu =0.858(4)`$ for $`N=16`$ and $`\nu =0.936(2)`$ for $`N=32`$, which compare reasonably with the estimates that one obtains from Eq. (8), i.e. $`\nu =0.885`$ for $`N=16`$ and $`\nu =0.946`$ for $`N=32`$.
In conclusion, the extension to six loops of the FT expansions solves the apparent contradictions between field theory and experiments. We find that new stable chiral fixed points exist for two- and three-component systems. The estimated exponents are in substantial agreement with experiments, whose conclusions on the nature of the phase transitions are thus confirmed. However, we note that first-order transitions are still possible for systems that are outside the attraction domain of the chiral fixed point. In this case, the RG flow runs away to a first-order transition.
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# Spatial Variability in the Ratio of Interstellar Atomic Deuterium to Hydrogen. II. Observations toward 𝛾² Velorum and 𝜁 Puppis by the Interstellar Medium Absorption Profile Spectrograph1footnote 11footnote 1This paper is dedicated in memory of Judith L. Tokel, wife of the first author, who passed away on 2000 June 10. Her enthusiastic support and encouragement were essential to its successful completion.
## 1 INTRODUCTION
The abundance ratio of atomic deuterium to hydrogen (D/H) in interstellar gas is widely regarded as an important tracer of galactic chemical evolution (Audouze & Tinsley 1974; Boesgaard & Steigman 1985; Tosi et al. 1998) and a key discriminant of the cosmic baryon-to-photon ratio $`(\eta )`$ in Big Bang Nucleosynthesis (BBN, Walker et al. 1991). A standard interpretation is that D should not be produced in significant quantities in astrophysical sites other than the Big Bang (Epstein, Latimer, & Schramm 1976). The generally accepted viewpoint is that while D is produced primordially, its destruction takes place when some of the gas is cycled through stars. The uncertainties surrounding this process represent a stumbling block in arriving at the primordial value. For this reason, a measurement of D/H is often regarded as a lower limit to the primordial ratio, and this may be translated into an upper limit to $`\eta `$ (for BBN, larger D/H implies lower $`\eta `$).
Rogerson & York (1973) made the first measurement of the atomic D/H abundance ratio in the interstellar medium (ISM) on the line of sight toward $`\beta `$ Cen. Measurements with the Copernicus satellite toward an additional 14 stars ($`100<d<500`$ pc) found that D/H values cluster around $`1.5\times 10^5`$, but with a dispersion that, in some cases, seemed to exceed the stated uncertainties (see review by Vidal-Madjar & Gry 1984). While a simple interpretation of the Copernicus data suggested that the D/H measurements revealed spatial variations, the reality of these differences has been difficult to substantiate due to concerns stemming from the somewhat inadequate resolution of Copernicus ($`15\mathrm{km}\mathrm{s}^1`$ FWHM) for this purpose.
Using the Goddard High Resolution Spectrograph (GHRS) and Space Telescope Imaging Spectrograph (STIS) on the Hubble Space Telescope (HST), several measurements of D/H in the local interstellar medium (LISM, $`d<100`$ pc) have been made with observations of Ly$`\alpha `$ (see Lemoine et al. 1999 for a review). Linsky et al. (1993, 1995) found D/H $`=1.60_{0.19}^{+0.14}\times 10^5`$ toward Capella. Vidal-Madjar et al. (1998) reported evidence for a factor of $`2`$ difference in D/H between two components on the line of sight toward the white dwarf G191-B2B ($`d=75`$ pc). Sahu et al. (1999) re-evaluated the GHRS data as well as new STIS echelle spectra of G191-B2B and concluded that, as a result of improved instrumental characterization (echelle scattered light correction near Ly$`\alpha `$), D/H values in both components appeared to be consistent with the usual LISM value. Questions remain, however. For example, Howk & Sembach (2000) found a different STIS echelle scattered light correction, one that agrees with the GHRS Ly$`\alpha `$ profile (Vidal-Madjar 2000). The resolution of the conflicting conclusions about the G191-B2B sightline may be provided by observations of higher Lyman series lines (Ly$`\beta `$ and Ly$`\gamma `$) by the Far Ultraviolet Spectroscopic Explorer (FUSE, Moos et al. 2000).
The Interstellar Medium Absorption Profile Spectrograph (IMAPS) provides a new window to study deuterium in the Galaxy by virtue of its wavelength coverage and very high spectral resolution ($`\lambda /\mathrm{\Delta }\lambda 80,000`$) which alleviate many of the previous obstacles to accurate D/H measurements. Prior to the IMAPS deuterium studies, D/H measurements in our galaxy after the Copernicus mission have been confined to the local ISM, since the saturated core of H I Ly$`\alpha `$ envelopes the D I Ly$`\alpha `$ feature for $`N`$(H I) $`1\times 10^{19}\mathrm{cm}^2`$ and higher Lyman lines lie beyond the reach of HST. The goal of the IMAPS deuterium program is to obtain high-quality measurements of D/H on sightlines toward bright OB stars beyond the LISM. In the first paper of this series, Jenkins et al. (1999a, hereafter Paper I) measured the D/H abundance ratio toward $`\delta `$ Ori A using IMAPS spectra. They found D/H $`=0.74_{0.13}^{+0.19}\times 10^5`$, a value that is much lower than that found in the LISM. Paper I also showed that the low abundance of D toward $`\delta `$ Ori A is not accompanied by an overabundance of N and O relative to H, as would be expected if the gas had been more thoroughly cycled through stellar interiors. We note that two other deuterium measurements using the Tübingen echelle spectrograph (resolution $`\lambda /\mathrm{\Delta }\lambda 10,000`$) on the ORFEUS-SPAS II mission have been recently reported: Gölz et al. (1998) found D/H $`=0.8_{0.4}^{+0.7}\times 10^5`$ toward BD +28 4211, an O-type subdwarf $``$100 pc away, and Bluhm et al. (1999) measured D/H $`=1.2_{0.4}^{+0.5}\times 10^5`$ toward BD +39 3226, a sdO star located $`270`$ pc from the Sun.
In this paper we analyze the lines of sight toward two well-known hot stars, $`\gamma ^2`$ Velorum (HD 68273) and $`\zeta `$ Puppis (HD 66811). They are among the brightest stars in the sky at 1000 Å, well studied by Copernicus, have H I column densities $`10^{20}\mathrm{cm}^2`$, and hence are prime D/H targets for IMAPS. These stars were observed under an ORFEUS-SPAS II Guest Investigator program.
The sightlines to $`\gamma ^2`$ Vel and $`\zeta `$ Pup have been studied extensively ($`\gamma ^2`$ Vel: Morton & Bhavsar 1979; Sahu 1992; Fitzpatrick & Spitzer 1994; $`\zeta `$ Pup: Morton & Dinerstein 1976; Morton 1978), including D/H measurements with observations from the Copernicus satellite. York & Rogerson (1976) found D/H$`=2.0_{0.7}^{+1.1}\times 10^5`$ toward $`\gamma ^2`$ Vel. Vidal-Madjar et al. (1977) studied the sightline toward $`\zeta `$ Pup and found that solutions consistent with the data permit a wide range of D/H values ($`1.7\times 10^5`$ D/H $`2.5\times 10^4`$). $`\gamma ^2`$ Vel and $`\zeta `$ Pup are the most luminous stars in the Vela-Puppis region. Their role in powering the Gum Nebula and their relationship to the Vela OB2 association, Vela supernova remnants, and various structures in the ISM in Vela- Puppis have been analyzed by Sahu (1992). Table 1 summarizes some basic properties of the two stars. The distances listed were derived from HIPPARCOS data by Schaerer, Schmutz & Grenon (1997), placing $`\gamma ^2`$ Vel significantly closer than earlier estimates.
$`\gamma ^2`$ Vel is the closest (and brightest) Wolf-Rayet star and consequently it has been studied intensively for more than a century (van der Hucht et al. 1981). Nevertheless, due to the complexity of the stellar system, its fundamental stellar parameters have been continuously debated and revised (as recently as 1997 when several papers were published, e.g., Schaerer et al. 1997; Schmutz et al. 1997). $`\gamma ^2`$ Vel is a double-lined spectroscopic binary star composed of a WR star (WC8) and a late O star (O8 III) of comparable brightness. The complex properties of this system are evident in the significant changes in the UV spectrum, and this behavior will be addressed in §4.3 when we discuss our determination of the interstellar hydrogen column density through measurements of the Ly$`\alpha `$ absorption feature.
$`\zeta `$ Pup is an extremely luminous, massive star (O4 Iaf) and the brightest O-type star in the sky. It is a prototype for O-star spectral properties and variability (Morton & Underhill 1977; Massa et al. 1995) and theoretical understanding of radiation-driven winds in massive stars (e.g. Puldrach et al. 1994). $`\zeta `$ Pup is believed to be a single star.
The high velocity resolution provided by IMAPS opens the way for more accurate line profile measurements in the far ultraviolet spectra of bright stars. Our observations of $`\gamma ^2`$ Vel and $`\zeta `$ Pup described in §2 build on our earlier results for $`\delta `$ Ori A described in Paper I. We follow with discussions of how we derived D I column densities (§3), H I column densities (§4), and values for D/H and N/H (§5). The paper concludes with a discussion in §6 of the possible implications of the differences in D/H for the three stars covered in this paper and Paper I.
## 2 OBSERVATIONS
The spectra analyzed in this paper were obtained with IMAPS during the ORFEUS-SPAS II mission (STS-80) in late 1996 (see Hurwitz et al. 1998). The design of IMAPS and its performance on earlier flights have been described in detail by Jenkins et al. (1988, 1996). In Paper I we described the in-flight performance of IMAPS during the 1996 mission. Here we only summarize the instrument’s principal characteristics.
IMAPS is an objective-grating echelle spectrograph designed to record the far-UV spectrum of bright stars at high spectral resolution. The optical design consists of a wire grid collimator to reject off-axis light from stars other than the target, an echelle grating with a 63° blaze angle, and a parabolic cross disperser. The spectral format is imaged on a solid KBr photocathode, whose electrons are magnetically focussed on a windowless, back-illuminated CCD with a 320x256 pixel format. The nominal wavelength range of 930–1150 Å is obtained in four selectable tilts that span the free spectral range of the echelle grating. The gratings were coated with LiF over aluminum, providing excellent throughput longward of 1000 Å. Although the reflectivity of LiF drops substantially shortward of 1000 Å, IMAPS achieved a useful throughput even in the 930 – 980 Å region. The spectral resolution in IMAPS spectra obtained during the 1996 flight was approximately $`\lambda /\mathrm{\Delta }\lambda 80,000`$, or $`4\mathrm{km}\mathrm{s}^1`$. The telluric O I lines (e.g. $`\lambda `$950.112 near Ly$`\delta `$ \- see Fig 1) have FWHM $`5\mathrm{km}\mathrm{s}^1`$, but these lines are probably partly resolved in IMAPS spectra (see Jenkins et al. 1999b).
$`\gamma ^2`$ Vel and $`\zeta `$ Pup were chosen for this IMAPS Guest Investigator program because they were available in 1996 November, have a flux near Ly$`\delta `$ $`>10^9\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1\mathrm{\AA }^1`$, have $`v\mathrm{sin}i100\mathrm{km}\mathrm{s}^1`$, and $`N`$(H I) $`10^{20}\mathrm{cm}^2`$.
$`\gamma ^2`$ Vel was observed by IMAPS on 1996 November 27 14:34–15:11 and 16:00–16:47 UT for a total of 4224 s. $`\zeta `$ Pup was observed on 1996 November 26 18:44–19:23 and 20:16–20:57 UT also for a total of 4224 s. The exposure time at each echelle grating position ranged from 817.6 s to 1226.4 s. These exposure times were chosen to obtain spectra with good signal-to-noise (S/N) ratios near the cores of the D I Ly$`\delta `$ (949.485 Å) and Ly$`ϵ`$ (937.548 Å) lines. Typically, we found S/N=10–15 in the local continuum near the Ly$`\delta `$ and Ly$`ϵ`$ interstellar features.
## 3 COLUMN DENSITY OF ATOMIC DEUTERIUM
### 3.1 General Considerations
The data reduction and analysis of the $`\gamma ^2`$ Vel and $`\zeta `$ Pup spectra and the determination of the D I column densities were identical to those described in Paper I for the line of sight toward $`\delta `$ Ori A. IMAPS spectra were analyzed to determine the total $`N`$(D I) (§3.3 and §3.4). A large number of International Ultraviolet Explorer (IUE) high resolution Ly$`\alpha `$ spectra were obtained from the National Space Science Data Center archive to determine the total H I column density ($`N`$(H I) ) in §4. We note that given that $`N`$(H I) $`10^{20}\mathrm{cm}^2`$ on these sightlines and that D/H $`10^5`$, the gas that gives rise to the H I Ly$`\alpha `$ damping wings is the same gas responsible for the D I Ly$`\delta `$ and Ly$`ϵ`$ features.
In §3.2 we show that N I is a good tracer of H I and D I on these sight lines. The line of sight column density per unit velocity \[$`N_a(v)`$, see Paper I\] for N I, defined by a range of N I lines recorded in the IMAPS spectra, provided a velocity template for modelling the D I profiles. The high S/N for the N I lines helped to constrain the model profiles that gave an acceptable fit to the lower S/N D I lines. This can prevent noise in the D I profiles from giving arbitrarily large or small $`N`$(D I) at specific velocities. We did not decompose the velocity profiles of N I or D I into separate (blended) Gaussian components – the D/H ratios determined here were based on total column densities for each sight line.
Using the method outlined in Paper I, we corrected for the effects of the weak lines Fe II $`\lambda `$937.652 and H<sub>2</sub> Lyman 14–0 $`P`$(2) $`\lambda `$949.351. The Fe II line is located at $`+51\mathrm{km}\mathrm{s}^1`$ in $`\gamma ^2`$ Vel and at $`+53\mathrm{km}\mathrm{s}^1`$ in $`\zeta `$ Pup on the heliocentric velocity scale whose zero point is the laboratory wavelength of D I Ly$`ϵ`$ (see Figs. 1 and 2). The H<sub>2</sub> line is located at $`+10\mathrm{km}\mathrm{s}^1`$ in $`\gamma ^2`$ Vel and at $`+12\mathrm{km}\mathrm{s}^1`$ in $`\zeta `$ Pup on the D I Ly$`\delta `$ heliocentric velocity scale. The strength and shape of these features were computed from other transitions of the same species recorded at longer wavelengths (and higher S/N) in the IMAPS spectra of $`\gamma ^2`$ Vel and $`\zeta `$ Pup. The computed central residual intensities of these features are nearly identical for the two stars: 0.66 for Fe II $`\lambda `$937.6 and 0.95 for H<sub>2</sub> $`\lambda `$949.3. The other potentially contaminating lines noted in Table 1 of Paper I can safely be ignored.
The possible contamination of the Ly$`\delta `$ order by scattered light from an adjacent order that contains the strong absorption lines of the N I $`\lambda \lambda `$954 multiplet was evaluated in the same manner as described in Paper I. The $`\chi ^2`$ analysis found that any residual contamination of the spectrum in the vicinity of Ly$`\delta `$ by the pattern of saturated N I features was $`<1`$% of the continuum in both $`\gamma ^2`$ Vel and $`\zeta `$ Pup. Amplitudes for this correction larger than $`1`$% caused unacceptably bad deviations in the bottom of the H I Ly$`\delta `$ profile. Given the general noise characteristics of the Ly$`\delta `$ profiles, this is a negligible effect.
For the case of $`\gamma ^2`$ Vel, there is another potential contamination source. A companion star, $`\gamma ^1`$ Vel, is located 42″ to the southwest of $`\gamma ^2`$ Vel. Its light should be accepted along with that from $`\gamma ^2`$ Vel, since IMAPS is an objective-grating spectrograph without an entrance slit to reject unwanted sources. (The grid collimator rejects light coming from more than 1° from the axis, however.) $`\gamma ^1`$ Vel is 2.5 magnitudes fainter than $`\gamma ^2`$ Vel in $`V`$ and has a spectral classification of B1 IV (Hoffleit & Jaschek, 1982). Even though $`\gamma ^1`$ Vel is cooler than $`\gamma ^2`$ Vel, its flux is not diminished much below the peak of the Planck distribution at our wavelengths of interest, so we expect its intensity at a given wavelength to be not more than about a factor of 10 fainter than $`\gamma ^2`$ Vel. Fortunately, when the observations were taken the roll angle of the spacecraft about the optical axis of IMAPS (governed by the position of the Sun in the sky) was such that the faint spectrum of $`\gamma ^1`$ Vel was displaced along the cross dispersion direction toward the long-wavelength part of the format. As a consequence, any spectral segment in the spectrum of $`\gamma ^2`$ Vel had light superposed on it from shorter wavelengths in $`\gamma ^1`$ Vel (the separation was slightly less than the distance between 3 echelle orders). The rapid decline in the sensitivity of IMAPS toward shorter wavelengths thus amplified the factor of 10 disparity in relative fluxes at any given position on the format. Therefore, the contamination of the $`\gamma ^2`$ Vel spectrum by light from $`\gamma ^1`$ Vel is negligible. As a final note, on visual inspection we are unable to see any ghost-like spectrum of $`\gamma ^1`$ Vel on top of the spectrum of $`\gamma ^2`$ Vel.
The Lyman line profiles for $`\gamma ^2`$ Vel and $`\zeta `$ Pup are shown in Figures 1 and 2. These profiles have been corrected for Fe II and H<sub>2</sub> line absorption, as described above. The background levels near the D I lines were determined from the broad, saturated cores of the adjacent H I profiles. As explained in Paper I, we compared resolution-degraded forms of the IMAPS profiles with those recorded by Copernicus to test the proposition that the cores of the H I lines indeed represented the zero-intensity levels in the vicinity of the adjacent D I features. This was done to check that we were not being misled by an effect from possible strong, broad wings extending away from the main peak in the instrumental profile of IMAPS. We concluded that indeed the cores of the H I profiles provided very good estimates of the zero levels in the vicinity of the D I lines.
The best answers for $`N`$(D I) and the deviations permitted by the data were evaluated by minimizing $`\chi ^2`$ (see Lampton, Margon & Bowyer 1976 and Bevington & Robinson 1992 for details) when all of our unknown parameters were allowed to vary. The ten free parameters for $`\gamma ^2`$ Vel are (1) $`N`$(D I) , (2) the gas temperature $`T`$, (3-8) the continuum slopes, Y-intercepts, and background levels for Ly$`\delta `$ and Ly$`ϵ`$, (9) a shift in the velocity zero point between the D I lines and the N I template, and (10) a coefficient for scaling the N I contamination signal in Ly$`\delta `$ (see §3.2 in Paper I). The gas temperature is an important parameter because the D I line can be broadened significantly compared to the N I template by thermal motions; see §4.1 in Paper I for details. There are three additional free parameters for $`\zeta `$ Pup because of the addition of Ly$`\gamma `$ to the analysis. Note that the linear continuum fitting (with specific velocity limits as given in §§3.3 and 3.4) is an integral part of the $`\chi ^2`$ evaluation, i.e., the deviations of the continuum levels are considered, in addition to the behavior inside the D I profiles. The same holds for the zero level as defined by the bottom of the adjacent H I profile. We used Powell’s method (Press et al. 1992, p. 406) to find the minimum $`\chi ^2`$. We then set confidence limits by increasing (or decreasing) $`N`$(D I) and $`T`$ with the other parameters freely varying until $`\chi ^2`$ increased by the appropriate amount for the confidence limit of interest with two useful parameters, $`N`$(D I) and $`T`$. We also used different initial values to establish that the minimum $`\chi ^2`$ is unique.
### 3.2 Velocity Profile Templates
In Paper I we discussed the benefits of obtaining a velocity profile based on high quality data for N I and O I, two species that are not significantly depleted in the ISM and which have very similar ionization balances to those of D I and H I (Ferlet 1981; York et al. 1983). This profile information is helpful in constraining the variety of possible interpretations that would produce acceptable fits with the D profiles. For N I, we used IMAPS spectra of the 10 lines in the multiplets at 952.4, 953.8, 954.1 and 1134.7 Å<sup>2</sup><sup>2</sup>2One of the lines in the 1134.7 Å multiplet, the line at 1134.165 Å, had to be omitted from consideration for $`\zeta `$ Pup because it was too close to the edge of the image format.. For all of the N I lines except those in the 952.4 Å multiplet, the background level was easily defined because the stronger lines were saturated. Of course, these lines were useful only for defining the behavior of the gas at velocities somewhat removed from the line core. Since the lines in the multiplet at 952.4 Å were not saturated, we had to determine the background level by a different method. For both $`\gamma ^2`$ Vel and $`\zeta `$ Pup, there are U1 scans of this weak multiplet available in the archive of spectra recorded by the Copernicus satellite (Rogerson et al. 1973). After comparing our IMAPS spectrum degraded to the resolution of Copernicus with the actual Copernicus scans, we determined the adjustments to the background levels that were needed for the IMAPS spectra of this N I multiplet. The backgrounds caused by grating scatter in the Copernicus spectra were taken from the observed count rates in the bottoms of the hydrogen Ly$`\gamma `$ and Ly$`\delta `$ lines, and these levels were subtracted off before the comparison was made.
As was done for our analysis in Paper I, we adopted a method developed by Jenkins & Peimbert (1997) to create a composite profile for the column density of N I as a function of velocity from the 10 lines in the four multiplets, with the weak lines in each case defining the main part of the profile and the strong lines outlining the exact shape of the profile’s wings, well away from the saturated part of the line. For places where there was overlap in the useful portions of the lines, there was satisfactory agreement. The $`f`$-values of Goldbach et al. (1992) were adopted for our analysis.
Figure 3 shows the $`N_a(v)`$ profiles derived for N I toward $`\gamma ^2`$ Vel and $`\zeta `$ Pup. In Paper I we showed that the profiles are unlikely to be contaminated by additional contributions from telluric absorption. $`N`$(N I) is $``$1% of $`N`$(O I) in the Earth’s atmosphere at the altitude of IMAPS ($``$305 km) during the ORFEUS-SPAS II mission. Even for our strongest N I transition used to derive the far wings of $`N_a(v)`$ for N I ($`\lambda `$1134.98), the telluric absorption should be only one third as strong as the O I\* $`\lambda `$950.112 feature present in Figures 1 and 2, and at a heliocentric velocity of $`+12.6\mathrm{km}\mathrm{s}^1`$ it would be buried in the saturated portion of the profile. Consequently, telluric N I has a negligible effect on the interstellar N I profiles.
O I is the best tracer of H I in the ISM since the ionization potential of O I (13.56 eV) is nearly identical to that of H I . The ionization potential of N I (14.53 eV) is only slightly greater that that of H I . Furthermore, both O and N are coupled to H by resonant charge exchange reactions. As discussed by Sofia & Jenkins (1998) and Jenkins et al. (2000), N I closely follows O I and H I unless $`n_\mathrm{e}n_{\mathrm{H}\mathrm{I}}`$.
In the IMAPS wavelength band there is no suitable set of O I transitions to completely define a velocity profile template. The available O I lines are highly saturated or are not detectable, either because they are buried in the core of a Lyman series line or because they are too weak. In Paper I we made use of an archival HST spectrum of the very weak intersystem transition of O I at 1355.6 Å for $`\delta `$ Ori A. There is no such spectrum available for $`\zeta `$ Pup, and the measurement of this line in the spectrum of $`\gamma ^2`$ Vel taken by Fitzpatrick & Spitzer (1994) is too noisy to adequately define the shape of the O I profile. Therefore, we were forced to use the next best option, N I, to define the velocity profile template.
Although we do not have a complete $`N_a(v)`$ profile for O I, we tested the assumption that N I traces H I by comparing selected portions of the O I $`N_a(v)`$ profile derived from available line profiles in $`\gamma ^2`$ Vel. First, we computed $`N_a(v)`$ from the wings of O I $`\lambda `$1039.230. Second, we used the GHRS spectrum of O I $`\lambda `$1355.6 (Fitzpatrick & Spitzer 1994) to define $`N_a(v)`$ in the core of the line. The computed $`N_a(v)`$ was scaled by the difference in the N and O solar abundances ($`0.90`$ dex, Anders & Grevesse 1989). The results, shown in Fig 3, demonstrate that the wings and core of $`N_a(v)`$ for N I and O I are in excellent agreement. Consequently, we have high confidence that the $`N_a(v)`$ profile for N I provides an accurate model for the velocity distribution of H I .
### 3.3 $`\gamma ^2`$ Velorum
The continuum near the Ly$`\delta `$ line shown in Figure 1 was determined in the velocity range $`(40,10)\mathrm{km}\mathrm{s}^1`$ on the blue side of the D I line and $`(+170,+205)\mathrm{km}\mathrm{s}^1`$ on the red side. For Ly$`ϵ`$, the continuum limits were $`(14,+8)\mathrm{km}\mathrm{s}^1`$ and $`(+170,+210)\mathrm{km}\mathrm{s}^1`$. The background level was determined from the core of the two H I lines in the velocity range $`(+60,+140)\mathrm{km}\mathrm{s}^1`$. The D I profiles were fit in the velocity range $`(10,+39)\mathrm{km}\mathrm{s}^1`$ for Ly$`\delta `$ and $`(+8,+39)\mathrm{km}\mathrm{s}^1`$ for Ly$`ϵ`$. With a spacing of independent velocity samples of $`1.25\mathrm{km}\mathrm{s}^1`$, this resulted in 276 degrees of freedom in the $`\chi ^2`$ analysis to fit the 10 free parameters described in §3.1.
By adjusting our estimate for the noise, which was only known to limited accuracy beforehand, we achieved a minimum value for $`\chi ^2`$ equal to 247.0 at a column density of $`1.12\times 10^{15}\mathrm{cm}^2`$. For 276 degrees of freedom, there was a 90% chance that we would have found this value of $`\chi ^2`$ or greater. With this 90% confidence level, we arrived at a conservatively high estimate for the noise. It then followed that this noise level, which is perhaps higher than the real noise, gave a conservatively large confidence interval for $`N`$(D I) , as determined by how rapidly the values of $`\chi ^2`$ deviated from the minimum value. With the noise level having been set in the manner just described, we explored for 90% and 99% confidence limits for $`N`$(D I) (i.e., $`1.65\sigma `$ and 2.58$`\sigma `$ deviations), which correspond to $`\chi ^2`$(min)+4.6 and $`\chi ^2`$(min)+9.2, respectively. Table 2 lists $`N`$(D I) and $`T`$ for the best value and these limits. The $`\pm `$90% confidence limits on the model D I profiles (cross-hatched regions) for Ly$`\delta `$ and Ly$`ϵ`$ are compared with the observed deuterium profiles in Figure 4.
In Figure 4 we also show with a dashed line the expected shape and depth of the D I profile for $`N`$(D I) =$`7.7\times 10^{14}\mathrm{cm}^2`$. This corresponds to a D/H ratio of $`1.5\times 10^5`$, assuming the H I column density derived below in §4. For this case, $`\chi ^2`$$`\chi ^2`$(min) = 52.4, which is clearly an unacceptable fit. The dot-dashed line in Figure 4 corresponds to the D I profile for $`N`$(D I) =$`3.8\times 10^{14}\mathrm{cm}^2`$, the D I column density necessary to achieve D/H=$`0.74\times 10^5`$, the value found toward $`\delta `$ Ori A in Paper I. A similar demonstration for the measurement of $`N`$(H I) is be given in §5.
### 3.4 $`\zeta `$ Puppis
The measurement of $`N`$(D I) toward $`\zeta `$ Pup used the same methodology described in the previous section with the following modifications. In addition to Ly$`\delta `$ and Ly$`ϵ`$ we were able to include the D I profile of Ly$`\gamma `$ in the $`\chi ^2`$ analysis. The blue and red regions for the linear continuum fits were $`(40,10)\mathrm{km}\mathrm{s}^1`$ and $`(+170,+190)\mathrm{km}\mathrm{s}^1`$ for Ly$`\gamma `$, $`(39,10)\mathrm{km}\mathrm{s}^1`$ and $`(+157,+190)\mathrm{km}\mathrm{s}^1`$ for Ly$`\delta `$, and $`(4,+3)\mathrm{km}\mathrm{s}^1`$ and $`(+180,+220)\mathrm{km}\mathrm{s}^1`$ for Ly$`ϵ`$. The background levels were determined from the saturated cores of the H I lines over these velocity limits: $`(+55,+130)`$ for Ly$`\gamma `$ and Ly$`\delta `$, $`(+65,+115)`$ for Ly$`ϵ`$. The model D I profiles were fit over the range $`(5,+40)\mathrm{km}\mathrm{s}^1`$ for Ly$`\gamma `$, $`(0,+40)\mathrm{km}\mathrm{s}^1`$ for Ly$`\delta `$, and $`(+3,+37)\mathrm{km}\mathrm{s}^1`$ for Ly$`ϵ`$. This resulted in 426 degrees of freedom to simultaneously fit the 13 free parameters described in §3.1. The minimum $`\chi ^2`$ is 388.5 and corresponds to the most probable value of $`N`$(D I) for $`\zeta `$ Pup, namely $`1.30\times 10^{15}\mathrm{cm}^2`$. The results of the $`\chi ^2`$ analysis are summarized in Table 3, including the 90% and 99% confidence limits. The same conservative treatment of the noise described for $`\gamma ^2`$ Vel was followed for $`\zeta `$ Pup. Figure 5 shows the observed D I Lyman line profiles and the range in the best fit model profiles and continua fits allowed by the 90% confidence limits.
## 4 COLUMN DENSITY OF ATOMIC HYDROGEN
### 4.1 Analysis Methodology
The primary objective of this IMAPS program is to determine D/H ratios of sufficient accuracy to test for spatial inhomogeneities. To achieve this goal, we must strive for a precision in $`N`$(H I) that is as good as (or preferably better than) that of the deuterium measurement. We must also understand the magnitude and nature of the uncertainties in $`N`$(H I) and $`N`$(D I) so that we can realistically assess the uncertainty in D/H for comparison to other measurements in the ISM. For these reasons, we decided to conduct our own investigation of the H I column densities toward $`\gamma ^2`$ Vel and $`\zeta `$ Pup, even though a number of measurements of $`N`$(H I) toward these stars are already available in the literature (e.g. Jenkins 1971; York & Rogerson 1976; Vidal-Madjar et al. 1977; Bohlin, Savage, & Drake 1978; Shull & Van Steenberg 1985; Diplas & Savage 1994).
We described our method for determining $`N`$(H I) in Paper I; here we provide a summary. In §4.2 and §4.3 we discuss details particular to the $`\zeta `$ Pup and $`\gamma ^2`$ Vel sight lines and present the results. We used the Ly$`\alpha `$ line to determine the total $`N`$(H I) .<sup>3</sup><sup>3</sup>3The Copernicus studies of D/H on the $`\gamma ^2`$ Vel and $`\zeta `$ Pup sight lines used the much weaker damping wings of Ly$`\beta `$ or Ly$`\gamma `$ to determine $`N`$(H I) . Since H I Ly$`\gamma `$ and Ly$`\delta `$ are on the flat part of the curve of growth, high velocity gas could influence their profiles. Ly$`\alpha `$ is immune from this potential problem. Due to the great breadth of the Lorentzian wings ($`\pm 1000\mathrm{km}\mathrm{s}^1`$ – see the velocity scale in Figure 9), Ly$`\alpha `$ absorption due to any high velocity interstellar gas ($`\pm 100\mathrm{km}\mathrm{s}^1`$) that may be present is confined to the saturated core of the profile. In principle, some high velocity H I , which is too widely dispersed in velocity to show up in the D I profiles, could contribute to the Ly$`\alpha `$ damping wings. However, this gas is not detected in the strong lines N I $`\lambda `$1134.98 and O I $`\lambda `$1039.23 observed with IMAPS. It is clear that such high velocity gas, if present, would have a low column density and a negligible effect on the Ly$`\alpha `$ damping wings. Since the Ly$`\alpha `$ line is not covered by IMAPS, was not observed by the HST spectrographs, and the available Copernicus Ly$`\alpha `$ scans suffer from several problems (see Paper I), we analyzed high dispersion IUE observations of Ly$`\alpha `$ to determine $`N`$(H I) . Furthermore, these stars were observed many times with IUE over the course of many years, and this offers an opportunity to evaluate potential sources of systematic error. For example, $`\gamma ^2`$ Vel is a spectroscopic binary, and the large IUE database allowed a search for systematic changes in the derived $`N`$(H I) as a function of orbital phase.
After screening and retrieving the IUE data<sup>4</sup><sup>4</sup>4We used the standard IUESIPS calibration of the IUE high dispersion spectra instead of the NEWSIPS calibration to avoid potential problems with background corrections and zero-level determination near Ly$`\alpha `$ found in some IUE SWP high dispersion spectra processed with NEWSIPS (see Massa et al. 1998), which could adversely affect our measurements. of interest and correcting for interorder scattered light (see Paper I), we used an approach first used by Jenkins (1971) to constrain the H I column density, i.e., we determined the $`N`$(H I) which provides the best fit to the Ly$`\alpha `$ profile with the optical depth $`\tau `$ at a given wavelength $`\lambda `$ calculated from the expression
$$\tau (\lambda )=N(\mathrm{H}\mathrm{I})\sigma (\lambda )=4.26\times 10^{20}N(\mathrm{H}\mathrm{I})(\lambda \lambda _0)^2$$
(1)
where $`\lambda _0`$ is the Ly$`\alpha `$ line center at the velocity centroid of the hydrogen (note that the useful portion of the Ly$`\alpha `$ profile is entirely due to the Lorentzian wings and the effects of instrumental and Doppler broadening can be neglected). As was done for D I , we also determined the important free parameters which can be adjusted to fit the H I Ly$`\alpha `$ absorption profile, then we found the set of parameters that minimized $`\chi ^2`$ using Powell’s method. We set confidence limits on the H I column as described in §3.1 with only one parameter of interest, $`N`$(H I) . The other (uninteresting) free parameters we selected for fitting the H I profile were three coefficients which specify a second-order polynomial fit to the continuum, and an additive correction to the flux zero point. The continuum was fit to the spectrum in several windows covering the range 1185 Å to 1276 Å. Despite our use of the Bianchi & Bohlin (1984) correction for interorder scattered light, in many cases inspection of the flat-bottomed, saturated portion of the Ly$`\alpha `$ profile showed that the zero intensity level was not quite correct. A zero point shift in the flux scale, one of the terms for evaluating $`\chi ^2`$, was determined from a region within the saturated core. For both stars, the zero point of the Ly$`\alpha `$ wavelength scale was set by placing the N I $`\lambda `$1200 multiplet in agreement with the IMAPS N I profiles. The Ly$`\alpha `$ profile was then fit only to the red wing because of the presence of strong stellar features superposed on the blue wing. However, the uncontaminated portion of the blue wing of Ly$`\alpha `$ was checked for consistency with the fit to the red wing and found to be in good agreement.
### 4.2 $`\zeta `$ Puppis
In Paper I, considerable attention was paid to the spectroscopic binary nature of $`\delta `$ Ori A and whether or not this could cause systematic errors in $`N`$(H I) . A similar analysis of $`\gamma ^2`$ Vel is presented in §4.3. The determination of $`N`$(H I) toward $`\zeta `$ Pup was less complex than the other stars observed by IMAPS to study D/H. As far as is known, $`\zeta `$ Pup is not a spectroscopic binary star. It is an O4 supergiant, so the stellar H I Ly$`\alpha `$ absorption line makes a negligible contribution to the H I absorption profile. The star is possibly a non-radial pulsator with very weakly variable stellar absorption lines with a period of 8.5 hr (Reid & Howarth 1996), and the variability of the stellar wind P-Cygni profiles is well known and studied (e.g., Prinja et al. 1992; Howarth, Prinja, & Massa 1995) and shows significant power at 19.2 hr and 5.2 day periods. Both of these sources of spectral variations are expected to have little impact on the interstellar H I column density derived from Ly$`\alpha `$, but this can be checked given the large number of observations and good temporal sampling.
There are more than 200 high-dispersion observations of $`\zeta `$ Pup in the IUE archive, primarily because the star was selected for intensive IUE observing programs to study stellar wind variability in massive stars. We concentrated on the spectra obtained in three specific multi-day observing runs in 1986, 1989, and 1995. In 1995 $`\zeta `$ Pup was observed continuously for 16 days (dubbed the ”MEGA“ campaign, Massa et al. 1995). We omitted three observations from these observing sessions because the archival data were corrupted or unavailable. This left 189 observations for measurement of $`N`$(H I) . All of the observations were obtained with the IUE large aperture and the signal-to-noise ratios are comparable; the $`\zeta `$ Pup data set is more uniform than the IUE Ly$`\alpha `$ data used for measuring $`N`$(H I) toward $`\delta `$ Ori A and $`\gamma ^2`$ Vel.
The $`N`$(H I) values derived from each individual observation are plotted in Figure 6, along with the 1$`\sigma `$ uncertainties, versus the observation date. Table 4 summarizes the mean H I column density $`<`$$`N`$(H I) $`>`$ and the root mean square (rms) dispersion $`\sigma `$ obtained from the 1986, 1989, and 1995 data sets analyzed separately, and $`<`$$`N`$(H I) $`>`$ for each data set is indicated with a heavy dashed line in Figure 6. We also list in Table 4 the formal error in the mean $`ϵ=\sigma /\sqrt{N}`$ (here $`N`$ is the number of measurements) and the reduced $`\chi ^2`$ for the three data sets.
In Paper I we discussed the potential sources of systematic error in the derivation of $`N`$(H I) from IUE spectra, and we suggested that the real uncertainty in the mean is likely to be greater than the formal estimate. The very large Ly$`\alpha `$ data set for $`\zeta `$ Pup shown in Figure 6 appears to confirm that this is indeed the case for this star. From this figure one can see that the column densities derived from the 1995 data are systematically lower than the column densities derived from the 1989 observing run. The mean of the 1995 data set is lower than the mean of the 1989 data by $`0.54\times 10^{19}\mathrm{cm}^2`$ (0.0256 dex), and this difference is substantially greater than the formal error estimates ($`ϵ`$).
This systematic error could be the result of secular changes in the stellar mass loss or circumstellar environment, a previously unrecognized stellar variability with a long period, or instrumental effects. Or perhaps some aspect of the interstellar sight line changed during this period. Since the IMAPS observations for $`N`$(D I) were made almost two years after the MEGA campaign, it is possible that a similar systematic error is present in any IUE estimate of $`N`$(H I) that we choose for comparison with the IMAPS $`N`$(D I) . The exact magnitude of this systematic uncertainty is difficult to estimate, but we find that $`N`$(H I) obtained from additional large aperture spectra of $`\zeta `$ Pup taken at other epochs are in good agreement with the estimates in Table 4. This indicates that the systematic error is not much larger than 0.0256 dex. For example, large aperture observations taken in 1988 May, 1989 December, and 1991 March, combined with the three $`<`$$`N`$(H I) $`>`$ values for the three data sets in Table 4, yield a mean $`N`$(H I) = $`9.34\pm 0.67\times 10^{19}\mathrm{cm}^2`$. We conclude that 0.0256 dex is a conservative estimate of the $`1\sigma `$ uncertainty in $`N`$(H I) . We therefore adopt the unweighted mean of the three $`<`$$`N`$(H I) $`>`$ values in Table 4 as the best estimate of $`N`$(H I) toward $`\zeta `$ Pup and 0.0256 dex as its $`1\sigma `$ uncertainty: $`N`$(H I) $`=9.18\pm 0.54\times 10^{19}\mathrm{cm}^2`$. We elected to use a mean of the $`<`$$`N`$(H I) $`>`$ values in Table 4 instead of a straight mean of all of the individual measurements shown in Fig. 6 because the latter would give excessive weight to the 1995 data set due to the much larger number of measurements obtained during that observing campaign.
### 4.3 $`\gamma ^2`$ Velorum
As noted in §1, $`\gamma ^2`$ Vel is a complex stellar system with a strongly variable UV stellar spectrum. Figure 7 shows two examples of the high-dispersion IUE spectra of $`\gamma ^2`$ Vel that we have used to derive $`N`$(H I) . These examples were selected to illustrate the variability of the P-Cygni emission lines in the vicinity of Ly$`\alpha `$, which are weakest at orbital phase $``$0.5. Figure 7 also shows the complex and variable stellar absorption superposed on much of the blue wing of the interstellar Ly$`\alpha `$ profile. This variability was a source of concern for this paper. Can we derive reliable H I column densities despite the complex spectral changes occurring in this star? Due to the stellar absorption, the blue wing of Ly$`\alpha `$ was not useful for constraining $`N`$(H I) . Fortunately, the red wing of Ly$`\alpha `$ was clean and stable as a function of orbital phase. This can be seen by comparing panel (a) to panel (b) in Figure 7. In this paper, we used only the red wing for fitting the Ly$`\alpha `$ profile, and the resultant fits generally look quite good (see Figures 7 and 9).
We examined $`N`$(H I) as a function of orbital phase of the spectroscopic binary to check for systematic errors in the derived interstellar H I column density due to the complicated variability of the stellar spectrum. There are considerable discrepancies in the orbital elements published by various groups (Stickland & Lloyd 1990; Schmutz et al. 1997, and references therein). We adopt the orbital elements derived by Schmutz et al. (1997) who find a period of $`78\stackrel{\mathrm{d}}{\mathrm{.}}53\pm 0\stackrel{\mathrm{d}}{\mathrm{.}}01`$ d with velocity semiamplitudes of $`K_{\mathrm{WR}}=122\pm 2\mathrm{km}\mathrm{s}^1`$ and $`K_\mathrm{O}=38.4\pm 2\mathrm{km}\mathrm{s}^1`$. In addition to the usual intrinsic variability of WR stars, it is believed that the $`\gamma ^2`$ Vel spectrum may also change as a result of periodic absorption of the O star component by the WR wind (Stickland & Lloyd 1990).
After screening the IUE high-dispersion observations of $`\gamma ^2`$ Vel, we were left with 42 spectra, 34 small aperture observations obtained early in the IUE mission and 8 later observations obtained with the large aperture. In Figure 8 we plot the H I column densities derived from all of the observations as a function of the phase of the spectroscopic binary, and Table 5 summarizes the $`<`$$`N`$(H I) $`>`$, $`\sigma `$, and $`ϵ`$ derived from the large aperture data only, the small aperture data only, and all of the data combined. While one might argue that a trend is apparent in Figure 8 with a minimum in the derived $`N`$(H I) values at orbital phase $``$ 0.5, this is a marginal result at best, and we do not believe that it should be taken too seriously. On the contrary, the generally good agreement of the H I column densities at various phases despite the dramatic changes in the stellar spectrum (see Figure 7) is encouraging, and we conclude from Figure 8 that the IUE data provide a good determination of $`<`$$`N`$(H I) $`>`$. Since the small and large aperture measurements are in good agreement, we adopted the results from the combined data set. Systematic errors probably cause the true uncertainty in $`<`$$`N`$(H I) $`>`$ to exceed the formal error estimate provided in Table 5. Drawing from our experience with the more extensive data set for $`\zeta `$ Pup, we conservatively adopt 0.0256 dex as the $`1\sigma `$ uncertainty in $`N`$(H I) toward $`\gamma ^2`$ Vel as well: $`N`$(H I) $`=5.13\pm 0.30\times 10^{19}\mathrm{cm}^2`$.
## 5 DEUTERIUM AND NITROGEN ABUNDANCE RATIOS
### 5.1 D/H
Combining the results for D I and H I column densities derived in the previous sections, we determined the atomic D/H abundance ratio for $`\gamma ^2`$ Vel and $`\zeta `$ Pup. Toward $`\gamma ^2`$ Vel we found that D/H $`=2.18_{0.31}^{+0.36}\times 10^5`$ (the errors are 90% confidence limits). This result is a large departure from the value of D/H usually attributed as “typical” in the ISM of the Galaxy ($`1.5\times 10^5`$). A large value for the average D/H ratio implies that even larger values could exist in individual components if other components have lower D/H values closer to the LISM ratio. Although the previous D/H measurement for $`\gamma ^2`$ Vel ($`=2.0_{0.7}^{+1.1}\times 10^5`$, York & Rogerson 1976) appears to be in close agreement with the IMAPS result, we consider this to be coincidental in view of the magnitude of the uncertainties in the Copernicus measurement.
For $`\zeta `$ Pup we find $`N`$(D I) $`=1.30_{0.17}^{+0.19}\times 10^{15}\mathrm{cm}^2`$, a value slightly below the lower limit derived from Copernicus spectra by Vidal-Madjar et al. (1977). They found that a large range in $`N`$(D I) ($`1.520.\times 10^{15}\mathrm{cm}^2`$) was possible for this sight line due to its complexity and a lack of adequate constraints. Since the value of $`N`$(H I) derived by Vidal-Madjar et al. (1977) for $`\zeta `$ Pup is consistent with ours, their large range in D/H must be a result of the $`N`$(D I) uncertainty. Much more is known now about the complexity of this sightline than at the time of the Copernicus study (e.g. Welty, Morton & Hobbs 1996). More importantly, the IMAPS spectra have sufficient spectral resolution that the structure of the $`\zeta `$ Pup sightline in neutral hydrogen can be defined by the $`N_a(v)`$ profile for N I and used to constrain the determination of $`N`$(D I) . Thus with higher spectral resolution and better knowledge of the velocity structure along the sightline, a more accurate value of $`N`$(D I) and D/H is derived. We find D/H $`=1.42_{0.23}^{+0.25}\times 10^5`$. The $`N`$(D I) , $`N`$(H I) , and D/H results for the three stars studied by IMAPS are summarized in Table 6, where all errors are presented as 90% confidence limits.
We now address a critical question, namely, is there sufficient uncertainty in the $`N`$(H I) and $`N`$(D I) measurements to reconcile the $`\gamma ^2`$ Vel D/H abundance ratio with the general result observed in the local ISM or with the D/H derived in Paper I for the sight line to $`\delta `$ Ori A? Figure 9a shows with dotted and dashed lines the expected appearance the Ly$`\alpha `$ profile would have if our measurement of $`N`$(D I) is correct but the $`N`$(H I) were made high enough so that D/H $`=1.5\times 10^5`$. In Figure 9, the dotted line represents the H I profile that would be needed for the most probable value for $`N`$(D I) , while the dashed lines represent the 90% confidence limits. For comparison, we also show in Figure 9a the best-fitting H I profile (solid line) derived in the previous section for this particular observation. Figure 9b is an analogous plot with $`N`$(H I) forced to take on a value so that D/H $`=0.74\times 10^5`$ as observed toward $`\delta `$ Ori A (Paper I). From Figure 9 it is clear that forcing $`N`$(H I) to make D/H $`=1.5\times 10^5`$, assuming $`N`$(D I) from §3.3, produces a very poor (and unacceptable) fit to the H I profile, and the fit becomes ridiculous if D/H $`=0.74\times 10^5`$. We have assumed the best-fit continuum (solid line) to compute the Ly$`\alpha `$ profiles shown in Figure 9. We have also explored whether or not the D/H ratios can be reconciled by adjusting the continuum placement, and we found that absurd continuum placements and shapes must be assumed to make D/H $`=1.5\times 10^5`$ or (even worse) $`0.74\times 10^5`$.
As demonstrated at the end of §3.3, the changes in $`N`$(D I) needed to make D/H $`=1.5\times 10^5`$ or $`=0.74\times 10^5`$, assuming our most probable value of $`N`$(H I) , gave a large mismatch with the data, which were clearly unacceptable. We conclude that spatial inhomogeneities in D/H in the ISM within $``$500 pc of the Sun have high statistical significance, especially when contrasting the D/H ratios in the directions of $`\gamma ^2`$ Vel and $`\delta `$ Ori A, which were measured with the same instrument and techniques.
### 5.2 Nitrogen Abundances
The N I column density $`N`$(N I) was computed by integrating the column density profiles shown in Figure 3. The $`N`$(N I) results for the three targets are listed in Table 6. This table also contains the resulting N/H and D/N abundance ratios. We examined the potential error in $`N`$(N I) by examining the errors in the portions of the various N I profiles used to construct $`N_a(v)`$. The $`1\sigma `$ uncertainty in $`N`$(N I) for $`\gamma ^2`$ Vel and $`\delta `$ Ori A is conservatively estimated to be 5% since the N I profiles were all of high quality. There may be some question about the structure in the core of $`N_a(v)`$ for $`\zeta `$ Pup. The uncertainty in $`N`$(N I) for $`\zeta `$ Pup was estimated by supposing the spike in $`N_a(v)`$ at $`v19\mathrm{km}\mathrm{s}^1`$ is due largely to noise fluctuations in the core of N I $`\lambda `$952.5, the weakest N I line we detect on this sightline. If we truncate the $`N_a(v)`$ profile at the level of the secondary maximum at $`v24\mathrm{km}\mathrm{s}^1`$ \[$`N_a(v)`$$`=5.6\times 10^{14}\mathrm{cm}^2(\mathrm{km}\mathrm{s}^1)^1`$\], we find that the area of the spike is $`3.6\times 10^{14}\mathrm{cm}^2`$, or 4.7% of the total. We then suppose that the total uncertainty in $`N`$(N I) for $`\zeta `$ Pup is $`\sqrt{2}`$ times this, or 1$`\sigma =5.1\times 10^{14}\mathrm{cm}^2`$, to allow for the possibility that there may be other such fluctuations.
## 6 DISCUSSION
There are two principal conclusions of the IMAPS D/H program. First, the atomic D/H ratio in the ISM, averaged over path lengths of 250 to 500 pc, exhibits significant spatial variability. Differences in the atomic D/H ratio on long pathlengths in the ISM have been suspected for many years (Vidal-Madjar et al. 1978; Vidal-Madjar & Gry 1984), but not until now have data of sufficient quality been available to evaluate and reduce statistical and systematic errors to levels where these differences are unequivocal. Second, we find no support for the simple picture that variations in D/H anticorrelate with those of N/H, i.e., one measure of how much the gas has been processed through stellar interiors. Figure 10 shows the relationship between D/H and N/H for the three sightlines studied by IMAPS plus the white dwarf G191$``$B2B (Vidal-Madjar et al. 1998; Sahu et al. 1999). We point out that some elements are systematically removed from the gas phase as they are incorporated into interstellar dust (Savage & Sembach 1996), but the abundance of N does not seem to be appreciably altered by this effect (Meyer et al. 1997).
Beyond the effects from depletions onto dust, spatial variations in interstellar gas abundances can arise as a natural consequence of galactic chemical evolution and the changing influences of different stellar populations. The lack of an anticorrelation between N/H and D/H depicted in Fig. 10 indicates that the variability of D/H is not just a consequence of different mixing ratios of material with differing levels of stellar processing, as we might anticipate, for instance, from the variable addition of metal-poor, infalling gas from the Galactic halo (Meyer et al. 1994). We may need to go further and draw a distinction between contributions from stars that simply destroy deuterium and those that both destroy deuterium and enrich the medium with additional nitrogen. That is, we could envision some stars cycling material only through their shallow layers that are only hot enough to burn deuterium, while others eject material from much deeper layers where the synthesis of heavier elements has taken place. This additional level of complexity could explain the behavior that we observed.
Global models of Galactic chemical evolution (Audouze & Tinsley 1974; Tosi 1988a,b; Dearborn, Steigman & Tosi 1996; Scully et al. 1996; Tosi et al. 1998) describe the destruction of D during stellar formation, evolution, and eventual mass loss. These models predict variations in D/H, N, and O abundances that are manifested as abundance gradients on a scale of $``$1 kpc. However, the predicted trends in galactic abundances may not accurately represent what is observable in the diffuse ISM. Tenorio-Tagle (1996) showed that the chemical enrichment of the diffuse ISM by OB associations, including supernovae from massive stars, is a slow process. Following a supernova explosion, chemically enriched ejecta remain clumpy and not well mixed with the diffuse ISM it encounters until it is incorporated into new star forming regions. This is a result of very long time scales for diffusion between different parcels of gas in the warm ($`10^4`$K) and cold ($`10^2`$K) phases of the ISM.
The diffusion time scale for enriched gas to thoroughly mix with the warm diffuse ISM can be very long ($`t_D>10^{10}`$ yr - Tenorio-Tagle 1996). On the spatial scale sampled by the IMAPS observations, different sight lines may encounter regions with very different dynamical and chemical histories. Differential galactic rotation and random cloud motions are expected to stir the diffuse ISM and chemically enriched parcels of gas, but these parcels retain their distinct chemical properties until they are disrupted through photoevaporation, most likely by the formation of new massive stars. Tenorio-Tagle finds that diffusion is efficient only for the hot phase of the ISM ($`t_D<10^6`$ yr), which accounts for a very small fraction of the total diffuse ISM. Only after the enriched gas and diffuse gas are highly ionized would the chemically enriched gas from the earlier generation of stars quickly diffuse into the ambient ISM. Thus, the time scale for mixing interstellar gases with different processing histories can be much longer than the chemical evolution time scale. However, it is possible that interstellar turbulence and its secondary phenomena may accelerate the mixing rate. Since the distribution of star forming regions (OB associations) shows large inhomogeneities on scales $`1`$ kpc, their corresponding chemical enrichment of the ISM may be expected to be nonuniform as well. This is perhaps revealed indirectly by variations in H II region abundances (Peimbert 1999) and solar-type stars at similar galactocentric radii (Edvardsson, et al. 1993).
Several processes unrelated to stellar nucleosynthesis may, under the right circumstances, alter the atomic D/H ratio of some parcels of interstellar gas (see Lemoine et al. 1999 for a review). D may be incorporated into HD (Watson 1973), but the fraction of molecular gas on our sight lines is very low (see Paper I). Differential radiation pressure on D and H (Vidal-Madjar et al. 1978; Bruston et al. 1981) may lead to a separation of D in some clouds near strong radiation fields. Adsorption of D onto dust grains (Jura 1982) may deplete D from the gas phase. Bauschlicher (1998) found that reactions of H and D with polycyclic aromatic hydrocarbon (PAH) cations might systematically provide some D enrichment in PAHs. A very different perspective has been offered by Mullan & Linsky (1999), who suggested that significant quantities of D may be formed in stellar flares from M dwarf stars and ejected into interstellar space. Some or all of these processes could be at work in the diffuse ISM and alter the atomic D/H ratio on individual sight lines, independent of the degree of chemical enrichment from stellar evolution. However, they have not yet been demonstrated to be quantitatively significant. Observational and theoretical tests of the efficiency and applicability of these processes are needed to better understand the mechanisms affecting the D/H ratio in the diffuse ISM.
While the D/H ratios derived from IMAPS spectra are in the general range expected from galactic chemical evolution models, a factor of three variation in the mean D/H ratio on path lengths of several hundred pc is unexpected. The apparent lack of an anti-correlation of the D/H abundance ratio with the metallicity of the gas and its variability on smaller than expected scales suggests that other processes in the Galaxy may be masking more general chemical evolution trends. This may pose a problem for deriving a “primordial” D/H by extrapolating back from D/H measurements in the Milky Way to extragalactic absorbers at higher and higher redshifts, or to even to evaluate “primordial” D/H directly from high redshift observations, until we understand the reasons for these differences. The spatial variations found in this study underline the importance of high-quality D/H determinations. D/H measurements in more distant regions of the Galaxy are needed to determine whether the properties of the gas within 500 pc of the Sun are representative of the Galactic disk. Observations with the FUSE satellite should probe such more distant environments and hopefully answer some of these questions.
The value of D/H for $`\gamma ^2`$ Vel is larger than that usually considered typical for the Milky Way. This robust result establishes a new lower limit to the primordial D/H ratio. Within the framework of standard Big Bang Nucleosynthesis (Walker et al. 1991), the large value of D/H found toward $`\gamma ^2`$ Vel is equivalent to a cosmic baryon density of $`\mathrm{\Omega }_Bh^2=0.023\pm 0.002`$. This error simply reflects the uncertainty in the D/H ratio toward $`\gamma ^2`$ Vel reported in this paper. We regard this value of $`\mathrm{\Omega }_Bh^2`$ as an upper limit since no correction has been applied for the destruction of deuterium in stars. This upper limit on $`\mathrm{\Omega }_Bh^2`$ is consistent with the preferred values of $`\mathrm{\Omega }_Bh^2`$ derived from recent analyses of the BOOMERANG and MAXIMA cosmic microwave background measurements (e.g., Lange et al. 2000; Tegmark & Zaldarriaga 2000; Hu et al. 2000). However, any lowering of this upper limit on $`\mathrm{\Omega }_Bh^2`$ to correct for astration will lead to a marginal disagreement with simple inflation models (see Figure 4 in Tegmark & Zaldarriaga 2000) and requires adjustments of other cosmological parameters. Alternatively, the $`\mathrm{\Omega }_Bh^2`$ upper limit from D/H may be taken as a prior assumption for the constraint of other cosmological parameters using the CMB data (e.g., see Hu et al. 2000).
We wish to thank the US and German space agencies, NASA and DARA, for their joint support of the ORFEUS-SPAS II mission that made these observations possible. The successful execution of our observations was the product of efforts over many years by engineering teams at Princeton University Observatory, Ball Aerospace Systems Group, and Daimler-Benz Aerospace. Important contributions to the success of IMAPS also came from the efforts of D. A. Content and R. A. Keski-Kuha and other members of the Optics Branch of the NASA Goddard Space Flight Center and from O. H. Siegmund and S. R. Jelinsky at the Berkeley Space Sciences Laboratory. We also thank Bruce Draine for insightful discussions on atomic interactions with dust grains and PAHs. This work was supported in part by NASA grant NAG5-616 to Princeton University. The IUE data were obtained from the National Space Science Data Center at NASA-Goddard.
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# 1 Introduction
## 1 Introduction
High Temperature Superconductors (HTSC) show a highly complex phenomenology, as it is evident from the extremely rich phase diagram; this issue, however, is not surprising since in this class of materials the electronic bands are very narrow. Kinetic energies are accordingly quite small, so that the system is sensitive to any kind of instability which may give rise to visible effects. In this situation, it is clear that, in order to provide an effective description of the properties of HTSC, electronic correlation may not be neglected.
Strongly correlated systems show some common characteristic features. On one hand the spectral weight associated with the coherent part of the one-particle spectral function is reduced with respect to the uncorrelated case and meanwhile a broad incoherent background arises. In addiction, the electronic dispersion of the coherent peak is narrowed by correlation effects. As a result, in a correlated system electronic kinetic energies, characterized by the Fermi energies $`E_F`$, may become comparable with phonon frequencies $`\omega _{ph}`$, leading to the breakdown of the adiabatic hypothesis (i.e. $`\omega _{ph}/E_F1`$) on which conventional superconductivity theory rests.
In addiction to the modification of one-particle properties, electronic correlation affects as well the two-particle ones of the system. In particular different analytical approaches point out to a predominance of the forward scattering (small q) in the charge density response function.<sup>?</sup> This feature, together with the above mentioned bandwidth reduction, makes Migdal’s theorem <sup>?</sup> unappliable.
In this paper we investigate the effects of electronic correlation on the electron-phonon coupling. The modulation of the electron-phonon scattering induced by correlation results to be particularly important in the Cooper channel leading to new feature of the superconductivity properties.
## 2 The Model
Systems of correlated electrons are usually described in terms of a Hubbard Hamiltonian, which can be considered a good paradigmatic model<sup>?</sup>:
$$\widehat{H}=\underset{i,j,\sigma }{}t_{i,j}\widehat{c}_{i,\sigma }^{}\widehat{c}_{j,\sigma }+U\underset{i}{}\widehat{n}_{i,}\widehat{n}_{i,}.$$
(1)
The behaviour of the system can easily be determined in the two limit cases $`(t/U)\mathrm{}`$ and $`(t/U)0`$. In the first case the Hubbard Hamiltonian reduces to the free-electron one: the system is a perfect Fermi gas, with energy levels $`\epsilon _𝐤`$ and Fermi energy $`E_F`$. In the second case the most favorable configuration is that in which the electrons are all “localized”, i.e. no hopping between sites is possible. In the intermediate region of parameters one expects a coexistence of the two behaviours.
In terms of the spectral function $`A(𝐤,\omega )`$ the situation can be described as follows: when $`(t/U)=\mathrm{}`$, $`A(𝐤,\omega )`$ is a delta-function centered at the energy $`\epsilon _𝐤`$, fulfilling the sum rule $`_𝐤𝑑\omega A(𝐤,\omega )e^{i\omega 0^+}=N`$, where $`N`$ is the total number of particles. As $`(t/U)`$ decreases, the delta-function, which represents itinerant electrons, will ‘lose weight’, i.e. the coherent spectral weight will be such that $`_𝐤𝑑\omega A_{co}(𝐤,\omega )e^{i\omega 0^+}=NZ`$, with $`Z<1`$. The spectral weight which is lost in the coherent part forms an incoherent (i.e. independent of $`𝐤`$) background of states which represents localized electrons and is such that $`_𝐤𝑑\omega A_{inc}(𝐤,\omega )e^{i\omega 0^+}=(1Z)N`$.<sup>?</sup> Even though Hubbard model cannot be solved exactly, there are many approximate methods of solution which can be used to determine the dependence of $`Z`$ on “physical” quantities, such as the doping $`\delta `$ or $`U`$. The above description retains its validity regardless of the particular approximation chosen.
In particular, without any loss of generality, the Green’s function of the system can be written as $`G=G_{co}+G_{inc}`$. Traditional analytical approaches usually focus on either of the two components, according to which properties they want to underline: itinerant properties (mean field slave-bosons, Gutzwiller) or insulating ones (Hubbard I approximation). Interplay between the two components can be recovered at a higher order level of approximation.<sup>?</sup>
In order to deal with this problems we have recently introduced an approximation based on a modified mean field solution within the slave boson technique, which allows us to treat both components on the same footing.
The Green’s function can be expressed in a particularly simple form:
$$G(𝐤,\omega )=Z\stackrel{~}{G}(𝐤,\omega )+(1Z)\underset{𝐤}{}\stackrel{~}{G}(𝐤,\omega ),$$
(2)
where $`\stackrel{~}{G}`$ is the standard mean-field slave boson solution describing only the coherent part. It is clear that $`Z`$ gives a parameter to estimate the degree of correlation of the system ($`Z=1`$ uncorrelated, $`Z=0`$ maximally correlated).
The above result can be quite easily proved by splitting the generic Green’s function in a local and non-local part:
$$G(i,j;t)=iT_tc_i(t)c_j^{}[1\delta _{i,j}]iT_tc_i(t)c_i^{}\delta _{i,j}$$
(3)
In the limit of infinite $`U`$ the slave bosons technique decomposes the real electron into two operators:
$$c_{i,\sigma }=b_i^{}f_{i,\sigma }.$$
The standard mean-field approximation \[$`b_i(t)bb_i(t)`$\] is now applied only to the non-local part of the Green’s function, while in the local part the no-double occupancy constraint can be exactly implemented by the relation
$$b_ib_i^{}f_i^{}(t)f_i=f_i^{}(t)f_i.$$
After some simple algebra, we obtain therefore:
$$G(i,j;t)=ib^2T_tf_{i,\sigma }(t)f_{j,\sigma }^{}i[1b^2]T_tf_{i,\sigma }(t)f_{i,\sigma }^{}\delta _{i,j},$$
(4)
which, written in $`𝐤`$, gives eq. (2) with $`Z=b^2`$.
It can be shown that the retarded Green’s function of this form preserves the total spectral weight ( $`_𝐤𝑑\omega \text{Im}G_R(𝐤,\omega )=N`$). As a consequence Luttinger’s theorem is naturally fulfilled, contrary to what happens with other approximations. This function can then be used as a basic ingredient to build up a diagramatic theory.
In this context we focus on the study of how charge-density response (namely electron-phonon interaction) is affected by electronic correlation.
## 3 Electron-phonon interaction: predominance of forward scattering
We wish to study how the effective electron-phonon coupling function $`g^2(𝐪,\omega )`$, is affected by electronic correlation. This issue has already been addressed with different analytical techniques focusing on charge fluctuations around mean-field solutions.<sup>?</sup> The results are in substantial agreement and show that electronic correlation induces a structure of the electron-phonon coupling which is strongly peaked around $`𝐪=0`$ leading to a predominance of forward scattering. This structure can be thought of as resulting by the poorer screening provided by the correlated electron system with respect to a normal metal, where highly itinerant electrons effectively screen out any charge modulation.
In this work we introduce an alternative and more intuitive way to describe correlation effects in an electron-phonon system based on the use of the Green’s function in eq. (2).
Let us consider the screening of the bare electron-phonon interaction $`g_0`$ by the electron-electron scattering. According the standard RPA approximation, the renormalized $`g^2`$(q) can be written as:
$$g^2(𝐪,\omega )=\frac{g_0^2(𝐪,\omega )}{1V(𝐪)\mathrm{\Pi }(𝐪,\omega )},$$
(5)
where V(q) is the Coulomb repulsive potential between electrons and $`\mathrm{\Pi }(𝐪,\omega )`$ is the charge density “bubble” of the correlated system:
$$\mathrm{\Pi }(𝐪,\omega )=\underset{𝐤}{}_{\mathrm{}}^+\mathrm{}\left(\frac{d\omega ^{}}{2\pi }\right)G(𝐤,\omega ^{})G(𝐤+𝐪,\omega +\omega ^{}).$$
(6)
In a normal metal $`\mathrm{\Pi }(0)=N(E_F)`$, where $`N(E_F)`$ is the density of states at the Fermi level, indicating that all the electrons at the Fermi surface contribute to the screening. The total response can be express in term of the “Thomas-Fermi” cut-off $`k_{TF}`$, which is usually larger than the Brillouin zone, and the resulting screened electron-phonon interaction is essentially $`𝐪`$-independent.
Things are different in a correlated system. By using the expression of Green’s function in eq. (2), $`\mathrm{\Pi }(𝐪,\omega )`$ can be rewritten in a coherent and incoherent part:
$$\mathrm{\Pi }(𝐪,\omega )=Z\mathrm{\Pi }_{co}(𝐪,\omega )+\frac{(1Z^2)}{Z}\mathrm{\Pi }_{inc}(\omega ),$$
(7)
where, just as for $`G`$, $`\mathrm{\Pi }_{inc}(\omega )=_𝐪\mathrm{\Pi }_{co}(𝐪,\omega )`$.
The two components contribute in a different way to the screening of the external charge. On a physical ground we expect that the part of the electrons which exhibits itinerant behaviour will be screening the external charge with a characteristic Thomas-Fermi cut-off which scales with $`Z`$. On the other hand the localized states will provide just a residual dynamical screening.
In this picture, the coherent part, which describes itinerant quasi-particles, can be described as a non-interacting renormalized system, by means of the standard Thomas-Fermi approximation: $`\mathrm{\Pi }_{co}(𝐪,\omega )=N(E_F)`$. The factor $`Z`$ which multiplies $`\mathrm{\Pi }_{co}(𝐪,\omega )`$ is due both to band renormalization and spectral weight reduction. The incoherent part, which describes localized states, exhibits a more complex behaviour as a function of the parameter $`Z`$ and of the exchanged frequency $`\omega `$. When $`\omega `$ is comparable with the effective kinetic energies of the electrons $`ZE_F`$, namely when the system is strongly correlated ($`Z0`$) , the charge response of the localized states is in counterphase, yielding a negative screening.
The resulting $`g^2(q)`$ can then be studied as a function of the degree of correlation of the system $`Z`$; in an intuitive picture $`Z`$ can be related to the hole-doping $`\delta `$ of the system through several approximations; in the infinite-$`U`$ limit it can be shown that $`Z\delta `$. In Fig. 1 is plotted the renormalized electron-phonon interaction as function of the exchanged momentum $`𝐪`$ for two different $`Z`$’s. The Thomas-Fermi constant $`k_{TF}`$ has been chosen $`k_{TF}=2k_F`$. As shown in the figure, $`g^2(q)`$ presents a sharp peak in the small-$`𝐪`$ region, which corresponds to forward scattering. A similar structure provides a cut-off $`q_c`$ for electron-phonon interactions in momentum-space which depends on the “degree of correlation” of the system.
Predominance of forward scattering has important consequences in the context of the electron-phonon theory of superconductivity. In particular, it is clear that in this situation Migdal’s theorem, which relies on the small parameter $`(\omega _{ph}/v_Fq)`$, cannot be justified.<sup>?</sup> Vertex corrections, usually neglected in virtue of Migdals’ theorem, need therefore to be explicitely included. In the past years, the extension of the theory of superconductivity in the so-called “nonadiabatic” regime, where first corrections beyond Migdal’s theorem are relevant, has been largely studied.<sup>?</sup> A major role is played by the momenta structure of the electron-phonon interaction: for small $`𝐪`$’s the resulting effects of the vertex corrections is mainly positive, leading to an enhancement of the electron-phonon coupling and to an increase of the superconducting critical temperature $`T_c`$ (see fig. 2).
## 4 Conclusions
In this work we present a rather simple and compact approach to the problem of electronic correlation in HTSC. It is based on the observation that the main features of strongly correlated electron systems described by Hubbard-type models are the reduction of the spectral weight associated with the itinerant part of the spectral function and the onset of an incoherent background of states. Through the introduction of a phenomenological parameter $`Z`$ which describes the “degree of correlation” of the system we were able to treat the two parts of the spectral function on the same footing. We applied this model to the computation of the electron-phonon coupling function $`g^2(q)`$ and found a structure which is strongly peaked in the small-$`𝐪`$ region and exhibits a strong dependence on $`Z`$
The peaked structure of $`g^2(q)`$ provides an upper cut-off in momentum-space for electron-phonon interactions, which leads to the breakdown of the adiabatic hypothesis in the momentum space $`(\omega _{ph}/v_Fq)`$. This points to the necessity of including electron-phonon vertex corrections in Migdal-Eliashberg equations, as predicted by the non-adiabatic theory of superconductivity, which shows that $`T_c`$’s of HTSC are enhanced by vertex corrections when the momentum $`𝐪`$ exchanged in electron-phonon interactions is small.
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# Ising films with surface defects
## I Introduction
Critical phenomena of magnetic films are of current interest, both experimentally and theoretically . In the limiting cases of one layer and of infinitely many layers, one deals with two-dimensional magnets and with standard bulk and surface magnetism , respectively. For systems consisting of a finite number of layers, interesting crossover phenomena between these limiting cases are expected.
In this article, we shall consider critical properties of ferromagnetic films of Ising magnets with various imperfections at the surface, motivated partly by possible experimental realizations of magnetic thin films with stripes of magnetic adatoms and stepped surfaces , partly by genuine theoretical interest. Imperfections may be due to regular or irregular changes in the surface couplings or due to additional structures on the surface. A simple example of the first case is a ladder of modified couplings in an otherwise uniform two-dimensional system as introduced by Bariev , see Fig. 1(a). We will study this briefly, since it can serve as a testing ground. Our main interest, however, is in additional structures, as depicted in Figs. 1(b)–(d). Thus we will investigate surfaces with magnetic adatoms in the form of
* one additional straight line, Fig. 1(b),
* two neighbouring lines, Fig. 1(c),
* a straight step of unit height, Fig. 1(d),
for various local couplings at the defects and for films of varying thickness.
Previous related work on Ising models includes the study of the step magnetization at the ordinary transition of rather thick films and the study of magnetism in thin films with rough surfaces .
In studying the influence of these imperfections especially on the critical behaviour, we use the density- matrix renormalization group technique (DMRG) , being most suitable in the case of merely one layer, and the Monte Carlo (MC) method , which allows to treat films of considerable thickness as well.
The article is organized as follows. In the next section, we present our findings on single layers with defects, applying DMRG. The MC results on single layers and on films with an additional line of magnetic adatoms and with a straight step on top of the surface are discussed in Section 3. A short summary concludes the article.
## II One layer: DMRG
The planar Ising model with line-like defects is a peculiar system, because it shows non-universal magnetic exponents. This is connected with the values $`\nu =1`$ and $`x_s=1/2`$ of the exponents for the correlation length and the surface magnetization of the pure system, respectively. A one-dimensional, energy-like perturbation then is marginal and can change the critical behaviour continuously. For this reason, the system has been the topic of various studies , with the focus most recently on a conformal treatment and on random systems . While the simple chain and ladder defects considered by Bariev are solvable free-fermion problems, the other cases we study are not integrable and one has to use numerical methods. In the following we dicuss the quantity of direct physical interest, the local magnetization at or near the defect lines.
To obtain it, we used the transfer matrix running along the direction of the defect, see Figs. 1(a)–(c), and determined its maximal eigenvector via the DMRG method . In this way one is treating an infinitely long strip of width $`M`$ with the defect located in the middle. Only the infinite-system algorithm was used, in which one enlarges the system step by step and always chooses an optimal reduced basis via the density matrix. This is very convenient, since one can insert different defects after the system has reached the desired size. No further sweeps to optimize the state were made, since tests on the ladder defect gave very good results without them. Most calculations were done with $`64`$ kept states and a truncation error around $`10^{15}`$. The local magnetization $`m(i)`$ was determined from the spin correlation function $`C(i)=\sigma _1\sigma _i`$ for free boundary conditions, or directly as $`\sigma _i`$ for fixed boundary spins. The width was always much larger than the correlation length and varied between $`M=100`$ and $`M=5000`$ for the temperature range studied ($`0.001<t<0.1`$, where $`t=1T/T_c`$ is the reduced temperature). The (absolute) error in $`m`$, determined by comparing with anaytical results was at most $`10^4`$ for a system at $`t=0.001`$, cut in the middle by a ladder defect. For less severe modifications and larger values of $`t`$ it was even smaller.
In Fig. 2 we show the correlation function $`C(i)`$ across the strip for ladder defects (Fig. 1(a)) and for an additional line (Fig. 1(b); in the DMRG study we considered the case $`J_n=J_s`$). The upper part gives an overall picture, while the lower one shows the defect region in more detail. For ladder defects the strength $`J_l`$ of the defect bonds was varied, whereas for an additional line it was the coupling $`J_a`$ between the line spins and the substrate. Since $`C(i)`$ factors for large distances, these curves also give the profile of the magnetization in the bulk. One can see how $`m`$ increases or decreases near the defect, depending on the sign of the perturbation (similar curves were obtained in for a random system). If one cuts the ladder bonds by choosing $`J_l=0`$, one obtains the boundary magnetization of the homogeneous model in the middle of the strip. On the upper side, the possible increase of $`m`$ depends on the details of the defect. It is limited if one varies $`J_a`$, because a line with infinite $`J_a`$ is equivalent to a chain defect in the plane with merely doubled bond strength.
The temperature dependence of $`m`$ is shown in Fig. 3 for the spins in the plane situated below one or two additional lines. One can see how it is increased over the Onsager value by increasing the coupling $`J_a`$. As expected, the effect is even stronger for two additional lines. In this case, $`m`$ has already twice the undisturbed value for the smallest shown $`t`$. Quantitatively, this enhancement is described by a decrease of the exponent $`\beta _l`$, the local critical exponent which describes the vanishing of the magnetization near the additional line of magnetic adatoms.
To investigate this, we have analyzed the temperature behaviour of $`m`$ in terms of an effective (critical) exponent $`\beta _{eff}`$, defined by
$$\beta _{eff}(t)=\mathrm{ln}(m(t_i)/m(t_{i+1}))/\mathrm{ln}(t_i/t_{i+1})$$
(1)
with $`t=(t_i+t_{i+1})/2`$ (alternatively, one could choose $`t`$ to be the geometric mean $`t=\sqrt{t_it_{i+1}}`$). As one approaches the critical point, $`t0`$, this quantity converges to the true local exponent $`\beta _l`$. It is also a very sensitive indicator for the numerical accuracy of a calculation.
Some typical results are given in Fig. 4 for one additional line and four values of the ratio $`J_l/J_s`$ of the couplings in the line. For $`J_l`$ = 0, one is treating the plane with independent attached spins and the Onsager result $`\beta =0.125`$ is recovered with high accuracy. In the other cases, the exponents both for the spin in the line and the one below it are shown and one sees that the two curves have different slopes, but a common limit for $`t0`$ which can be determined very precisely. The values for $`\beta _l`$ found in this way are accurate to at least three digits. For the case $`J_a/J_s1`$ which, as mentioned, is equivalent to a line defect in the plane, this was checked explicitely by comparing with the analytical result. In the figure, also a negative $`J_l`$ is shown, which leads to a reduction of $`m`$ and an increase of $`\beta _l`$ over the Onsager value . In this case, a limiting value $`0.142`$ is approached rapidly for $`J_l/J_s<1`$. This is the same effect as for a chain defect in the plane with strong antiferromagnetic couplings . In that case, the exponent is increased up to the value 0.5 of the free surface. The sign of $`J_a`$, on the other hand, has no influence on the exponent.
The results for $`\beta _l`$ are collected in Table 1 and in Fig. 5, where the exponent is plotted as a function of the varied couplings (keeping the other couplings fixed and equal to $`J_s`$). For comparison also the analytical results , for simple chain and ladder defects are shown in Fig. 5. One notes that, for a single line and small modifications, it does not matter much whether one changes $`J_a`$ or $`J_l`$. A large $`J_l/J_s`$, however, has a much more pronounced effect than $`J_a/J_s`$, since it corresponds to additional spins which are almost rigidly locked together. For the double line, the exponent drops much faster, reaching $`10^2`$ already around $`J_a/J_s1`$. For more additional lines, i.e. for a terrace on the surface as in Fig. 1(d), this effect would be even stronger. In this case, the magnetization would practically jump as in a first-order transition.
## III Films: Monte Carlo simulations
### A One additional line of spins
Extending the DMRG calculations on an Ising layer with one additional line of spins, we did Monte Carlo simulations on the corresponding Ising films, consisting of $`L`$ layers with one line of magnetic adatoms on top of the surface, see Fig. 1(b). We set $`J_l=J_s`$, with $`J_a=J_n=J_s`$ (variant A, treating the spins directly below the additional line as surface spins, as it was done in the DMRG study) or $`J_a=J_n=J_b`$ (variant B, treating the spins directly below the additional line as bulk spins).– In the layers, periodic boundary conditions are used.
Let $`S_{i,j,k}=\pm 1`$ be the spin on site $`(i,j)`$ in the $`k`$th layer. Taking layers of $`M\times N`$ spins, the spins in the additional line on top of the surface through the center are located at $`(i=(M+1)/2,j,k=0)`$, $`M`$ odd, with $`j`$ running from 1 to $`N`$. We computed, among others, the line magnetization $`m(i,k;L)`$ defined by
$$m(i,k;L)=\frac{1}{N}\left|\underset{j}{}S_{i,j,k}\right|$$
(2)
The magnetization of the line on top of the surface, $`m_l`$, is given by $`m_l(L)`$= $`m((M+1)/2,k=0;L)`$.
In the simulations, the film thickness $`L`$ ranged from 1 to 40, with layer sizes being sufficiently large to circumvent finite–size effects (up to $`161\times 320`$). To speed up computations, the single–cluster–flip algorithm was implemented. We studied the cases (i) $`J_s=J_b`$ as well as (ii) $`J_s=2J_b`$ (variants A and B), which lead to the two characteristic scenarios of surface critical phenomena for $`L,M,N\mathrm{}`$ (semi–infinite case). In the first case, bulk and surface spins order simultaneously at temperature $`T_c`$ (ordinary transition), while in the second one the surface spins order at a higher temperature, $`T_s`$ (surface transition) .
(i) At the ordinary transition of the semi–infinite Ising model, $`L\mathrm{}`$, the magnetization deep in the bulk vanishes like $`mt^\beta `$, with $`t=|TT_c|/T_c`$, where $`\beta =0.31\mathrm{}`$ . At the perfect, flat surface, one finds $`mt^{\beta _1}`$, with $`\beta _10.80`$ . The vanishing of the magnetization in the additional line of spins on top of the surface is expected to be governed by $`\beta _1`$ as well, i.e. $`\beta _l(L\mathrm{})=\beta _1`$ . On the other hand, for a single perfect layer, $`L=1`$, it is well known that $`mt^{\beta _{2d}}`$, $`\beta _{2d}=1/8`$. Adding a row of spins, we obtain, from the DMRG calculations, $`m_lt^{\beta _l(L=1)}`$ with $`\beta _l(L=1)0.084`$, see Table 1.
To monitor the influence of the layer thickness $`L`$ at the ordinary transition, we computed magnetization profiles $`m(i,k;L)`$, the critical temperature $`T_c(L)`$, and the critical exponent $`\beta _l(L)`$. The dependence of the transition temperature on the thickness $`L`$ has been studied before for flat films , and it is, certainly, not affected by the presence of the additional line. In Fig. 6, the magnetization in the defect line, $`m_l(L)`$, is depicted as a function of temperature for $`L`$ ranging from 1 to 5, illustrating the increase of the transition temperature with $`L`$. In the ordered phase, $`T<T_c`$, the line magnetization $`m(i,k;L)`$ is, in each layer, maximal for the center line, $`m((M+1)/2,k;L)`$, see also Fig. 2. The maximum is most pronounced at $`k=1`$ (we shall denote the magnetization in that line beneath the additional row of spins by $`m_{lb}=m((M+1)/2,1;L)`$), due to the increased coordination number, compared to the other surface lines. The magnetization in the additional line, $`m_l`$, is suppressed compared to $`m_{lb}`$, because of missing neighbouring spins.
Various crossover effects show up in the effective exponent $`\beta _{eff}(i,k;L)`$, defined by $`m(i,k;L)t^{\beta _{eff}(i,k;L)}`$, corresponding to the slope in a standard log–log–plot of the temperature dependence of the magnetization , see eq. (1). On approach to $`T_c`$, one expects to observe the limiting cases $`\beta _{eff}(i,k;L)\beta `$ for $`k`$ and $`L`$ large, $`\beta _1`$ for $`k`$ small and $`L`$ large, $`\beta _l(L=1)`$ for $`L=1`$ and $`i=(M+1)/2`$, $`k=0`$ or 1, and $`\beta _{2d}`$ for $`L=1`$ and sufficiently far away from the additional line in the centre.
The crossover behaviour is illustrated in Fig. 7, showing $`\beta _{eff}((M+1)/2,k;L)`$ for $`m_l(L)`$, $`k=0`$, and $`m_{lb}(L)`$, $`k=1`$, with the film thickness ranging from $`L=1`$ to 10. $`\beta _{eff}((M+1)/2,0;L)`$ decreases monotonically, except for $`L=1`$, over a wide range of temperatures on lowering $`t`$, but with the effective exponent, at fixed $`t`$, increasing clearly with the film thickness, as depicted in Fig. 7(a). The data seem to indicate that the asymptotic critical exponent $`\beta _l(L)`$, as $`t0`$, of the magnetization in the additional line of magnetic adatoms increases, however, only weakly with $`L`$, being quite small, around 0.1, for $`L`$ going up to 10 (the increase itself may be argued to reflect the diminishing role of the defect line on the two–dimensional critical fluctuations in thicker films; of course, $`\beta (L)`$ is bounded by 1/8 for finite $`L`$). For the magnetization beneath the additional line, $`m_{lb}`$, corrections to the asymptotics are rather large as well, see Fig. 7(b). Here the effective exponent $`\beta _{eff}((M+1)/2,1;L)`$ changes with temperature in a non–monotonic fashion, except for $`L=1`$. In agreement with the observations for $`m_l`$, the true critical exponent $`\beta _l(L)`$ is rather small, around 0.1, increasing only weakly with $`L`$. The location of the maximum in $`\beta _{eff}((M+1)/2,1;L)`$ indicates the temperature, at which one crosses over from the regime dominated by two–dimensional critical fluctuations, close to the phase transition, to the regime, further away from the critical point, where the fluctuations are (nearly) isotropic and three–dimensional. Thence, at the maximum the corresponding correlation length is argued to be about the thickness of the film $`L`$. In the thermodynamic limit, $`L\mathrm{}`$, the maximum is believed to shift towards $`t=0`$, with its height being $`\beta _l=\beta _10.80`$.
Note that the strong corrections to scaling, as seen by the deviations of the effective exponents from their asymptotic values, may cause severe difficulties in extracting the true critical exponents in simulations as well as in experiments.
Similar crossover phenomena, now between $`\beta _{2d}`$, $`\beta _1`$ and $`\beta `$, are expected to occur for the magnetization far away from the defect line, when varying the film thickness. This aspect, however, is of minor importance in the context of this study.
(ii) At the surface transition of flat Ising films, the surface magnetization vanishes, on approach to $`T_s`$, like $`mt^{\beta _{2d}}`$, independent of $`L`$. The critical exponent for the magnetization at the additional line of magnetic atoms on top of the surface, $`\beta _l`$, with $`m_{l(lb)}t^{\beta _l}`$, depends on the local couplings at that line, as seen from our DMRG results for $`L=1`$. Indeed, the situation is similar to that of the edge magnetization at the surface transition, the edge corresponding to an extended defect line , where non–universality holds as well.
For $`J_s=2J_b`$ and $`J_a=J_n=J_s`$, variant A, one obtains for a single layer, from the DMRG method, $`\beta _l(L=1)0.084`$, i.e. the value is below that of the perfect two–dimensional Ising model because of the increase in $`m_l`$ due to the additional line of spins, see Fig. 2. The value increases weakly with layer thickness, becoming in the limit of the semi–infinite system $`\beta _l(L=\mathrm{})=0.091\pm 0.002`$, as inferred from MC data for films with thickness $`L`$ up to 40, and reasonable extrapolations. Because the critical fluctuations in a film of finite thickness are ultimatively of two–dimensional nature, one expects a non–universal critical behaviour at the defect line with $`\beta _l`$ depending on $`L`$. Actually, the slight increase of the critical exponent with $`L`$ reflects the impact of the bulk spins, which now tend to lower the magnetization in the defect line.
For $`J_s=2J_b`$ and $`J_a=J_n=J_b`$, variant B, both for single layers, $`L=1`$, and films, the magnetization profile $`m(i,k;L)`$ close to $`T_s`$ is non–monotonic exhibiting a minimum at the center line $`i=(M+1)/2`$, see Fig. 8. This minimum is due to the reduction of the couplings $`J_n`$ at the defect below the value $`J_s`$ elsewhere in the surface. As one goes deeper into the bulk, the magnetization profile smoothens, which can be readily understood.
The critical exponent describing the vanishing of $`m_l`$ (or $`m_{lb}`$) depends rather weakly on the thickness $`L`$ of the film. For $`L=1`$, we estimate from the MC data $`\beta _l(L=1)=0.38\pm 0.01`$, i.e. a value above the Onsager value of the perfect two–dimensional Ising model resulting from the decrease of the magnetization at the defect line . The effective exponent decreases on approach to criticality, $`t0`$, when considering $`m_{lb}`$, while it increases when considering $`m_l`$, allowing to estimate $`\beta _l`$ accurately.–From data for fairly thick films, $`L`$ up to 40, we estimate $`\beta _l`$ of the semi–infinite system to be $`\beta _l(L=\mathrm{})=0.34\pm 0.02`$. The slight change of $`\beta _l(L)`$ with $`L`$ for films of finite thickness is, again, believed to be due to the correlations of the spins at the defect line with the bulk spins, which affect $`\beta _l`$ in such a way that it is non–universal when the critical fluctuations are of two–dimensional character.
### B Step
Finally, we briefly report our findings for the critical properties of the step magnetization. A straight step is introduced (actually two steps, to allow for periodic boundary conditions) by adding half a layer of magnetic adatoms to the surface of the magnetic film , see Fig. 1(d). We discrimate two couplings, $`J_s`$ if both neighbouring spins are surface spins, and $`J_b`$ otherwise. We consider the line magnetization of the spins at the step edge, $`m_{se}`$, vanishing on approach to the transition as $`m_{se}t^{\beta _{se}}`$.
For $`J_s=J_b`$, i.e. at the ordinary transition, one obtains $`\beta _{se}0.80`$ in semi–infinite Ising models, $`L\mathrm{}`$, i.e. the same value as for the critical exponent of the surface magnetization, as had been shown in a previous Monte Carlo study on thick Ising films with a step , in agreement with analytical considerations . However, in thin films, the critical behaviour is quite different. In the simulations, for a single layer $`L=1`$ plus half a layer, we find a critical exponent close to 1/2 (its concrete value depends rather sensitively on a very accurate determination of $`T_c`$), i.e. a value close to that of the surface critical exponent $`\beta _1`$ of the two-dimensional Ising model (note also its robustness against randomness in the couplings ). This observation can be understood in the following way. One is dealing with a composite system displaying, as the layer size goes to infinity, two distinct phase transitions, one at the critical temperature of the Ising plane, $`k_BT_c(L=1)/J_s=2.269\mathrm{}`$, and one at the critical temperature of the double layer, $`k_BT_c(L=2)/J_s3.2`$. Related composite Ising models have been investigated before , showing that on approach to the upper critical temperature, where half of the system is disordered, the critical behaviour of the magnetization at the interface (i.e., here, at the step) is governed by the surface critical exponent. The same scenario is expected to hold for finite films with $`T_c(L+1)>T_c(L)`$. However, the temperature region where this behaviour can be observed, will become smaller and smaller as $`L`$ increases.
At the surface transition, the same considerations are believed to be valid. Indeed, in the case $`J_s=2J_b`$, we found $`\beta _{se}`$ to be quite close to 1/2 for a single layer, $`L=1`$, plus half a layer. For trivial reasons, $`\beta _{se}(L=1)=1/2`$ holds for $`J_sJ_b`$, when the bottom layer and the extra half layer decouple with the step edge being the surface of a two–dimensional Ising model. In the thermodynamic limit, where $`T_c(L+1)=T_c(L)`$, so that the above decoupling considerations do not apply, we estimated from MC data for films with up to 40 layers, a value of $`\beta _{se}=0.33\pm 0.02`$. Presumably, in that limit, $`\beta _{se}`$ is non–universal at the surface transition, depending on the ratio $`J_s/J_b`$.
## IV Summary
Using density-matrix renormalization group and Monte Carlo techniques, we studied critical properties of magnetic Ising films with various surface defects.
In particular, the effect of the local couplings at one or two additional lines of magnetic adatoms on the surface as well as at straight steps of monoatomic height has been investigated, especially in the limiting cases of films consisting of merely one layer and rather thick films.
In the case of a single layer, $`L=1`$, with additional lines of magnetic adatoms, the critical exponent of the magnetization at the surface defect is non–universal. The dependence of its value on the local couplings, as compared to that of the perfect two–dimensional situation, follows the trends observed for the exactly soluble two–dimensional Ising model with ladder and chain like bond–defects. The value may be lower or larger than in the perfect situation, 1/8, corresponding to an increase or decrease in the magnetization at the defect line. Adding half a layer of spins, one recovers, at the step, the surface critical exponent, 1/2, of the two–dimensional Ising model.
In the limit $`L\mathrm{}`$, varying the strength of the surface couplings may lead either to a surface or an ordinary phase transition. The change of the critical exponent of the magnetization at the defect has been found to depend only fairly weakly, for both types of transition, on the film thickness $`L`$ in the case of one additional line of spins. At steps, the critical exponent is argued to be 1/2, for films of finite thickness and both kinds of transition, in agreement with the simulations.
In the paper, we have not only presented the results for the exponents, but also shown various magnetization curves directly, so as to give an impression of the size of the effects. This is meant to encourage further experimental work on such surface structures and their magnetic properties.
###### Acknowledgements.
We would like to thank K. Baberschke and J. Kirschner for discussions of the experimental situation. M.C.C. thanks the Deutscher Akademischer Austauschdienst (DAAD) for financial support.
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# Non-Deterministic Learning Dynamics in Large Neural Networks due to Structural Data Bias
## 1 Introduction
Rosenblatt first introduced the perceptron and proved the famous perceptron convergence theorem in 1962. It is an indicator of the richness of the perceptron as a dynamical system that almost 40 years later it continues to yield fascinating results which have hitherto remained hidden. Especially during the last decade, considerable progress has been made in understanding the dynamics of learning in artificial neural networks through the application of the methods of statistical mechanics. The dynamics of on-line learning in perceptrons has been analysed intensively, but for the most part such studies have been carried out in the idealised scenario of so-called complete training sets (in which the number of training examples is large compared with $`N`$, the number of degrees of freedom), and have also assumed a homogeneous input data distribution. A recent review of work in this field is contained in . A general theory of learning in the context of restricted training sets (where the size of the training set is proportional to $`N`$) is generally much more difficult, although an exact solution of the dynamical equations for the more elementary problem of unbiased on-line Hebbian learning with restricted training sets and noisy teachers has been found . Nevertheless, substantial progress has been made towards a general theory of learning with restricted training sets and the reader may refer, for example, to , for details.
In this paper we consider complete training sets, but we admit the possibility of a structural bias of the input vectors. This is a significant issue since in real-world situations a training sample will generally have a non-zero average; this is especially important in the case of on-line learning, where examples are not available prior to learning, so that one cannot correct for any bias prior to processing. This in itself would be sufficient motivation for the present study. However, it turns out that the introduction of structurally biased input data leads to qualitative (rather than only quantitative) modifications of the actual learning curves observed in numerical simulations and the mathematical theories required for their description. Various authors have studied so-called clustered examples, in which examples are drawn from two Gaussian distributions situated close to each other, with an input bias of order $`N^{\frac{1}{2}}`$ (i.e. in magnitude similar to finite-size effects). Learning with input bias has also been considered in in the context of linear networks ; the linear theory was then used to construct an approximation for a class of non-linear models, and it was shown that on-line learning is more robust to input bias and out-performs batch learning when such bias is present.
Here we consider a situation which is more natural and less restrictive than the one considered in , and which does not require the linearity of : we study the familiar (non-linear) perceptron, with the perceptron learning rule and with a structural, i.e. $`𝒪(N^0)`$, bias in the input data. Using $`𝑩`$, $`𝑱`$, and $`𝑨`$ to denote the teacher weights, the student weights and the bias vector (precise definitions follow), we develop our theory in terms of three macroscopic observables: the standard observables $`Q=𝑱^2,R=𝑱𝑩,`$ and a new observable $`S=𝑱𝑨`$ (the overlap between student weights and bias vector). In contrast to the the dynamics of the bias-free case, we find that in the presence of an $`𝒪(N^0)`$ input bias the system passes though three phases, characterised by different scaling of typical times and of macroscopic observables. This could already have been anticipated on the basis of numerical simulations, see e.g. figure 1. We obtain a closed system of equations in which the evolution of $`\{Q,R\}`$ is deterministic in the limit $`N\mathrm{}`$, as in the bias-free case, but where $`S`$ is (generally) a stochastic variable, whose conditional probability distribution $`P_t(S|Q,R)`$ becomes non-trivial. Phase I is a short phase, in which the system reduces the alignment of the student weight vector $`𝑱`$ relative to the bias vector $`𝑨`$. During phase I, the observable $`S`$ is deterministic, is rapidly driven towards zero, and no learning takes place. Before the state $`S=0`$ is reached, however, the system enters phase II, a very short phase in which $`S`$ evolves stochastically to a quasi-stationary probability distribution (which we calculate) and in which both $`Q`$ and $`R`$ are frozen. In phase III, where most of the learning takes place, the $`S`$ distribution is modified by a non-negligible random walk element, which generates a diffusion term in the equation controlling the evolution of $`P_t(S|Q,R)`$, whereas $`Q`$ and $`R`$ satisfy coupled differential equations which involve averages over $`P_t(S|Q,R)`$. The stochastic nature of $`S`$ is reflected in the fact that the generalisation error also exhibits fluctuations (see figure 1). The (exact) equations describing phase III cannot be further simplified, but we introduce an approximation yielding more tractable equations for $`\{Q,R\}`$, which still have the merit of reducing to the more familiar equations when no-bias is present. Moreover, they are found to be in excellent agreement with the results obtained from numerical simulations. Compared to the unbiased case, having a finite bias is found to change the pre-factor in the asymptotic power law of the asymptotic time-dependence of the generalisation error, but not the exponent. A preliminary and more intuitive presentation of some of the present results can be found in .
## 2 Definitions
We study on-line learning in a student perceptron $`\mathrm{\Sigma }:\{1,1\}^N\{1,1\}`$, which learns a task defined by a teacher perceptron $`T:\{1,1\}^N\{1,1\}`$ whose fixed weight vector is $`𝑩\mathrm{}^N`$. Teacher and student output are given by the familiar recipes
$$T(𝝃)=\mathrm{sgn}[𝑩𝝃],\mathrm{\Sigma }(𝝃)=\mathrm{sgn}[𝑱𝝃],$$
We assume that $`𝑩`$ is normalised such that $`𝑩^2=1`$, the components being drawn randomly with mean zero and standard deviation of order $`𝒪(N^{1/2})`$, and statistically independent of the input data. In order to model the bias in the input sample we assume that for $`𝝃=(\xi _1,\mathrm{},\xi _N)\{1,1\}^N`$ all $`\xi _i`$ are independent, with $`\xi _i=a`$, so that the probability of drawing $`𝝃`$ is given by
$$p(𝝃)=\underset{i}{}\frac{1}{2}[1+a\xi _i]$$
(1)
We define $`\xi _i=a+v_i`$, such that the (independent) $`v_i`$ have mean zero and variance $`\sigma ^2=1a^2`$, and the short-hand $`𝑨=a(1,\mathrm{},1)`$ (i.e. a vector with all $`N`$ entries equal to $`a`$, to be referred to as the ‘bias vector’). The teacher-bias overlap $`𝑩𝑨`$ is now a random parameter which is $`𝒪(1)`$, since $`(𝑩𝑨)^2=_{i,j}a^2B_iB_j=a^2`$, whose distribution will be Gaussian for $`N\mathrm{}`$, with mean 0 and standard deviation $`a`$.
The student perceptron $`\mathrm{\Sigma }`$ is being trained according to an on-line learning rule of the form $`𝑱_{m+1}=𝑱_m+\mathrm{\Delta }𝑱_m`$, where at each iteration step an input vector $`𝝃_m`$ is drawn independently according to $`(\text{1})`$, and where
$$\mathrm{\Delta }𝑱_m=\frac{\eta }{N}𝝃_m\mathrm{sgn}(𝑩𝝃_m)[|𝑱_m|,𝑱_m𝝃_m,\mathrm{sgn}(𝑩𝝃_m)]$$
For Hebbian learning, for instance, we have
$$[J,u,T]=1:\mathrm{\Delta }𝑱_m=\frac{\eta }{N}𝝃_m\mathrm{sgn}(𝑩𝝃_m)$$
whilst the familiar perceptron learning rule is defined by
$$[J,u,T]=\theta [uT]:\mathrm{\Delta }𝑱_m=\frac{\eta }{2N}𝝃_m[\mathrm{sgn}(𝑩𝝃_m)\mathrm{sgn}(𝑱_m𝝃_m)]$$
(2)
We will derive, from the microscopic stochastic process for the weight vector $`𝑱`$, a macroscopic dynamical theory in terms of the familiar observables $`Q=𝑱^2`$ and $`R=𝑱𝑩`$, as well as (in order to obtain closure) a new observable $`S=𝑱𝑨`$ measuring the overlap between the vector $`𝑱`$ and the bias vector. The teacher and student output can then be written in the form
$$\mathrm{\Sigma }(𝝃)=\mathrm{sgn}[\lambda _1+x],T(𝝃)=\mathrm{sgn}[\lambda _2+y]\mathrm{with}\lambda _1=\widehat{𝑱}𝑨,\lambda _2=𝑩𝑨,$$
with $`\widehat{𝑱}=𝑱/|𝑱|`$, and where the local fields $`\{x,y,z\}`$ are defined by $`x=\widehat{𝑱}𝒗`$, $`y=𝑩𝒗`$ and $`z=\widehat{𝑨}𝒗`$ (the latter field $`z`$ will also enter our calculation in due course). Note that $`\lambda _1=S/\sqrt{Q}`$. For large $`N`$, the three fields $`\{x,y,z\}`$ are zero-average Gaussian random variables, each with variance $`\sigma ^2=1a^2`$, and with correlation coefficients given by
$$xy=\omega \sigma ^2,xz=\sigma ^2S/|𝑨|,yz=\sigma ^2\lambda _2/|𝑨|.$$
(3)
We note that equation (3) implies that $`z`$ will be independent of $`(x,y)`$ for large $`N`$ so that
$$p(x,y,z)=[\sigma \sqrt{2\pi }]^1e^{z^2/2\sigma ^2}p(x,y),p(x,y)=\left[2\pi \sigma ^2\sqrt{1\omega ^2}\right]^1e^{\frac{1}{2}[x^22\omega xy+y^2]/\sigma ^2(1\omega ^2)}$$
(4)
with $`\omega =\widehat{𝑱}𝑩=R/\sqrt{Q}`$. It will turn out that most of the averages to appear in this paper, involving (4) (to be written as $`\mathrm{}`$), may be expressed in terms of the function $`K(x)=\mathrm{erf}(x/\sqrt{2})`$. The generalisation error $`E_g=\theta [(\widehat{𝑱}𝝃)(𝑩𝝃)]`$, for example, can be written as
$$E_g=𝑑x𝑑yp(x,y)\theta [(\lambda _1+x)(\lambda _2+y)]=I_1(\lambda _1,\lambda _2,\omega )+I_1(\lambda _1,\lambda _2,\omega )$$
(5)
where
$$I_1(\lambda _1,\lambda _2,\omega )=_{\lambda _1}^{\mathrm{}}𝑑x_{\lambda _2}^{\mathrm{}}𝑑yp(x,y)=\frac{1}{4}\left[1K(\frac{\lambda _2}{\sigma })\right]\frac{1}{2}_{\frac{\lambda _2}{\sigma }}^{\mathrm{}}DyK(\frac{\lambda _1\omega \sigma y}{\sigma \sqrt{1\omega ^2}})$$
with the Gaussian measure $`Dy=(2\pi )^{\frac{1}{2}}e^{\frac{1}{2}y^2}dy`$ (see appendix A for details). This then gives
$$E_g=\frac{1}{2}\frac{1}{2}_{\frac{\lambda _2}{\sigma }}^{\mathrm{}}DyK\left(\frac{\lambda _1+\omega \sigma y}{\sigma \sqrt{1\omega ^2}}\right)+\frac{1}{2}_{\frac{\lambda _2}{\sigma }}^{\mathrm{}}DyK\left(\frac{\lambda _1\omega \sigma y}{\sigma \sqrt{1\omega ^2}}\right)$$
(6)
Note that, due to the identity $`_0^{\mathrm{}}DyK(\omega y/\sqrt{1\omega ^2})=\frac{1}{2}\frac{1}{\pi }\mathrm{arccos}\omega `$, formula (6) reduces, as it should, to the well known expression $`E_g=\pi ^1\mathrm{arccos}\omega `$ in the case where the input bias is zero (i.e. for $`a0`$).
## 3 From Microscopic to Macroscopic Laws
We now consider the dynamics of the macroscopic observables $`\{Q,R,S\}`$ in the limit of large $`N.`$ In the bias-free case, where for large $`N`$ the fluctuations in the macroscopic observables are insignificant, this can be done in a direct and simple way. Here, for $`a0`$, the situation is qualitatively different, since (as will turn out) the fluctuations in $`S`$ will no longer vanish, and their distribution will have a strong impact on the macroscopic laws. In order to provide a setting for our theory we briefly review a well known procedure which enables us to pass from a discrete to a continuous time description. We suppose that at time $`t`$ the probability that the perceptron has undergone precisely $`m`$ updates is given by the Poisson distribution $`\pi _m(t)=\frac{1}{m!}(Nt)^me^{Nt}`$. For large $`N`$ this will give us $`t=\frac{m}{N}+𝒪(N^{1/2})`$, the usual real-valued time unit, and the uncertainty as to where we are on the time axis vanishes as $`N\mathrm{}`$. It is not hard to show that the probability density $`p_t(𝑱)`$ of finding the vector $`𝑱`$ at time $`t`$ satisfies
$$\frac{d}{dt}p_t(𝑱)=N𝑑𝑱^{}\left\{\text{}\delta [𝑱𝑱^{}\mathrm{\Delta }𝑱]_𝝃\delta [𝑱𝑱^{}]\right\}p_t(𝑱^{})$$
where, for the perceptron learning rule (2), the single-step modification $`\mathrm{\Delta }𝑱`$ is given by
$$\mathrm{\Delta }𝑱=\frac{\eta }{2N}𝝃[\mathrm{sgn}(𝑩𝝃)\mathrm{sgn}(𝑱𝝃)]$$
and where and $`\mathrm{}_𝝃`$ denotes the average over all questions $`𝝃`$ in the training set $`\{1,1\}^N`$. The macroscopic observables $`𝛀=(Q,R,S)`$, in turn, have the probability density $`P_t(𝛀)=𝑑𝑱p_t(𝑱)\delta [𝛀𝛀(𝑱)]`$, which satisfies the macroscopic stochastic equation
$$\frac{d}{dt}P_t(𝛀)=𝑑𝛀^{}𝒲_t[𝛀,𝛀^{}]P_t(𝛀^{})$$
where
$$𝒲_t[𝛀,𝛀^{}]=N\delta [𝛀𝛀(𝑱+\mathrm{\Delta }𝑱)]_𝝃\delta [𝛀𝛀(𝑱)]_{𝛀^{},t}$$
with the so-called sub-shell (or conditional) average $`\mathrm{}_{𝛀^{},t}`$, defined as
$$f(𝑱)_{𝛀,t}=\frac{𝑑𝑱p_t(𝑱)\delta [𝛀𝛀(𝑱)]f(𝑱)}{𝑑𝑱p_t(𝑱)\delta [𝛀𝛀(𝑱)]}.$$
It is possible to make various assumptions regarding the scaling behaviour of our observables at time $`t=0`$, but once this has been specified the scaling at subsequent times is determined by the dynamics. We make the natural assumption that $`Q(0)=𝒪(1)`$ so that, in accordance with our assumptions regarding the statistics of $`𝑩`$, we have $`R(0)=𝒪(N^{1/2})`$. We suppose that $`S(0)=𝒪(N^{1/2})`$, the maximum permitted by the Schwarz inequality.
In this context it is worth remarking that in the idealised case of zero bias, Hebbian learning is known to out-perform the perceptron learning rule; but in the more realistic situation of even moderately biased data the Hebbian rule fails miserably. For example, if we assume that $`S(0)`$ is $`𝒪(1)`$,and that $`Q,R`$ are initially $`𝒪(1),`$ it follows from the learning rule (or from the methods which we apply below to the perceptron learning rule) that in the initial evolution of the Hebbian system $`dS/d\tau =\eta a^2K(\lambda _2/\sigma ),`$ where $`\tau =Nt,`$ so that $`S`$ rapidly diverges and no learning takes place; the student vector $`𝑱`$ cannot break away from its alignment to the bias vector. We shall show, however, that the perceptron has no problem coping with extreme initial conditions such as $`S(0)=𝒪(N^{1/2}),`$ and that in due course effective learning occurs. The Hebbian example also serves to show that, even if we were to choose the weaker initial scaling $`S(0)=𝒪(N^0)`$, dependent on the specific choice we make for the learning rule, the order parameter $`S`$ might well be driven towards $`S=𝒪(N^{\frac{1}{2}})`$ states.
A systematic exploration of the possible scaling scenarios reveals the following.<sup>1</sup><sup>1</sup>1For brevity we will in this paper only describe the resulting self-consistent solution, which is indeed perfectly consistent with the observations in numerical simulations such as in figure 1. For the perceptron learning rule and for the initial scaling conditions as specified above, the only self-consistent solution of the macroscopic equations is one describing a situation where the system passes through three phases $`\{\mathrm{I},\mathrm{I}\mathrm{I},\mathrm{I}\mathrm{I}\mathrm{I}\}`$ defined by time scales $`t=\{\tau N^{1/2},\tau N^1,\tau \}`$, in which our observables are $`𝒪(1)`$ quantities in all three phases, with the exception of $`S`$ which is $`𝒪(N^{1/2})`$ in phase I. We will write $`S=\stackrel{~}{S}N^{1/2}`$ in Phase I, with $`\stackrel{~}{S}=𝒪(N^0)`$, and formulate our Phase I equations in terms of $`\stackrel{~}{S}`$ rather than $`S`$. The number of iterations $`m`$ is related to the original time $`t`$ by $`m=Nt`$ so that the number of iterations up to time $`\tau `$, in each of the three phases, is given by $`m=\{\tau N^{1/2},\tau N^0,\tau N\}`$. We incorporate these scaling properties into our equations in each of the three phases, by working henceforth only with $`𝒪(N^0)`$ time units $`\tau `$ and $`𝒪(N^0)`$ observables $`𝛀`$, which satisfy
$$\frac{d}{d\tau }P_\tau (𝛀)=𝑑𝛀^{}𝒲_\tau [𝛀,𝛀^{}]P_\tau (𝛀^{})$$
(7)
with
$$𝒲_\tau [𝛀,𝛀^{}]=F_{\mathrm{I},\mathrm{I}\mathrm{I},\mathrm{I}\mathrm{I}\mathrm{I}}\delta [𝛀𝛀(𝑱+\mathrm{\Delta }𝑱)]_𝝃\delta [𝛀𝛀(𝑱)]_{𝛀^{},t}$$
$$=\frac{F_{\mathrm{I},\mathrm{I}\mathrm{I},\mathrm{I}\mathrm{I}\mathrm{I}}}{(2\pi )^3}𝑑\widehat{𝛀}e^{i\widehat{𝛀}𝛀}\left\{e^{i\widehat{𝛀}𝛀(𝑱+\mathrm{\Delta }𝑱)}_𝝃e^{i\widehat{𝛀}𝛀(𝑱)}\right\}_{𝛀^{},t}$$
(8)
and $`F_\mathrm{I}=N^{1/2},F_{\mathrm{I}\mathrm{I}}=N^0,F_{\mathrm{I}\mathrm{I}\mathrm{I}}=N.`$ In a subsequent stage it will be convenient to write $`\mathrm{\Delta }𝑱=𝒌+𝒌^{}`$, where
$$𝒌=\frac{\eta }{2N}𝑨[\mathrm{sgn}(𝑩𝝃)\mathrm{sgn}(𝑱𝝃)],𝒌^{}=\frac{\eta }{2N}𝒗[\mathrm{sgn}(𝑩𝝃)\mathrm{sgn}(𝑱𝝃)]$$
(9)
so that
$$\mathrm{\Delta }𝑱𝑨=\frac{1}{2}\eta a^2[\mathrm{sgn}(\lambda _2+y)\mathrm{sgn}(\lambda _1+x))]+\eta a\frac{z}{2\sqrt{N}}[\mathrm{sgn}(\lambda _2+y)\mathrm{sgn}(\lambda _1+x))]$$
(10)
We are now in a position to discuss the dynamics in each of the three phases in which different scaling laws apply.
## 4 Phase I: Elimination of Bias-Induced Activation
In Phase I we define the $`𝒪(N^0)`$ observables $`𝛀=(\stackrel{~}{S},Q,R)=(𝑱𝑨/\sqrt{N},Q,R)`$ and $`F_I=\sqrt{N}.`$ Upon expanding the exponential $`e^{i\widehat{𝛀}𝛀(𝑱+\mathrm{\Delta }𝑱)}`$ in powers of $`\mathrm{\Delta }𝑱`$ we obtain from equation (8)
$$𝒲_\tau [𝛀,𝛀^{}]=\frac{1}{(2\pi )^3}𝑑\widehat{𝛀}e^{i\widehat{𝛀}[𝛀𝛀(𝑱)]}\left\{\mathrm{}\right\}_\mathrm{I}_{𝛀^{},\tau }$$
where
$$\left\{\mathrm{}\right\}_\mathrm{I}=i\underset{i\mu }{}N^{1/2}\mathrm{\Delta }J_i\frac{\mathrm{\Omega }_\mu }{J_i}_𝝃\widehat{\mathrm{\Omega }}_\mu +\frac{1}{2}i\underset{ij\mu }{}N^{1/2}\mathrm{\Delta }J_i\mathrm{\Delta }J_j\frac{^2\mathrm{\Omega }_\mu }{J_iJ_j}_𝝃\widehat{\mathrm{\Omega }}_\mu $$
$$+\frac{1}{2}\underset{ij\mu \nu }{}N^{1/2}\mathrm{\Delta }J_i\mathrm{\Delta }J_j\frac{\mathrm{\Omega }_\mu }{J_i}\frac{\mathrm{\Omega }_\nu }{J_j}_𝝃\widehat{\mathrm{\Omega }}_\mu \widehat{\mathrm{\Omega }}_\nu +𝒪(N^1).$$
A straightforward calculation using equation (9) and the two averages $`\mathrm{sgn}(\lambda _2+y)=K(\frac{\lambda _2}{\sigma })`$ and $`\mathrm{sgn}(\lambda _1+x)=K(\frac{\lambda _1}{\sigma })=\mathrm{sgn}(\stackrel{~}{S})`$ (which is valid for large $`N`$ in phase I) now gives
$$\left\{\mathrm{}\right\}_\mathrm{I}=\frac{1}{2}i\eta a^2[K(\frac{\lambda _2}{\sigma })\mathrm{sgn}(\stackrel{~}{S})]\widehat{\mathrm{\Omega }}_1+i\eta \stackrel{~}{S}[K(\frac{\lambda _2}{\sigma })\mathrm{sgn}(\stackrel{~}{S})]\widehat{\mathrm{\Omega }}_2.$$
We can now apply equation (7) to compute the time derivative of the probability density $`P_\tau (𝛀)`$. Note that the sub-shell average $`\mathrm{}_{𝛀^{},\tau }`$ involves an integration over all $`𝑱`$ for which $`𝛀(𝑱)=𝛀^{}`$ (in a distributional sense) so in calculating the relevant integrals we may effectively replace $`𝛀(𝑱)`$ by $`𝛀^{}`$ at appropriate stages. For example, $`𝑑\widehat{𝛀}\widehat{\mathrm{\Omega }}_je^{i\widehat{𝛀}[𝛀𝛀(𝑱)]}=i(2\pi )^3_j^{}\delta [𝛀𝛀^{}]`$, where $`_j^{}`$ denotes differentiation with respect to $`\mathrm{\Omega }_j^{}.`$ We now find for $`P_\tau (\stackrel{~}{S},Q,R)`$ a Liouville equation
$$\frac{d}{d\tau }P_\tau (\stackrel{~}{S},Q,R)=\frac{}{\stackrel{~}{S}}\left[\frac{\eta a^2}{2}[K(\frac{\lambda _2}{\sigma })\mathrm{sgn}(\stackrel{~}{S})]P_\tau (\stackrel{~}{S},Q,R)\right]\frac{}{Q}\left[\eta \stackrel{~}{S}[K(\frac{\lambda _2}{\sigma })\mathrm{sgn}(\stackrel{~}{S})]P_\tau (\stackrel{~}{S},Q,R)\right]$$
with the deterministic solution $`P_\tau (\stackrel{~}{S},Q,R)=\delta [\stackrel{~}{S}\stackrel{~}{S}(\tau )]\delta [QQ(\tau )]\delta [RR(\tau )]`$, where the actual deterministic trajectory $`\{\stackrel{~}{S}(\tau ),Q(\tau ),R(\tau )\}`$ is the solution of the coupled flow equations
$$\frac{d}{d\tau }\stackrel{~}{S}=\frac{1}{2}\eta a^2[K(\frac{\lambda _2}{\sigma })\mathrm{sgn}(\stackrel{~}{S})],\frac{d}{d\tau }Q=\eta \stackrel{~}{S}[K(\frac{\lambda _2}{\sigma })\mathrm{sgn}(\stackrel{~}{S})],\frac{d}{d\tau }R=0.$$
It follows that $`\stackrel{~}{S}(\tau )=\stackrel{~}{S}(0)+\frac{1}{2}\eta a^2\tau [K(\lambda _2/\sigma )\mathrm{sgn}(\stackrel{~}{S})]`$. We see that $`\stackrel{~}{S}`$ is driven to zero in times $`\tau =\tau _\pm `$ (with $`\pm `$ referring to the cases $`\stackrel{~}{S}_0>0`$ and $`\stackrel{~}{S}_0<0`$, respectively), which are given by
$$\tau _\pm =\frac{2|\stackrel{~}{S}_0|}{\eta a^2(1K(\frac{\lambda _2}{\sigma }))}.$$
Irrespective of the value of $`\stackrel{~}{S}_0`$, the system seeks to eliminate any strong alignment of the learning vector $`𝑱`$ relative to the bias vector $`𝑨`$. This is clearly confirmed by numerical simulations. Our equation for $`Q`$ also readily integrates to give $`Q=Q_0+[\stackrel{~}{S}^2\stackrel{~}{S}_0^2]/a^2`$. We see that the length $`J=\sqrt{Q}`$ of the student weight vector decreases and that $`J[J_0^2\stackrel{~}{S}_0^2/a^2]^{\frac{1}{2}}`$ as $`\tau \tau _\pm `$. Again, this is confirmed by numerical simulations. The equation $`dR/d\tau =0`$ implies that $`\omega J`$ is constant in Phase I. As can be clearly seen in figure 1, no learning takes place in this phase, since expression (6) for $`E_g`$ reduces to $`E_g=\frac{1}{2}[1\mathrm{sgn}(S)K(\lambda _2/\sigma )]`$ in the limit $`|\lambda _1|\mathrm{}`$ (note: $`\lambda _1=S/J`$). However, at times $`\tau `$ approaching $`\tau _\pm `$ it is no longer valid to argue that $`S`$ is $`𝒪(\sqrt{N})`$; it is now $`𝒪(N^0)`$ and we enter the scaling regime of Phase II.
## 5 Phase II: Transition to Error Correction
As shown in the previous section, $`S`$ is an $`𝒪(N^0)`$ quantity in phase II, and it is also clear that $`\{Q,R\}`$ are $`𝒪(N^0)`$ at the start of phase II. In phase II (and, as we will see, also in phase III) we have to consider the observables $`𝛀=(S,𝚽)`$, with $`𝚽=(Q,R)`$; the reason for this slight departure from our phase I terminology will soon become clear. We can now express equation (8) as
$$𝒲_\tau [𝛀,𝛀^{}]=F_{\mathrm{I}\mathrm{I},\mathrm{I}\mathrm{I}\mathrm{I}}\frac{d\widehat{𝛀}}{(2\pi )^3}e^{i\widehat{𝛀}𝛀}e^{i\widehat{S}S(𝑱+\mathrm{\Delta }𝑱)i_{\mu =1}^2\widehat{\mathrm{\Phi }}_\mu \mathrm{\Phi }_\mu (𝑱+\mathrm{\Delta }𝑱)}_𝝃e^{i\widehat{𝛀}𝛀(𝑱)}_{𝛀^{},\tau }$$
Here
$$\mathrm{\Phi }_\mu (𝑱+\mathrm{\Delta }𝑱)=\mathrm{\Phi }_\mu (𝑱)+\underset{i}{}\mathrm{\Delta }J_i\frac{\mathrm{\Phi }_\mu }{J_i}+\frac{1}{2}\underset{ij}{}\mathrm{\Delta }J_i\mathrm{\Delta }J_j\frac{^2\mathrm{\Phi }_\mu }{J_iJ_j}$$
(this expansion is exact, since $`\{Q,R\}`$ are quadratic and linear functions, respectively). Substituting and expanding the exponential gives
$$𝒲_\tau [𝛀,𝛀^{}]=F_{\mathrm{I}\mathrm{I},\mathrm{I}\mathrm{I}\mathrm{I}}\frac{d\widehat{𝛀}}{(2\pi )^3}e^{i\widehat{𝛀}[𝛀𝛀(𝑱)]}e^{i\widehat{S}\mathrm{\Delta }𝑱𝑨}1_𝝃\frac{d\widehat{𝛀}}{(2\pi )^3}e^{i\widehat{𝛀}[𝛀𝛀(𝑱)]}\left\{\mathrm{}\right\}_{\mathrm{I}\mathrm{I},\mathrm{I}\mathrm{I}\mathrm{I}}_𝛀^{}$$
where
$$\left\{\mathrm{}\right\}_{\mathrm{I}\mathrm{I},\mathrm{I}\mathrm{I}\mathrm{I}}=iF_{\mathrm{I}\mathrm{I},\mathrm{I}\mathrm{I}\mathrm{I}}\underset{i\mu }{}\mathrm{\Delta }J_i\frac{\mathrm{\Phi }_\mu }{J_i}e^{i\widehat{S}\mathrm{\Delta }𝑱𝑨}_𝝃\widehat{\mathrm{\Phi }}_\mu +\frac{1}{2}iF_{\mathrm{I}\mathrm{I},\mathrm{I}\mathrm{I}\mathrm{I}}\underset{ij\mu }{}\mathrm{\Delta }J_i\mathrm{\Delta }J_j\frac{^2\mathrm{\Phi }_\mu }{J_iJ_j}e^{i\widehat{S}𝒌𝑨}_𝝃\widehat{\mathrm{\Phi }}_\mu $$
$$+\frac{1}{2}F_{\mathrm{I}\mathrm{I},\mathrm{I}\mathrm{I}\mathrm{I}}\underset{ij\mu \nu }{}\mathrm{\Delta }J_i\mathrm{\Delta }J_j\frac{\mathrm{\Phi }_\mu }{J_i}\frac{\mathrm{\Phi }_\nu }{J_j}e^{i\widehat{S}\mathrm{\Delta }𝑱𝑨}_𝝃\widehat{\mathrm{\Phi }}_\mu \widehat{\mathrm{\Phi }}_\nu +𝒪(N^{3/2})$$
(11)
Note: whereas it is valid to expand $`e^{i\widehat{𝚽}𝚽(𝑱+\mathrm{\Delta }𝑱)}`$ in the manner just described, we cannot treat $`e^{i\widehat{S}S(𝑱+\mathrm{\Delta }𝑱)}`$ in the same way since $`_i\mathrm{\Delta }J_i(S/J_i)=𝑨\mathrm{\Delta }𝑱=𝒪(N^0)`$ in phases II and III. Equations (7,11) form the basis for our study of Phases II and III.
The time scale $`\tau `$ in Phase II is related to $`t`$ via $`t=\tau N^1`$, so that $`F_{\mathrm{I}\mathrm{I}}=N^0`$, but although this phase is of short duration it has an important role as regards the stochastic evolution of the bias overlap parameter $`S`$. It is straightforward to show that the third term in equation (11) makes no contribution in the limit of large $`N`$. Moreover, in the very short phase II we may approximate $`\mathrm{\Delta }𝑱𝑨`$ by $`𝒌𝑨`$ (10). Referring to the details and notation in the appendix we have
$$e^{i\widehat{S}𝒌𝑨}=1E_g+e^{i\eta a^2\widehat{S}}I_1(\lambda _1,\lambda _2,\omega )+e^{i\eta a^2\widehat{S}}I_1(\lambda _1,\lambda _2,\omega ).$$
(12)
and we then find that in phase II
$$𝒲_\tau [𝛀,𝛀^{}]=I_1(\lambda _1,\lambda _2,\omega )\delta [SS^{}+\eta a^2]\delta [𝚽𝚽^{}]+I_1(\lambda _1,\lambda _2,\omega )\delta [SS^{}\eta a^2]\delta [𝚽𝚽^{}]E_g\delta [𝛀𝛀^{}].$$
Substitution into (7) and repetition of the arguments used for phase I we find that $`Q`$ and $`R`$ remain constant in phase II, whilst the conditional distribution $`P_\tau (S|Q,R)`$ satisfies
$$\frac{d}{d\tau }P_\tau (S|Q,R)=I_1(\lambda _1(S^{}),\lambda _2,\omega )P_\tau (S^{}|Q,R)+I_1(\lambda _1(S^+),\lambda _2,\omega )P_\tau (S^+|Q,R)$$
$$E_g(S,Q,R)P_\tau (S|Q,R)$$
where $`S^\pm =S\pm \eta a^2`$. The distribution equilibrates, on the relevant time-scale, to a stationary distribution $`P(S|Q,R)`$ given as the solution of
$$E_g(S,Q,R)P(S|Q,R)=I_1(\lambda _1(S^{}),\lambda _2,\omega )P(S^{}|Q,R)+I_1(\lambda _1(S^+),\lambda _2,\omega )P(S^+|Q,R).$$
Using relation (5) we find that this equilibrium condition can be written as $`A(S)+B(S)=A(S^+)+B(S^{})`$, where $`A(S)=I_1(\lambda _1(S),\lambda _2,\omega )P(S|Q,R)`$ and $`B(S)=I_1(\lambda _1(S),\lambda _2,\omega )P(S|Q,R)`$. One can easily show by taking Fourier transforms that it is satisfied by $`B(S)=A(S^+)`$, the correctness of which is evident by substitution. In this phase the permissible values of $`S`$ are those which differ from some initial value $`S(0)`$ by an integral multiple of $`\eta a^2`$. Upon writing the allowed values of $`S`$ as $`S_n=S(0)+n\eta a^2`$, we immediately obtain $`P(S|Q,R)=_{n=\mathrm{}}^{\mathrm{}}w(S_{n+1}|Q,R)\delta [SS_n]`$, where
$$w(S_{n+1}|Q,R)=\frac{I_1(\lambda _1(n),\lambda _2,\omega )}{I_1(\lambda _1(n+1),\lambda _2,\omega )}w(S_n|Q,R),$$
(13)
with $`I_1`$ as given in (22). Equation (13) fully determines the quasi-stationary distribution $`P(S|Q,R)`$. Comparison with numerical simulations shows very satisfactory agreement, see e.g. figure 2. The above picture is also in line with our intuition, since in a single step the change in $`S`$ is given by
$$\mathrm{\Delta }S=\mathrm{\Delta }𝑱𝑨=\frac{1}{2}\eta a^2[\mathrm{sgn}(\lambda _2+y)\mathrm{sgn}(\lambda _1+x)]+\frac{1}{2}\eta a\frac{z}{\sqrt{N}}[\mathrm{sgn}(\lambda _2+y)\mathrm{sgn}(\lambda _1+x)].$$
Provided we can neglect the $`N^{\frac{1}{2}}`$ term in this expression, which is true on the time scale of phase II, we see that in a single update $`\mathrm{\Delta }S\{0,\pm \eta a^2\}`$. However, if the $`N^{\frac{1}{2}}`$ term could be neglected indefinitely this would imply that, far into the future, the system would retain a memory of its initial conditions. In fact the term $`\frac{1}{2}\eta az[\mathrm{sgn}(\lambda _2+y)\mathrm{sgn}(\lambda _1+x)]/\sqrt{N}`$ represents a random walk superposed on the quasi-stationary distribution found for $`S`$ in phase II.
## 6 Phase III: Error Correction
As we enter phase III, where $`F_{\mathrm{I}\mathrm{I}\mathrm{I}}=N`$, the above ‘random walk’ term will come to have a significant role after about $`N`$ iterations<sup>2</sup><sup>2</sup>2We are grateful to Peter Sollich for pointing this out., leading to a modified probability distribution which contains a diffusion term: $`S_nS_n+s(t)`$. The walk is given by
$$s(t)=\frac{\eta a}{2\sqrt{N}}\underset{\mu =1}{\overset{Nt}{}}z(\mu )[\mathrm{sgn}(\lambda _2+y(\mu ))\mathrm{sgn}(\lambda _1(\mu )+x(\mu ))]$$
in an obvious notation, where the fields $`z(\mu )`$ are, as we have seen earlier, independent of $`(x,y)`$. The random walk addition $`s(t)`$ has mean zero, and variance given by
$$s^2(t)=\frac{\eta ^2a^2\sigma ^2}{2N}\underset{\mu =1}{\overset{Nt}{}}[1\mathrm{sgn}(\lambda _2+y(\mu ))\mathrm{sgn}(\lambda _1(\mu )+x(\mu ))]=t(\eta a\sigma )^2E_g$$
(14)
where $`E_g`$ is to be interpreted as a time average of $`E_g`$ over phase III, up to time $`t`$.
In order to extract the macroscopic laws in phase III we will now have to analyse this diffusion effect carefully, starting from equation (11). The details of this analysis are given in appendix B, where we show that for large $`N`$ the macroscopic distribution in phase III will again be of the form $`P_\tau (S,Q,R)=P_\tau (S|Q,R)\delta [QQ(\tau )]\delta [RR(\tau )]`$, but now with the deterministic values $`\{Q(\tau ),R(\tau )\}`$ given as the solution of the coupled equations
$$\frac{d}{d\tau }Q=\eta \sqrt{Q}𝑑S(K_1+L_1+M_1)P_\tau (S|Q,R)+\frac{1}{2}\eta ^2𝑑S(K_3+L_3+M_3)P_\tau (S|Q,R)$$
(15)
$$\frac{d}{d\tau }R=\frac{1}{2}\eta 𝑑S(K_2+L_2+M_2)P_\tau (S|Q,R)$$
(16)
The factors $`\{K_i,L_i,M_i\}`$, defined in appendix B, are indeed functions of $`S`$ (via $`\lambda _1`$) and of $`\{Q,R\}`$. The origin and meaning of these two equations can be appreciated more clearly by writing them in the following, somewhat more appealing, form (without as yet specifying the learning rule $`[J,u,T])`$:
$$\frac{d}{d\tau }Q=2\eta J𝑑SP_\tau (S|Q,R)(\lambda _1+x)\mathrm{sgn}(\lambda _2+y)[\sqrt{Q},\lambda _1+x,\mathrm{sgn}(\lambda _2+y)]$$
$$+\eta ^2𝑑SP_\tau (S|Q,R)^2[\sqrt{Q},\lambda _1+x,\mathrm{sgn}(\lambda _2+y)]$$
$$\frac{d}{d\tau }R=\eta 𝑑SP_\tau (S|Q,R)|y|[\sqrt{Q},\lambda _1+x,\mathrm{sgn}(\lambda _2+y)]$$
(see for details). Although equations (15,16) are superficially similar to the equations which we derived in phase I, we now have a situation in which functions of $`S`$ are weighted with respect to the probability distribution $`P_\tau (S|Q,R)`$ which satisfies a partial differential equation derived from equation (26) (in appendix B) by integration over $`Q`$ and $`R`$, namely
$$\frac{d}{d\tau }P_\tau (S|Q,R)=$$
$$N[\text{}I_1(\lambda _1(S^+,\lambda _2,\omega )P_\tau (S^+|Q,R)+I_1(\lambda _1(S^{},\lambda _2,\omega )P_\tau (S^{}|Q,R)E_g(S,Q,R)P_\tau (S|Q,R)]$$
$$+\frac{1}{2}\eta ^2a^2\sigma ^2\left[\frac{^2}{S^2}[I_1(\lambda _1(S^+),\lambda _2,\omega )P_\tau (S^+|Q,R)]+\frac{^2}{S^2}[I_1(\lambda _1(S^{}),\lambda _2,\omega )P_\tau (S^{}|Q,R)]\right]$$
(17)
Equations (15,16,17), together with the definitions of the short-hands $`\{K_i,L_i,M_i\}`$ as given in appendix B, provide an exact and closed set of equations for the macroscopic dynamics in phase III, in terms of the observables $`\{S,Q,R\}`$. In the large $`N`$ limit, $`Q`$ and $`R`$ satisfy deterministic equations, as in conventional no-bias theories, but $`S`$ remains stochastic throughout phase III. Furthermore, the persistent appearance of the factor $`\lambda _2`$ (which depends on the actual realisation of the teacher weights) induces sample-to-sample fluctuations. An example of the result of solving the coupled equations (15,16,17) numerically (via a numerical realisation, i.e. Monte Carlo, of the conditional stochastic process (17) for $`S`$) is shown in figure 3, and compared with numerical simulations of the underlying microscopic perceptron learning process. The agreement between theory and experiment is quite satisfactory.
## 7 Asymptotics of the Generalisation Error
A full numerical study of our equations (15,16,17) would be difficult, but these equations undergo a great simplification, permitting further analysis, if we make the approximation $`P_\tau (S|Q,R)=\delta [SS]`$, and assume that $`\lambda _1(S)=\lambda _2`$; numerical simulations confirm the validity of the replacement of $`\lambda _1`$ by $`\lambda _2`$ on average in phase III. In this approximation equations (15,16) become
$$\frac{d}{d\tau }Q=\eta \sqrt{Q}[K_1+L_1+M_1]+\frac{1}{2}\eta ^2[K_3+L_3+M_3]\frac{d}{d\tau }R=\frac{1}{2}[K_2+L_2+M_2]$$
Note that $`K_1+L_1+M_1=\lambda _1[(A_1+B_1+C_1)(A_2+B_2+C_2)]+(A_3+B_3+C_3)(A_4+B_4+C_4)`$. Referring to appendix A for the relevant expressions for $`\{A_i,B_i,C_i\}`$ in terms of the integrals $`I_1(\lambda _1,\lambda _2,\omega )`$ and $`I_2(\lambda _1,\lambda _2,\omega )`$, and using the identity $`K(\alpha )=_{\mathrm{}}^{\mathrm{}}DyK((\alpha \omega y)/\sqrt{1\omega ^2})`$, we find that in the approximation $`\lambda _1=\lambda _2`$ the following identities hold:
$$A_1+B_1+C_1=K(\lambda _2/\sigma ),A_2+B_2+C_2=K(\lambda _2/\sigma )$$
$$A_3+B_3+C_3=\sqrt{\frac{2}{\pi }}\omega \sigma e^{\frac{\lambda _2^2}{2\sigma ^2}},A_4+B_4+C_4=\sqrt{\frac{2}{\pi }}\sigma e^{\frac{\lambda _2^2}{2\sigma ^2}},K_1+L_1+M_1=\sqrt{\frac{2}{\pi }}\sigma (1\omega )e^{\frac{\lambda _2^2}{2\sigma ^2}}.$$
$$K_2+L_2+M_2=(A_5+B_5+C_5)(A_6+B_6+C_6)=\sqrt{\frac{2}{\pi }}\sigma (1\omega )e^{\frac{\lambda _2^2}{2\sigma ^2}}$$
$$K_3+L_3+M_3=1_{\frac{\lambda _2}{\sigma }}^{\mathrm{}}DyK\left(\frac{\lambda _2+\omega \sigma y}{\sigma \sqrt{1\omega ^2}}\right)+_{\frac{\lambda _2}{\sigma }}^{\mathrm{}}DyK\left(\frac{\lambda _2\omega \sigma y}{\sigma \sqrt{1\omega ^2}}\right)=2E_g$$
and equations (15,16) therefore become (upon rewriting the equation for $`Q`$ in terms of $`J=\sqrt{Q}`$):
$$\frac{d}{d\tau }J=\frac{\eta }{\sqrt{2\pi }}\sigma (1\omega )e^{\frac{\lambda _2^2}{2\sigma ^2}}+\frac{\eta ^2}{2J}E_g\frac{d}{d\tau }R=\frac{\eta }{\sqrt{2\pi }}\sigma (1\omega )e^{\frac{\lambda _2^2}{2\sigma ^2}}$$
(18)
The corresponding equation for $`\omega =R/J`$ is
$$\frac{d}{d\tau }\omega =\frac{\eta }{J\sqrt{2\pi }}\sigma (1\omega ^2)e^{\frac{\lambda _2^2}{2\sigma ^2}}\frac{\omega \eta ^2}{2J}E_g$$
(19)
which is to be solved in combination with (6). Numerical solution of these equations is found to be in very good agreement with the results of numerical simulations, even for finite times; however, it is relevant to consider what basis exists for making the approximation $`\lambda _1=\lambda _2`$, other than the fact that it works. We have already observed that the probability distribution for $`S`$ in phase III is a random walk superposed on the underlying discrete distribution which emerged in phase II. Equation (14) indicates that the random walk, reflected in the diffusion terms in equation (17), could in principle lead to a large variance for $`S`$, were this random walk not coupled to the underlying discrete distribution via equation (17). The discrete distribution and the random walk, however, are found to interact in such a way that the fluctuations actually tend to zero in the limit $`\tau \mathrm{}`$; this is confirmed by the results of numerical simulations which show that the fluctuations in $`\lambda _1=S/J`$ decrease with time and that on average $`\lambda _1`$ tends to $`\lambda _2`$, see figure 4. In a single step the average change in $`S`$ is equal to
$$\frac{1}{2}\eta a^2[\mathrm{sgn}(\lambda _2+y)\mathrm{sgn}(\lambda _1+x)]+\frac{1}{2}\eta a\frac{z}{\sqrt{N}}[\mathrm{sgn}(\lambda _2+y)\mathrm{sgn}(\lambda _1+x)]=\frac{1}{2}\eta a^2[K(\frac{\lambda _2}{\sigma })K(\frac{\lambda _1}{\sigma })]$$
so, as the fluctuations in $`S`$ diminish, we do indeed expect that $`\lambda _1`$ will tend to $`\lambda _2.`$
We will now use the coupled equations (6,19) to derive an asymptotic expression for the generalisation error $`E_g`$. Differentiation of (6) with respect to $`\omega `$ gives
$$\frac{E_g}{\omega }=\frac{e^{\frac{1}{2}\beta ^2}}{\pi \sqrt{1\omega ^2}}e^{\frac{1}{2}\beta ^2(1\omega )/(1+\omega )}$$
with the constant $`\beta =\lambda _2/\sigma `$. Changing the variable to $`\omega =\mathrm{cos}\theta `$, and expanding for $`\theta 0`$ gives
$$E_g=\pi ^1e^{\frac{1}{2}\beta ^2}_0^\theta 𝑑ue^{\frac{1}{2}\beta ^2\mathrm{tan}^2(u/2)}=\pi ^1e^{\frac{1}{2}\beta ^2}[\theta \frac{\beta ^2\theta ^3}{24}+𝒪(\theta ^5)]$$
(20)
Equation (18) for $`J`$ and equation (19) for $`\omega `$ can now be written
$$\frac{dJ}{d\tau }=\frac{\eta \sigma }{\sqrt{2\pi }}(1\mathrm{cos}\theta )e^{\frac{1}{2}\beta ^2}+\frac{\eta ^2E_g}{2J}J\mathrm{sin}\theta \frac{d\theta }{d\tau }=\frac{\eta \sigma }{\sqrt{2\pi }}\mathrm{sin}^2\theta e^{\frac{1}{2}\beta ^2}\frac{\eta ^2E_g\mathrm{cos}\theta }{2J}$$
Using the expansion $`\mathrm{tan}\theta =\theta +\frac{1}{3}\theta ^3+𝒪(\theta ^5)`$ we then expand our previous equations for the evolution of $`J`$ and $`\theta `$, giving
$$\frac{d}{d\tau }\theta =e^{\frac{1}{2}\beta ^2}\left\{\frac{\eta \sigma \theta }{J\sqrt{2\pi }}+\frac{\eta ^2}{2\pi J^2}\frac{\eta ^2\theta ^2\rho }{2\pi J^2}\right\}+𝒪(\theta ^4),\mathrm{with}\rho =\frac{1}{24}\beta ^2+\frac{1}{3}$$
$$\frac{d}{d\tau }J=e^{\frac{1}{2}\beta ^2}\left\{\frac{\eta \sigma \theta ^2}{2\sqrt{2\pi }}+\frac{\eta ^2}{2\pi J}[\theta \frac{1}{24}\beta ^2\theta ^3]\right\}+𝒪(\theta ^5).$$
Upon making the asymptotic ansatz $`J=A/\theta `$, the equation for $`J`$ can now be expressed so as to give a second equation for $`\theta `$. The two resulting equations for $`d\theta /d\tau `$ are
$$\frac{d}{d\tau }\theta =\frac{\eta \sigma \theta ^4}{2A\sqrt{2\pi }}e^{\frac{1}{2}\beta ^2}\frac{\eta ^2\theta ^4}{2A^2\pi }\left[1\frac{\beta ^2\theta ^2}{24}+𝒪(\theta ^4)\right]e^{\frac{1}{2}\beta ^2}$$
and
$$\frac{d}{d\tau }\theta =\frac{\eta \sigma \theta ^2}{A\sqrt{2\pi }}e^{\frac{1}{2}\beta ^2}+\frac{\eta ^2\theta ^2}{2\pi A^2}e^{\frac{1}{2}\beta ^2}+𝒪(\theta ^4)$$
Consistency requires that $`A`$ be given by $`A=\eta /\sigma \sqrt{2\pi }`$. The asymptotic equation for $`\theta `$ subsequently becomes $`d\theta /d\tau =\frac{1}{2}\sigma ^2\theta ^4e^{\frac{1}{2}\beta ^2}`$, from which we obtain the asymptotic power law $`\theta =k\tau ^\alpha ,`$ where $`\alpha =\frac{1}{3}`$ and $`k^3=2e^{\frac{1}{2}\beta ^2}/3\sigma ^2`$. Combining this, finally, with (20) we then obtain, recalling that in phase III one simply has $`\tau =m/N=t`$:
$$E_g(t)=\rho (a)e^{\lambda _2^2/3\sigma ^2}t^{\frac{1}{3}}(t\mathrm{}),\rho (a)=\left[\frac{2}{3\pi ^3}\right]^{\frac{1}{3}}(1a^2)^{\frac{1}{3}}$$
(21)
Note that the power of $`\tau `$ occurring in this expression is the same as the power which appears in the asymptotic form of the generalisation error in the conventional no-bias theory; the coefficient is however different, but reduces to the familiar form in the case of zero bias, where $`a=\lambda _2=0`$ and $`\sigma =1`$. Moreover, our prediction of the asymptotic form of $`E_g`$ is in excellent agreement with the results of numerical simulations. This is evident from figure (5), where we show the observed function $`\rho (a)`$, defined as $`\rho (a)=lim_t\mathrm{}E_g(t)t^{\frac{1}{3}}e^{\lambda _2^2/3\sigma ^2}`$, versus the theoretical prediction as given in (21). Note that the dependence of (21) on the teacher-bias overlap $`\lambda _2=𝑩𝑨`$ implies sample-to-sample fluctuations.
## 8 Discussion
We have studied analytically the dynamics of on-line learning in non-linear perceptrons, trained according to the perceptron rule, for the scenario of having structurally biased, i.e. $`𝒪(N^0)`$, input data. The bias changes qualitatively the learning process, inducing three distinct phases (with different scaling properties) and persistent stochastic as well as sample-to-sample fluctuations in the generalisation error, even for $`N\mathrm{}`$. At a theoretical level, the need to introduce an extra order parameter $`S`$ (the projection of the student weight vector in the direction of the bias) which is neither deterministic nor self-averaging makes the analysis considerably more involved than that of the idealised bias free case. In the third and final phase, in which most of the learning takes place, we have obtained a set of exact closed equations which involve the conditional probability density of $`S`$. However, because of their complicated nature, an exact analytic solution of these equations appears to be out of the question, as is also generally the case in the more familiar no-bias scenarios. Nevertheless we have found that an approximate (and much simpler) version of our equations yields results which are in excellent agreement with numerical simulations. We show that the asymptotic power law for the generalisation error is largely preserved, with the bias showing up only in the pre-factor. At various stages throughout out calculations we have compared the predictions of our macroscopic dynamic equations with the results of numerical simulations of the underlying (microscopic) learning process, which consistently showed excellent agreement.
Although in this paper we have confined ourselves to the perceptron learning rule, it is clear that our analysis is in no way restricted to this particular rule, and can be applied to other rules such as the AdaTron learning rule, where $`\mathrm{\Delta }𝑱=\frac{\eta }{2N}𝝃[\mathrm{sgn}(𝑩𝝃)\mathrm{sgn}(𝑱𝝃)]|𝑱𝝃|`$; one could even study optimal learning rates and optimal learning rules, generalising to the case of having $`a0`$. Preliminary studies of the AdaTron learning rule with structurally biased data show, for instance, that the simple result (13), describing the phase II distribution in the case of the perceptron, is replaced by the integral equation
$$E_gP_\tau (\stackrel{~}{S}|Q,R)=\eta a^2J_{\lambda _1}^{\mathrm{}}d\rho G(\rho ,\lambda _2)P_\tau (\stackrel{~}{S}+(\lambda _1+\rho )\eta a^2J)|Q,R)$$
$$+\eta a^2J_{\lambda _1}^{\mathrm{}}d\rho G(\rho ,\lambda _2)P_\tau (\stackrel{~}{S}+(\lambda _1\rho )\eta a^2J)|Q,R).$$
where $`G`$ is defined by
$$G(x,\lambda _2)=\frac{e^{\frac{x^2}{2\sigma ^2}}}{2\sigma \sqrt{2\pi }}\left[1K\left(\frac{\lambda _2+\omega x}{\sigma \sqrt{1\omega ^2}}\right)\right],$$
The discrete distribution which in the present paper we found for the perceptron in phase II no longer applies in the Adatron case, and is replaced by a continuous distribution which satisfies the above integral equation. The analysis of the AdaTron in the case of biased data is more complicated than for the perceptron, as might have been expected from the nature of the AdaTron learning rule, but much of the work which we have presented for the perceptron can be carried through and the results will be published in . There is also scope for a more detailed mathematical investigation of the partial differential equation which we derived to describe the conditional probability distribution $`P_\tau (S|Q,R)`$ for the perceptron, but this is likely to be difficult, and beyond the scope of the present paper.
### Acknowledgements
It is a pleasure to thank Peter Sollich and Michael Biehl for helpful comments and discussions. JAFH wishes to thank King’s College London Association for financial support.
## Appendix A Integrals and Averages
We recall that the function $`K`$ is defined by $`K(x)=\mathrm{erf}(x/\sqrt{2})`$. In terms of this definition note that $`_\tau ^{\mathrm{}}𝑑\zeta e^{\frac{1}{2}\zeta ^2}=\frac{1}{2}\sqrt{2\pi }[1K(\tau )]`$. We now proceed to list various integrals which occur in our calculations, or are referred to in the text, and where appropriate outline a brief derivation. Recall that the joint distribution of $`(x,y)=(\widehat{𝑱}𝒗,𝑩𝒗)`$ is given by
$$p(x,y)=[2\pi \sigma ^2\sqrt{1\omega ^2}]^1e^{\frac{1}{2}\frac{[x^22\omega xy+y^2]}{\sigma ^2(1\omega ^2)}}$$
where $`\sigma ^2=1a^2`$ and $`\omega =𝑩\widehat{𝑱}`$. We then find that
$$I_1(\lambda _1,\lambda _2,\omega )=_{\lambda _1}^{\mathrm{}}𝑑x_{\lambda _2}^{\mathrm{}}𝑑yp(x,y)=_{\lambda _2}^{\mathrm{}}\frac{dy}{2\pi \sigma }e^{\frac{y^2}{2\sigma ^2}}_{\frac{\lambda _1\omega y}{\sigma \sqrt{1\omega ^2}}}^{\mathrm{}}𝑑\zeta e^{\frac{1}{2}\zeta ^2}$$
$$=\frac{1}{4}\left[1K(\frac{\lambda _2}{\sigma })\right]\frac{1}{2}_{\frac{\lambda _2}{\sigma }}^{\mathrm{}}DyK(\frac{\lambda _1\omega \sigma y}{\sigma \sqrt{1\omega ^2}})$$
(22)
and similarly
$$I_2(\lambda _1,\lambda _2,\omega )=_{\lambda _1}^{\mathrm{}}𝑑xx_{\lambda _2}^{\mathrm{}}𝑑yp(x,y)=_{\lambda _1}^{\mathrm{}}\frac{dxx}{2\sigma \sqrt{2\pi }}e^{\frac{x^2}{2\sigma ^2}}_{\lambda _1}^{\mathrm{}}\frac{dxx}{2\sqrt{2\pi }\sigma }e^{\frac{x^2}{2\sigma ^2}}K(\frac{\lambda _2\omega x}{\sigma \sqrt{1\omega ^2}})$$
$$=\frac{\sigma }{2\sqrt{2\pi }}e^{\frac{\lambda _1^2}{2\sigma ^2}}\left[1K\left(\frac{\lambda _2\omega \lambda _1}{\sigma \sqrt{1\omega ^2}}\right)\right]+\frac{\omega \sigma }{2\sqrt{2\pi }}e^{\frac{\lambda _2^2}{2\sigma ^2}}\left[1K\left(\frac{\lambda _1\omega \lambda _2}{\sigma \sqrt{1\omega ^2}}\right)\right]$$
The following averages with respect to the distribution $`p(x,y)`$ are easily calculated:
$$\mathrm{sgn}(\lambda _1+x)=K\left(\frac{\lambda _1}{\sigma }\right),x\mathrm{sgn}(\lambda _2+y)=\sqrt{\frac{2}{\pi }}\omega \sigma e^{\frac{\lambda _2^2}{2\sigma ^2}},x\mathrm{sgn}(\lambda _1+x)=\sqrt{\frac{2}{\pi }}\sigma e^{\frac{\lambda _1^2}{2\sigma ^2}}$$
$$\mathrm{sgn}(\lambda _1+x)\mathrm{sgn}(\lambda _2+y)=I_1(\lambda _1,\lambda _2,\omega )I_1(\lambda _1,\lambda _2,\omega )I_1(\lambda _1,\lambda _2,\omega )+I_1(\lambda _1,\lambda _2,\omega )$$
$$=_{\frac{\lambda _2}{\sigma }}^{\mathrm{}}DyK\left(\frac{\lambda _1+\omega \sigma y}{\sigma \sqrt{1\omega ^2}}\right)_{\frac{\lambda _2}{\sigma }}^{\mathrm{}}DyK\left(\frac{\lambda _1\omega \sigma y}{\sigma \sqrt{1\omega ^2}}\right)$$
Finally, in studying phases II and III we require the following averages:
$$\mathrm{sgn}(\lambda _2+y)e^{i\widehat{S}𝒌𝑨}=A_1+B_1e^{i\widehat{S}\eta a^2}+C_1e^{i\widehat{S}\eta a^2}$$
$$\mathrm{sgn}(\lambda _1+x)e^{i\widehat{S}𝒌𝑨}=A_2+B_2e^{i\widehat{S}\eta a^2}+C_2e^{i\widehat{S}\eta a^2}$$
$$x\mathrm{sgn}(\lambda _2+y)e^{i\widehat{S}𝒌𝑨}=A_3+B_3e^{i\widehat{S}\eta a^2}+C_3e^{i\widehat{S}\eta a^2}$$
$$x\mathrm{sgn}(\lambda _1+x)e^{i\widehat{S}𝒌𝑨}=A_4+B_4e^{i\widehat{S}\eta a^2}+C_4e^{i\widehat{S}\eta a^2}$$
$$y\mathrm{sgn}(\lambda _2+y)e^{i\widehat{S}𝒌𝑨}=A_5+B_5e^{i\widehat{S}\eta a^2}+C_5e^{i\widehat{S}\eta a^2}$$
$$y\mathrm{sgn}(\lambda _1+x)e^{i\widehat{S}𝒌𝑨}=A_6+B_6e^{i\widehat{S}\eta a^2}+C_6e^{i\widehat{S}\eta a^2}$$
$$e^{i\widehat{S}𝒌𝑨}=A_7+B_7e^{i\widehat{S}\eta a^2}+C_7e^{i\widehat{S}\eta a^2}$$
$$\mathrm{sgn}(\lambda _1+x)\mathrm{sgn}(\lambda _2+y)e^{i\widehat{S}𝒌𝑨}=A_8+B_8e^{i\widehat{S}\eta a^2}+C_8e^{i\widehat{S}\eta a^2}$$
where
$$\begin{array}{ccc}A_1=[I_1(\lambda _1,\lambda _2,\omega )I_1(\lambda _1,\lambda _2,\omega )],\hfill & B_1=I_1(\lambda _1,\lambda _2\omega ),\hfill & C_1=I_1(\lambda _1,\lambda _2,\omega ),\hfill \\ A_2=[I_1(\lambda _1,\lambda _2,\omega )I_1(\lambda _1,\lambda _2,\omega )],\hfill & B_2=I_1(\lambda _1,\lambda _2,\omega ),\hfill & C_2=I_1(\lambda _1,\lambda _2,\omega ),\hfill \\ A_3=[I_2(\lambda _1,\lambda _2,\omega )+I_2(\lambda _1,\lambda _2,\omega )],\hfill & B_3=I_2(\lambda _1,\lambda _2,\omega ),\hfill & C_3=I_2(\lambda _1,\lambda _2,\omega ),\hfill \\ A_4=[I_2(\lambda _1,\lambda _2,\omega )+I_2(\lambda _1,\lambda _2,\omega )],\hfill & B_4=I_2(\lambda _1,\lambda _2,\omega ),\hfill & C_4=I_2(\lambda _1,\lambda _2,\omega ),\hfill \\ A_5=[I_2(\lambda _2,\lambda _1,\omega )+I_2(\lambda _2,\lambda _1,\omega )],\hfill & B_5=I_2(\lambda _2,\lambda _1,\omega ),\hfill & C_5=I_2(\lambda _2,\lambda _1,\omega ),\hfill \\ A_6=[I_2(\lambda _2,\lambda _1,\omega )+I_2(\lambda _2,\lambda _1,\omega )],\hfill & B_6=I_2(\lambda _2,\lambda _1,\omega ),\hfill & C_6=I_2(\lambda _2,\lambda _1,\omega ),\hfill \\ A_7=1E_g,\hfill & B_7=I_1(\lambda _1,\lambda _2,\omega ),\hfill & C_7=I_1(\lambda _1,\lambda _2,\omega ),\hfill \\ A_8=[I_1(\lambda _1,\lambda _2,\omega )+I_1(\lambda _1,\lambda _2,\omega )],\hfill & B_8=I_1(\lambda _1,\lambda _2,\omega ),\hfill & C_8=I_1(\lambda _1,\lambda _2,\omega ).\hfill \end{array}$$
All these formulae may be established by elementary methods. For example,
$$x\mathrm{sgn}(\lambda _2+y)e^{i\widehat{S}𝒌𝑨}=𝑑x𝑑yxp(x,y)\mathrm{sgn}(\lambda _2+y)e^{\frac{1}{2}i\widehat{S}\eta a^2[\mathrm{sgn}(\lambda _2+y)\mathrm{sgn}(\lambda _1+x)]}$$
$$=\left[_{\mathrm{}}^{\lambda _1}𝑑xx+_{\lambda _1}^{\mathrm{}}𝑑xx\right]\left[_{\lambda _2}^{\mathrm{}}𝑑ye^{\frac{1}{2}i\widehat{S}\eta a^2(1\mathrm{sgn}(\lambda _1+x))}p(x,y)\right]$$
$$\left[_{\mathrm{}}^{\lambda _1}𝑑xx+_{\lambda _1}^{\mathrm{}}𝑑xx\right]\left[_{\lambda _2}^{\mathrm{}}𝑑ye^{\frac{1}{2}i\widehat{S}\eta a^2(1+\mathrm{sgn}(\lambda _1+x))}p(x,y)\right]$$
$$=_{\lambda _1}^{\mathrm{}}𝑑xx_{\lambda _2}^{\mathrm{}}𝑑yp(x,y)_{\lambda _1}^{\mathrm{}}𝑑xx_{\lambda _2}^{\mathrm{}}𝑑ye^{i\widehat{S}\eta a^2}p(x,y)$$
$$_{\lambda _1}^{\mathrm{}}𝑑xx_{\lambda _2}^{\mathrm{}}𝑑ye^{i\widehat{S}\eta a^2}p(x,y)+_{\lambda _1}^{\mathrm{}}𝑑xx_{\lambda _2}^{\mathrm{}}𝑑yp(x,y)$$
$$=I_2(\lambda _1,\lambda _2,\omega )e^{i\widehat{S}\eta a^2}I_2(\lambda _1,\lambda _2,\omega )e^{i\widehat{S}\eta a^2}I_2(\lambda _1,\lambda _2,\omega )+I_2(\lambda _1,\lambda _2,\omega )$$
$$=A_3+B_3e^{i\widehat{S}\eta a^2}+C_3e^{i\widehat{S}\eta a^2},$$
as required.
## Appendix B Analysis of Macroscopic Distribution in Phase III
Here we give the details of our analysis of the macroscopic distribution $`P_\tau (S,Q,R)`$ in phase III, starting from equation (11). We note that, in phase III:
$$e^{i\widehat{S}\mathrm{\Delta }𝑱𝑨}=e^{i\widehat{S}𝒌𝑨}\left\{1\frac{i\eta az\widehat{S}}{2\sqrt{N}}[\mathrm{sgn}(\lambda _2+y)\mathrm{sgn}(\lambda _1+x)]\frac{(\eta az\widehat{S})^2}{4N}[1\mathrm{sgn}(\lambda _1+x)\mathrm{sgn}(\lambda _2+y)]+\mathrm{}\right\}$$
The terms which we neglected are $`𝒪(N^2)`$, since when performing averages over the training set the average of the $`z^3`$ term is zero. Equation (11) now yields
$$𝒲_\tau [𝛀,𝛀^{}]=N\frac{d\widehat{𝛀}}{(2\pi )^3}e^{i\widehat{𝛀}[𝛀𝛀(𝑱)]}e^{i\widehat{S}𝒌𝑨}1\frac{(\eta a\sigma \widehat{S})^2}{4N}e^{i\widehat{S}𝒌𝑨}[1\mathrm{sgn}(\lambda _1+x)\mathrm{sgn}(\lambda _2+y)]_𝝃$$
$$\frac{d\widehat{𝛀}}{(2\pi )^3}e^{i\widehat{𝛀}[𝛀𝛀(𝑱)]}\left\{\mathrm{}\right\}_{\mathrm{I}\mathrm{I}\mathrm{I}}_𝛀^{}$$
(23)
where
$$\left\{\mathrm{}\right\}_{\mathrm{I}\mathrm{I}\mathrm{I}}=iN\underset{i\mu }{}k_i\frac{\mathrm{\Phi }_\mu }{J_i}e^{i\widehat{S}𝒌𝑨}_𝝃\widehat{\mathrm{\Phi }}_\mu +\frac{1}{2}iN\underset{ij\mu }{}k_ik_j\frac{^2\mathrm{\Phi }_\mu }{J_iJ_j}e^{i\widehat{S}𝒌𝑨}_𝝃\widehat{\mathrm{\Phi }}_\mu $$
$$+\frac{1}{2}N\underset{ij\mu \nu }{}k_ik_j\frac{\mathrm{\Phi }_\mu }{J_i}\frac{\mathrm{\Phi }_\nu }{J_j}e^{i\widehat{S}𝒌𝑨}_𝝃\widehat{\mathrm{\Phi }}_\mu \widehat{\mathrm{\Phi }}_\nu $$
(24)
We showed in appendix A that
$$𝑑x𝑑yp(x,y)\mathrm{sgn}(\lambda _1+x)\mathrm{sgn}(\lambda _2+y)e^{i\widehat{S}𝒌𝑨}=I_1(\lambda _1,\lambda _2,\omega )I_1(\lambda _1,\lambda _2,\omega )e^{i\widehat{S}\eta a^2}$$
$$I_1(\lambda _1,\lambda _2,\omega )e^{i\widehat{S}\eta a^2}+I_1(\lambda _1,\lambda _2,\omega )$$
so that
$$𝑑x𝑑yp(x,y)e^{i\widehat{S}𝒌𝑨}[1\mathrm{sgn}(\lambda _1+x)\mathrm{sgn}(\lambda _2+y)]=2[e^{i\widehat{S}\eta a^2}I_1(\lambda _1,\lambda _2,\omega )+e^{i\widehat{S}\eta a^2}I_1(\lambda _1,\lambda _2,\omega )],$$
by virtue of equation (12) and the fact that $`I_1(\lambda _1,\lambda _2,\omega )+I_1(\lambda _1,\lambda _2,\omega )=1E_g`$. Bearing in mind the sub-shell average we may write
$$𝑑\widehat{S}\widehat{S}^2e^{i\widehat{S}[S\pm \eta a^2S(𝑱)]}=2\pi \frac{^2}{S_{}^{}{}_{}{}^{2}}\delta [S\pm \eta a^2S^{}]$$
Upon combining equations (7,12,23) we find that in phase III the joint probability density $`P_\tau (S,Q,R)`$ satisfies
$$\frac{d}{d\tau }P_\tau (S,Q,R)=N[\text{}I_1(\lambda _1(S^+),\lambda _2,\omega )P_\tau (S^+,Q,R)+I_1(\lambda _1(S^{}),\lambda _2,\omega )P_\tau (S^{},Q,R)$$
$$E_g(S,Q,R)P_\tau (S,Q,R)\text{}]+\frac{1}{2}\eta ^2a^2\sigma ^2d𝛀^{}P_\tau (𝛀^{})[\text{}I_1(\lambda _1(S^{}),\lambda _2,\omega )\frac{^2}{S_{}^{}{}_{}{}^{2}}\delta [S+\eta a^2S^{}]$$
$$+I_1(\lambda _1(S^{}),\lambda _2,\omega )\frac{^2}{S_{}^{}{}_{}{}^{2}}\delta [S\eta a^2S^{}]]\delta [𝚽𝚽^{}]\frac{d𝛀^{}}{(2\pi )^3}P_\tau (𝛀^{})d\widehat{𝛀}e^{i\widehat{𝛀}[𝛀𝛀(𝑱)]}\left\{\mathrm{}\right\}_{\mathrm{I}\mathrm{I}\mathrm{I}}_𝛀^{}$$
where $`\left\{\mathrm{}\right\}_{\mathrm{I}\mathrm{I}\mathrm{I}}`$ is given by equation (24), and hence
$$\frac{d}{d\tau }P_\tau (S,Q,R)=N[\text{}I_1(\lambda _1(S^+),\lambda _2,\omega )P_\tau (S^+,Q,R)+I_1(\lambda _1(S^{}),\lambda _2,\omega )P_\tau (S^{},Q,R)$$
$$E_g(S,Q,R)P_\tau (S,Q,R)\text{}]+\frac{1}{2}\eta ^2a^2\sigma ^2[\frac{^2}{S^2}[I_1(\lambda _1(S^+),\lambda _2,\omega )P_\tau (S^+,Q,R)]$$
$$+\frac{^2}{S^2}[I_1(\lambda _1(S^{}),\lambda _2,\omega )P_\tau (S^{},Q,R)]]\frac{d𝛀^{}}{(2\pi )^3}P_\tau (𝛀^{})d\widehat{𝛀}e^{i\widehat{𝛀}[𝛀𝛀(𝑱)]}\left\{\mathrm{}\right\}_{\mathrm{I}\mathrm{I}\mathrm{I}}_𝛀^{}$$
(25)
As regards the evaluation of $`\left\{\mathrm{}\right\}_{\mathrm{I}\mathrm{I}\mathrm{I}}`$ we note that
$$N\underset{i\mu }{}k_i\frac{\mathrm{\Phi }_\mu }{J_i}e^{i\widehat{S}𝒌𝑨}_𝝃\widehat{\mathrm{\Phi }}_\mu =\eta J(\lambda _1+x)[\mathrm{sgn}(\lambda _2+y)\mathrm{sgn}(\lambda _1+x)]e^{i\widehat{S}𝒌𝑨}\widehat{\mathrm{\Phi }}_1$$
$$+\frac{1}{2}\eta (\lambda _2+x)[\mathrm{sgn}(\lambda _2+y)\mathrm{sgn}(\lambda _1+x)]e^{i\widehat{s}𝒌𝑨}\widehat{\mathrm{\Phi }}_2$$
$$=\eta J[K_1+L_1e^{i\widehat{s}\eta a^2}+M_1e^{i\widehat{s}\eta a^2}]\widehat{\mathrm{\Phi }}_1+\frac{1}{2}\eta [K_2+L_2e^{i\widehat{s}\eta a^2}+M_2e^{i\widehat{S}\eta a^2}]\widehat{\mathrm{\Phi }}_2$$
in which
$$K_1=\lambda _1(A_1A_2)+(A_3A_4),K_2=\lambda _2(A_1A_2)+(A_5A_6),$$
$$L_1=\lambda _1(B_1B_2)+(B_3B_4),L_2=\lambda _2(B_1B_2)+(B_5B_6),$$
$$M_1=\lambda _1(C_1C_2)+(C_3C_4),M_2=\lambda _2(C_1C_2)+(C_5C_6)$$
and $`A_i,B_i,C_i`$ are functions defined in appendix A and expressed in terms of the integrals $`I_1(\lambda _1,\lambda _2,\omega )`$ and $`I_2(\lambda _1,\lambda _2,\omega )`$. In a similar way we find that
$$N\underset{ij\mu }{}k_ik_j\frac{^2\mathrm{\Phi }_\mu }{J_iJ_j}e^{i\widehat{S}𝒌𝑨}_𝝃\widehat{\mathrm{\Phi }}_\mu =\eta ^2[1\mathrm{sgn}(\lambda _1+x)\mathrm{sgn}(\lambda _2+y)]e^{i\widehat{s}𝒌𝑨}\widehat{\mathrm{\Phi }}_1$$
$$=\eta ^2[K_3+L_3e^{i\widehat{s}\eta a^2}+M_3^{i\widehat{s}\eta a^2}]\widehat{\mathrm{\Phi }}_1$$
where $`K_3=A_7A_8`$, $`L_3=B_7B_8`$, and $`M_3=C_7C_8`$. Note that the term
$$N\underset{ij\mu \nu }{}k_ik_j\frac{\mathrm{\Phi }_\mu }{J_i}\frac{\mathrm{\Phi }_\nu }{J_j}e^{i\widehat{S}𝒌𝑨}_𝝃\widehat{\mathrm{\Phi }}_\mu \widehat{\mathrm{\Phi }}_\nu $$
makes no contribution to $`𝒲_\tau [𝛀,𝛀^{}]`$ in the limit of large $`N`$. Using equations (24) and (25) we can now carry out the remaining integrations using standard formulae from distribution theory, as described for earlier phases, and find that
$$\frac{d}{d\tau }P_\tau (S,Q,R)=$$
$$N\left[\text{}I_1(\lambda _1(S^+),\lambda _2,\omega )P_\tau (S^+,Q,R)+I_1(\lambda _1(S^{}),\lambda _2,\omega )P_\tau (S^{},Q,R)E_g(S,Q,R)P_\tau (S,Q,R)\right]$$
$$+\frac{1}{2}\eta ^2a^2\sigma ^2\left[\frac{^2}{S^2}[I_1(\lambda _1(S^+),\lambda _2,\omega )P_\tau (S^+,Q,R)]+\frac{^2}{S^2}[I_1(\lambda _1(S^{}),\lambda _2,\omega )P_\tau (S^{},Q,R)]\right]$$
$$\frac{}{Q}[\text{}\eta J[K_1P_\tau (S,Q,R)+L_1P_\tau (S^+,Q,R)+M_1P_\tau (S^{},Q,R)]$$
$$+\frac{1}{2}\eta ^2[K_3P_\tau (S,Q,R)+L_3P_\tau (S^+,Q,R)+M_3P_\tau (S^{},Q,R)]\text{}]$$
$$\frac{}{R}\left[\frac{1}{2}\eta [K_2P_\tau (S,Q,R)+L_2P_\tau (S^+,Q,R)+M_2P_\tau (S^{},Q,R)]\right]$$
(26)
Integration over $`S`$ now gives, in combination with the relation $`E_g=I_1(\lambda _1,\lambda _2,\omega )+I_1(\lambda _1,\lambda _2,\omega )`$:
$$\frac{d}{d\tau }P_\tau (Q,R)=\frac{}{Q}\{\text{}P_\tau (Q,R)[\eta JdS(K_1+L_1+M_1)P_\tau (S|Q,R)$$
$$+\frac{1}{2}\eta ^2dS(K_3+L_3+M_3)P_\tau (S|Q,R)]\}\frac{}{R}\{\text{}P_\tau (Q,R)\left[\frac{1}{2}\eta dS(K_2+L_2+M_2)P_\tau (S|Q,R)\right]\}$$
which is a Liouville equation with solution $`P_\tau (Q,R)=\delta [QQ(\tau )]\delta [RR(\tau )]`$, where the deterministic flow trajectories $`(Q(\tau ),R(\tau ))`$ are given as the solutions of (15,16), as claimed.
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# Acceleration of UHE Cosmic Rays in Gamma-Ray Bursts
## 1 Introduction
The origin of the cosmic ray population beyond $`10^{15}eV`$ is an important enigma of modern astrophysics which will hopefully be clarified by the Pierre Auger instrument. Beyond this energy, the galactic magnetic field is unable to scatter cosmic rays and moreover supernovae remnants cannot accelerate protons. Pulsars could accelerate protons up to $`10^{17}eV`$. This population, that has an isotropic energy spectrum with a powerlaw index $`3.1`$, probably comes from extragalactic sources. The “knee” of the spectrum around $`10^{15}eV`$ is so smooth that the extragalactic contribution could be important at lower energy also.
The photoproduction of pions by cosmic ray protons on the Cosmic Microwave Background, the GZK-effect, occurs beyond the energy threshold of $`3\times 10^{19}eV`$ and higher energy protons cannot come from sources located farther than $`100Mpc`$ (Aharonian & Cronin aha (1994)). Few events beyond the GZK-threshold have recently been detected by AGASA (Hayashida et al. aga (1998)) and the “Fly’s eyes” (Sokolsky fly (1998)), and the Auger experiment will considerably improve the statistics of these events. These preliminary results suggest that a new population of cosmic rays is observed at energies higher than $`10^{18}eV`$, because the spectrum becomes harder. The origin of these few events, although some pointing ability of the instruments, remains very uncertain. The identified extragalactic sources within $`100Mpc`$ are very few.
Vietri (vie (1995)) and Waxman (wax (1995)) showed that the rate of GRBs in the Universe ($`10^8yr^1Mpc^3`$) and their energy ($`10^{51}10^{53}erg`$) make them very good candidates as sources of the UHECRs, if $`10\%`$ of their energy is converted into UHECR energy. The most considered GRB model, namely the “fireball” model (Rees & Mészáros ree (1992)), is based on a relativistic blast wave having a Lorentz factor $`\mathrm{\Gamma }_s`$ larger than $`10^2`$, this high value being necessary for gamma photons to escape from pair production. A Fermi cycle at a relativistic shock can amplify the energy of cosmic ray by a factor $`\mathrm{\Gamma }_s^2`$. Thus two Fermi cycles could be sufficient to reach the expected energy of the UHECRs. However, Gallant & Achterberg (gal (1999)) recently showed that only the first cycle can produce such amplification, the next cycles amplifying by a factor of 2 only. Moreover the escape probability is large ($`0.30.5`$).
The GRB light curves indicate strong internal disturbances and they have been considered as the main cause of gamma-ray emission and a possible cause of cosmic ray generation (Waxman & Bahcall wax2 (1999) ; Daigne & Mochkovitch daim (1998)). In this paper we propose an interpretation of these disturbances in terms of relativistic hydromagnetic fronts, calculate their generation through the streaming instability caused by baryon loading, show that there is an efficient nonlinear generation of both forward and backward waves, estimate the amount of energy they extract from the fireball and calculate the acceleration of cosmic rays by their crossings. The baryon loading occurs first during the “primary” GRB when the relativistic pair wind entrains the debris of the progenitor, and then during the interaction of the fireball with the interstellar medium. Furthermore, the first loading is likely to occur in the “hypernova” scenario, however, this is less obvious in the “merging” scenario.
## 2 Loading baryons
The canonical fireball model (Rees & Mészáros ree (1992)) presents a primary explosion during which the plasma is optically thick until $`t_{}`$ ($`R_{}=ct_{}`$), then an adiabatic and free expansion until the ambient medium starts to exert sufficient ram pressure; this defines a deceleration time $`t_d`$ and a deceleration radius $`R_d`$ (typically $`t_{}=3\times 10^3t_d`$). The third stage is thus the deceleration one during which most of the energy of the fireball is dissipated in the interstellar medium, giving rise to the afterglow. We assume that the fireball is initially composed of $`e^+e^{}`$ pairs and that this ultra-relativistic wind of bulk Lorentz factor $`\mathrm{\Gamma }`$, flowing along a strong poloidal magnetic field, is perturbed by ambient baryons of mass density $`\rho _b`$. There are two stages of baryon loading, during the “primary” GRB and during the interaction with the interstellar medium.
### 2.1 Hydrodynamics of the baryon loading stage
The main features of the hydrodynamics of the relativistic fireball can be derived from the energy invariant as long as the radiation losses are negligible (assuming radial magnetic field lines):
$$E=(e+\beta ^2P)\gamma ^2d^3v,$$
(1)
where $`e`$ is the comoving energy-mass density, $`P`$ the comoving pressure, $`\beta `$ the flow motion (in unit of light velocity), and $`\gamma `$ the Lorentz factor of the flow ($`\gamma (1\beta ^2)^{1/2}`$), distributed around the typical bulk Lorentz factor $`\mathrm{\Gamma }`$. We assume that the plasma shell is composed of a low pressure baryonic component of mass $`M`$ and a high pressure relativistic component containing pairs and cosmic rays of average energy density in comoving frame $`e_{}=3P`$. We approximate the energy by the following expression, assuming an appropriate self-similar evolution of the shell (which allows us to properly define the bulk Lorentz factor):
$$E=M\mathrm{\Gamma }c^2+4P\mathrm{\Gamma }V_0,$$
(2)
where $`V_0=4\pi R^2\delta R_0`$ is the covolume of the shell, $`R(t)`$ the shell radius, $`\delta R_0`$ its comoving thickness. During the expansion, the energy in the shell is redistributed between kinetic and internal energies through the competition of adiabatic cooling and heating by the flux of incoming matter. Part of the mass flux increases the low pressure baryonic mass: $`\dot{M}=(1\alpha )\rho _b\mathrm{\Gamma }c4\pi R^2`$; the other part contributes to the high energy component: $`\dot{e}_{}_{heating}=\alpha \rho _b\mathrm{\Gamma }c^34\pi R^2/V_0`$, assuming a heating length shorter than the fireball width. The adiabatic cooling is such that $`\dot{e}_{}_{adiab}=\frac{4}{3}\frac{\dot{V}_0}{V_0}e_{}`$. Thus setting $`\dot{E}=0`$, after some simple manipulations, we obtain (with $`\alpha _01+\alpha /3`$):
$$\alpha _0\rho _b4\pi R^2c^3\mathrm{\Gamma }^2+\frac{1}{3}\frac{\dot{V}_0}{V_0}M\mathrm{\Gamma }c^2+(\frac{\dot{\mathrm{\Gamma }}}{\mathrm{\Gamma }}\frac{1}{3}\frac{\dot{V}_0}{V_0})E=0.$$
(3)
It can easily be checked that, when there is no mass input, an adiabatic expansion occurs such that $`\mathrm{\Gamma }=\mathrm{\Gamma }_0(V_0(t)/V_0(t_0))^{1/3}`$ as long as $`M\mathrm{\Gamma }c^2E`$, followed by a free expansion $`\mathrm{\Gamma }E/Mc^2`$ when the internal energy has become smaller than the kinetic energy. Now to analyze the effect of the mass sweeping, we assume a self-similar expansion such that $`d\mathrm{ln}V_0/d\mathrm{ln}t=\chi `$ and define the dimensionless quantity $`\eta (t)`$ (within a coefficient of order unity) as the ratio of the fireball energy over the swept energy-mass of the ambient medium, measured in the comoving frame, such that
$$E=\eta (t)\rho _b\frac{4}{3}\pi R^3\mathrm{\Gamma }^2c^2.$$
(4)
Then equation (3) leads to the differential equation:
$$\frac{d\eta }{d\mathrm{ln}t}=2\eta (\eta _{\mathrm{}}\eta ),$$
(5)
where $`\eta _{\mathrm{}}`$ is the asymptotic value taken by $`\eta (t)`$, namely $`\eta _{\mathrm{}}=\frac{3}{2}+\frac{\chi }{3}`$. With $`\eta (t_0)=\eta _0`$, the solution is
$$\eta (t)=\eta _{\mathrm{}}\frac{\eta _0(t/t_0)^{2\eta _{\mathrm{}}}}{\eta _{\mathrm{}}\eta _0+\eta _0(t/t_0)^{2\eta _{\mathrm{}}}}.$$
(6)
The evolution starts in the adiabatic regime, $`\eta (t)\eta _0(t/t_0)^{2\eta _{\mathrm{}}}`$ and $`\mathrm{\Gamma }t^{\chi /3}`$. If $`\eta `$ reaches its asymptotic value, then the bulk Lorentz factor decreases according to the canonic law $`\mathrm{\Gamma }t^{3/2}`$. However, the interest of this derivation is to stress that this Blandford & Mc Kee (blan (1976)) law does not mean that a free expansion has been reached and a strong relativistic shock set up. Indeed the energy of the fireball can still be dominated by the internal energy, so that a shock does not form, and the incoming baryons interact with the fireball plasma through a streaming instability. This is probably what happens during the primary stage of baryon loading where a high entropy pair plasma entrains some baryonic mass, expected to be of the order $`10^6M_{}`$, coming from the debris of the progenitor. This interaction is probably more important than the final stage of interaction with the interstellar medium in generating perturbations in the flow as described in subsection 2.3 and section 3.
Thus there are two stages of baryon loading (see Fig. 1), the primary stage when the high entropy pair wind entrains a mass $`M_010^6M_{}`$ coming from the debris of the massive star giving rise to the hypernova or of the merging compact objects, and a secondary stage corresponding to the sweeping of the ambient interstellar medium but that loads typically $`10^3M_0`$ only.
### 2.2 The magnetic field and some scales
We assume that this baryonic plasma flows along a strong poloidal magnetic field. This is the major assumption of this paper: $`B>B_{eq}`$ (we note $`\beta _p(B_{eq}/B)^2<1`$) and $`B>B_m`$ where $`B_{eq}`$ is the equipartition field corresponding to the plasma pressure in the fireball, and $`B_m`$ is the equipartition field corresponding to the ram pressure exerted by the incoming baryons.
$$B_m10^2(\frac{n_b}{1cm^3})^{1/2}\frac{\mathrm{\Gamma }}{10^3}Gauss.$$
(7)
Indeed, the poloidal magnetic pressure follows the same decrease as the relativistic pressure with distance; for conical expansion, $`B_p^2/8\pi PR^4`$. Moreover, a toroidal magnetic field is likely significant in the edge of the collimated flow and can play a confining role. The decay of the toroidal field is weaker, because $`B_tR^1`$ in a conical expansion. If the magnetic field dominates at the beginning, it remains so with the expansion. Moreover, a tiny initial toroidal component of the magnetic field can become dominant in the expansion and favors collimation. Like in extragalactic jets, the poloidal field can be dominant along the axis of the flow and the toroidal field dominant in the edge of the flow. Thus we assume an internal poloidal field such that $`B_p=B_d(R/R_d)^2`$ with $`B_dB_m`$ at $`R_d10^{16}cm`$ and an external toroidal field $`B_tB_d(R/R_d)^1`$; in particular at $`R_{}`$ (typically $`10^3R_d`$), $`B10^8G`$. Of course these estimates are sensitive to the model of magnetic field distribution, however the numbers we proposed are reasonable and help to explain out what happens.
In the early stage the proton energy is severely limited by synchrotron losses. Assuming an acceleration time $`t_{acc}=\kappa t_L`$, where $`t_L`$ is the Larmor time, the maximum proton energy is such that the acceleration time is equal to the synchrotron time and gives
$$ϵ_{max}=2\frac{10^{11}}{\sqrt{\kappa }}\left(\frac{B}{1G}\right)^{1/2}GeV.$$
(8)
Taking $`\kappa 10`$, the energy of $`10^{17}eV`$ can be achieved only in the region where $`B`$ is smaller than $`B_s10^6G`$, which requires $`R>R_s10^2R_d`$. Bearing in mind that $`\mathrm{\Gamma }`$ reaches a value between $`10^2`$$`10^3`$, the UHECRs ($`10^{20}eV`$ in the observer frame) are necessarily produced in the region $`R>R_s`$.
A neutrino radiation is produced during the early stage through pp-collisions thanks to Fermi acceleration. Indeed, $`ϵ_{max}`$ given by (8) exceeds $`1GeV`$ when $`B<10^{12}G`$, therefore where $`R>R_h10^2R_{}`$. This low energy neutrino emission ends when the pp mean free path becomes larger than $`\delta R`$, thus when $`R=R_{pp}10^2R_hR_{}`$. This stage of Fermi acceleration supports the predictions made by Paczynski & Xu (pac (1994)), namely a neutrino emission between $`30GeV`$$`TeV`$ of global energy of a few percent of the fireball energy.
### 2.3 Relaxation of the baryon stream in the fireball
In this subsection, we examine how the backstream of baryons in the fireball, initially composed of $`e^+e^{}`$ pairs, undergoes a relaxation by triggering a beam instability in the pair plasma.
The hydromagnetic perturbations so produced scatter the baryons that are then entrained and some of them accelerate to high energy. We presume that these perturbations are those revealed by the light curve and which also accelerate the electrons responsible for the synchrotron and inverse Compton gamma-emission. Indeed, initial perturbations do not amplify in the expansion, they even decay (this differs from an expanding universe that is self-gravitating). Therefore we stress that these perturbations must be generated during the expansion by the appropriate instability.
At the presumed energy of the particles, the Coulomb interactions are negligible (indeed, for relativistic electrons, the Coulomb cross section is as low as the Thomson cross section; as long as the shell is optically thin to Compton scattering, the mean free path is larger than the shell width; this is of course more obvious when protons are involved in the collision). Therefore the incoming baryons interact only with the magnetic field carried by the fireball. The magnetic field represents a more efficient obstacle if it is perpendicular to the flow, and the interaction starts in the form of an intense backward fast magneto-sonic wave. In a confined relativistic plasma, the fast mode propagates with a Lorentz factor $`\gamma _F=\gamma _S\gamma _{}`$, where the sound Lorentz factor (corresponding to the parallel slow mode actually) $`\gamma _S=\sqrt{3/2}`$ and the Lorentz factor of the generalized Alfvén waves (see Pelletier & Marcowith pelmar (1998)) $`\gamma _{}=(1+\frac{1}{2\beta _p})^{1/2}`$. However, since we consider a collimated expansion, the most natural circumstance is the interaction along the parallel magnetic field within a wide solid angle around the axis.
In this scope, the interaction scenario is the following. In the comoving frame, the proton back-stream, of velocity $`v_b=\beta _bc`$ and Lorentz factor $`\mathrm{\Gamma }_b`$, along the field line generates backward Alfvén waves of velocity $`V_{}=\beta _{}c`$ at a rate $`g`$ given by (Marcowith et al. mar (1997) ; Pelletier & Marcowith pelmar (1998)):
$$g=g_0\frac{\mathrm{\Gamma }}{\mathrm{\Gamma }_b}\sqrt{\beta _b\beta _{}},$$
(9)
where $`g_0\frac{\omega _{pi}}{\mathrm{\Gamma }}(2\beta _{})^{1/2}`$, $`\omega _{pi}`$ being the plasma frequency of the ions in the stream. This maximum growth occurs for the wave number
$$k_0=\frac{\omega _{ci}}{\mathrm{\Gamma }_b(v_bV_{})}\frac{\gamma _{}^2}{\mathrm{\Gamma }_b}\frac{\omega _{ci}}{c}.$$
(10)
When these backward waves have reached a high level, they scatter the baryons that are then entrained by the pair flow. Thus, we presume that the large internal disturbances observed in the light curve result from loading baryons in the fireball either at the early stage of the explosion or during the “free” expansion stage. At the beginning of the fireball expansion in the interstellar medium, the diffusion length is not short enough compared to the narrow width of the shell. Then baryon entrainment starts when the level of the resonant Alfvén waves (hereafter A-waves) is strong enough to get a diffusion length shorter than the shell thickness. Then the kinetic flux of the incoming baryons is transformed into a backward flux of intense A-waves. As long as these waves do not interact with the plasma shell (their wavelengthes are larger than the dissipation scale), the shell does not decelerate and just an electromagnetic wake is generated propagating towards to shell center but more slowly than the ultrarelativistic expansion of the shell so that it seems to advance. Thus for a baryon mass loading rate $`\dot{M}_b=4\pi R^2\rho _b\mathrm{\Gamma }c`$, the energy flux, measured in the comoving frame, generated in the form of backward disturbances in the shell is:
$$S=V_{}W_{}=\frac{\dot{M}_b}{4\pi R^2}\mathrm{\Gamma }c^2.$$
(11)
For a fireball of initial energy $`E_0`$, this baryon loading produces only a weak shell deceleration as long as
$$t<t_d=\frac{3E_0}{\dot{M}_b\mathrm{\Gamma }c^2},$$
(12)
where $`t_d`$ is the deceleration time. In other words, the shell is still weakly perturbed as long as the mass loading rate is smaller than $`\dot{M}_d=\rho _b\mathrm{\Gamma }4\pi R_d^2c`$ where $`R_d=ct_d`$ is the deceleration radius. There are in fact two deceleration radii corresponding to the two stages of baryon loading.
The relaxation of the baryon stream in the fireball is described by the following simplified model coupling the relative level of Alfvén energy density $`u`$ with the relative motion of the stream $`\beta _b\beta _{}`$ with respect to the backward A-waves so-generated. The coordinate $`r`$ denotes the radial distance from this external sheet, increasing $`r`$ means approaching the shell center and the “initial” condition is $`\mathrm{\Gamma }_b(0)=\mathrm{\Gamma }`$.
$`V_{}{\displaystyle \frac{du}{dr}}`$ $`=`$ $`2gu`$ (13)
$`c{\displaystyle \frac{d\mathrm{\Gamma }_b}{dr}}`$ $`=`$ $`\nu \beta \mathrm{\Gamma }_b`$ (14)
where in a quasi-linear approximation the slowing rate $`\nu =\nu _0u`$, the more intense the waves the more efficient the scattering. The energy flux conservation implies that
$$\nu _0=\frac{2g\mathrm{\Gamma }}{u_{\mathrm{}}\beta _{}\beta \mathrm{\Gamma }_b}$$
(15)
where $`u_{\mathrm{}}`$ is the asymptotic value of $`u`$, namely $`W_{}/(B^2/8\pi )`$ given by (11) and thus $`u_{\mathrm{}}=(B_m/B)^2<1`$.
The relaxation of the stream can be characterized by a single typical length $`l_r=V_{}/g_0`$. Fig. 2 illustrates the stream relaxation in the shell that occurs in few $`l_r`$.
Let us consider the interaction with the interstellar medium. The typical wavelength of the excited waves, measured in light seconds ($`l.sec.`$), is such that
$$\lambda _010^3\frac{\mathrm{\Gamma }}{10^3}(\frac{B}{10^3G})^1l.sec.,$$
(16)
whereas the relaxation length is
$$l_r\frac{\mathrm{\Gamma }}{10^3}(\frac{n_b}{1cm^3})^{1/2}l.sec..$$
(17)
This has to be compared with the shell width in the comoving frame; at the deceleration time its estimate is (Rees & Mészáros ree (1992))
$$\begin{array}{c}\delta R_d3\times 10^2(\frac{n_b}{1cm^3})^{1/3}(\frac{\mathrm{\Gamma }}{10^3})^{5/3}\hfill \\ \hfill \times (\frac{E}{10^{51}erg})^{1/3}(\frac{\mathrm{\Omega }}{4\pi })^{2/3}l.sec..\end{array}$$
(18)
Let us consider now the interaction in the “primary” GRB. Expanding $`10^6M_{}`$ within a sphere of one light-second radius (comoving) leads to a particle density of $`10^{19}cm^3`$. The plasma is now collisional, but the collective process previously developed provides a faster relaxation than both Coulomb and pp-collision ones. The typical wavelength $`\lambda _010^{10}\frac{\mathrm{\Gamma }}{10^3}(\frac{B}{10^{10}G})^1l.sec.`$ and the relaxation length $`l_r10^9\frac{\mathrm{\Gamma }}{10^3}(\frac{n_b}{10^{18}cm^3})^{1/2}l.sec.`$. Incidently, we note that although the time is quite short, the density is so high that the Lawson criterium for fusion is largely satisfied; moreover the thermal energy of the particles is high ($`\overline{ϵ}10^{12}eV`$). Thus, alpha particles and neutrons should be produced and moreover, as mentioned previously, a significant neutrino emission.
Once the relaxation is achieved, there are two major nonlinear interactions with the plasma shell through magnetosonic compression. First, intense backward A-waves that propagate almost parallel to the magnetic field produce fast magnetosonic compression governed by the Hada’s system generalized to relativistic plasma by Pelletier & Marcowith (pelmar (1998)). The transverse magnetic perturbation $`b`$ (reduced to the averaged field) exerts a pressure that produces fast parallel perturbed motion $`u_{}`$ (specific momentum) that is proportional to the parallel electric field (see Pelletier pel (1999)). This, parallel electric field is responsible for particle acceleration or stochastic heating. For delocalized waves, particles that resonate with the parallel electric field produce the nonlinear Landau damping of the A-waves. The so-produced $`E_{}`$ efficiently injects electrons and positrons in the high energy population.
There is also another strong nonlinear effect with the slow magnetosonic mode (S-mode) that efficiently backscatters the primary flux of A-waves. This is presented in the next section.
## 3 Brillouin backscattering
As seen in the previous section, the flux of backward A-waves becomes stronger and stronger as long as the shell is still entraining baryons. If these waves are not be absorbed in the plasma shell, no deceleration nor heating of the shell would stem from the relaxation of the incoming baryons. Moreover, Fermi acceleration with these waves can work only if there are both forward and backward waves. One could expect some reflection of these waves in the internal edge of the shell; however, because they propagate at a velocity close to the light velocity, no significant change of impedance and thus no significant reflection occurs at the edge. Both absorption and backscattering efficiently develop by excitation of sound or slow magnetosonic modes.
### 3.1 Parametric Brillouin instability
A large fraction of this flux can be backscattered (thus producing a forward A-flux) by exciting the slow mode of the shell plasma. This is analogous to the Brillouin backscattering, but with A-waves instead of ordinary electromagnetic waves. The analogous Raman scattering does not work with A-waves (three mode couplings does not work with A-waves). A mother A-wave ($`\omega _0,k_0`$) of relative amplitude $`b_0`$ spontaneously gives rise to a slow magnetosonic mode ($`\omega _s,k_s`$), but this S-mode is unable to carry the whole flux of energy-momentum and a backward (forward with respect of the shell expansion) A-wave is generated with a lower energy-momentum ($`\omega _{},k_{}`$) such that:
$$\omega _0=\omega _s+\omega _{}andk_0=k_s+k_{}.$$
(19)
For parallel propagating resonant waves,
$$k_s=\frac{2k_0}{1+V_s/V_{}}andk_{}=k_0\frac{1V_s/V_{}}{1+V_s/V_{}},$$
(20)
one gets the most efficient rate of coherent wave decay:
$$G_{decay}=\frac{1}{2}\sqrt{\omega _s\omega _{}}(\frac{b_0^2}{2\beta _p})^{1/2}.$$
(21)
This result is a particular solution that can be derived from the complete analysis made by Champeaux et al. (cham (1999)). For incoherent scattering, the decay rate is lower, proportional to $`b_0^2/2\beta _p`$ instead of its square root.
The S-modes are in turn absorbed by resonant interaction with “thermal” particles. This process is clearly the most efficient to transmit the momentum and energy of the baryon stream to the plasma shell.
Now intense forward and backward flux of long A-waves are still remaining. There is no other way to absorb them than cosmic ray acceleration through Fermi processes that involve resonant interaction between high energy particles and long A-waves.
### 3.2 A toy model of the backscattering process
In principle, solving the nonlinear system that governs the backscattering process allows us to predict how the dissipated energy is shared between heating the thermal particles of the shell and acceleration of cosmic rays. The detailed theory involves the numbers of quanta for each mode ($`N_+,N_{},N_s`$) for the mother A-waves, the backscattered A-waves and the generated slow waves respectively), a probability rate of scattering $`w`$ and absorption rate ($`\gamma _+,\gamma _{},\gamma _s`$), all these triplets depending of the three wave vectors. In this paper we merely address a very simplified version of the system, similar to the system obtained with random phase approximation, but averaged over the spectral bands. We think it would be useful to solve such a toy system to illustrate the process and to see whether reasonable estimates can be obtained by developing this method. For the sake of simplicity, we assume that backscattering occurs behind the stream relaxation sheet (in fact, it could even start in the relaxation sheet). The system is thus the following (in this notation a backscattered wave propagates outwards):
$`V_{}{\displaystyle \frac{}{r}}N_+`$ $`=`$ $`w(N_+N_s+N_+N_{}N_{}N_s)\gamma _+N_+`$ (22)
$`V_{}{\displaystyle \frac{}{r}}N_{}`$ $`=`$ $`w(N_+N_s+N_+N_{}N_{}N_s)\gamma _{}N_{}`$ (23)
$`V_s{\displaystyle \frac{}{r}}N_s`$ $`=`$ $`w(N_+N_s+N_+N_{}N_{}N_s)\gamma _sN_s`$ (24)
The absorption rate $`\gamma _s`$ of the S-waves by the “thermal” plasma is always larger than the absorption rates $`\gamma _+`$ and $`\gamma _{}`$ of the A-waves by cosmic rays (gyro-synchrotron absorption); moreover the S-waves have shorter wavelengths than the A-waves, the backscattered waves have even larger wavelengthes than their mother wave and thus even less absorbed. There are two scales in the system: a short characteristic scale associated to the growth of the S-waves $`l_s=V_s/(\gamma _sG)`$, where $`G=w(N_+(0)N_{}(0))`$ is the nonlinear scaterring rate, and a longer scale associated with the decay of the mother waves $`l_{}=V_{}(\gamma _sG)/\gamma _sG`$; $`l_sl_{}`$ when $`G\gamma _s`$. This decay length $`l_{}\beta \lambda _0/u`$. A ratio $`T`$ of the incident flux is transmitted, a ratio $`R`$ is reflected, a ratio $`A`$ is absorbed and one has:
$$T+R+A=1;$$
(25)
and the absorption ratio is divided into a thermal one through S-waves, $`A_s`$, and a nonthermal one through A-waves, $`A_{}`$. Thus when $`V_{}/\gamma _+<\delta R`$, $`T0`$. By solving the system, one obtains the backscattering rate: $`R=\frac{\omega _{}N_{}(0)}{\omega _0N_+(0)}`$. One calculates the thermal absorption ratio:
$$A_s=\frac{_0^{\mathrm{}}\gamma _s\omega _sN_s𝑑r}{V_{}\omega _0N_+(0)}.$$
(26)
The absorption rate into cosmic rays is deduced by using (25). Such a calculation is illustrated in Fig. 3. For reasonable absorption rates, we obtain $`R0.17`$, $`A_s0.53`$ and $`A_{}0.30`$. Of course, a more realistic calculation should be made to predict these ratios, but this simple model indicates that, through the backscattering process, a sizable fraction of incoming energy can be converted into cosmic rays. In particular, a more detailed calculation would take into account that the longer wavelength modes are less damped and thus can be transmitted behind the shell.
The combined process of stream relaxation and Brillouin scattering followed by energy absorption by particles occurs first during the “primary” stage of the relativistic expansion when the pair wind entrains ambient baryons coming from the progenitor debris. As argued previously, this interaction is likely to occur without shock formation. Then after a new adiabatic and free expansion, a new stage of interaction occurs with the interstellar medium. The ratio of the energy dissipated by the cosmic rays acceleration over the initial fireball energy is $`E_{cr}/E=A_{}R(t)^3/R_d^3`$ before the deceleration time; thus, at $`R=0.3R_d`$, about $`10\%`$ of the fireball energy in converted into cosmic rays. When the fireball radius reaches the deceleration radius, a strong collisionless shock is set up. However it can be considered as a hardening of the previous process. Indeed, incoming protons are reflected in the relaxation front giving rise of the external shock, whereas the backward magnetosonic wavefronts turn to internal shocks. The previous description applies during the predeceleration stage. However, although the theory of beam-plasma instability is no longer relevant when the shock is set up, the Alfvén fronts are still behind the shock and still accelerate cosmic rays efficiently.
## 4 Fermi acceleration of UHE Cosmic Rays
As stressed by Pelletier (pel (1999)), intense forward and backward A-waves of long wavelength (delocalized or localized) are necessary to accelerate UHE Cosmic Rays. The intense localized fronts propagate with a nonlinearly modified velocity such that $`\gamma _{}\gamma _{nl}=\gamma _{}+\delta \gamma `$ with $`\delta \gamma /\gamma _{}\beta _{}b_m^2`$ ($`b_m`$ being the maximum reduced amplitude of the perturbed magnetic field of the front measured in its comoving frame). These intense wave packets scatter particles of Larmor radius smaller than their width or wavelength in their frame in a few gyro-periods. In a scattering time, the particles energy gains a factor $`\gamma _{}^2`$ and the momenta are concentrated in the front cone of half-angle $`1/\gamma _{}`$. Further acceleration requires other fronts propagating in the opposite direction, because wavefronts propagating in the same direction at almost the same velocity $`V_{}`$ tend to roughly isotropize the distribution of interacting particles with respect to their comoving frame. Thus it is crucial to get both forward and backward fronts to accelerate cosmic rays and the Brillouin backscattering process of section 3, that generates longwavelength perturbations in a time $`G^1`$, is the appropriate and efficient solution. The relativistic regime of Fermi acceleration is deeply different from the nonrelativistic regime, not only because of the anisotropy effect but also because the acceleration time scale can become shorter than the scattering time scale. This is a crucial advantage that is illustrated by the following result (Pelletier pel (1999)) obtained in the case of a permanent flux of both forward and backward waves, the energy gain during a time $`\mathrm{\Delta }t`$, smaller than the scattering time, is such that
$$<\mathrm{\Delta }p^2>=p^2\beta _{}^2\gamma _{}^2\nu _s\mathrm{\Delta }t,$$
(27)
where $`\nu _s`$ is the scattering frequency. Moreover, because $`B>B_{eq}`$, $`\gamma _{}>2`$, and the energy jump at each scattering is large and the statistic evolution cannot be treated by a Fokker-Planck description.
Let us give some more details about the localized fronts. Intense longwaves tend to self-organize in forward and backward relativistic fronts of amplitude $`b_m(r_{}/\xi )^{1/2}`$, where $`\xi `$ is the front width and $`r_{}=(<\gamma ^2>/\overline{\gamma })V_{}/\omega _c`$ is the minimum scale of relativistic MHD (Pelletier & Marcowith pelmar (1998)); it is worth noting that the perturbation amplitudes appear $`\mathrm{\Gamma }^2`$ times larger in the observation frame. At each crossing of these fronts, the gain is by a factor $`\gamma _{}^4`$ within few gyro-periods. Such localized fronts can be approximately described as solitons (Pelletier & Marcowith pelmar (1998); Pelletier pel (1999)), they differ from ideal solitons because of the Fermi acceleration that produces a kind of Landau-synchrotron damping of them. When ideal solitons cross each other they do not destroy, whereas damping destroys them and their complete absorption corresponds to the strongest efficiency of the Fermi process. A similar idea was used by Daigne & Mochkovitch (daim (1998)) with internal shocks to explain the gamma emission. Here we interpret the internal shocks as internal Alfvén fronts which can contribute to short scale variations in the light curve, and possibly more easily than pure hydrodynamic shocks.
For a bulk Lorentz factor of the shell of $`10^3`$, particles in GRBs have to reach $`10^{17}eV`$ in the comoving frame in order to supply the UHE Cosmic Rays population beyond the GZK threshold. We consider the interaction with the interstellar medium only because, as previously showed, synchrotron losses prevent UHE cosmic ray production, inside a radius $`R_s`$ (subsection 2.2). Assuming that the largest wavelengthes are one tenth of the effective shell width (limited by causal connection), a magnetic field just larger than the equipartition value at deceleration radius ($`B_{eq}10^3G`$ say) is enough to get Larmor radii of that size (i.e. $`10^5pc`$), which corresponds to an energy of order $`10^{17}eV`$. With $`\overline{\gamma }10^8`$, the scale $`r_{}10l.sec.`$. Therefore, few fronts having a width of $`110r_{}`$ that cross each other within a shell width of $`3\times 10^2l.sec.`$ can produce protons of $`10^{21}eV`$ within one second with respect of an observer frame during the interaction with the interstellar medium. When the strong relativistic shock has formed ($`R>R_d`$), external particles gain energy by a first Fermi half-cycle by a factor $`\mathrm{\Gamma }_s^2`$. They are then injected in the relativistic A-fronts behind the shock and suffer further acceleration up to the ultimate energy.
The local distribution is more likely quasi monoenergetic rather than a powerlaw. However the local characteristic energy is proportional to the product $`BR`$ and thus the global distribution reflects the distribution of the product $`BR`$; as stated by Pelletier (pel (1999)), a distribution close to $`ϵ^2`$ is obtained for $`BR^m`$ with $`m`$ close to $`2`$.
Our main conclusion is that Gamma-Ray Bursts are capable of producing the UHE Cosmic Rays through a Fermi process with relativistic Alfvén waves. These intense waves are generated by the two stages of baryon entrainment, first during the primary GRB stage, second during the interaction with the interstellar medium. The Fermi process works because of the Brillouin backscattering process that turns out to be appropriate and efficient, moreover it allows the heating and deceleration of the shell plasma by the incoming flux. Although the primary stage does not produce UHE cosmic rays because of synchrotron losses, the Fermi acceleration allows the maintenance of a high energy proton population between $`R_h`$ and $`R_{pp}`$ that can emit a significant low energy neutrino flux of the order $`10^2E`$. A high energy neutrino flux can be generated through the p$`\gamma `$-process by UHE cosmic rays at the end of the free expansion, typically a sizable fraction of the UHE cosmic ray energy in the form of $`10^{14}`$$`10^{16}eV`$ neutrinos (Rachen & Mészáros rac (1998)).
###### Acknowledgements.
The authors are grateful to Y. Gallant, G. Henri, A. Marcowith, R. Mochkovitch and J. Rachen for fruitful discussions.
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# 1. Introduction.
## 1. Introduction.
S. Donaldson in \[Do96\] has adapted the concept of very ample bundle to the symplectic setting. Following this idea some results of complex projective geometry has been generalized to symplectic geometry, as Bertini’s theorem and Lefschetz’s hyperplane theorem in \[Do96\], conectedness of the space of “good” sections of a very ample bundle in \[Au97\], divisors on projective fibrations in \[Pa98\], special position theorems in \[Pa99\], existence of Lefschetz pencils in \[Do99\], existence of branched coverings of symplectic 4-manifolds over $`^2`$ and associated invariants \[Au99, Au99b\], and Kodaira’s embeddings theorems and symplectic determinantal submanifolds in \[MPS99\]. These ideas have opened a new insight in symplectic geometry allowing to understand symplectic manifolds through the study of the linear systems associated to a “very ample” vector bundle.
In \[IMP99\] the idea of \[Do96\] of working with approximately $`J`$-holomorphic sections was partially traslated to the contact case. With these sections and a generalization of the local estimated transversality result of \[Do96\], which is the key of the symplectic approach, the ideas of \[Do96, Au97\] were translated word by word to the contact setting. The only important loss was the isotopy results which have been developed in the symplectic theory starting with the ideas of D. Auroux in \[Au97\]. In this paper we show how to develop a contact geometry of linear systems analogous to the symplectic case. We also show how to recover partially the isotopy in the theory.
We will prove a theorem analogue to that of \[Do99\]. In fact, we will show the existence of a certain class of pencils on a contact manifold. The main tool in the proof will be a generalization of the local transversality theorem proved in \[Do99\]. As in the symplectic case this result is not strictly necessary for the proof, but it has interest in its own. It could be used to simplify the constructions of other linear systems, as for instance, the one in \[MPS99\] in the symplectic case.
A compatible chart in a contact manifold $`(C,D)`$ at a point $`x`$ will be a chart $`\varphi :U_xC^n\times `$, where $`U_x`$ is neighborhood of $`x`$, verifying $`\varphi (x)=(0,0)`$, $`(\varphi _{})(D(x))=^n\times \{0\}`$ and moreover verifying that the presymplectic form $`(\varphi _{})d\theta (x)`$, when restricted to $`^n\times \{0\}`$, is a positive form of type $`(1,1)`$ at the origin of coordinates. If the contact manifold is exact we will impose also that $`\varphi _{}(R)(x)),\frac{}{s}>0`$ and say that the chart is oriented compatible, where $`R`$ is the Reeb vector field and $`s`$ is the real coordinate.
###### Definition 1.1.
A (oriented) contact pencil on a closed (exact) contact manifold $`C`$ consists of the following data:
1. a codimension $`4`$ contact submanifold $`AC`$,
2. a finite set of smooth contact curves $`\mathrm{\Delta }=_{iI}\gamma _iCA`$,
3. a smooth map $`f:VA^1`$, whose restriction to the complementary of the set $`\mathrm{\Delta }`$ is a submersion, satisfying also that $`P_\mathrm{\Delta }=f(\mathrm{\Delta })`$ is a set of locally smooth curves with transversal selfintersections.
Also the data have to admit the following standard local models:
* At any point $`aA`$, there are (oriented) compatible coordinates $`(z_1,\mathrm{},z_n,s)^n\times `$ such that $`A`$ is locally given by $`\{z_1,z_2=0\}`$. And the function $`f`$ has the expression $`f(z_1,\mathrm{},z_n,s)=\frac{z_1}{z_2}^1`$ near $`a`$.
* At a point $`b_i\gamma _i`$ there are (oriented) compatible coordinates in which $`f`$ is written as $`f(b_i)+\phi (s)+z_1^2+\mathrm{}+z_n^2`$, where $`\phi :`$ verifies $`\phi (0)=0`$ and $`\phi ^{}(0)0`$.
It is clear from the local model that the counterimage $`f^1(p)`$ is a subset of $`CA`$, whose closure in $`C`$ is smooth at $`f^1(p)A`$. Abusing language, we will call fiber over $`b`$ to the closure of the counterimage. It is a smooth submanifold if $`b`$ is a regular value of $`f`$. In other case we will have one or two singularities locally modelled by:
(1) $`\phi (s)+z_1^2+\mathrm{}+z_n^2.`$
In case $`\text{dim}C=3`$, the smooth fibers will be (oriented) links on $`C`$. The link operation that is performed when the image crosses a circle $`\gamma _i`$ of the sphere $`^1`$ looks, after general projection to a plane, as
The main result of this article is
###### Theorem 1.2.
Given a contact closed manifold $`(C,D)`$ (resp. exact) and $`\alpha H_{2n1}(C,)`$ which is reduction of an integer class, there exists a contact pencil on $`C`$ (resp. oriented) whose fibers are contact submanifolds, homologous to $`\alpha `$.
With this result at hand, it is easy to understand the possibility of finding general isotopic constructions for the contact submanifolds constructed in \[IMP99\]. Once fixed a compatible complex structure in the distribution, in Section 6 we will show how to define a sequence of contact fibrations $`f_k`$, satisfying $`f_k(0)N_k`$ and $`f_k(\mathrm{})N_k^{}`$, where $`N_k`$ and $`N_k^{}`$ are sequences of codimension $`2`$ contact submanifolds constructed with the method developed in \[IMP99\]. In order to assure the isotopy between $`N_k`$ and $`N_k^{}`$ we have only to construct a path between $`0`$ and $`\mathrm{}`$ in $`^1`$ which does not intersect $`\mathrm{\Delta }`$, because in this case $`f`$ restricted to the path is surjective. But this is only possible if $`0`$ and $`\mathrm{}`$ are in the same connected component of $`^1\mathrm{\Delta }`$, which is not true in general (contrarely to the symplectic case where $`\mathrm{\Delta }`$ consists of isolated points). We will study in Subsection 6.1 the topological relationship between the counterimages of the points of a path crossing $`\mathrm{\Delta }`$. This will prove that
$$H_i(N_k)=H_i(N_k^{}),$$
$$\pi _i(N_k)=\pi _i(N_k^{}),$$
for $`k=0,\mathrm{},n2`$. Moreover we will show how to make the result independent of the chosen complex structure and of the contact form, thus providing a comprobation of the “contact Lefschetz hyperplane theorem” proved in \[IMP99\]. The result brings us back to the Lefschetz original ideas to prove the hyperplane theorem (see \[Le24\]). It would be interesting to complete the proof in Lefschetz’s way, since this would be a more geometrical proof than the one contained in \[IMP99\], the latter being an adaptation of the modern Morse-theoretic argument due to Bott, Andreotti and Frankel (and in the symplectic case to Donaldson).
In Section 2 we will give the basic results in contact geometry needed to develop the proof. In Section 3 we will give the proof of the exact case assuming transversality results of approximately holomorphic contact geometry. Afterwards, in Sections 4 and 5, we develop the local approximately holomorphic techniques needed to achieve this transversality. In Section 6 we will study the relationship between approximately holomorphic sections transverse to 0. Finally, we will adapt all the study to the non-exact case in Section 7.
Acknowledments:
I am very grateful to S. Donaldson by his kindness passing me the preprint \[Do99\] which is absolutely fundamental to develop this work and also by his useful observations about some of the parts of this paper. Also I want to mention D. Auroux and R. Paoletti for passing me their useful preprints. Finally we want to thank to the members of the GESTA <sup>1</sup><sup>1</sup>1Geometría simpléctica con técnicas algebraicas. seminar in Madrid their support and interest through the elaboration of this work. I want to mention especially Vicente Muñoz by his useful commentaries to the previous versions of this paper.
## 2. Definitions and results.
We will assume along the proofs that $`(C,D)`$ is an exact contact manifold, where we have fixed a global contact form $`\theta `$. In Section 7 we will precise the changes needed to extend the results to the non-exact setting. We start recalling with the basic definitions and results of \[IMP99\].
### 2.1. Basic concepts
A complex structure on $`C`$ is a complex structure $`J`$ defined on $`D`$, interpreted as a symplectic vector bundle, compatible with the symplectic form $`d\theta `$. Recall that the distribution $`D`$ is contact, if and only if the restriction of $`d\theta `$ to $`D`$ is symplectic. The contact $`1`$-form defines a vector field $`R`$ by the condition
$$i_R\theta =1,i_Rd\theta =0,$$
It is called the Reeb vector field. Given a complex structure we obtain a metric for the manifold
$$g_J(u,v)=d\theta (u,Jv)+\theta \theta .$$
This metric is called a contact metric. It depends on the fixed contact form. The $`k`$-rescaled contact metric will be defined as $`g_{J,k}=kg_J`$. If there is not risk of confusion we will denote simply $`g_k`$, supposing $`J`$ fixed along the proofs. (Observe that the $`k`$-rescaled contact metric $`g_k`$ is not the contact metric associated to $`k^\alpha \theta `$, for $`\alpha `$)
Now we give some definitions to control contact structures in $`^{2n+1}`$.
###### Definition 2.1.
The maximum angle between two subspaces $`U,V\text{Gr}_{}(r,n)`$ is defined as:
$$\mathrm{}_M(U,V)=\underset{uU}{\mathrm{max}}\mathrm{}(u,V).$$
This angle defines a distance in the topological space $`\text{Gr}_{}(r,n)`$ (for details see \[MPS99\]).
###### Definition 2.2.
Let $`\theta _k`$ be a sequence of contact forms in $`^{2n+1}`$, with associated distributions $`D_k=\text{Ker}\theta _k`$. The sequence is called $`c`$-asymptotically flat in the set $`U^{2n+1}`$ if
$$\mathrm{}_M(D_k(0),D_k(x))ck^{1/2},\text{for all }xU.$$
The sequence is called asymptotically flat if there exists some $`c>0`$ for which it is $`c`$-asymptotically flat.
The standard contact structure in $`^{2n+1}`$ is defined as $`\theta _0=ds+_{j=1}^nx_jdy_j`$, where $`(x_j,y_j,s)^{2n+1}`$. The $`k`$-rescaled contact metric is the contact metric associated to $`\theta _{k^{1/2}}=k^{1/2}ds+x_jdy_j`$, which is obtained from $`\theta _0`$ scaling the coordinates by a factor $`k^{1/2}`$. It is obvious from the definition that $`\theta _{k^{1/2}}`$ is a sequence of asymptotically flat contact forms on any bounded set of $`^{2n+1}`$.
If a contact distribution $`D`$ is defined at the origin by the horizontal subspace $`^{2n}\times \{0\}^{2n+1}`$, we can define at a neighborhood of this point a canonical complex structure in $`D`$ by means of vertical projection of the canonical one defined on $`^{2n}`$. This complex structure will be denoted $`J_0`$. In the case of the standard contact structure we can extend $`J_0`$ over all $`^{2n+1}`$.
Following \[IMP99\] we can define $``$ and $`\overline{}`$ operators in any kind of function, morphism or section defined on a contact manifold, restricting ourselves to the contact distribution, or equivalently projecting along the Reeb direction.
### 2.2. Approximately holomorphic geometry
We will talk about uniform constant, polynomial, when this constant, polynomial, etc. does not depend on the chosen point $`xC`$, nor in the integer $`k`$ appearing in the context. (However these constants can depend on the modulus of the given sections, on the nature of its derivatives, on the step of a recurrent reasoning, etc. but always independently of $`k`$).
We will use from \[IMP99\] the following
###### Lemma 2.3.
Given $`(C,\theta )`$ a closed contact manifold and $`J`$ a compatible complex structure on $`\text{Ker}\theta `$, there exists a uniform constant $`c>0`$ and a contact Darboux chart $`\psi :(C,\theta )(^{2n+1},\theta _0)`$ satisfying
$$\frac{1}{2}g(v,w)(\psi _{})_xv,(\psi _{})_xw2g(v,w),xB_g(x,c),v,wT_xC.$$
This implies that $`|^r\psi |=O(1)`$ and $`|^r\psi ^1|=O(1)`$, for $`r=1,2,3`$. Moreover $`|\overline{}\psi (y)|c^{}d(x,y)`$, for a uniform constant $`c^{}`$.
A $`k^{1/2}`$-Darboux chart is a chart $`\varphi _k:B(x,ϵ)^{2n+1}`$ such that $`(\varphi _k)_{}\theta =\theta _{k^{1/2}}`$. The following result is a direct corollary of Lemma 2.3:
###### Corollary 2.4.
Given $`(C,\theta )`$ a closed contact manifold and $`J`$ a compatible complex structure on $`\text{Ker}\theta `$, there exists a uniform constant $`c>0`$ such that there exists a contact chart $`\psi :(C,\theta )(^{2n+1},\theta _{k^{1/2}})`$ satisfying
$$\frac{1}{2}d_{g_k}(v,w)(\psi _{})_xv,(\psi _{})_xw2d_{g_k}(v,w),xB_{g_k}(x,c),v,wT_xC.$$
This implies that $`|^r\psi |=O(1)`$ and $`|^r\psi ^1|=O(1)`$, for $`r=1,2,3`$. Moreover $`|^r\overline{}\psi (y)|c^{}k^{1/2}`$, for $`r=0,1,2`$, for a uniform constant $`c^{}`$.
These results allows to trivialize locally contact manifolds in an approximately holomorphic way. The analogous notion of the $`c`$-bounds of \[Do99\] in the contact case is the following
###### Definition 2.5.
A sequence of sections $`s_k`$ of hermitian bundles $`E_k`$ over the contact manifold $`(C,\theta )`$ has mixed $`C^r`$-bounds $`(c_D,c_R)`$ at the point $`xC`$ if it satisfies
$`|s_k(x)|<c_D,`$
$`|_D^js_k(x)|<c_D,j=1,\mathrm{},r.`$
$`|^js_k(x)|<c_R,j=1,\mathrm{},r.`$
$`|^j\overline{}s_k(x)|<c_Rk^{1/2},j=0,\mathrm{},r1.`$
The sequence has uniform mixed $`C^r`$-bounds $`(c_D,c_R)`$ if it satisfies these bounds at every point.
There we denote by $`_D`$ the restriction of the operator $``$ to the subspace $`D`$. The metric used in the manifold $`C`$ to measure the norms in the precedent definition is the rescaled contact metric $`g_k`$. Recall that if $`s_k^j`$ has $`(c_D^j,c_R^j)`$ mixed $`C^r`$-bounds ($`j=1,2`$), then $`s_k^1+s_k^2`$ has $`(c_D^1+c_D^2,c_R^1+c_R^2)`$ mixed $`C^r`$-bounds.
The following results are used to trivialize bundles over contact manifolds.
###### Definition 2.6.
A sequence of sections $`s_k`$ of bundles $`E_k`$ has mixed Gaussian decay in $`C^r`$-norm away from a point $`xC`$ if there exist a uniform polynomial $`P`$ and uniform constants $`\lambda >0`$, $`c_a`$ such that for all $`yX`$, the sequence $`s_k(y)`$ has mixed $`C^r`$-bounds
$$(P(d_k(x,y))\mathrm{exp}(\lambda d_k(x,y)^2),c_aP(d_k(x,y))\mathrm{exp}(\lambda d_k(x,y)^2)).$$
Recall that the uniform constant $`c_a`$ is necessary because we will check afterwards that for different sequences of sections this constant cannot be fixed.
The prequantizable line bundle $`L`$ over an exact contact manifold $`C`$ is defined as the complex line bundle with connection such that $`curv(L)=id\theta `$.
###### Lemma 2.7 (Lemma 5 from \[IMP99\]).
Let $`(C,\theta )`$ be a closed contact manifold. There exists a uniform constant $`c_s>0`$, such that given any point $`xC`$, there exists a sequence of sections $`\sigma _{k,x}`$ of $`L^k`$ satisfying $`|\sigma _{k,x}|c_s`$ at every $`y`$ in a ball of $`g_k`$-radius $`10`$ centered at $`x`$ and the sections $`\sigma _{k,x}`$ have uniform mixed Gaussian decay away from $`x`$ in $`C^3`$-norm (in this case $`c_a=1`$).
### 2.3. Transversality results.
Following Donaldson \[Do99\] and Auroux \[Au99\] we set up the following definitions. A linear application $`f:^n^r`$ is $`\eta `$-transverse if it has a right inverse $`\nu :^r^n`$ such that $`|\nu |\eta ^1`$. In the non-linear case we will say that $`f:U^n^r`$ is $`\eta `$-tranverse to $`y^r`$ over $`U`$ if $`xU`$, such that $`|f(x)y|<\eta `$, then $`df`$ is $`\eta `$-transverse. Recall that this is an open condition. In fact, if $`|fg|_{C^1,U}<ϵ/10`$ and $`f`$ is $`ϵ`$-transverse to $`y`$ over $`U`$ then $`g`$ is, say, $`ϵ/2`$-transverse to $`y`$ in $`U`$.
The definition of transversality to 0 for sections of hermitian bundles over riemannian manifolds is totally analogous. In the case of contact manifolds we have to strengthen the conditions.
###### Definition 2.8.
A section $`s_k`$ of the hermitian vector bundle $`E_k`$ over the contact manifold $`(C,\theta )`$ is $`\eta `$-transverse to 0 on $`UC`$ if for all $`xU`$ such that $`|s_k(x)|<\eta `$ then $`_Ds_k(x)`$ is $`\eta `$-transverse to 0 (with respect to the $`g_k`$-metric in $`C`$).
From the discussion of \[IMP99\] it follows that a sequence of sections $`s_k`$ of the bundles $`E_k`$ over the contact manifold $`(C,\theta )`$, which has uniform $`(c_D,c_R)`$ bounds and which is $`\eta `$-transverse to 0, has as zero set a contact submanifold, for $`k`$ large enough. The precise result of \[IMP99\] is
###### Theorem 2.9.
Given a closed exact contact manifold $`(C,\theta )`$. Let $`ϵ>0`$ and let $`s_k`$ be a sequence of sections of the bundles $`EL^k`$, for a fixed hermitian bundle $`E`$, with uniform mixed $`C^r`$-bounds $`(c_D,c_R)`$. Then there exists a real number $`\eta >0`$ (depending on $`ϵ`$, $`c_D`$ and $`c_R`$), and a sequence $`\sigma _k`$ such that:
1. $`\sigma _ks_k`$ has mixed $`C^r`$-bounds $`(ϵ,c_R^{})`$.
2. $`\sigma _k`$ is $`\eta `$-transverse to 0.
In \[IMP99\] the proof is developed for sequences with mixed $`C^2`$-bounds, but there is not any problem in generalizing it to the mixed $`C^r`$-bounds case. We need in this article mixed $`C^3`$ bounds. Given a section $`s_k=(s_k^0,s_k^1)`$ of the bundle $`^2SL^k`$, for a fixed hermitian line bundle $`S`$, whose zero set is $`Z(s_k)`$, we denote $`F^{s_k}=\frac{s_k^1}{s_k^0}:CZ(s_k)^1`$ the projectivization of the section. The holomorphic part of the differential of this application will be denoted by $`F^{s_k}`$. Now we state a generalization of Theorem 2.9 which will be proved in Section 4:
###### Theorem 2.10.
Given a closed exact contact manifold $`(C,\theta )`$. Let $`ϵ>0`$ and let $`s_k=(s_k^0,s_k^1)`$ be a sequence of sections of the bundles $`^2SL^k`$ with mixed $`C^3`$-bounds $`(c_D,c_R)`$. Suppose that $`s_k^0`$ and $`s_k`$ are both transverse to 0. Then there exists $`\eta >0`$ (depending on $`ϵ`$, $`c_D`$ and $`c_R`$), and a sequence $`\sigma _k`$ satisfying
1. $`\sigma _ks_k`$ has mixed $`C^3`$-bounds $`(ϵ,c_R^{})`$,
2. $`\sigma _k^0`$ and $`\sigma _k^0\sigma _k^1`$ are $`\eta `$-transverse to 0,
3. $`F_k^\sigma `$ is $`\eta `$-transverse to 0 away from $`Z(\sigma _k^0)`$.
The techniques used in Section 4 improve slightly the ones of \[IMP99\] and thus could allow a simpler proof of Theorem 2.9 avoiding some of the complications of the globalization process in that article.
## 3. Proof of the main result.
Take a complex line bundle $`S`$ with connection $``$ satisfying that $`curv()=PD(\alpha )`$. This is possible since $`\alpha `$ is an integer class. Starting with any sequence of sections $`s_k^{}`$ of $`^2SL^k`$ with mixed $`C^3`$-bounds, we can perturb it using Theorems 2.9 and 2.10 to achieve a sequence of sections $`s_k`$ verifying properties 2 and 3 of Theorem 2.10 and with mixed $`C^3`$-bounds $`(c_D,c_R)`$. We consider this sequence as starting datum and will use it to construct the oriented contact pencil.
From the $`\eta `$-transversality of $`s_k`$ the zero set $`A=Z(s_k)`$ is a codimension $`4`$ contact manifold where $`F^{s_k}`$ is not well defined. We will write $`F`$ instead of $`F^{s_k}`$ whenever it causes no confusion. Now we will study, as in \[Do99\], the shape of the “bad set” $`\mathrm{\Gamma }=\{xC:|F||\overline{}F|\}`$. From \[IMP99\] we know that if we prove that $`\mathrm{\Gamma }=\{xC:d_DF=0\}`$ we will have obtained that the fibres of $`F`$ are contact at all smooth points. This will be proved in several steps.
###### Lemma 3.1.
There is a constant $`\xi >0`$, depending only on $`c_D`$, $`c_R`$ and $`ϵ`$, such that if $`k`$ is large enough then $`|s_0|\xi `$ on $`\mathrm{\Gamma }`$.
Proof: The proof is analogous to that of Lemma 7 in \[Do99\]. $`\mathrm{}`$
Define $`\mathrm{\Delta }`$ as the set of critical points of $`F`$. The connected components of $`\mathrm{\Delta }`$ form a discrete set of smooth curves by the transversality condition imposed to $`F`$. Also we can assure that this set of components is finite because it is contained in $`\mathrm{\Gamma }`$, which by Lemma 3.1 is contained in the complementary of a $`\gamma `$-neighborhood of $`W_{\mathrm{}}=Z(s_k^0)`$, for $`\gamma >0`$ a uniform constant small enough. We define:
$$\mathrm{\Omega }_\xi =\{pC,|s_0(p)|>\xi /2\}.$$
The following step, adapting again Donaldson’s argument, is to estimate the shape of the set $`\mathrm{\Delta }`$. This is the content of the following:
###### Proposition 3.2.
There is a uniform constant $`\rho _0>0`$, such that the $`\rho _0`$-neighborhoods of each connected component $`\gamma _i`$ of $`\mathrm{\Delta }`$ are disjoint and are contained in $`\mathrm{\Omega }_\xi `$. Moreover for any $`\rho <\rho _0`$, for $`k`$ large enough (only depending on $`\rho `$), the set $`\mathrm{\Gamma }`$ is contained in a $`\rho `$-neighborhood of the set $`\mathrm{\Delta }`$.
The proof of this result is absolutely analogous to the proof of Proposition 9 in \[Do99\] and depends strongly on Lemma 8 in that paper. We refer the reader to \[Do99\] for the argument.
Finally, we need to perturb the sequence $`F^{s_k}`$ in arbitrarily small neighborhoods of $`A`$ and $`\mathrm{\Delta }`$ to achieve the local models required in Definition 1.1. The perturbation required in $`\mathrm{\Delta }`$ needs a careful analysis, but again the situation in $`A`$ is a straightforward generalization of \[Do99\]. We need only to define
$$LD_x=s_0^ks_1^k:TC_xL_x^kL_x^k.$$
With this notation the result we need in our case is
###### Lemma 3.3.
For a point $`xA`$, $`F`$ can be represented in the standard model of Definition 1.1 at $`x`$ if and only if $`TA_xD`$ is a symplectic subspace and the restriction of $`d\theta `$ to the symplectic orthogonal $`CA_x`$ (in $`D`$) of $`TA_xD`$ is a positive form of type $`(1,1)`$ with respect to the complex structure on $`CA_x`$ induced by $`LD_x`$.
We do not provide a proof of this Lemma, which follows word by word the proof of Lemma 11 in \[Do99\]. Thanks to the $`(c_D^{},c_R^{})`$ mixed $`C^3`$-bounds and the transversality of the sequence $`s_k`$, it is easy to check that a small $`C^3`$-perturbation of the sequence satisfies the hypothesis of Lemma 3.3, thus completing the study in the neighborhood of $`A`$.
Now we study the map $`F`$ near $`\mathrm{\Delta }`$. Again Donaldson’s ideas work in this case, however the adaptation of the proof needs some changes. Select a smooth connected curve $`\gamma _i`$ in $`\mathrm{\Delta }`$. We are going to perturb $`F`$ in a $`\gamma `$-neighborhood of $`\mathrm{\Delta }`$. By Lemma 3.2, the perturbations can be made in each connected component $`\gamma _i`$ in an independent way.
Recall that, for $`k`$ large enough, the curve $`\gamma _i`$ is contact, i.e. $`x\gamma _i,T_xC=T_x\gamma _iD_x`$. Moreover, the angle between $`T_x\gamma _i`$ and $`D_x`$ is bounded below by a uniform constant because of the transversality of the sequence. Using the contact metric $`g_J`$ associated to the fixed complex structure $`J`$ to define a geodesic flow, we can obtain a diffeomorphism:
$$\varphi _i:U_\rho V_\rho S^1\times ^n,$$
where $`U_\rho `$ is the $`\rho `$-neighborhood of $`\gamma _i`$ (in $`g_k`$ metric) and $`V_\rho `$ is its image by the flow, which is an open neighborhood of $`S^1\times \{0\}`$. We can construct a metric in $`S^1`$ by imposing the condition that $`(\varphi _i)_{|\gamma _i}`$ is an isometry with respect to the rescaled contact metric $`g_k`$. In $`^n`$ we will fix the standard metric. The product metric will be denoted by $`g_0^k`$. We can select with this choice a uniform $`\rho >0`$ such that
(2)
$$\lambda _mg_k(d\varphi _i(v),d\varphi _i(w))|g_0^k(v,w)|\lambda _Mg_k(d\varphi _i(v),d\varphi _i(w)),v,wT_xC,xU_\rho ,$$
where $`\lambda _m,\lambda _M>0`$ are uniform constants. Once we have fixed $`\varphi _i`$, we obtain a distribution $`D_k`$ in $`\varphi _i(U_\rho )`$ constructed as the image of the distribution $`D`$. Denote by $`D_h`$ the integrable distribution given as $`\{p\}\times ^n`$ defined in $`S^1\times ^n`$. Perhaps after shrinking $`\rho `$ uniformly, we can check that
(3)
$$\mathrm{}_M(D_k(s,z),D_h(s,z))<c_u|z|k^{1/2},(s,z)V_\rho S^1\times ^n,$$
where $`c_u>0`$ is a uniform constant. Moreover we can impose without loss of generality that $`(\varphi _i)_{}J(\gamma _i)=J_0`$. So we can project orthogonally $`(\varphi _i)_{}J`$ (defined on $`D_k`$) to $`D_h`$ obtaining a new almost complex structure $`\widehat{J}`$. In fact it is easy to check, as in \[Do96\], that
(4)
$$=_0+\overline{\mu }\overline{}_0,\overline{}=\overline{}_0+\mu _0,$$
where $``$ and $`_0`$ are the operators defined by the structures $`\widehat{J}`$ and $`J_0`$ in $`S^1\times ^n`$, and $`|\mu (z)|c|z|k^{1/2}`$, where $`c`$ is a uniform constant. In order to finish the proof we follow these steps: we will define the perturbation, afterwards we will prove it satisfies the conditions for the distribution $`D_h`$ with the almost complex structure $`\widehat{J}`$, and finally we will check the result for the distribution $`D_k`$.
Given any differentiable function $`f:C`$, we denote $`f_0=f\varphi _i^1`$. By the inequalities (2) we can use $`F_0`$ instead of $`F`$ for all the computations using the induced distribution $`D_k`$. To construct the perturbation we define the complex Hessian $`H=\frac{1}{2}F`$. Using the trivialization $`\varphi _i`$ we may regard it as
$$H_0(s,z)=H_{\alpha \beta }(s)z_\alpha z_\beta ,$$
on $`\varphi _i(U_\rho )`$. Also we take a cut-off function $`\beta _\rho :S^1\times ^n[0,1]`$ satisfying
1. $`\beta _\rho (\varphi _i(p))=1`$, if $`d_k(p,\gamma _i)\frac{\rho }{2}`$.
2. $`\beta _\rho (\varphi _i(p))=0`$, if $`d_k(p,\gamma _i)\rho `$.
3. $`|\beta _\rho |=O(\rho ^1)`$.
We can adjust $`\beta `$ to assure condition 3 because of equation (2). The constant $`\rho <\rho _0`$ will be fixed along the proof to assure that the conditions are satisfied (namely we will have to shrink $`\rho `$ in a uniform way). A modification of $`F`$ will be
$$f_0(s,z)=\beta _\rho (w^{}(s)+H_0(s,z))+(1\beta _\rho )F_0(s,z),$$
where $`w^{}:S^1`$ is any smooth function.
We denote by $``$ and $`\overline{}`$ the operators associated to the almost complex structure $`\widehat{J}`$ acting on $`D_h`$. Respectively we denote by $`_k`$ and $`\overline{}_k`$ the operators associated to the action of $`(\varphi _i)_{}J`$ in $`D_k`$. The equivalent to Lemma 10 of \[Do99\] is
###### Lemma 3.4.
If $`\rho >0`$ is small enough, $`k`$ is sufficiently large, $`|w^{}(s)F_0(s,0)|`$ is sufficiently small and $`|\frac{dw^{}(s)}{ds}|=O(1)`$, then the inequality $`|_kf_0||\overline{}_kf_0|`$ is only satisfied in $`\gamma _i`$.
Proof: First, assume that we are at a point where $`\beta _\rho =1`$. Then $`f_0=w^{}+H_0`$ and
$$f_0=H_0,\overline{}f_0=\overline{}H_0.$$
The $`\eta `$-transversality of $`F_0`$ yields the bound
$$|H_0(s,z)|\eta |z||\overline{}(F_0)_{(z=0)}||z|.$$
Now recall that $`\overline{}+\overline{}=0`$ on functions and that the norm of $`\overline{}F_0`$ is controled because $`F`$ has uniform mixed $`C^3`$-bounds. Then we can write
$$|H_0(s,z)|\eta |z|c_uk^{1/2}|z|.$$
Using the inequality (3) and that $`|\frac{}{s}H_{\alpha \beta }(s)|=O(1)`$ (as follows from the mixed $`C^3`$-bounds of $`F`$) and recalling the bound $`|\frac{dw^{}}{ds}|=O(1)`$, we rewrite the inequality as
(5)
$$|_kH_0(s,z)|\eta |z|c_u^{}k^{1/2}|z|,$$
where $`c_u^{}>0`$ is a uniform constant. On the other hand,
$$|\overline{}H_0|c_u|z|^2k^{1/2},$$
and thus, by an analogous argument,
(6)
$$|\overline{}_kH_0|c_u^{\prime \prime }|z|^2k^{1/2}+c_u^{\prime \prime }|z|k^{1/2}.$$
Now, if we impose that $`|_kH_0||\overline{}_kH_0|`$, we obtain by comparing (5) and (6), $`z=0`$ for $`k`$ large enough.
Now we study points in the annulus containing the support of $`\beta _\rho `$. In this case,
$$\overline{}f_0=\overline{}\beta _\rho (w^{}+H_0F_0)+\beta _\rho \overline{}H_0+(1\beta _\rho )\overline{}F_0.$$
Bounding the right hand side as in \[Do99\] we obtain an expression for the value $`|\overline{}f_0|`$. Using (3) and the bounds $`|\frac{dw^{}}{ds}|=O(1)`$ and $`|\overline{}_kF_0|=O(k^{1/2})`$, we conclude
$$|\overline{}_kf_0|c(\rho ^2+k^{1/2}+|F_0(s,0)w^{}|\rho ^1).$$
In the same way we know that
$$_kf_0=_k\beta _\rho (w^{}+H_0F_0)+\beta _\rho _kH_0+(1\beta _\rho )_kF_0.$$
Using the transversality of $`F`$ we can obtain a lower bound for $`|_kf_0|`$. The argument follows the one of Donaldson arriving to the final expression
$$|_kf_0||\overline{}_kf_0|\frac{\eta \rho }{2}c(\rho ^2+k^{1/2}+|w^{}F_0(s,0)|\rho ^1).$$
Obviously, once fixed a sufficiently small $`\rho `$, for $`k`$ large enough and $`|F_0(s,0)w^{}|`$ small enough compared with $`\rho `$ the inequality is strictly positive for any point in the annulus. $`\mathrm{}`$
To finish the proof we only have to check that the function $`f`$ conforms the local models at any point of $`\mathrm{\Delta }`$. We use the perturbation $`w^{}`$ to assure that the curves $`\gamma _i`$ project smoothly into $`^1`$, this is equivalent to impose
$$_vf(x)0,$$
for any point $`x\mathrm{\Delta }`$ and any nonzero $`vT_x\mathrm{\Delta }`$. This can be achieved by a generic perturbation, and so $`w^{}`$ can be selected to get it. Also we can assure, using this perturbation, that the intersections of branches of $`f(\mathrm{\Delta })`$ are transverse, again by a genericity argument.
Finally, given any $`x\mathrm{\Delta }`$ there exists coordinates $`(s,x_1,y_1,\mathrm{},x_n,y_n)`$ such that $`f`$ is locally written as
$$f(s,z_1,\mathrm{},z_n)=\phi (s)+H_{\alpha \beta }(s)z_\alpha z_\beta ,$$
where $`\phi ^{}(0)0`$ (this condition is equivalent to the smoothness of the branches of $`f(\mathrm{\Delta })`$). Now at a neighborhood of $`s=0`$ the $`1`$-parametric family $`(H_{\alpha \beta }(s))`$ of bilinear complex forms can be diagonalized by a smooth family of complex changes of coordinates if the eigenvalues of $`(H_{\alpha \beta })`$ are all distint. This is a genericity condition that can be achieved at all the points of $`\mathrm{\Delta }`$ by a generic perturbation of $`O(k^{1/2})`$ in $`F`$, before starting the perturbation process which we have developed along this section. With this condition, we obtain a smooth family of invertible complex matrix $`P(s)`$ such that
$$H_{\alpha \beta }(s)z_\alpha z_\beta =z^TP(s)^TP(s)z,$$
where $`z=(z_1,\mathrm{},z_n)`$. Therefore the change of coordinates
$`\times ^n`$ $``$ $`\times ^n`$
$`(s,z)`$ $``$ $`(s,P(s)z)`$
gives us the required local model. This finishes the proof of the main theorem.
## 4. Transversality results.
In this section we will prove Theorem 2.10.
### 4.1. The globalization scheme
First we recall how the globalization process developed in \[Do96\] adapts to the contact setting. This adaptation was carried out in \[IMP99\]. Now, we set up the process in a functorial way in the style of \[Au99\]. As always, we denote by $`(C,\theta )`$ an exact contact manifold.
###### Definition 4.1.
A family of properties $`P(ϵ,x)_{xC,ϵ>0}`$ of sections of bundles over $`C`$ is local and mixed $`C^r`$-open if, given a section $`s`$ satisfying $`P(ϵ,x)`$, any section $`\sigma `$ such that $`s\sigma `$ has $`(\eta ,c_R)`$ mixed $`C^r`$-bounds satisfies $`P(ϵc_u\eta ,x)`$, for some constant $`c_u`$.
###### Proposition 4.2 (\[IMP99\]).
Let $`P(ϵ,x)_{xC,ϵ>0}`$ be a local and mixed $`C^r`$-open family of properties of sections of vector bundles $`E_k`$ over $`C`$. Assume there exist uniform constants $`c`$, $`c^{}`$, $`c^{\prime \prime }`$, $`p`$ and a function $`f:^3^+`$ such that, given any $`xC`$, any small enough $`\delta >0`$, and mixed $`C^r`$-bounded sections $`s_k`$ of $`E_k`$ with uniform mixed $`C^r`$-bounds, say $`(c_D,c_R)`$, there exist, for all large enough $`k`$, mixed $`C^r`$-bounded sections $`\tau _{k,x}`$ of $`E_k`$ with the following properties:
1. $`\tau _{k,x}`$ has mixed $`C^r`$-bounds $`(c^{\prime \prime }\delta ,f(c^{\prime \prime }\delta ,c_D,c_R))`$,
2. the sections $`\frac{1}{\delta }\tau _{k,x}`$ have mixed Gaussian decay away from $`x`$ in $`C^r`$-norm,
3. $`s_k+\tau _{k,x}`$ satisfy the property $`P(\eta ,y)`$ for all $`yB_{g_k}(x,c)`$, with $`\eta =c^{}\delta (\mathrm{log}(\delta ^1))^p`$.
Then, given any $`\alpha >0`$ and mixed $`C^r`$-bounded sections $`s_k`$ of $`E_k`$, there exist, for all large enough $`k`$, mixed $`C^r`$-bounded sections $`\sigma _k`$ of $`E_k`$, such that $`s_k\sigma _k`$ has mixed $`C^r`$-bounds $`(\alpha ,c_R)`$ for some $`c_R>0`$. Also, the sections $`\sigma _k`$ satisfy $`P(ϵ,x)`$ for some uniform $`ϵ>0`$ at any $`xC`$.
Sketch of the proof: Although this is just a slight variation of the globalization argument in \[IMP99\] we provide the main lines by completeness. Let $`S`$ be a finite set of points in $`C`$ verifying the following properties:
1. $`_{xS}B_{g_k}(x,c)C`$.
2. There exists a partition $`S=_{jJ}S_j`$ verifying that $`d_{g_k}(x,y)>N`$ if $`x,yS_j`$. $`N`$ will be fixed along the proof.
3. The cardinal of $`J`$ is $`O(N^{2n+1})`$.
The idea is to achieve the property $`P(ϵ,x)`$ in all the balls $`B_{g_k}(x,c)`$, $`xS_j`$ at once. Using the hypothesis at each $`xS_j`$ we can build a local section $`\tau _{k,x}`$ which satisfies $`P(\eta ,y)`$ for all $`yB_{g_k}(x,c)`$, where $`\eta =c^{\prime \prime }\delta (\mathrm{log}(\delta ^1))^p`$. We select $`N`$ large enough, such that at a given $`xS_j`$ the components of the perturbations due to the rest of the points of the set $`S_j`$ through the directions defined by the distribution $`D`$ do not destroy the transversality obtained at $`B_{g_k}(x,c)`$. The condition we impose to $`N`$ is
(7)
$$c^{\prime \prime }\delta (\mathrm{log}(\delta ^1))^p2c_uc^{}\delta \mathrm{exp}(\lambda N^2),$$
where $`c_u`$ is a uniform constant. The left hand side of the inequality is the amount of transversality obtained by the perturbation at a point $`xS_j`$. The right hand side is a value less than the double of the value of the norm of the sum of all the other perturbations in the set $`S_j`$ (different from the selected and multiplied by the uniform constant provided by the local an $`C^r`$-open property). So the inequality assures us that, after adding the perturbation, we have obtained $`P(\frac{1}{2}c^{\prime \prime }\delta (\mathrm{log}(\delta ^1))^p,y)`$ at a $`c`$-neighborhood of $`S_j`$.
Now the process is clear. We perturb in $`S_1`$ to obtain $`P(\eta _1,x)`$ at a $`c`$-neighborhood of it. Afterwards we add again a perturbation with mixed $`C^r`$-bounds $`(\frac{1}{c_u}\eta _1/2,c_R^2)`$ to achieve $`P(min(\eta _2,\eta _1/2),y)`$ for all $`y`$ at a $`c`$-neighborhood of $`S_1S_2`$. The number of steps is independent of $`k`$, so the achieved final property is $`P(\eta ,x)`$ where $`\eta `$ does not depend on $`k`$. The added section has also mixed $`C^r`$-bounds $`(ϵ,c_R)`$. The constant $`ϵ`$ is obtained by adjusting the $`\eta _i`$ in the iteration, which is possible by property 1. The constant $`c_R`$ is uniform because at each stage does depend only on the precedent values of $`c_R`$ and $`\eta _j`$, so it is independent of $`k`$, and thus the uniformity is obvious since the number of stages is independent of $`k`$ (i.e. $`O(N^{2n+1})`$). Remark that $`c_R`$ cannot be bounded since property 1 provides uniformity but not control of the constant!
To end we have to check the inequality (7) in all the steps of the process. Following the asymptotic analysis of \[Do96\] we conclude that it is possible to obtain an integer $`N`$ independently of the step $`j`$. This ends the proof. $`\mathrm{}`$
### 4.2. The local perturbation.
We are going to achieve the transversality property for $`F`$, required in Theorem 2.10, by using Proposition 4.2. The first observation is the following
###### Lemma 4.3.
There exists a uniform constant $`\xi >0`$, such that, for $`k`$ large enough, any point $`x`$ verifying $`|F(x)|<\xi `$ lies in the set $`\mathrm{\Omega }_{2\xi }`$.
The proof is identical to that of Lemma 3.1 and we refer again to \[Do99\] for details.
We are going to perturb the sequence of sections $`s_k^1`$ on order to obtain transversality. The open mixed $`C^2`$-open property $`P(ϵ,x)`$ which we have to obtain in $`C`$ is that $`F`$ is $`ϵ`$-transverse to 0 at $`x`$. We are in the hypothesis of Definition 4.1. Now we are going to construct a local section satisfying the hypothesis of Proposition 4.2, thus concluding the proof. We can choose $`c>0`$ uniformly to assure that $`B_{g_k}(x,8c)\mathrm{\Omega }_\xi `$ for any $`x\mathrm{\Omega }_{\frac{3}{2}\xi }`$ and also that $`B_{g_k}(x,8c)`$ is in the complementary set of $`\mathrm{\Omega }_\xi `$ for any $`x`$ in the complementary set of $`\mathrm{\Omega }_{\frac{3}{2}\xi }`$.
If we are in the complementary set of $`\mathrm{\Omega }_{\frac{3}{2}\xi }`$ we can choose the section $`\tau _{k,x}`$ to be zero, because of Lemma 4.3.
Choose a point $`x\mathrm{\Omega }_{\frac{3}{2}\xi }`$. Take an approximately contact holomorphic chart $`(s,z_k^1,\mathrm{},z_k^n)`$ satisfying the properties of Corollary 2.4. Generalizing \[Au99, MPS99\] we define the $`1`$-forms
$$\mu _k^j=(\frac{z_k^js_{k,x}^{ref}}{s_k^0}).$$
In $`\mathrm{\Omega }_\xi `$ these forms have uniform mixed $`C^3`$-bounds because the term $`^r\overline{}z_k^j`$ can be bounded near $`x`$, and furthermore both $`s_{k,x}^{ref}`$ and $`s_k^0`$ have mixed $`C^3`$-bounds, as well as $`s_k^0`$ is also bounded below. Recall that $`(\mu _k^j)_{j=1,\mathrm{},n}`$ is a unitary basis of $`D^{}`$ at $`x`$, and it is almost unitary in the ball $`B_{g_k}(x,8c)`$ (i.e. the basis $`\mu _k^j`$ is arbitrarily close to a unitary basis), after eventually shrinking $`c`$ uniformly. We define an application $`v:B_{g_k}(x,8c)^n`$ by the formula:
$$F=\underset{j=1}{\overset{n}{}}v_j\mu _k^j.$$
In matrix notation,
(8)
$$F=v^T\mu _k,$$
where $`\mu _k^T=(\mu _k^1,\mathrm{},\mu _k^n)`$. Thus it is possible to understand $`\mu _k`$ as a linear map $`\mu _k(y):^nT_y^{}M`$. Multiplying by $`\mu _k^1`$ in (8) we obtain
$$v^T=F\mu _k^1.$$
This implies that $`|v|=O(1)`$, since $`\mu _k`$ is approximately unitary. We compute now the derivatives of $`v`$ using equation (8). To do this, recall that $`F`$ has mixed $`C^2`$-bounds in the ball $`B_{g_k}(x,8c)`$ (because $`F`$ has mixed $`C^3`$-bounds). Also, in the same way, $`\mu _k`$ has mixed $`C^2`$-bounds. Differentiating (8)
$$F=v^T\mu _k+v^T\mu _k,$$
which implies that $`|v^T\mu _k|=O(1)`$ and using again that $`\mu _k`$ is approximately unitary, we finally get
$$|v|=O(1)$$
Differentiating respect to $`\overline{}`$ we find that $`|\overline{}v|=O(k^{1/2})`$, and iterating the process
$$|v|=O(1),|\overline{}v|=O(k^{1/2}).$$
Now we use the approximately contact-holomorphic chart $`\mathrm{\Psi }`$ defined in Corollary 2.4, after eventual uniform shrinking of $`c`$. We construct the function $`\widehat{v}=v\mathrm{\Psi }^1`$. By the properties of $`\mathrm{\Psi }`$, it is easy to check that:
$$|\widehat{v}|=O(1),|^r\widehat{v}|=O(1),|^{r1}\overline{}\widehat{v}|=O(k^{1/2}),r=1,2.$$
Scaling the coordinates by a uniform constant we can assume that $`\mathrm{\Psi }(B_{g_k}(x,2c))B_{2n+1}(0,2)\mathrm{\Psi }(B_{g_k}(x,8c))`$. Now, we are in the hypothesis of the following
###### Proposition 4.4.
Let $`f_k:B\times [0,1]^m`$ be a sequence of functions where $`B`$ is the ball of radius $`1`$ in $`^n`$ and $`B\times [0,1]`$ is equipped with a sequence of contact forms $`\theta (k)`$ whose distributions are asymptotically flat. Let $`0<\delta <1/2`$ be a constant and let $`\sigma =\delta (\mathrm{log}(\delta ^1))^p`$, where $`p`$ is a integer depending only on the dimensions. Assume that $`f_k`$ satisfies over $`B\times [0,1]`$ the following bounds
$$|f_k|1,|\overline{}_0f_k|\sigma ,|\overline{}_0f_k|\sigma ,$$
for $`k`$ large enough, where $`\overline{}_0`$ is the $`(0,1)`$ operator defined in $`D(k)=\mathrm{ker}\theta (k)`$ by vertical projection of the standard complex structure $`J_0`$. Then for $`k`$ large enough there exists a smooth curve $`w_k:[0,1]^m`$ such that $`|w_k|<\delta `$ and the function $`f_kw_k`$ is $`\sigma `$-transverse to zero on $`B(0,1/2)\times [0,1]`$. Moreover, if $`|f_k/s|<1`$ and $`|f_k/s|<1`$, we can choose $`w_k`$ such that $`|d^iw_k/ds^i|<\mathrm{\Phi }(\delta )`$, $`(i=1,2)`$; $`d^jw_k/ds^j(0)=0`$ and $`d^jw_k/ds^j(1)=0`$, for all $`j`$, where $`c`$ is a uniform constant and $`\mathrm{\Phi }:^+^+`$ is a function depending only on the dimensions.
This proposition will be proved in Section 5. This is the analogous in the contact case to Theorem 12 in \[Do99\]. For the particular value of $`m=1`$ it has been proved in \[IMP99\].
Now we apply this proposition to the map $`\widehat{v}`$ over $`[1,1]\times B`$, because for $`k`$ large enough, it satisfies the hypothesis (without loss of generality we can choose $`[1,1]\times B`$ instead of $`[0,1]\times B`$, also we suppose $`|\widehat{v}|_{C^2}1`$, multiplying by uniform constants). The obtained path $`\widehat{w}`$ is extended to $`\times ^n`$ as
$$\widehat{w}(s,z)=\{\begin{array}{cc}\widehat{w}(1)\hfill & ,\text{for }s>1,\hfill \\ \widehat{w}(s)\hfill & ,\text{for }s[1,1],\hfill \\ \widehat{w}(1)\hfill & ,\text{for }s<1.\hfill \end{array}$$
We keep the same notation for this extended map. The map $`\widehat{v}\widehat{w}`$ is transverse to 0 in $`B(0,1)\times [1,1]`$, thus $`\widehat{v}\widehat{w}`$ is also transverse to 0 in $`B_{2n+1}(0,1)`$. Remark that the chart $`\mathrm{\Psi }`$ is defined in a ball of radius $`O(k^{1/2})`$ . So the pull-back of $`\widehat{w}`$ by $`\mathrm{\Psi }`$, denoted $`w`$, is well defined in this ball, and it will be enough for our purposes. By the properties of $`\mathrm{\Psi }`$ we have that $`v+w`$ is $`\eta `$-transverse to 0 over the ball $`B_{g_k}(x,c)`$. The constant $`\eta `$ does not coincide with the transversality obtained by $`\widehat{v}+\widehat{w}`$ by a uniform factor, so $`\eta =c_u\delta (\mathrm{log}(\delta ^1))^p`$, the constants $`c_u`$ and $`p`$ are uniform and $`\delta `$ is the norm of $`w`$ except by a uniform factor $`c_u^{}`$, i.e. $`|w|<c_u^{}\delta `$, $`c_u^{}0`$.
The needed perturbation is
$$\tau _{k,x}=\underset{j=1}{\overset{n}{}}w_jz_js_{k,x}^{ref}.$$
The first question is whether $`\tau _{k,x}`$ has the desired mixed $`C^3`$-bounds. The verification is mere routine and it follows the lines of \[IMP99\]. As an example, we compute the bound for $`\overline{}\tau _{k,x}`$:
$$\overline{}\tau _{k,x}=\overline{}w_jz_js_{k,x}^{ref}+w_j\overline{}z_js_{k,x}^{ref}+w_jz_j\overline{}s_{k,x}^{ref}.$$
The third term is easily bounded by
$`|w_jz_j\overline{}s_{k,x}^{ref}|`$ $``$ $`O(1)2d_k(x,y)c_ak^{1/2}P(d_k(x,y))\mathrm{exp}(\lambda d_k(x,y)^2)=`$
$`=`$ $`k^{1/2}Q(d_k(x,y))\mathrm{exp}(\lambda d_k(x,y)^2).`$
The second one is bounded in the same way recalling that $`\overline{}z=c_0k^{1/2}d_k(x,y)`$, where $`c_0>0`$ is uniform. For the first term we proceed as follows:
$$|\overline{}w_jz_js_{k,x}^{ref}||\overline{}w_j|P(d_k(x,y))exp(\lambda d_k(x,y)^2).$$
Recall that $`\theta _{k^{1/2}}`$ is asymptotically flat. Using this fact we obtain that $`|w_j|O(k^{1/2})d_k(x,y)`$, thus in particular $`|\overline{}w_j|O(k^{1/2})d_k(x,y)`$. This concludes the bounding.
Now we have to compare $`\tau _{k,x}`$ with $`w`$. If $`\frac{\tau _{k,x}}{s_k^0}`$ were $`w^T\mu _k`$, the proof would be rapidly concluded. But we are almost in this situation because
$$\frac{\tau _{k,x}}{s_k^0}=\mathrm{\Sigma }w_j\frac{z_js_{k,x}^{ref}}{s_k^0}+\mathrm{\Sigma }w_j\mu _j.$$
Using (3) again, we obtain
$$\frac{\tau _{k,x}}{s_k^0}=O(k^{1/2})+w^T\mu _k.$$
Therefore, we can find $`\stackrel{~}{w}`$ such that $`\frac{\tau _{k,x}}{s_k^0}=\stackrel{~}{w}^T\mu _k`$ and $`|w\stackrel{~}{w}|=0(k^{1/2})`$. So, for $`k`$ large enough $`v\stackrel{~}{w}`$ is, say, $`0.9\eta `$-transverse to 0 in $`B_{g_k}(x,c)`$. We end by remarking that $`\frac{s_k^1\tau _{k,x}}{s_k^0}`$ is $`c_b\eta `$-transverse to 0, $`c_b>0`$ a uniform constant, if $`v\stackrel{~}{w}`$ is $`0.9\eta `$-transverse to 0. This is obvious since both just differ by the application of the almost unitary matrix $`\mu _k`$. $`\mathrm{}`$
## 5. Local results.
The aim of this section is to prove Theorem 4.4. This is the generalization of the local results in the symplectic setting, needed to achieve controled transversality in the contact case.
### 5.1. Reduction to integral distributions.
We can easily reduce the proof to the following
###### Proposition 5.1.
Let $`f:B\times [0,1]^m`$ be a complex valued function, where $`B`$ is the ball of radius $`1`$ in $`^n`$. Let $`0<\delta <1/2`$ be a constant and let $`\sigma =\delta (\mathrm{log}(\delta ^1))^p`$, where $`p`$ is a suitable fixed integer depending only on the dimensions $`n,m`$. Assume that $`f_s`$ satisfies the following bounds over $`B\times [0,1]`$
$$|f_s|1,|\overline{}f_s|\sigma ,|\overline{}f_s|\sigma .$$
Then there exists a smooth curve $`w:[0,1]^m`$ such that $`|w|<\delta `$ and the function $`f_sw(s)`$ is $`\sigma `$-transverse to zero over the ball $`B(0,1/2)`$. Moreover, if $`|f_s/s|<1`$ and $`|f_s/s|<1`$, we can choose $`w`$ such that $`|d^iw/ds^i|<\mathrm{\Phi }(\delta )`$ ($`i=1,2`$), $`d^jw/ds^j(0)=0`$ and $`d^jw/ds^j(1)=0`$ for all $`j`$, where $`\mathrm{\Phi }`$ is a function depending only on the dimensions $`n,m`$.
The proof of this Proposition is a generalization of Lemma 10 in \[IMP99\] to the case $`m>1`$.
Proof of Theorem 4.4: We have to obtain the transversality of $`f_kw_k`$ when we restrict $`(f_kw_k)`$ to the distribution defined by $`\theta (k)`$. Recall that $`\theta (k)`$ is asymptotically flat, so $`\mathrm{}_M(D_h,D_k)=O(k^{1/2})`$, where $`D_h=^n\times \{p\}`$ and $`D_k=\text{Ker}\theta (k)`$. The key idea is that $`(f_kw_k)=O(1)`$, so if $`f_kw_k`$ is $`\eta `$-transverse to $`D_h`$, then it is, say, $`0.9\eta `$-transverse to $`D_k`$ for $`k`$ large enough. The factor $`0.9`$ can be eliminated by increasing $`p`$ uniformly. $`\mathrm{}`$
### 5.2. Proof of Proposition 5.1.
Let $`f:^n^m`$ be a smooth application. We define the subset $`U(f,w,\delta ,\sigma )`$ of the ball $`B(w,\delta )`$ of radius $`\delta `$ as the set of points such that $`f`$ is $`\sigma `$-transverse to $`w`$ in $`\frac{1}{2}B^{2n}`$, the ball of radius $`1/2`$ in $`^n`$. First we prove the following
###### Theorem 5.2 (Extension of Theorem 12 in \[Do99\]).
For any $`n`$, $`m`$, $`\delta >0`$ and $`0<\gamma <1`$ there is a $`p=p(n,m,\gamma )`$ such that if we define $`\sigma =\delta (\mathrm{log}(\delta ^1))^p`$ then for all the maps $`f:B^{2n}^m`$ verifying that
$$|f|1,|\overline{}f|\sigma ,|\overline{}f|\sigma ,$$
there exists a connected component of $`U(f,w,\delta ,\sigma )`$ containing another path-connected set $`U^{}(f,w,\delta ,\sigma )`$ whose volume is at least $`\gamma `$ times the total volume of $`B(w,\delta )`$ and such that given two points $`x,yU^{}(f,w,\delta ,\sigma )`$ is possible to find a smooth curve $`\gamma `$ in $`U(f,w,\delta ,\sigma )`$ joining $`x`$ and $`y`$ with curvature at each point and length bounded by $`\mathrm{\Phi }(\delta )`$, where $`\mathrm{\Phi }`$ is a function depending only on the dimensions $`n,m`$.
We repeat the the proof of Theorem 12 of \[Do99\] taking care of some details.
Proof: The first step in the proof is to approximate $`f_s`$ by a holomorphic function $`\widehat{f}_s`$ such that $`|f_s\widehat{f}_s|_{B,C^1}<c\sigma `$ (see Lemma 28 in \[Do96\]). This process does not hold over all of the unit ball. This is the reason why we restrict ourselves to the ball $`B^{}=\frac{1}{2}B^{2n}`$ in the sequel. Then we approximate $`\widehat{f}_s`$ by a polynomial. We can obtain polynomials $`g_s`$ such that $`|g_sf_s|_{B^{},C^1}c\sigma `$ and their degree $`d`$ can be stimated by $`O(\mathrm{log}(\sigma ^1))`$.
Adapting notations of \[Au97, IMP99\] we denote by $`Z_{h_s,ϵ}`$ the images of the set of points of $`B^{}`$ which are not $`ϵ`$-transverse to 0 for $`h_s`$. We want to prove that one component of the complementary set of $`Z_{f_s,\sigma }`$ satisfies the required properties. The first observation is that the $`C^1`$-closedness of $`f_s`$ and $`g_s`$ assures us that $`Z_{f_s,\sigma }Z_{g_s,(c+1)\sigma }`$.
We use now the following
###### Theorem 5.3 (Theorem 26 of \[Do99\]).
Given a polynomial $`g:^n^m`$ of degree $`d`$ and $`ϵ>0`$ there is a real-algebraic subvariety $`A(g)^m`$ of codimension $`2`$ and degree $`D`$ such that $`Z_{g,ϵ}`$ is contained in the $`Kϵ`$-neighborhood of $`A(g)`$, where $`K,D(d+1)^p`$, for some integer $`p`$ depending only on the dimensions $`n,m`$.
Also we use the following Donaldson’s result:
###### Proposition 5.4 (Proposition 31 in \[Do99\]).
For each integer $`N`$ and real number $`\theta >0`$, there is a $`\mu =\mu (\theta ,N)`$ with the following property. For any real-algebraic hypersurface $`A^n`$ of degree $`D`$ and $`ϵ(D+1)^\mu `$,
$$Vol(B^NA_ϵ)\theta .$$
Denote $`\sigma ^{}=(c+1)\sigma `$. With these two results and following the discussion in \[Do96\] p. 689 about the behaviour of the function $`\delta (\mathrm{log}\delta ^1)^p`$ we can assure the $`\sigma ^{}=\delta \mathrm{log}(\delta ^1)^p`$-neighborhood, $`p`$ a fixed integer, of the bad set $`A(g_s)`$ has volume arbitrarily small. Also we can assure the same condition for the $`3\sigma ^{}`$-neighborhood (changing $`p`$ slightly). We take a covering of $`B(0,\delta )`$ by balls $`B(x_i,\sigma ^{}/2)`$ of centers $`x_i`$ and radius $`\sigma ^{}/2`$ and assuring that the covering of balls with radius $`\sigma ^{}`$ centered in the same points cover each point of $`B(0,\delta )`$ only a finite uniform number of times, for instance less than $`\nu `$ times. Denote by $`C`$ the set of centers $`x_i`$ of the balls of the covering contained in the $`2\sigma ^{}`$-neighborhood of $`A(g_s)`$. Recall that the union of these balls is contained in $`A_{3\sigma ^{}}(g_s)`$, so we can conclude that
$$\underset{x_iCr_C}{}vol(B(x_i,\sigma ^{}))\nu vol(A_{3\sigma ^{}}(g_s)).$$
And from this expression we can easily obtain a bound for the number of balls in $`C`$ as
$$card(C)c_u\delta /(\sigma ^{})^2,$$
where $`c_u`$ is a uniform constant and so this number only depends on $`\delta `$. But observe that
$$A_\sigma ^{}(g_s)\underset{x_iC}{}B(x_i,\sigma ^{}/2)\underset{x_iC}{}B(x_i,\sigma ^{})=W_sA_{3\sigma ^{}}(g_s)$$
.
Recalling that $`W_s`$ has an arbitrary small volume, a standard isoperimetric inequality gives us that in the complementary of $`W_s`$ we can find a connected component of arbitrarily big volume (see, for instance, \[Au99\] for more details). This component will be the set $`U^{}(f,w,\delta ,\sigma )=U^{}`$ in the statement of the Theorem. Obviously in the complementary of $`A_\sigma (f_s)A_\sigma ^{}(g_s)`$ there will be a connected component containing $`U^{}`$, this component will be the set $`U(f,w,\delta ,\sigma )=U`$. To finish we have only to check that these two sets satisfy the required properties.
We adapt the ideas of \[IMP99\]. We call $`N=card(C)`$. Observe that we have fixed $`\sigma `$, and then $`N=f(\delta )`$ and $`2\pi \sigma N=g(\delta )`$, for some functions $`f`$ and $`g`$. Now take $`y,z`$ points in the large connected component $`U`$ of the complementary of $`Z_s^+`$. Denote $`L(y,z)`$ the straight segment joining them. This segment cuts at most at $`2N`$ points $`y_0,z_0,y_1,\mathrm{}`$ to the border $`Bor=(_{x_iC}B(x_i,\sigma ^{}))`$. Obviously $`L(z_i,y_{i+1})U`$ and $`L(y_i,z_i)_{x_iC}B(x_i,\sigma ^{})`$. We replace the lines $`L(y_i,z_i)`$ by curves $`C(y_i,z_i)`$ contained in $`Bor`$ connecting $`y_i`$ and $`z_i`$. We construct the curves following maximal diameters of the spheres which define the border and so $`\text{length}(C(y_i,z_i))g(\delta )`$. Therefore the curve
$$\gamma ^{}=L(y,y_0)C(y_0,z_0)L(z_0,y_1)\mathrm{}$$
satisfies
$`\text{length}(\gamma ^{})`$ $`=`$ $`L(y,y_0)+L(z_q,z)+\mathrm{\Sigma }L(y_i,y_{i+1})+\mathrm{\Sigma }C(y_i,z_i)2\delta +`$
$`+f(\delta )g(\delta )=\mathrm{\Phi }(\delta )/2,`$
where $`\mathrm{\Phi }`$ is some function depending only on the dimensions. Perturbing slightly $`\gamma ^{}`$ to make it differentiable and removing it from the border we obtain $`\gamma `$ which, bounding enough the perturbation, satisfies $`\text{length}(\gamma )\mathrm{\Phi }(\delta )`$. So the length between two points can be bounded by a function of $`\delta `$. Moreover, if we translate the diameters in the border $`Bor`$ till the border of $`_{x_iC}B(x_i,\sigma ^{}/2)`$, we can assure that the curvature of the path can be bounded by $`O(1/\sigma ^{})`$, again a function of $`\delta `$.
$`\mathrm{}`$
With the precedent result we can easily finish the proof of Proposition 5.1. Recall that $`f_s`$ satisfies
$$|\frac{f_s}{s}|1,|\frac{f_s}{s}|1.$$
This implies by a symple application of the Mean Value Theorem of the Differential Calculus (for vectorial applications) that:
$$|f_sf_{s+ϵ}|<ϵ,|f_sf_{s+ϵ}|<ϵ.$$
Recall that given two applications $`f,g:^n^m`$ such that $`|fg|_{C^1}<ϵ`$, if $`f`$ has a left inverse with norm less than $`\eta ^1`$, for a large enough uniform constant $`c_0`$, $`g`$ has a left inverse of norm less than $`(\eta c_0ϵ)^1`$. This implies that
(9)
$$U(f_s,0,\delta ,\sigma )U(f_{s+ϵ},0,\delta ,\sigma c_0ϵ).$$
We choose a number $`\sigma `$, which satisfies the hypothesis of Theorem 5.2, with $`\gamma >\frac{1}{2}`$. Take an integer $`q`$ satisfying that
(10)
$$\frac{1}{q}<\sigma /c_0.$$
We can choose $`x_iU(g_{i/q},0,\delta ,(c+1)\sigma )U(g_{(i+1)/q},0,\delta ,(c+1)\sigma )`$, for $`i=0,\mathrm{},q1`$. Consider the path $`H_i=x_i\times [i/q,(i+1)/q]`$. There exists a smooth path $`V_i`$ connecting $`x_i`$ with $`x_{i+1}`$ in $`U^{}(g_{(i+1)/q},0,\delta ,\sigma )`$. Its length is bounded above by $`\mathrm{\Phi }(\delta )`$, as well as its curvature at any point. We parametrize $`V_i`$ by its arch-length.
We choose a smooth function $`\beta :[0,1][0,1]`$ satisfying:
$$\beta (x)=\{\begin{array}{cc}0,\hfill & x[0,1/4],\hfill \\ 0<\beta (x)<1,\hfill & 1/4<x<3/4,\hfill \\ 1,\hfill & x[3/4,1].\hfill \end{array}$$
and compute $`|\beta |_{C_2}=c_b`$. Denote $`\beta _i(x)=\beta (x)\text{length}(V_i)`$. This function has norm $`|\beta _i|_{C_2}<c_b\mathrm{\Phi }(\delta )`$. Define a path
$$w((i/q+ϵ)=V_i\left(\beta _i\left(\frac{1}{q}(ϵ+\frac{q}{2})\right)\right),|ϵ|q/2.$$
The map $`w`$ is smooth and we obtain $`|w|_{C_2}<c_b\mathrm{\Phi }(\delta )`$, a bound which is a function of $`\delta `$. Remark that the first derivative depends on the length of the path $`V_i`$ and the second on its curvature. Finally we observe that $`|V_i(s)w(s)|<\frac{1}{2q}`$. It implies, by (9) and (10), that $`w(s)`$ is $`\sigma /2`$ transverse to 0, for all $`s[0,1]`$. We can find an integer $`p^{}`$ such that $`\sigma ^{}=\delta (\mathrm{log}(\delta ^1))^p^{}<\sigma /2`$ and the proof is finished.
## 6. Topological considerations.
In this section we will do some simple topological remarks about the contact Lefschetz pencils. In the first subsection we will study the topological relationship between the smooth fibres of a contact fibration.
### 6.1. Crossing critical curves.
Over each of the connected components of $`^1f_k(\mathrm{\Delta })`$ the fibres of $`f_k`$ are isotopic. Now, we are interested in analyzing the topological behaviour of the fibres when crossing through $`\mathrm{\Delta }`$. We will prove in this subsection the following
###### Proposition 6.1.
Given a contact fibration $`(f,\mathrm{\Delta },A)`$, choose a path $`\gamma :[0,1]^1`$ such that $`\gamma (1/2)f(\mathrm{\Delta })`$ and the rest of the points of $`\gamma `$ are regular values of $`f`$. Then $`N=f^1(\gamma (1))`$ is built up from $`N^{}=f^1(\gamma (0))`$ by adding a $`n`$-dimensional handle and removing another $`n`$-dimensional handle. Therefore $`H_j(N)=H_j(N^{})`$ (resp. $`\pi _j(N)=\pi _j(N^{})`$) for $`j=0,\mathrm{},n2`$.
Proof: We restrict ourselves without loss of generality to neighborhoods of $`\gamma (1/2)`$ and of the critical point of $`f`$, where we can define a compatible chart. So, with the usual identifications, we can suppose that $`f:\times ^n`$. Moreover, for simplicity, we will assume that $`f(s,z_1,\mathrm{},z_n)=s+z_1^2+\mathrm{}+z_n^2`$ (being the general case a straightforward generalization) and the path will be $`\gamma (t)=2(t1/2)i`$ with a critical value for $`f`$ at $`t=1/2`$.
The proof of the result reduces to show that $`B=f^1(\gamma ([0,1]))`$ is a cobordism between $`N=f^1(\gamma (0))`$ and $`N^{}=f^1(\gamma (1))`$ with only one surgery of index $`n`$. And this follows if we find a Morse function with a critical point of index $`n`$. Choose $`h=im(f)=2x_1y_1+2x_2y_2+\mathrm{}+2x_ny_n`$. We can assume at a neighborhood $`U`$ of the critical point $`\mathrm{𝟎}=(0,\mathrm{},0)`$ that $`BUg^1(0)`$, where $`g(s,x_1,\mathrm{})=Re(f)=s+x_1^2+y_1^2+\mathrm{}+x_n^2+y_n^2`$. To compute the index of $`g`$ we have to restrict ourselves to $`\mathrm{ker}g(\mathrm{𝟎})`$ which is
$$\mathrm{ker}g(\mathrm{𝟎})=\{(0,z_1,\mathrm{},z_n):(z_1,\mathrm{},z_n)^n.\}$$
Finally, recall that
$$h=\left(\begin{array}{ccc}0& 0& 0\\ 0& 0& 2I_n\\ 0& 2I_n& 0\end{array}\right).$$
Restricting to $`\mathrm{ker}g(\mathrm{𝟎})`$ we obtain
$$h_{|\mathrm{ker}g(\mathrm{𝟎})}=\left(\begin{array}{cc}0& 2I_n\\ 2I_n& 0\end{array}\right),$$
which has index $`n`$. $`\mathrm{}`$
Now we can obtain a geometrical relationship between the contact submanifolds obtained in Theorem 2.9 as zero sets of transverse sequences with mixed $`C^3`$-bounds. The result will be
###### Proposition 6.2.
Let $`S`$ be a line bundle with connection over a closed exact contact manifold $`(C,D)`$ Given $`N_k`$ and $`N_k^{}`$ sequences of contact submanifolds obtained as zero sets of transverse sections with mixed $`C^3`$-bounds of the bundles $`SL^k`$, then $`N_k`$ and $`N_k^{}`$ are cobordant through a cobordism defined by surgeries of index $`n`$.
Proof: Initially we will suppose that the contact form $`\theta `$ and the almost complex structure $`J`$ used in the definition and construction of the sequences of sections are coincident in the two cases.
Choose two of such sequences $`\sigma _k^1`$ and $`\sigma _k^2`$ which are $`\eta `$-transverse to 0, for some $`\eta >0`$. We claim that we can assume without loss of generality that $`\sigma _k^1\sigma _k^2`$ is $`\eta ^{}`$-transverse to 0, for some $`\eta >0`$. If we suppose this to be false, we can apply Theorem 2.9 and find a sequence of sections $`\tau _k^1\tau _k^2`$ with mixed $`C^3`$-bounds $`(\eta /2,c_R)`$, such that they are $`\eta ^{}`$-transverse to 0. Define
$$\sigma _{k,t}^j=\sigma _k^j+t\tau _k^j,j=1,2.$$
It is easy to check that the sections $`\sigma _k^j(t)`$ are $`\eta /2`$-transverse to 0, therefore its zero sets are isotopic. So we can impose directly that $`\sigma _k^0\sigma _k^1`$ is transverse to 0.
Define the function
$`F_k:CZ(\sigma _k^1\sigma _k^2)`$ $``$ $`^1`$
$`p`$ $``$ $`{\displaystyle \frac{\sigma _k^2(p)}{\sigma _k^1(p)}}.`$
By an analogous argument to that in the precedent lines we assume that $`F_k`$ is $`\eta `$-transverse to 0 away from $`Z(\sigma _k^1\sigma _k^2)`$ Recall that $`F_k^1(0)=Z(\sigma _k^1)`$ and $`F_k^1(\mathrm{})=Z(\sigma _k^2)`$. Following the proof of Section 3 we obtain a contact fibration $`f_k`$, for $`k`$ large enough. But, we do not impose the local model at a neighborhood of $`A`$, therefore we do not perturb any neighborhood of $`A`$. We observe that the perturbation at a neighborhood of the set of critical points $`\mathrm{\Delta }`$ does not change the counterimages of the points $`0`$ and $`\mathrm{}`$. It follows since there exists a $`\rho _0`$-neighborhood, being $`\rho _0>0`$ a uniform constant, of these two fibres which does not intersect the “bad set” $`\mathrm{\Gamma }`$, because of the mixed $`C^3`$-bounds of the sections $`\sigma _k^1`$ and $`\sigma _k^2`$ (which assure contactness at a neighborhood of the zero set of the section). So we perturb the sequence $`F_k`$ in a $`\rho `$-neighborhood of $`\mathrm{\Delta }`$, with $`0<\rho <\rho _0`$. Therefore, even after performing the needed perturbations, $`f_k^1(0)=N_k`$ and $`f_k^1(\mathrm{})=N_k^{}`$.
By Proposition 6.1 we obtain, following a path between $`0`$ and $`\mathrm{}`$, that $`f_k^1(0)=N_k`$ and $`f_k^1(\mathrm{})=N_k^{}`$ are related by a sequence of operations of index $`n`$. This finishes the proof if we admit that the complex structures coincide for the two sequences of sections.
Suppose now that the sequence of submanifolds $`N_k`$ is the zero set of a sequence of sections $`\sigma _k`$ with mixed $`C^3`$-bounds respect to an almost complex structure $`J_0`$. Furthermore $`N_k^{}`$ comes from $`s_k^{}`$ with mixed $`C^3`$-bounds respect to $`J_1`$. Fix a continuous path $`J_t`$ joining $`J_0`$ and $`J_1`$ in the moduli of compatible almost complex structures. We use the following
###### Lemma 6.3.
Let $`J_0`$ be a compatible almost complex structure in $`C`$. There exists a uniform $`ϵ>0`$ satisfying that for any compatible almost complex structure $`J`$ such that $`|JJ_0|<ϵ`$, there exist two sequences of sections $`s_k`$ and $`s_k^{}`$ with mixed $`C^3`$-bounds, respect to $`J_0`$ and $`J`$, which are $`\eta `$-transverse to 0, for some $`\eta >0`$. Moreover the zero sets of $`s_k`$ and $`s_k^{}`$ are isotopic for $`k`$ large enough.
Find the uniform constants $`ϵ_t=ϵ`$ provided by Lemma 6.3 for the almost complex structures $`J_t`$. This implies, by the continuity of $`J_t`$, that there exists $`ϵ_t^{}>0`$ such that we can find two sequences of sections $`s_k^t`$ and $`s_k^t^{}`$ with mixed $`C^3`$-bounds respect to $`J_t`$ and $`J_{t+ϵ_t^{}}`$ whose zero sets are isotopic, for $`k`$ large enough.
Now we cover the segment $`[0,1]`$ by open segments $`(t,t+ϵ_t^{})`$. The segment is compact, so we can find a finite subset $`(0,ϵ_0^{}),(t_1,t_1+ϵ_{t_1}^{}),\mathrm{}(t_N,1)`$ of the precedent set of segments which covers $`[0,1]`$. Obviously, without loss of generality, we can choose $`t_j+ϵ_{t_j}^{}=t_{j+1}`$.
Denote $`Z(s)`$ the zero set of a given section $`s`$. Observe that the sets $`Z(\sigma _k)`$ and $`Z(s_k^0)`$ are related through a cobordism of index $`n`$, for $`k`$ large enough. But, recall now that $`Z(s_k^0)`$ and $`Z(s_k^0^{})`$ are isotopic for $`k`$ large enough. Again, $`Z(s_k^0^{})`$ and $`Z(s_k^{t_1})`$ are related through surgeries of index $`n`$… Following this argument we find that $`Z(\sigma _k^0)`$ and $`Z(\sigma _k^1)`$ are related through surgeries of index $`n`$.
Finally recall that if $`s_k`$ is $`\eta `$-tranverse to 0 with mixed $`C^r`$-bounds $`(c_D,c_R)`$ with respect to a complex structure $`J`$ and to a contact form $`\theta `$, then $`s_k`$ is $`\eta /c`$-transverse to 0 with mixed $`C^r`$-bounds $`(cc_D,cc_R)`$ with respect to the same contact structure $`J`$ and to a contact form $`\theta ^{}`$. The constant $`c>0`$ only depends on the contact forms $`\theta `$ and $`\theta ^{}`$. This proves that the result does not depend on the chosen contact form. $`\mathrm{}`$
Proof of Lemma 6.3: Construct a sequence of sections $`\sigma _k`$ $`\eta `$-transverse to 0 using the globalization argument of Proposition 4.2, for some uniform constant $`\eta >0`$. Thus, we obtain
$$\sigma _k=\underset{jJ}{}w_ks_{k,x_j}^{ref},$$
where the points $`x_j`$ are elements of the set $`S`$ with the properties described in the proof of that Proposition. Recall that the definition of $`\sigma _k`$ makes sense if we change the almost complex structure $`J_0`$ by another one $`J`$. The obtained section $`\sigma _k^{}`$ will be different because the reference sections $`s_{k,x_j}^{ref}`$ depend on the complex structure. However, we claim that $`\sigma _k^{}`$ has mixed $`C^3`$-bounds and that $`\sigma _k\sigma _k^{}`$ has mixed $`C^3`$-bounds $`(c_uϵ,c_R^{})`$, where $`c_u`$ and $`c_R^{}`$ are uniform constants. Impose
$$c_uϵ<\frac{\eta }{2}.$$
Obviously the constant $`ϵ`$ can be chosen uniformly and also $`\sigma _k`$ and $`\sigma _k^{}`$ have isotopic zero sets. This finishes the proof.
We have only to check that $`\sigma _k\sigma _k^{}`$ have the desired mixed $`C^3`$-bounds. We must follow the proof in \[IMP99\] to check this condition. We only give the key ideas:
1. Check that $`|s_{k,x,J_0}^{ref}s_{k,x,J}|_{C^2}=O(ϵ).`$
2. Compute the integer $`N`$ in the globalization process for the sequence in $`J_0`$. Check that this constant is enough also for the globalization process with $`J`$ if we shrink $`ϵ`$ uniformly in this globalization process. Why uniformly? Because $`N`$ is independent of $`k`$.
$`\mathrm{}`$
Proposition 6.2 is weaker than the “contact Lefschetz hyperplane theorem” proved in \[IMP99\], but it is more geometrical and enligthens the behaviour of the generic sections of the bundles $`L^k`$. It would be interesting to study the different connected components of $`^1\mathrm{\Delta }`$ to control the topology of all the “approximately holomorphic” contact zeroes of $`L^k`$.
## 7. The non-exact case.
To conclude the discussion we consider now the non-exact case. The important point is the following standard result
###### Proposition 7.1.
Given a non-exact contact manifold $`(C,D)`$. There exists an exact contact manifold $`(\widehat{C},\widehat{D})`$ which is a non-trivial double covering of $`C`$. The projection is a contactomorphism and it can be found a contact form $`\widehat{\theta }`$ in $`\widehat{C}`$ such that the structure $`_2`$-action of the covering is a strong anticontactomorphism, i.e. for $`\alpha _2`$, it is $`\alpha _{}\widehat{\theta }=\widehat{\theta }`$.
For a simple proof, see \[IMP99\]. Now, we follow the ideas in \[IMP99\]. We lift $`S`$ to the double covering and denote it again by $`S`$. It is easy to find an almost complex compatible structure $`J`$ satisfying that $`\alpha _{}J=J`$, this implies that $`\alpha _{}g_k=g_k`$. The application $`\alpha `$ lifs to a morphism of bundles $`\stackrel{~}{\alpha }:L^k\overline{L}^k`$ (in fact recalling from \[IMP99\] that $`L`$ is trivialized by construction, we take the identity in each fiber). This morphism preserves the connection. So it is easy to check that if $`s_k`$ is a sequence of sections of the bundles $`SL^k`$ with $`C^r`$-bounds $`(c_D,c_R)`$ with respect to the contact form $`\theta `$ and the almost complex structure $`J`$, then $`\alpha _{}s_k`$ is a sequence of sections of the bundle $`S\overline{L}^k`$ with respect to the contact form $`\theta `$ and to the compatible almost complex structure $`J`$ with the same mixed $`C^r`$-bounds $`(c_D,c_R)`$. But, now we identify the bundle $`L^k`$ with $`\overline{L}^k`$ using the anticomplex isomorphism provided by the identity. Then a simple computation shows that $`\alpha _{}s_k`$ is a sequence os sections of the bundles $`SL^k`$ with the same mixed $`C^3`$-bounds.
Our only task is to assure that all the objects in the construction are $`G`$-invariant. The precedent considerations assure this condition in the local constructions. Now we are going to study the globalization process. The first point is to achieve the set of points $`S`$ invariantly. Moreover each $`S_j`$ has to be $`_2`$-invariant. This is only true if the action of $`_2`$ is free, because in any other case we would have problems to assure the second property of the set. The perturbation is performed in a similar way that in the standard case. The key idea is that the perturbation term $`\tau _{k,x}`$ constructed to obtain transversality close to a point $`x`$ can be transported to the point $`\alpha (x)`$ by means of $`\tau _{k,\alpha (x)}=\alpha _{}\tau _{k,x}`$. There is no interference between $`x`$ and $`\alpha (x)`$, because
$$d_k(x,\alpha (x))=O(k^{1/2}).$$
The term $`\tau _{k,\alpha (x)}`$ produces the same transversality as $`\tau _{k,x}`$ because of the $`_2`$-invariance of the construction.
The perturbation of the function $`F^{s_k}`$ can be made invariantly. In the set $`A`$ it is clear. In $`\mathrm{\Delta }`$, the trivialization is invariant, as well as the function $`H`$ and we can choose $`w^{}`$ invariant without loss of generality.
So we have constructed a pencil in $`\widehat{C}`$ which is $`_2`$-invariant. This $`_2`$-invariance property allows us to quotient by the group $`_2`$ obtaining new data $`A=\widehat{A}/_2`$, $`\mathrm{\Delta }=\widehat{\mathrm{\Delta }}/_2`$ and $`f=f/_2`$. It is a trivial exercise to check that the object so defined is a contact pencil on $`C`$. $`\mathrm{}`$
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# PROSPECTS FOR THE LENSING OF SUPERNOVAE
## 1 Introduction
There are still many outstanding problems in the fields of structure formation and dark matter which gravitational lensing has not yet been able to address. Outside of the visible extent of galaxies not a great deal is known about the distribution of dark matter (DM) on scales smaller than galaxy clusters. Galaxy–galaxy lensing has put some constraints on the mass of galaxy halos, but their size scale and the degree to which they are smooth mass distributions or collections of subclumps are not well determined.$`^{\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?}}`$ Even the concept of galaxy halos, with a single galaxy within each of them and well defined sizes, may not be correct. Nor is it clear how the observable properties of galaxies relate to the DM distribution around them. Recently dark matter simulations have revealed several problems with the cold dark matter (CDM) model. One of these is that CDM predicts a large amount of small scale structure within halos which is unaccounted for in the observed distribution of light.$`^\mathrm{?}`$ In addition, the composition of DM is still a mystery: it could be large compact objects like black holes or it could be microscopic particles like WIMPS. Even with primordial nucleosynthesis bounds roughly half the baryons are in some form that has yet to be directly detected – conceivably in condensed objects. I will describe here how the gravitational lensing of type Ia SNe can help to answer some of these questions. Lack of space dictates that this only be an outline and that the reader be referred to more detailed papers.
High redshift type Ia SNe have recently been used by two collaborations to measure cosmological parameters.$`^{\mathrm{?},\mathrm{?}}`$ These measurements and the proposed application of Ia’s to gravitational lensing are made possible by the discovery of a tight correlation between the light curve shapes and peak luminosities of these SNe.$`^{\mathrm{?},\mathrm{?}}`$ The measured standard deviation of corrected peak luminosities in local SNe is now $`\stackrel{<}{}0.12\text{ mag}`$. The successes of these collaborations have inspired plans for more aggressive, larger searches for high redshift SNe in the near future. The volume and quality of data is likely to increase dramatically. Besides constraining $`\mathrm{\Omega }_m`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }`$ these data, amongst other things, would be ideal for gravitational lensing studies.
In the next section I will discuss how SNe can be used to determine if dark matter is composed of macroscopic compact objects and in section 3 the case of DM composed of microscopic particles is taken up.
## 2 Lensing by macroscopic compact objects
Some magnification probability distributions for a point source at $`z=1`$ are plotted in figure 1. The means of all these distributions are zero which corresponds to the usual Friedman-Robertson-Walker (FRW) luminosity distance. When DM is made entirely of macroscopic compact objects the distribution peaks well below the mean making most SNe under-luminous. The peak of the distribution is just slightly brighter than the solution corresponding to the empty beam or Dyer–Roeder luminosity distance $`^\mathrm{?}`$. This is the formal solution of Sachs optical scalar equations$`^\mathrm{?}`$ for a beam that passes through only empty space (no Ricci focusing) and has no shear on it. The long, high magnification tail to these distributions correspond to rare cases where a DM object is very close to the line of sight. At $`z<2`$ it is unlikely that multiple lens lie so close to the line of sight that their lensing effects become nonlinearly coupled which significantly simplifies calculations.
The distributions for microscopic DM shown in figure 1 assume that the DM is clustered into halos surrounding galaxies. The exact form of the clustering is not important for this section. As the figure shows, these distributions are significantly more centrally concentrated about the FRW solution. With macroscopic DM the magnification is dominated by Ricci focusing – the isotropic expansion or contraction of the image caused by matter within the beam – while in the macroscopic case (with true point sources) it is entirely due to shear caused by matter outside of the beam. This is the essential difference between the two cases and why they can be differentiated using SNe observations.
To differentiate between DM candidates we construct the statistic
$`_p{\displaystyle \frac{1}{N_{sn}}}\mathrm{ln}\left[{\displaystyle \frac{P\left(\{\delta \mu \}|\text{ macro DM,noise}\right)}{P\left(\{\delta \mu \}|\text{ micro DM,noise}\right)}}\right].`$ (1)
where the $`P`$’s are the probability of getting the $`N_{sn}`$ observed SN magnifications given that DM is of the specified type and given the expected noise. $`_p`$ is close to normally distributed for a modest number of SNe. Some example distributions of $`_p`$ are plotted in figure 2. If DM is microscopic $`_p`$ is expected to be small while if DM is macroscopic it will be large. Figure 2 shows that with 100 SNe at $`z1`$ one is unlikely to confuse the two cases.
This technique works for DM objects with masses $`\stackrel{>}{}10^3\text{ M}_{}`$. For smaller masses it is probable that the expanding SN photosphere will make the magnification time dependent which is an interesting subject in itself. For more details on differentiating DM candidates see Metcalf & Silk.$`^\mathrm{?}`$
## 3 Correlations between light and magnification
When the DM is made of microscopic particles it can be treated as a transparent, massive fluid clumped into structures that are much larger than the beam size. Because of lensing the variance of SN brightnesses will increase with redshift$`^\mathrm{?}`$ which introduces noise into the cosmological parameter estimates. The increases in variance is not the best way of detecting the lensing however. Using the correlation between foreground light and SN brightness reduces possible systematic errors associated with the evolution of the SN population and makes the measurement more directly sensitive to the mass and size scale of dark matter halos.
Lets define the weighted foreground flux for each SN as
$`={\displaystyle \underset{y_i<R}{}}w(z_i,z_s)f_i`$ (2)
where $`f_i`$ the observed flux from the $`i`$th galaxy and $`y_i`$ is its distance (angular or proper) from the line of sight to the SN. The weight function depends on the galaxy and SN redshifts. Figure 3 shows the expected correlation of $``$ with SN brightness, $`\delta \delta b`$, for a SN at $`z_s=1`$. This quantity is proportional to the average surface density within distance $`R`$ of a galaxy.
By measuring this correlation the size scale of galactic halos could be constrained significantly better than with galaxy–galaxy lensing using a much larger data set. This is a result of the magnification being a steeper function of galactic radius than shear when the density profile is steeper than isothermal. In addition, by selecting or searching for SNe behind galaxy clusters the structure of clusters and the smaller halos within them can be probed. For instance the tidal truncation of galactic halos could be investigated. Higher order correlations can also be used to look for substructure in galaxy halos. This subject is treated more thoroughly in a complete paper by the author.$`^\mathrm{?}`$
## 4 Discussion
The number of high redshift SNe required for the studies proposed here are well within the projected numbers for future SN searches – the VISTA telescope is expected to find hundreds at $`z\stackrel{>}{}1`$.<sup>a</sup><sup>a</sup>ahome page: http://www-star.qmw.ac.uk/$``$jpe/vista/ and the proposed SNAPSAT satellite could find thousands up to $`z1.7`$.<sup>b</sup><sup>b</sup>bhome page: http://snap.lbl.gov/proposal/ The errors assumed in this paper are quite conservative compared to those expected for satellite observations. The future looks bright for this new field in gravitational lensing.
## Acknowledgments
Many thanks to all the friends that made this such an enjoyable conference. Thanks also to J. Silk for all his help while some of this work was being done.
## References
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# 1 Introduction
## 1 Introduction
Six-dimensional conformal field theories (CFT<sub>6</sub>) have attracted some interest in view of recent advances in string/M-theory. In particular, the low energy dynamics of the collective coordinates of $`N`$ coinciding M5 branes of M-theory realize a quite interesting class of non-trivial CFT<sub>6</sub> with maximal supersymmetry, the so-called $`𝒩=(0,2)`$ interacting theories originally described in . While they still lack a lagrangian formulation, the AdS/CFT conjecture has provided concrete tools to extract some informations on these rather mysterious theories in their large $`N`$ limit, such as their trace anomalies , the spectrum of their operators and some of their 2- and 3-point correlation functions . To understand better the structure of these interacting theories, in refs. some of their more accessible properties were compared with those of another maximally supersymmetric CFT<sub>6</sub>: the non-interacting one made up by $`N`$ copies of the free $`𝒩=(0,2)`$ tensor multiplet containing 5 scalars, 1 two-form with selfdual field strength and 2 Weyl fermions. In particular, in the trace anomalies for the free theory were computed and compared with the corresponding ones for the interacting theory obtained through the AdS/CFT conjecture . In this comparison it is crucial to disentangle the coefficients of the universal part of the anomalies by separating out a trivial sector. The latter can always be cancelled by the variation of a local counterterm which can be added to the effective action. This disentanglement will be the subject of the present paper and we will start addressing it after a brief introduction to the topic of trace anomalies.
Trace anomalies can be characterized by the anomalous Weyl variation of a general coordinate invariant effective action depending on a background metric. Discovered originally in (and reviewed in ) they can be computed using Feynman graphs, as in the original papers, or more efficiently using the heat kernel methods of De Witt (as employed e.g. in ) or by a quantum mechanical representation as proposed in . Their structure has been analyzed in by cohomological methods which encode the information on the Wess–Zumino consistency conditions specialized to the Weyl symmetry. Finally, a useful classification was described in where trace anomalies are divided into three classes: type A (always proportional to the topological Euler density), type B (made up by independent Weyl invariants) and trivial anomalies (obtainable as Weyl variations of local functionals and in general expressible as total derivatives). The number of type B and trivial anomalies grows quite rapidly with the number of dimensions. Already in six dimensions there are 1 type A, 3 type B and 6 trivial anomalies satisfying the consistency conditions on a set of 17 independent terms with the property of being cubic in the curvature (and thus with the correct dimensions to constitute possible trace anomalies). Usually a concrete calculation on a specific model delivers the anomaly as a linear combination of these 17 terms, and one is left with the problem of disentangling the correct universal part of the anomaly (type A and B).
This problem was solved pragmatically in by expressing the $`6d`$ trace anomalies in a special basis for the curvature invariants (a basis which employs the Weyl tensor and traceless Ricci tensor instead of the Riemann and Ricci tensors). This special basis makes it easier to cast the anomaly into a form with the expected type A and B contributions plus a combination of total derivatives. The latter part was interpreted as a trivial anomaly since on general grounds one expects trivial anomalies to be total derivatives. This procedure worked for the four cases of free scalars, fermions, two-forms and interacting $`(0,2)`$ theory considered in . Now, one may object that the calculations made there used a basis of 7 independent total derivatives while the cohomological analysis of predicts only 6 trivial anomalies. Thus, to make sure of the correct identification of the various anomalies, we have decided to perform again a cohomological analysis to derive a basis for the trivial anomalies. This allows us to check that indeed there are only 6 trivial anomalies, for which we find the explicit expressions and identify the local counterterms that can cancel them. Then we verify that one specific linear combination of the 7 total derivative terms used in doesn’t solve the consistency conditions, but never appears in the results for the trace anomalies of the various cases treated there, thus confirming the correctness of those results.
The final output of our analysis is a systematic classification of the type A, type B and trivial anomalies for trace anomalies in six dimensions. This knowledge can be useful to put new calculations of such anomalies in a preferred basis and extract unambiguously the universal coefficients of the type A and B parts. It is presumably the difficulties related to the proper factorization of trivial anomalies that has caused a miscalculation of the trace anomalies for a scalar field in .
Thus in sect. 2 we review and solve the consistency conditions. In sect. 3 we cast those solutions into a more useful basis by taking into account their character as anomalies of type A, B and trivial. In sect. 4 we present our conclusions. Finally, we leave appendices A, B, C and D for more technical parts where we list useful results of our calculations.
## 2 Consistency conditions
Let’s consider the effective action $`W[g]`$ for a CFT coupled to a background metric $`g_{ab}`$. We will use an euclidean signature. For simplicity we assume absence of chiral gravitational anomalies and thus can consider $`W[g]`$ to be general coordinate invariant. This assumption is not necessary . Under an infinitesimal Weyl transformation depending on an infinitesimal arbitrary function $`\sigma (x)`$
$`\delta _\sigma g_{ab}(x)=2\sigma (x)g_{ab}(x),`$ (1)
the effective action generically suffers an anomalous variation
$`\delta _\sigma W[g]={\displaystyle d^6x\sqrt{g}\sigma (x)A(x)}.`$ (2)
It is well known that by functional differentiation with respect to $`\frac{2}{\sqrt{g}}\frac{\delta }{\delta g_{ab}}`$ the effective action $`W[g]`$ generates correlation functions of the stress tensor $`T^{ab}`$. Thus eq. (2) produces an anomalous trace to the stress tensor
$`T^a{}_{a}{}^{}(x)=A(x)`$ (3)
which depends on the background curvature. General coordinate invariance guarantees that the anomaly $`A(x)`$ is a scalar and dimensional considerations in $`d=6`$ fix it to be cubic in the curvature (two covariant derivatives count as one curvature). However, those particular anomalies that can be obtained also from the Weyl variation of local functionals of the metric are considered trivial since they can be cancelled by subtracting the same local functionals from the original non-local effective action.
Since the anomaly $`A(x)`$ is obtained by varying a functional, there are integrability conditions that can be identified by applying the commutator algebra of the Weyl symmetry, $`[\delta _{\sigma _1},\delta _{\sigma _2}]=0`$, to the effective action. Such integrability conditions are generically known as Wess–Zumino consistency conditions, originally derived for chiral anomalies in .
Now, we follow the work of Bonora et al. as a guideline to study the consistency conditions
$`[\delta _{\sigma _1},\delta _{\sigma _2}]W[g]=0`$ (4)
and derive all of their solutions in $`d=6`$. We use the following conventions for the curvature tensors
$`[_a,_b]V^c=R_{ab}{}_{}{}^{c}{}_{d}{}^{}V_{}^{d},R_{ab}=R_{ca}{}_{}{}^{c}{}_{b}{}^{},R=R^a{}_{a}{}^{},`$ (5)
so that the scalar curvature of a sphere is positive, and use the same basis of 17 independent curvature invariants as in <sup>1</sup><sup>1</sup>1 However we differ in the definitions of the various curvature tensors, so there are some sign differences with respect to ref. . In particular, we agree on the sign of the Riemann tensor $`R_{abcd}`$ but have opposite sign for the Ricci tensor $`R_{ab}`$ and scalar curvature $`R`$.
$$\begin{array}{ccc}K_1=R^3\hfill & K_2=RR_{ab}^2\hfill & K_3=RR_{abmn}^2\hfill \\ K_4=R_a{}_{}{}^{m}R_{m}^{}{}_{}{}^{i}R_{i}^{}^a\hfill & K_5=R_{ab}R_{mn}R^{mabn}\hfill & K_6=R_{ab}R^{amnl}R^b_{mnl}\hfill \\ K_7=R_{ab}{}_{}{}^{mn}R_{mn}^{}{}_{}{}^{ij}R_{ij}^{}^{ab}\hfill & K_8=R_{amnb}R^{mijn}R_i{}_{}{}^{ab}_j\hfill & K_9=R^2R\hfill \\ K_{10}=R_{ab}^2R^{ab}\hfill & K_{11}=R_{abmn}^2R^{abmn}\hfill & K_{12}=R^{ab}_a_bR\hfill \\ K_{13}=(_aR_{mn})^2\hfill & K_{14}=_aR_{bm}^bR^{am}\hfill & K_{15}=(_iR_{abmn})^2\hfill \\ K_{16}=^2R^2\hfill & K_{17}=^4R.\hfill & \end{array}$$
(6)
All other terms cubic in the curvature are linear combinations of the above invariants after taking into account the symmetry properties and the Bianchi identities of the Riemann tensor.
Any trace anomaly can be expanded in the above basis
$`\delta _\sigma W[g]={\displaystyle d^6x\sqrt{g}\sigma (x)\underset{i=1}{\overset{17}{}}a^iK_i}`$ (7)
and after computing a second Weyl variation one obtains
$`[\delta _{\sigma _2},\delta _{\sigma _1}]W[g]={\displaystyle d^6x\sqrt{g}\underset{i=1}{\overset{17}{}}\underset{\alpha =1}{\overset{9}{}}f^\alpha {}_{i}{}^{}a_{}^{i}H_\alpha }`$ (8)
where the rectangular matrix of coefficients $`f^\alpha _i`$ can be constructed using the variations of the terms entering eq. (7) and reported in appendix A, and where the 9 independent (unintegrated) 2-cochains $`H_\alpha `$ are given by<sup>2</sup><sup>2</sup>2The symbol $`[]`$ denotes antisymmetrization, namely $`a_{[1}b_{2]}=a_1b_2a_2b_1`$.
$$\begin{array}{cc}H_1=R^2\sigma _{[1}^2\sigma _{2]}\hfill & H_2=R_{ab}^2\sigma _{[1}^2\sigma _{2]}\hfill \\ H_3=RR_{ab}\sigma _{[1}^a^b\sigma _{2]}\hfill & H_4=R_{am}R^m{}_{b}{}^{}\sigma _{[1}^{}^a^b\sigma _{2]}\hfill \\ H_5=R_{abmn}^2\sigma _{[1}^2\sigma _{2]}\hfill & H_6=(^2R)\sigma _{[1}^2\sigma _{2]}\hfill \\ H_7=R\sigma _{[1}^4\sigma _{2]}\hfill & H_8=R_{ab}\sigma _{[1}^a^b^2\sigma _{2]}\hfill \\ H_9=R^{mn}R_{amnb}\sigma _{[1}^a^b\sigma _{2]}.\hfill & \end{array}$$
(9)
Now the consistency condition eq. (4) applied to eq. (8) requires that
$`{\displaystyle \underset{i=1}{\overset{17}{}}}f^\alpha {}_{i}{}^{}a_{}^{i}=0,\alpha =1,\mathrm{},9.`$ (10)
The $`9\times 17`$ matrix $`f^\alpha _i`$ has rank 7, so there are 7 independent constraints for the coefficients $`a^i`$ to form a consistent anomaly. The resulting 10 independent anomalies can be presented as
$`M_I(x)={\displaystyle \underset{i=1}{\overset{17}{}}}a_I^iK_i(x)`$ (11)
where the 10 vectors $`a_I^i`$, $`I=1,\mathrm{},10`$, form a basis for the solutions of eq. (10). These vectors are constructed in appendices A (where one can read off the matrix $`f^\alpha _i`$) and C. We list them later on in eqs. (1322).
Trivial anomalies are those that can be obtained by varying a local functional. To recognize them we compute the Weyl variation of the most general local functional obtained as a linear combination with coefficients $`c^i`$ of the integrated curvature invariants $`K_i`$ (note that we can restrict the index $`i10`$ since by partial integration the remaining terms are not linearly independent)
$`\delta _\sigma {\displaystyle d^6x\sqrt{g}\underset{i=1}{\overset{10}{}}c^iK_i}={\displaystyle d^6x\sqrt{g}\sigma (x)\underset{i=1}{\overset{10}{}}\underset{j=1}{\overset{17}{}}g^j{}_{i}{}^{}c_{}^{i}K_j}.`$ (12)
The matrix of coefficients $`g^j_i`$ has rank 6 and therefore identifies 6 trivial anomalies. These are constructed in appendices B (where one can read off the matrix $`g^j_i`$) and C, and reported here below in eqs. (1722). In appendix C one may also find the local functionals which generate the trivial anomalies (see eq. (42)).
Now we present the solutions of the consistency conditions just described. A basis for the non-trivial anomalies is given by
$`M_1`$ $`=`$ $`{\displaystyle \frac{19}{800}}K_1{\displaystyle \frac{57}{160}}K_2+{\displaystyle \frac{3}{40}}K_3+{\displaystyle \frac{7}{16}}K_4{\displaystyle \frac{9}{8}}K_5{\displaystyle \frac{3}{4}}K_6+K_8`$ (13)
$`M_2`$ $`=`$ $`{\displaystyle \frac{9}{200}}K_1{\displaystyle \frac{27}{40}}K_2+{\displaystyle \frac{3}{10}}K_3+{\displaystyle \frac{5}{4}}K_4{\displaystyle \frac{3}{2}}K_53K_6+K_7`$ (14)
$`M_3`$ $`=`$ $`K_1+8K_2+2K_310K_4+10K_5{\displaystyle \frac{1}{2}}K_9+5K_{10}5K_{11}`$ (15)
$`M_4`$ $`=`$ $`K_1+12K_23K_316K_4+24K_5+24K_64K_78K_8`$ (16)
where we have chosen to agree with the ones reported in ref. . Instead a suitable basis for the remaining trivial anomalies is given by
$`M_5`$ $`=`$ $`6K_63K_7+12K_8+K_{10}7K_{11}11K_{13}+12K_{14}4K_{15}`$ (17)
$`M_6`$ $`=`$ $`{\displaystyle \frac{1}{5}}K_9+K_{10}+{\displaystyle \frac{2}{5}}K_{12}+K_{13}`$ (18)
$`M_7`$ $`=`$ $`K_4+K_5{\displaystyle \frac{3}{20}}K_9+{\displaystyle \frac{4}{5}}K_{12}+K_{14}`$ (19)
$`M_8`$ $`=`$ $`{\displaystyle \frac{1}{5}}K_9+K_{11}+{\displaystyle \frac{2}{5}}K_{12}+K_{15}`$ (20)
$`M_9`$ $`=`$ $`K_{16}`$ (21)
$`M_{10}`$ $`=`$ $`K_{17}.`$ (22)
This is the main result we were searching for.
## 3 A useful basis for six dimensional trace anomalies
In the previous section we have derived the solutions to the consistency conditions. We now put those solutions into a more useful basis by taking into account their character as type A, B or trivial anomalies, as classified in .
The type A anomaly is unique and proportional to the six dimensional topological Euler density and can be written as
$`E_6`$ $`=`$ $`ϵ_{m_1n_1m_2n_2m_3n_3}ϵ^{a_1b_1a_2b_2a_3b_3}R^{m_1n_1}{}_{a_1b_1}{}^{}R_{}^{m_2n_2}{}_{a_2b_2}{}^{}R_{}^{m_3n_3}{}_{a_3b_3}{}^{}=8M_4.`$ (23)
The anomalies of type B are given instead by the three following Weyl invariants
$`I_1`$ $`=`$ $`C_{amnb}C^{mijn}C_i{}_{}{}^{ab}{}_{j}{}^{}=M_1`$ (24)
$`I_2`$ $`=`$ $`C_{ab}{}_{}{}^{mn}C_{mn}^{}{}_{}{}^{ij}C_{ij}^{}{}_{}{}^{ab}=M_2`$ (25)
$`I_3`$ $`=`$ $`C_{mabc}\left(^2\delta _n^m+4R_n^m{\displaystyle \frac{6}{5}}R\delta _n^m\right)C^{nabc}+_iJ^i`$ (26)
$`=`$ $`{\displaystyle \frac{16}{3}}M_1+{\displaystyle \frac{8}{3}}M_2{\displaystyle \frac{1}{5}}M_3+{\displaystyle \frac{2}{3}}M_4+_iJ^i`$ (27)
where
$$C_{abcd}=R_{abcd}\frac{1}{4}(g_{ac}R_{bd}+g_{bd}R_{ac}g_{ad}R_{bc}g_{bc}R_{ad})+\frac{1}{20}(g_{ac}g_{bd}g_{ad}g_{bc})R$$
(28)
is the Weyl tensor in 6 dimensions and
$`_iJ^i={\displaystyle \frac{2}{3}}M_5{\displaystyle \frac{13}{3}}M_6+2M_7+{\displaystyle \frac{1}{3}}M_8`$ (29)
is a trivial anomaly that make $`I_3`$ locally Weyl invariant once multiplied by the measure $`\sqrt{g}`$ . Finally the independent six trivial anomalies can be identified by $`M_5`$, $`M_6`$, $`M_7`$, $`M_8`$, $`M_9`$, $`M_{10}`$ as listed in eqs. (1722).
To summarize, a preferred basis for the trace anomalies which takes into account the classification of ref. is given by $`(E_6;I_1,I_2,I_3;M_5,M_6,M_7,M_8,M_9,M_{10})`$. The first four elements make up a basis for the true trace anomalies and in that form have been used in the calculations of ref. .
It may be useful to recall that in the basis of Bonora et al. $`M_4`$ gives the type A anomaly, while $`M_1`$ and $`M_2`$ are type B anomalies. On the other hand, $`M_3`$ contains a spurious contribution from the Euler density and it is not classifiable as the remaining type B anomaly: it is preferable to use $`I_3`$ instead.
Anselmi has introduced in ref. the notion of pondered Euler density $`\stackrel{~}{E}_6`$ by adding a suitable trivial anomaly to $`E_6`$ to make it linear in the conformal factor once evaluated on conformally flat metrics
$`\stackrel{~}{E}_6=E_6+\left({\displaystyle \frac{288}{5}}20\zeta \right)M_6+\left(20\zeta {\displaystyle \frac{408}{5}}\right)M_7+\left({\displaystyle \frac{\zeta }{2}}{\displaystyle \frac{9}{25}}\right)M_9{\displaystyle \frac{24}{5}}M_{10}.`$ (30)
This is an equivalent way of presenting the type A anomaly which may be useful for various applications. Note that $`\zeta `$ labels a 1-parameter family of trivial anomalies. It can be chosen at will showing that the definition of a pondered density is not unique. This construction can be extended to any even dimension .
A further characterization of type A anomalies has been proposed in by studying the AdS/CFT holographic correspondence, while their systematic computation for free models in arbitrary dimensions has been carried out recently in . Finally, it is worth mentioning that CFTs with a special linear relation between the type A and B anomalies, the so-called $`c=a`$ theories, have been identified in as an interesting subclass of conformal theories with special properties.
## 4 Conclusions
We have presented a systematic derivation and classification of trace anomalies in six dimensions. We have solved the consistency conditions and listed the 10 independent solutions as type A ($`E_6`$), type B ($`I_1,I_2,I_3`$) and trivial ($`M_5,M_6,M_7,M_8,M_9,M_{10}`$) trace anomalies. We summarize them by using the $`K_i`$ basis in Table 1.
Our motivation to perform this analysis was to make sure that the identifications of type A and B anomalies made in for various models was correct. As we show explicitly in appendix D that is the case: a spurious term, which doesn’t solve the consistency conditions but enters the basis of total derivatives used to identify trivial anomalies, always drops out in the relevant cases. On the other hand, the calculation of the trace anomalies for a scalar field performed in didn’t produce the same result as in because trivial anomalies were not properly factorized: in the coefficients of $`K_4`$, $`K_5`$, $`K_7`$ and $`K_8`$ were interpreted as the coefficients of true trace anomalies but those structures appear in our list of trivial anomalies, so their coefficients do not have a universal meaning since they can be modified by adding local counterterms to the effective action. Nevertheless, the method of can still be used to compute trace anomalies. By inspecting Table 1 one can recognize the terms that are not corrupted by trivial anomalies: a simple set of 4 independent structures is given by $`K_1`$, $`K_2`$, $`K_3`$, $`K_4K_5`$.
With the present knowledge of trivial anomalies at hand, an interesting investigation could be to study their flows in CFT<sub>6</sub> deformed by relevant operators. This would provide a test in six dimensions of some of the properties studied in ref. .
## Appendix A
We define the 1-cochains
$`𝒲_i^{(a)}={\displaystyle d^6x\sqrt{g}\sigma _a(x)K_i},i=1,\mathrm{},17`$ (31)
where $`\sigma _a`$ denote infinitesimal parameters of Weyl transformations and $`K_i`$ belong to the list of curvature invariants given in eq. (LABEL:inv), and the 2-cochains
$`_\alpha ={\displaystyle d^6x\sqrt{g}H_\alpha },\alpha =1,\mathrm{},9`$ (32)
with the list of $`H_\alpha `$ reported in eq. (9). Now, we compute the variations
$`\mathrm{\Delta }𝒲_i\delta _{\sigma _2}𝒲_i^{(1)}\delta _{\sigma _1}𝒲_i^{(2)}`$ (33)
which can be expanded in the basis of the functionals $`_\alpha `$ . We find:
$$\begin{array}{ccc}\mathrm{\Delta }𝒲_1\hfill & =\hfill & 30_1\hfill \\ \mathrm{\Delta }𝒲_2\hfill & =\hfill & 2_110_28_3\hfill \\ \mathrm{\Delta }𝒲_3\hfill & =\hfill & 8_310_5\hfill \\ \mathrm{\Delta }𝒲_4\hfill & =\hfill & 3_212_4\hfill \\ \mathrm{\Delta }𝒲_5\hfill & =\hfill & 2_2+2_32_48_9\hfill \\ \mathrm{\Delta }𝒲_6\hfill & =\hfill & _1+4_2+4_312_42_5+12_9\hfill \\ \mathrm{\Delta }𝒲_7\hfill & =\hfill & 3_1+12_2+12_324_43_5+24_9\hfill \\ \mathrm{\Delta }𝒲_8\hfill & =\hfill & \frac{3}{4}_1+3_2+3_36_4\frac{3}{4}_5\hfill \\ \mathrm{\Delta }𝒲_9\hfill & =\hfill & 2_110_610_7\hfill \\ \mathrm{\Delta }𝒲_{10}\hfill & =\hfill & _16_24_34_43_6_74_816_9\hfill \\ \mathrm{\Delta }𝒲_{11}\hfill & =\hfill & 4_112_216_3+16_44_52_64_832_9\hfill \\ \mathrm{\Delta }𝒲_{12}\hfill & =\hfill & _15_610_8\hfill \\ \mathrm{\Delta }𝒲_{13}\hfill & =\hfill & _1+6_2+4_3+4_4+3_6_7+8_8+16_9\hfill \\ \mathrm{\Delta }𝒲_{14}\hfill & =\hfill & \frac{1}{2}_1+_22_3+14_4+\frac{5}{2}_6\frac{3}{2}_7+8_8+8_9\hfill \\ \mathrm{\Delta }𝒲_{15}\hfill & =\hfill & 4_1+12_2+16_316_4+4_5+2_62_7+8_8+32_9\hfill \\ \mathrm{\Delta }𝒲_{16}\hfill & =\hfill & 0\hfill \\ \mathrm{\Delta }𝒲_{17}\hfill & =\hfill & 0.\hfill \end{array}$$
(34)
From this list one can read off the matrix $`f^\alpha _i`$ of eq. (10).
## Appendix B
To find all trivial anomalies, we compute the Weyl variation of the most general dimensionless local functional of the metric
$$𝒦=\underset{i=1}{\overset{10}{}}c^i𝒦_i\text{where}𝒦_i=d^6x\sqrt{g}K_i.$$
(35)
Defining as before $`𝒲_i=d^6x\sqrt{g}\sigma (x)K_i`$ we have
$$\begin{array}{ccc}\delta _\sigma 𝒦_1\hfill & =\hfill & 30𝒲_{16}\hfill \\ \delta _\sigma 𝒦_2\hfill & =\hfill & 4𝒲_920𝒲_{10}8𝒲_{12}20𝒲_{13}6𝒲_{16}\hfill \\ \delta _\sigma 𝒦_3\hfill & =\hfill & 4𝒲_920𝒲_{11}8𝒲_{12}20𝒲_{15}4𝒲_{16}\hfill \\ \delta _\sigma 𝒦_4\hfill & =\hfill & 12𝒲_412𝒲_5+3𝒲_96𝒲_{10}12𝒲_{12}6𝒲_{13}12𝒲_{14}\frac{3}{2}𝒲_{16}\hfill \\ \delta _\sigma 𝒦_5\hfill & =\hfill & 10𝒲_410𝒲_54𝒲_6+2𝒲_78𝒲_8\frac{1}{2}𝒲_9+12𝒲_{10}+2𝒲_{11}4𝒲_{12}\hfill \\ & & +20𝒲_{13}18𝒲_{14}+\frac{3}{4}𝒲_{16}\hfill \\ \delta _\sigma 𝒦_6\hfill & =\hfill & 6𝒲_63𝒲_7+12𝒲_8+𝒲_94𝒲_{10}7𝒲_{11}2𝒲_{12}16𝒲_{13}+12𝒲_{14}\hfill \\ & & 4𝒲_{15}\frac{1}{2}𝒲_{16}\hfill \\ \delta _\sigma 𝒦_7\hfill & =\hfill & 12𝒲_66𝒲_7+24𝒲_812𝒲_{11}24𝒲_{13}+24𝒲_{14}6𝒲_{15}\hfill \\ \delta _\sigma 𝒦_8\hfill & =\hfill & 6𝒲_46𝒲_5+6𝒲_{10}\frac{3}{2}𝒲_{11}3𝒲_{12}+6𝒲_{13}6𝒲_{14}\frac{3}{2}𝒲_{15}\hfill \\ \delta _\sigma 𝒦_9\hfill & =\hfill & 2𝒲_{16}20𝒲_{17}\hfill \\ \delta _\sigma 𝒦_{10}\hfill & =\hfill & 20𝒲_420𝒲_58𝒲_6+4𝒲_716𝒲_8+3𝒲_9+4𝒲_{10}+4𝒲_{11}16𝒲_{12}\hfill \\ & & +20𝒲_{13}36𝒲_{14}\frac{3}{2}𝒲_{16}6𝒲_{17}.\hfill \end{array}$$
(36)
Other variations are not necessary since by partial integration we find
$$\begin{array}{ccc}𝒦_{11}\hfill & =\hfill & 4𝒦_44𝒦_5+2𝒦_6𝒦_7+4𝒦_8𝒦_9+4𝒦_{10}\hfill \\ 𝒦_{12}\hfill & =\hfill & \frac{1}{2}𝒦_9\hfill \\ 𝒦_{13}\hfill & =\hfill & 𝒦_{10}\hfill \\ 𝒦_{14}\hfill & =\hfill & 𝒦_4𝒦_5\frac{1}{4}𝒦_9\hfill \\ 𝒦_{15}\hfill & =\hfill & 𝒦_{11}\hfill \\ 𝒦_{16}\hfill & =\hfill & 𝒦_{17}=0.\hfill \end{array}$$
(37)
From the set of Weyl variations written above one can easily constructs the matrix $`g^j_i`$ of eq. (12). This matrix has rank 6 and so there are 6 independent trivial anomalies, which can be identified as the variations of $`𝒦_1`$, $`𝒦_2`$, $`𝒦_3`$, $`𝒦_5`$, $`𝒦_6`$, $`𝒦_9`$.
## Appendix C
Here we solve the consistency conditions in eq. (4). They consist in the following homogeneous system of linear equations (same as eq. (10) of the main text)
$`{\displaystyle \underset{i=1}{\overset{17}{}}}f^\alpha {}_{i}{}^{}a_{}^{i}=0,\alpha =1,\mathrm{},9.`$ (38)
The $`9\times 17`$ matrix $`f^\alpha _i`$ can be constructed form the calculations reported in appendix A. It has rank 7 and so one has 10 independent solutions for the $`a^i`$. A possible choice for the parameters of the solution is: $`a^6`$, $`a^7`$, $`a^8`$, $`a^{10}`$, $`a^{11}`$, $`a^{13}`$, $`a^{14}`$, $`a^{15}`$, $`a^{16}`$, $`a^{17}`$. Thus, one can straightforwardly derive a suitable basis for the anomalies:
$$\begin{array}{ccc}A_1\hfill & =\hfill & \frac{21}{200}K_1+\frac{43}{40}K_2\frac{1}{5}K_3\frac{5}{4}K_4+\frac{3}{2}K_5+K_6\hfill \\ A_2\hfill & =\hfill & \frac{27}{100}K_1+\frac{51}{20}K_2\frac{3}{10}K_3\frac{5}{2}K_4+3K_5+K_7\hfill \\ A_3\hfill & =\hfill & \frac{11}{200}K_1+\frac{9}{20}K_2\frac{3}{40}K_3\frac{1}{2}K_4+K_8\hfill \\ A_4\hfill & =\hfill & \frac{3}{25}K_1K_22K_5\frac{1}{10}K_9+K_{10}\frac{2}{5}K_{12}\hfill \\ A_5\hfill & =\hfill & \frac{8}{25}K_1\frac{13}{5}K_2\frac{2}{5}K_3+2K_44K_5+K_{11}\frac{2}{5}K_{12}\hfill \\ A_6\hfill & =\hfill & \frac{3}{25}K_1+K_2+2K_5\frac{1}{10}K_9+\frac{4}{5}K_{12}+K_{13}\hfill \\ A_7\hfill & =\hfill & K_4+K_5\frac{3}{20}K_9+\frac{4}{5}K_{12}+K_{14}\hfill \\ A_8\hfill & =\hfill & \frac{8}{25}K_1+\frac{13}{5}K_2+\frac{2}{5}K_32K_4+4K_5\frac{1}{5}K_9+\frac{4}{5}K_{12}+K_{15}\hfill \\ A_9\hfill & =\hfill & K_{16}\hfill \\ A_{10}\hfill & =\hfill & K_{17}.\hfill \end{array}$$
(39)
Now, in appendix B we have identified a basis of trivial anomalies. One can use that knowledge to make a change of basis and separate the 6 trivial anomalies from the non-trivial ones. We have chosen the latter to agree with those of ref. obtaining
$$\begin{array}{c}M_1=\frac{3}{4}A_1+A_3,M_2=3A_1+A_2,M_3=5A_45A_5,\hfill \\ M_4=24A_14A_28A_3,\hfill \\ M_5=6A_13A_2+12A_3+A_47A_511A_6+12A_74A_8,\hfill \\ M_6=A_4+A_6,M_7=A_7,M_8=A_8+A_5,M_9=A_9,M_{10}=A_{10}.\hfill \end{array}$$
(40)
This is the basis reported in eqs. (1322). The last six elements $`M_5,\mathrm{},M_{10}`$ are trivial and one can easily identify the local functionals generating them. Defining
$$_i=d^6x\sqrt{g}\sigma (x)M_i,$$
(41)
we have
$$\begin{array}{c}_5=\delta _\sigma \left(\frac{1}{30}𝒦_1\frac{1}{4}𝒦_2+𝒦_6\right),_6=\delta _\sigma \left(\frac{1}{100}𝒦_1\frac{1}{20}𝒦_2\right),\hfill \\ _7=\delta _\sigma \left(\frac{37}{6000}𝒦_1\frac{7}{150}𝒦_2+\frac{1}{75}𝒦_3\frac{1}{10}𝒦_5\frac{1}{15}𝒦_6\right),_8=\delta _\sigma \left(\frac{1}{150}𝒦_1\frac{1}{20}𝒦_3\right),\hfill \\ _9=\delta _\sigma \left(\frac{1}{30}𝒦_1\right),_{10}=\delta _\sigma \left(\frac{1}{300}𝒦_1\frac{1}{20}𝒦_9\right).\hfill \end{array}$$
(42)
## Appendix D
In (and in ) the following basis of invariants was also used
$$\begin{array}{ccc}B_1=^4R\hfill & B_2=(_aR)^2\hfill & B_3=(_aB_{mn})^2\hfill \\ B_4=_aB_{bm}^bB^{am}\hfill & B_5=(_iC_{abmn})^2\hfill & B_6=R^2R\hfill \\ B_7=B_{ab}^2B^{ab}\hfill & B_8=B_{ab}_m^bB^{am}\hfill & B_9=C_{abmn}^2C^{abmn}\hfill \\ B_{10}=R^3\hfill & B_{11}=RB_{ab}^2\hfill & B_{12}=RC_{abmn}^2\hfill \\ B_{13}=B_a{}_{}{}^{m}B_{m}^{}{}_{}{}^{i}B_{i}^{}^a\hfill & B_{14}=B_{ab}B_{mn}C^{ambn}\hfill & B_{15}=B_{ab}C^{amnl}C^b_{mnl}\hfill \\ B_{16}=C_{ab}{}_{}{}^{mn}C_{mn}^{}{}_{}{}^{ij}C_{ij}^{}^{ab}\hfill & B_{17}=C_{ambn}C^{aibj}C^m{}_{i}{}^{}{}_{}{}^{n}_j\hfill & \end{array}$$
(43)
where $`C_{abcd}`$ is the Weyl tensor defined in (28) and $`B_{ab}`$ is the traceless part of the Ricci tensor
$$B_{ab}=R_{ab}\frac{1}{6}Rg_{ab}.$$
(44)
The trivial anomalies in are expressed using the following set of total derivatives
$$\begin{array}{c}C_1=B_1,C_2=B_2+B_6,C_3=B_3+B_7,C_4=B_4+B_8,C_5=B_5+B_9,\hfill \\ C_6=\frac{1}{9}B_2B_4\frac{1}{5}B_{11}\frac{3}{2}B_{13}+B_{14},\hfill \\ C_7=\frac{1}{60}B_2\frac{3}{4}B_3+\frac{3}{4}B_4+\frac{1}{4}B_5+\frac{1}{12}B_{12}+\frac{1}{2}B_{15}\frac{1}{4}B_{16}B_{17}.\hfill \end{array}$$
(45)
However, these total derivatives do not form a subset of the space of trivial anomalies. In fact we find
$$\begin{array}{c}C_1=M_{10},C_2=\frac{1}{2}M_9,C_3=\frac{1}{12}M_9+(K_{10}+K_{13}),\hfill \\ C_4=\frac{7}{6}M_6+M_7\frac{5}{72}M_9+\frac{7}{6}(K_{10}+K_{13}),C_5=M_6+M_8+\frac{1}{20}M_9,\hfill \\ C_6=2M_6M_7+\frac{1}{8}M_92(K_{10}+K_{13}),C_7=\frac{1}{12}M_5\frac{1}{12}M_6\frac{1}{4}M_7+\frac{7}{12}M_8+\frac{1}{32}M_9.\hfill \end{array}$$
(46)
Because of the term $`(K_{10}+K_{13})`$ the space spanned by the $`C_i`$ is larger than the space of trivial anomalies. The latter is expressed by
$`{\displaystyle \underset{i=1}{\overset{7}{}}}a^iC_i\text{is a trivial anomaly iff}6a_3+7a_412a_6=0.`$ (47)
In all the linear combinations of $`C_i`$ satisfy the above relation, i.e. are trivial anomalies. This is a good check on the correctness of the anomaly computations performed there.
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# Resistivity of Mixed-Phase Manganites
\[
## Abstract
The resistivity $`\rho _{dc}`$ of manganites is studied using a random-resistor-network, based on phase-separation between metallic and insulating domains. When percolation occurs, both as chemical composition and temperature vary, results in good agreement with experiments are obtained. Similar conclusions are reached using quantum calculations and microscopic considerations. Above the Curie temperature, it is argued that ferromagnetic clusters should exist in Mn-oxides. Small magnetic fields induce large $`\rho _{dc}`$ changes and a bad-metal state with (disconnected) insulating domains. PACS numbers: 71.10.-w, 75.10.-b, 75.30.Kz
\]
The study of manganites is one of the main areas of research in Strongly Correlated Electrons. Three main reasons have triggered this wide interest: (1) The low-bandwidth materials have unexplained transport properties. They are insulators at room temperature, changing into bad metals at low temperatures. A sharp peak in the resistivity $`\rho _{dc}`$ appears at the ferromagnetic (FM) transition. Small magnetic fields turn the insulator into a metal, with a “Colossal” Magneto-Resistance (CMR). (2) The phase diagram T-x (T=temperature, x=hole density) is rich, with complex spin, charge, and orbital order. (3) Mn-oxides have intrinsic inhomogeneities in most of the T-x plane even in single crystals.
This challenging behavior has been addressed by previous theoretical studies. Regarding item (2), the various FM, antiferromagnetic (AF), orbital-ordered, and charge-ordered (CO) phases have already emerged from simulations and mean-field approximations. Regarding item (3), phase separation (PS) has been proposed to explain the inhomogeneities. PS can be: (a) electronic, with nanometer-size clusters, or (b) structural, where disorder can induce up to micrometer-size clusters and percolation, when influencing on first-order transitions. However, the explanation of transport, item (1), is more complicated since estimations of $`\rho _{dc}`$ are notoriously difficult. In addition, the prominent inhomogeneities of Mn-oxides have not been incorporated into realistic $`\rho _{dc}`$ calculations. The behavior of $`\rho _{dc}`$ remains unexplained, although it is central to manganite physics.
Our goal in this paper is to present a rationalization of the $`\rho _{dc}`$ vs T curves of Mn-oxides based on the currently prevailing phase-separated/percolative framework for these compounds. In this context, items (1) and (3) above are closely related. The analysis necessarily involves phenomenological considerations, since percolation cannot be addressed on sufficiently large lattices using accurate microscopic models. However, the $`\mu `$-meter clusters in experiments strongly suggest that a coarse-grain approach should be sufficient. In addition, results of microscopic calculations presented below are consistent with those of the macroscopic approach.
The main concept introduced here is summarized in Figs.1a-b. The manganite state in the CMR regime is assumed to be percolative, with metallic filaments across the sample (Fig.1a). Percolation indeed occurs in models, and in many experiments. The insulating and metallic (percolative) regions are assumed to have resistances $`\mathrm{R}_\mathrm{I}`$(T) and $`\mathrm{R}_\mathrm{M}^{\mathrm{per}}`$(T), respectively, as sketched in Fig.1b. $`\mathrm{R}_\mathrm{M}^{\mathrm{per}}`$ is large at T=0 due to the complex shape of the conducting paths, and grows with T as in any metal, eventually diverging when the percolative path melts with increasing T. Note that at room temperature $`\mathrm{R}_\mathrm{I}`$$`<`$$`\mathrm{R}_\mathrm{M}^{\mathrm{per}}`$ and, thus, most of the conduction in this regime occurs through the insulator. On the other hand, $`\mathrm{R}_\mathrm{I}`$ is so large at low-T that current can only flow through the percolative paths. This suggests a simple two parallel resistances description (Fig.1b), where a peak in the effective resistance at intermediate T is natural.
To substantiate this idea, first consider results obtained using a random-resistor-network that mimics the prominent FM-CO mixtures in Mn-oxides in the CMR regime. Two-dimensional (2D) and three-dimensional (3D) square and cubic clusters are used, with link resistivities randomly selected as metallic ($`\rho _M`$) or insulating ($`\rho _I`$), with a fixed metallic fraction $`p`$ (which in, e.g., $`(\mathrm{La}_{5/8\mathrm{y}}\mathrm{Pr}_\mathrm{y})\mathrm{Ca}_{3/8}\mathrm{MnO}_3`$ (LPCMO) it is proportional to the amount of La ). The lattice spacing of this effective network is comparable to the FM or CO domain size, much larger than the Mn-Mn distance. The actual values of $`\rho _M`$ and $`\rho _I`$ vs T were directly taken from LPCMO data (y=0.00 and 0.42, respectively), and for simplicity they are used in both 3D and 2D clusters. Other materials were tried and the analysis below does not depend qualitatively on the reference compounds.
The Kirchoff equations for the network were solved iteratively on large clusters using well-known techniques, and the net resistivity was found. Typical results are shown in Fig.1c. Only the limiting cases $`p`$=0 (all insulator) and 1 (all metal) are taken from experiments. As expected, a percolative regime exists between $`p`$=0.4 and 0.5, where $`\rho _{dc}`$(T=0) is as large as in LPCMO and other materials. $`\rho _{dc}`$ has insulating behavior at room-T, even for $`p`$ as high as 0.65, while at low-T a (bad) metallic behavior is observed. A broad peak appears at intermediate T’s and $`p`$’s. Similar results exist in 3D (inset of Fig.1c).
It is remarkable that Fig.1c is already in good qualitative agreement with some Mn-oxide experiments, such as for $`\mathrm{La}_{0.96\mathrm{y}}\mathrm{Nd}_\mathrm{y}\mathrm{K}_{0.04}\mathrm{MnO}_3`$ , or even non-manganite materials, such as $`\mathrm{CaFe}_{1\mathrm{x}}\mathrm{Co}_\mathrm{x}\mathrm{O}_3`$ where an AF-FM competition occurs. However, many manganites present a more pronounced $`\rho _{dc}`$ peak at intermediate T’s. To reproduce this feature, it is necessary to introduce a percolative process not only as $`p`$ (or x) varies, but also as T changes. This is reasonable since the metallic component triggered by ferromagnetism is sensitive to T, and the FM clusters shrink in size as T increases. This proposal was tested qualitatively using two models: (i) the Random Field Ising model (RFIM), that describes the disorder-induced PS, with spin up and down crudely representing the competing metal and insulator, and (ii) the well-known one-orbital model (with parameters t, $`\mathrm{J}_\mathrm{H}`$, and $`\mathrm{J}^{}`$ representing the $`\mathrm{e}_\mathrm{g}`$-hopping, Hund coupling, and Heisenberg exchange among $`\mathrm{t}_{2\mathrm{g}}`$-spins, respectively). The latter is supplemented by a term $`_i`$$`\varphi _i`$n<sub>i</sub>, with $`\varphi _i`$ randomly taken from \[-W,W\] and n<sub>i</sub> the on-site density at site $`i`$. This disorder generates coexisting clusters near first-order FM-AF transitions, as explained in Ref..
In Fig.2a, it is shown a portion of a typical Monte Carlo (MC) simulation of the RFIM at T=1.6 on a 500$`\times `$500 cluster, with fixed random fields taken from \[-1.0,1.0\] (J=1 is the FM Ising coupling). These parameters are the same as in Ref.. Three main domains were found: spins-up (black), spins-down (white), and regions with a small spin expectation value (grey). The generation with T of these paramagnetic (PM) areas in the surface of the up and down domains weakens the percolative tendencies of the RFIM: for example, in Fig.2a domains connected at T=0, become disconnected at finite-T. Similar behavior occurs in the microscopic one-orbital case, as shown in Fig.2b for a disordered configuration with percolative features. Increasing T from $``$0.0 to 0.05t decouples the two FM regions. FM, AF, and PM regimes dominate at finite T, as in the RFIM. The T=0.05t Drude weight (not shown) is much smaller than at T$``$0.0.
To incorporate the indications of T-induced percolation (Figs.2a-b) in the phenomenological approach, a T-dependent metallic fraction $`p(\mathrm{T})`$ is needed. $`p`$ must decrease as T grows, should vary rapidly near the Curie temperature $`\mathrm{T}_\mathrm{C}`$ as the magnetization does, but otherwise its T-dependence is unknown. Fortunately, the qualitative results using several functions are similar, and a typical case is shown in Fig.2c. The $`\rho _{dc}`$’s obtained by this simple procedure now clearly resemble those found in experiments, with a robust peak at intermediate T’s, moving toward lower T’s as the system becomes more insulating. This agreement with experiments is unlikely to be accidental, and justifies a posteriori our assumptions. Note that if our approach is correct, consistency requires that in Mn-oxides above $`\mathrm{T}_\mathrm{C}`$ there should exist (disconnected) FM clusters on an insulating matrix, since $`p`$ does not drop abruptly to 0 at $`\mathrm{T}_\mathrm{C}`$. A new temperature scale T is predicted, with those FM clusters existing in the range $`\mathrm{T}_\mathrm{C}`$$``$T$``$T. The density of states in this regime likely has a $`pseudogap`$, based on previous investigations of mixed-phase states.
Consider now nanometer-scale clusters. Here quantum effects cannot be neglected. However, the problem is still too difficult to be treated microscopically, and an effective description is needed. For this purpose, instead of a resistor network, a 3D lattice model with NN electron hopping (and zero chemical potential) is here used, with link hopping amplitudes randomly selected to be either “metallic” ($`\mathrm{t}_\mathrm{M}`$) or “insulating” ($`\mathrm{t}_\mathrm{I}`$), representing effective hoppings through the nanoclusters with 1 (0) corresponding to the FM (AF) regions of the microscopic model at large $`\mathrm{J}_\mathrm{H}`$. Such “lattice of quantum wires” has been used before to study quantum percolation. The cluster conductance $`C`$ (in $`\mathrm{e}^2/\mathrm{h}`$ units) is calculated using the Kubo formula within a Landauer setup, connecting the finite clusters to semi-infinite leads using an infinitesimal voltage drop. The self-energy of these leads is known exactly, and the lead on the, e.g., right is used for an iterative calculation of the self-energy from right to left through the cluster, supplemented by a self-energy matching at the left end. With the Green function obtained by this procedure, $`C`$ is evaluated using known formulas. Since this is $`quantum`$ percolation, the study is only in 3D (2D localization leads to zero conductivity).
The hoppings $`\mathrm{t}_\mathrm{M}`$ and $`\mathrm{t}_\mathrm{I}`$ are not available from experiments. However, $`\mathrm{t}_\mathrm{M}`$ should decrease with increasing T following the FM-phase magnetization, while $`\mathrm{t}_\mathrm{I}`$ increases with T (since, e.g., the zero-conductivity T=0 AF configuration disorders as T grows). $`C`$ was obtained on up to 20<sup>3</sup> clusters using the $`\mathrm{t}_\mathrm{M}`$(T) and $`\mathrm{t}_\mathrm{I}`$(T) in the inset of Fig.3a, but the results do not depend qualitatively on the particular functions used, as long as $`\mathrm{t}_\mathrm{I}`$ changes rapidly with T near room-T, as $`\rho _I`$ does in experiments. The metallic fraction $`p`$ in Fig.3a was made T-dependent as in Fig.2c, and the critical percolation at T=0 is expected to be located near $`p_c`$$``$0.45 . Results are shown in Fig.3a. Once again a remarkable qualitative agreement with experiments is obtained, suggesting that both nano- and micro-meter clusters lead to similar results.
For completeness, $`C`$ was also calculated using microscopic models on small lattices. A MC simulation of the one-orbital model on an 8$`\times `$8 PBC cluster and density x=0.5 was performed. From previous work, it is known that a metal-insulator first-order transition occurs at $`\mathrm{J}_\mathrm{c}^{}`$$``$0.07 (if $`\mathrm{J}_\mathrm{H}`$=$`\mathrm{}`$). Disorder is introduced such that in the NN-sites link $``$ij$``$ the hopping is $`\mathrm{t}_{\mathrm{ij}}`$=1+$`\delta _{ij}`$ and Heisenberg coupling is $`\mathrm{J}_{\mathrm{ij}}^{}`$=$`\mathrm{J}^{}`$(1+$`\delta _{ij}`$), where $`\delta _{ij}`$ is randomly taken from \[-$`\mathrm{\Delta }`$,$`\mathrm{\Delta }`$\], and $`\mathrm{J}^{}`$ is uniform. This disorder makes the transition continuous. The MC procedure generates $`\mathrm{t}_{2\mathrm{g}}`$-spin configurations from which NN-sites effective hoppings can be calculated (as in double exchange models). These hoppings are used to evaluate $`C`$ (Fig.3b). Due to the disorder, $`C`$ interpolates smoothly from metal to insulator varying $`\mathrm{J}^{}`$ (otherwise a discontinuous transition occurs). $`C^1`$(T$``$0) can be very large, but finite, if the appropriate value of $`\mathrm{J}^{}`$ is selected. In this respect the result has clear similarities with those of the macroscopic approach. However, the full shape of the experimental $`\rho _{dc}`$ curves is difficult to reproduce with microscopic models on small clusters where percolation cannot be studied. Nevertheless, for the (few) disorder configurations with percolative-like characteristics found on small systems (as in Fig.2b), the associated $`C`$ vs T has a broad maximum at a finite T.
Consider now nonzero magnetic fields ($`h`$). In the random-network model mimicking coexisting FM-AF regions, a small $`h`$ will increase the FM fraction $`p`$ by a concomitant small amount. However, near percolation tiny modifications in $`p`$ can induce large $`\rho _{dc}`$ changes, as shown in Fig.2b where a 5% modification in $`p`$ at low-T can alter $`\rho _{dc}`$ by two orders of magnitude. In the percolative regime, “small” perturbations can drastically change the conductivity. This analysis predicts that the metallic state reached from the insulator with magnetic fields is $`not`$ homogeneous but must still have a substantial fraction of insulating clusters. This is consistent with the experimental large $`\rho _{dc}`$(T=0) of such a state.
Another effect contributes to this phenomenon. It exists even on chains where percolation does not occur, and it is illustrated in Fig.4a where $`C^1`$ is shown using the microscopic half-doped one-orbital model at $`\mathrm{J}^{}`$=0.14 where the system is at a FM(metal)-AF(insulator) transition, the latter with the periodic spin structure $``$, as shown in Ref.. The field is introduced as $`h`$$`_\mathrm{i}\mathrm{M}_\mathrm{i}^\mathrm{z}`$, where $`\mathrm{M}_\mathrm{i}^\mathrm{z}`$=$`\mathrm{s}_\mathrm{i}^\mathrm{z}+(3/2)\mathrm{S}_\mathrm{i}^\mathrm{z}`$, with $`\mathrm{S}_\mathrm{i}^\mathrm{z}`$ the $`\mathrm{z}`$-component of the classical $`\mathrm{t}_{2\mathrm{g}}`$-spin at site $`\mathrm{i}`$ with norm 1, and $`\mathrm{s}_\mathrm{i}^\mathrm{z}`$ the spin of the mobile electron at the same site. Disorder in the hopping and $`\mathrm{J}^{}`$ (as in Fig.3b) of strength $`\mathrm{\Delta }`$=0.03 produces coexisting FM-AF clusters. The nearly perfect AF-links at low-T induce a huge $`C^1`$ at $`h`$=0. However, field modifications of just 0.01t ($``$17 Tesla, if $`\mathrm{t}`$=0.2eV) produce dramatic changes in $`C^1`$ at low-T (Fig.4a). The resulting $`C^1`$ at $`h`$$``$0.0 is still large, but compared with $`C^1`$($`h`$=0) the effect is notorious. An analysis of the spin correlations vs $`h`$ shows that these large resistance changes mainly originate in the $`AF`$ regions, since small fields produce a small spin canting and concomitant small conductivity, creating a $`valve`$ effect between metallic domains. In real manganites, a relatively modest $`\rho _{dc}`$ change in the insulating regions could contribute appreciably to the large MR. To simulate this effect, $`\rho _I`$ of the 2D network of Fig.2c was slightly modified ($`\mathrm{\Delta }`$$`\sigma _{dc}`$=$`0.1`$$`(\mathrm{\Omega }cm)^1`$), with $`\rho _M`$ untouched. The resulting $`\rho _{dc}`$ changes (Fig.4b) are indeed large at low-T, comparable to those obtained changing $`p`$ by a few percent.
Summarizing, the manganite $`\rho _{dc}`$ was studied within the PS-framework using a semi-phenomenological approach. At room-T, conduction predominantly through the insulating regions leads to d$`\rho _{dc}`$/dT$`<`$0, while at low-T the metallic filaments carry the current. The large MR induced by magnetic fields is caused by small changes in the metallic fraction $`p`$ and/or in the conductivity of the insulator, effects which severely affect transport near percolation. Our approach provides a simple explanation of the CMR effect, without invoking polaronic or Anderson localization concepts, and independently of the origin (Coulomb vs Jahn-Teller) of the competing phases.
The authors thank S. Yunoki for the programs used for Figs.2b-3b-4a, and NSF (DMR-9814350), Comision Interministerial de Ciencia y Tecnologia de España (PB96-0085), MARTECH, CDCH (Univ. Central de Venezuela), and Fundación Antorchas for support.
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# Noise-driven evolution in stellar systems: Theory
## 1 Introduction
Both numerical experiments and observations suggest a universality in galaxy profiles resulting from dissipationless collapse. Remnants from merger simulations (e.g. van Albada 1982, Barnes 1989) typically manifest the often observed $`r^{1/4}`$-type profile in elliptical galaxies. More recently, the same feature has arisen in the context of cosmological large-scale structure simulations. Researchers beginning with Navarro, Frenk & White (1997) have noted that CDM collapse results in a haloes with a similar overall profile. Debate remains about the precise functional form of these profiles, but these and the $`r^{1/4}`$-law are all quite similar, suggesting that some general violent-relaxation-like dynamical mechanism is at work.
A precise explanation of the dynamics of violent relaxation, the process driving the convergent evolution in by large fluctuations in potential, remains elusive. There have been several attacks on this problem. Tremaine, Henon & Lynden-Bell (1986) follow a statistical mechanics approach and explore extremizing functionals (see their paper for a review of prior work). Spergel & Hernquist (1992) extended this approach with the additional ansatz that the orbit–perturbation interaction can be approximated by kicks near perigalacticon. Recently, this has been followed by Mangalam, Nityananda & Sridhar (1999) who incorporate these physical mechanisms in a semi-analytic evolution equation based on diffusion in action space. Over the same period, Kandrup has pursued a mathematically rigorous exploration into the geometry of phase space resulting from the details of Hamiltonian flows (e.g. see Kandrup 1998 and references therein).
In this paper, I take a different approach to the problem. Rather than trying to identify an appropriate form for the phase-space distribution to start, I treat the problem as an initial value problem for a stochastic process and develop an evolution equation for exploring its consequence on stellar systems. A companion paper (Paper 2) will apply this to the evolution of galactic haloes in particular. The underlying motivation is as follows. Previous work by the author has shown that fluctuations in stellar systems on the largest scales can be strongly amplified by their own self gravity (e.g. Weinberg 1993). This means that large-scale fluctuations will greatly exceed their Poisson amplitudes. In particular, Weinberg (1993) explored the idealized case of periodic cube; the fluctuations become very large as the system size approaches its Jeans’ length. Similarly for stable galaxies, the fluctuations in a system will be largest at its discrete modes. In addition, Weinberg (1994) argues that galaxies will often have very weakly damped $`m=1`$ (sloshing or seich) modes and these result in large excitations when excited (see Vesperini & Weinberg 2000 for another example). Putting this together, one might ask: if noise preferentially excites particular modes independent of the noise source, is it possible the repetitive stochastic response of the galaxy will lead to some characteristic features, independent of its initial conditions?
In order to answer this question, this paper concentrates on the theoretical framework for describing the evolution due to stochastic fluctuations. Beginning with a description of the linear response of a galaxy to excitation, and assuming that the process is Markovian, one may expand the Boltzmann collision term in a series. Analogous to two-body collisions, only the first two terms contribute and the resulting evolution equation has a Fokker-Planck form. This equation is very far from being analytically tractable because the diffusion coefficients depend on integrals over all of phase space under the stochastic perturbation. It does not solve the classic violent relaxation problem because this approach is perturbative and assumes that the stellar system remains near an equilibrium, Nonetheless it does incorporate the same underlying processes, the self-gravitating response to noise, and yields profiles with many of the same features that are found in the numerical experiments as Paper 2 will demonstrate.
The plan for this paper is as follows. The overall approach is outlined in §2.1 followed by an explicit derivation for spherical haloes in §2.2. This could be straightforwardly extended to many standard geometries geometry (e.g. disk or disk and halo together). We describe the character of the noise for two general cases, transient and quasi-periodic perturbers (a halo of black holes, for example) in §3. These two cases result in qualitatively different behavior and represent the most plausible astronomical scenarios. Transient noise is probably the most relevant and important. Finally, the main features are summarized in §4. Paper 2 will review the basic physics and apply these methods to understanding the evolution of halo in a noisy environment, e.g. just after formation, and may be a better place to begin for those interested in astronomical consequences rather than the kinetic theory.
## 2 Derivation of the evolution equation
### 2.1 Overview
The outer halo in large galaxies is less than 10 dynamical times old so primordial inhomogeneities will not have had time to phase mix and continued disturbances from mergers will not have relaxed (see Tremaine 1992 for additional discussion). These, as well as any intrinsic sources of noise, e.g. a population of $`10^6\mathrm{M}_{}`$ massive black holes, dwarf galaxies, debris streams (Johnston 1998, Morrison et al. 2000) or dark clusters, are amplified by the self-gravity of the halo. In the current epoch, these distortions create potentially observable asymmetries in the stellar and gaseous Galactic disk. Just after galaxy formation, this noise may be sufficient to drive the evolution of the halo (Paper 2). Other possible applications include the evolution of proto-stellar clusters and the evolution of the proto-stellar binary distribution in the noisy star forming environment (work in progress). All of these problems motivated the development of a stellar dynamical kinetic equation that can handle general stochastic processes.
The general problem will be familiar to most dynamicists and could be solved following the standard two-body approach: beginning with the collisional Boltzmann equation, one writes the collision term in Master equation form and expands in a Taylor series to derive Fokker-Planck equation (e.g. Binney & Tremaine 1987, Spitzer 1987). For a spherically symmetric system, the phase-space distribution is a function of two action and and two angle variables. Averaging over times short compared to the relaxation time but long compared to the dynamical time (orbit-averaging), one obtains a Fokker-Planck equation.
Alternatively, recent work in statistical mechanics and noise theory in general has developed a body of methods for treating stochastic differential equations based directly on transition probabilities. If $`P(x^{},t+\tau |x,t)`$ is the transition probability to some new state $`x^{}`$ at time $`t+\tau `$ from the initial state $`x`$ at time $`t`$, then the following integral equation determines all subsequent evolution of the distribution $`f(x)`$:
$$f(x,t+\tau )=𝑑x^{}P(x,t+\tau |x^{},t)f(x^{},t).$$
By expanding the transition probability in its moments of $`xx^{}`$ for small $`\tau `$, one may derive a differential equation of the form:
$$\frac{f(x,t)}{t}=\underset{n=1}{\overset{\mathrm{}}{}}\left(\frac{}{x}\right)^nD^{(n)}(x,t)f(x,t)$$
(1)
where the coefficients $`D^{(n)}`$ are proportional to the time-derivative of the moments in transition probability. This is known as the Kramers-Moyal expansion.
For galaxy evolution, the state variable $`x`$ in the equations above is replaced with the phase-space six vector. Since we are interested in long-term evolution, we may orbit average the evolution equation. Writing the phase space as actions and angles turns the orbit average into an angle average. For any given noise process, we can solve for the change in actions of any orbit in the galaxy using a perturbative approach and similarly derive the change in the phase-space distribution function as described in Weinberg (1998). Evaluation of these quantities lead directly to the moments needed for the Fokker-Planck-type evolution equation. The most subtle aspect of the development below is enforcing a consistent time ordering. Implicit in the angle averaging is a short dynamical time scale and a long evolutionary time scale. The stochastic perturbations are on short time scales and instantaneous from the evolutionary point of view. Although this approximation may seem restrictive and only marginally true for some scenarios of interest, it yields results in good agreement with n-body simulation in the case of globular cluster evolution and the closely related case of self-gravitating fluctuations explored in Weinberg (1998).
After deriving the evolution equation below, we work out the details of the Fokker-Planck coefficients for two examples in §3: transient distortions such as fly-by encounters and quasi-periodic distortions such as orbiting super-massive black holes. The consequences of both of these will be explored in Paper 2.
### 2.2 Derivation of the Fokker-Planck equation
The main advantage of the Kramers-Moyal expansion is that the noise process appears explicitly in an initial-value form. Because the response of the galaxy to a stochastic event can be straightforwardly computed in perturbation theory or by n-body simulation, we can compute the noise-driven evolution for a wide variety of scenarios. For this reason, it is better suited to treating general stochastic noise than the Master equation, even though the two approaches are formally equivalent.
Although the Kramers-Moyal expansion has an infinite number of terms in general (cf. eq. 1), the Pawula Theorem (Pawula 1967) shows that consistency demands that the expansion either stops after two terms and takes the standard Fokker-Planck form or must have an infinite number of terms. For the stellar dynamical case considered here, the series truncates after two terms. The approach approach sketched here is clearly described by Risken (1989, Chap. 4).
As outlined in §2.1, one begins the derivation with the definition of transition probability:
$$f(𝐈,t+\tau )=𝑑𝐈^{}P(𝐈,t+\tau |𝐈^{},t)f(𝐈^{},t)$$
(2)
where an average over the rapidly oscillating angles is implied and $`P`$ is the conditional probability that a state has $`𝐈`$ at time $`t+\tau `$ if it has $`𝐈^{}`$ at time $`t`$ initially. A Taylor series expansion of the integrand in $`\mathrm{\Delta }𝐈^{}𝐈`$ followed by a change of variables and integration over $`\mathrm{\Delta }`$. In the limit $`\tau 0`$ this expansion leads directly to
$$\frac{f(𝐈,t+\tau )}{t}=\underset{n=1}{\overset{\mathrm{}}{}}\left(\frac{}{𝐈}\right)^nD^{(n)}(𝐈,t)f(𝐈,t).$$
(3)
Note that $`P`$ is the probability of a change in $`𝐈`$ due to stochastic events. Therefore, the formal time derivatives in the expansion and phase space integral is better considered as the limit for small $`\tau `$ (but for $`\tau `$ greater than a dynamical time) of the ensemble average of stochastic events. If we let $`\xi `$ be the stochastic value of $`𝐈`$, the expansion coefficients describing the stochastic variables $`\xi `$ is:
$$D^{(n)}(x,t)=\frac{1}{n!}\underset{\tau 0}{lim}\frac{1}{\tau }[\xi (t+\tau )𝐈]^n|_{\xi (t)=𝐈}.$$
(4)
In $`N`$ dimensions, this becomes:
$$D_{j_1,j_2,\mathrm{},j_n}^{(n)}(𝐈,t)=\frac{1}{n!}\underset{\tau 0}{lim}\frac{1}{\tau }M_{j_1,j_2,\mathrm{},j_n}^{(n)}(𝐈,t,\tau )$$
(5)
where the moments $`M^{(n)}`$ are
$$M_{j_1,j_2,\mathrm{},j_n}^{(n)}(𝐱,t,\tau )=(\overline{I}_{j_1}I_{j_1})(\overline{I}_{j_2}I_{j_2})\mathrm{}(\overline{I}_{j_n}I_{j_n}).$$
(6)
The angle brackets denote the integration over the conditional probability of obtaining the state variable $`\overline{I}_j`$ at time $`t+\tau `$ given $`I_j`$ at time $`t`$.
We will see below that the response function guarantees the that the Kramers-Moyal expansion terminates after two terms in the limit of weak perturbations and will assume this here. This is consistent with our intuition that repetitive weak, stochastic excitation in a galaxy is a Markov process. The evolution equation is then the Fokker-Planck equation in standard form:
$$\frac{f(𝐈,t)}{t}=𝐋_{FP}(𝐈,t)f(𝐈,t)$$
(7)
where
$$𝐋_{FP}(𝐈,t)=\frac{}{I_i}D_i^{(1)}(𝐈,t)+\frac{^2}{I_iI_j}D_{ij}^{(2)}(𝐈,t).$$
(8)
We see that the natural coordinates for the Boltzmann equation are action-angle variables. For a collisionless equilibrium, the actions are constant and the angles advance at constant rate. In deriving this Fokker-Planck equation, one averages over orbital periods in favor of following the evolution over longer time scales. The distribution function is then a function of actions alone, $`f=f(𝐈)`$.
However, we will concentrate on the evolution of haloes modeled as a collisionless spherical distribution for application in Paper 2 and adopt the traditional $`E,J,J_z`$ or $`E,\kappa J/J_{max}(E),\mathrm{cos}\beta J_z/J`$ phase-space variables for computational purposes. In addition, we will not consider any processes with a preferred axis so we can average equation (7) over $`\beta `$ with no loss of information. If we are willing to restrict ourselves to an isotropic distribution, we can average over $`\kappa `$ to yield a $`1+1`$ dimensional Fokker-Planck equation in time and energy, $`E`$; this is described in §2.3 below. Because the diffusion coefficients depend on the distribution function, the Fokker-Planck equation in (7) is non-linear, just as for globular cluster evolution. However, it is straightforward to solve this numerically by iteration.
To derive the coefficients $`D^{(1)}`$ and $`D^{(2)}`$, we will use the method described in Weinberg (1998, hereafter Paper 1). To summarize, we represent distortions in the structure of halo in a biorthogonal basis. Any distortion can be summarized then by a set of coefficients in three indices. Because large spatial scales are most important in understanding global evolution, we can truncate this expansion and still recover most of the power. Moreover, we can analytically compute the self-gravitating response of the halo to some arbitrary perturbation as previously described. This development gives us $`\xi (t)`$ for all phase-space variables (cf. eq. 4) and the appropriate ensemble averages give us the required diffusion coefficients. For example, if one uses the same biorthogonal basis in an n-body simulation of a desired transient process, the time series of coefficients can be used directly to derive the coefficients $`D^{(1)}`$ and $`D^{(2)}`$ after removing the time-invariant (DC) component which corresponds to the equilibrium background. We will consider point mass perturbers in §3.1 and transient perturbers (dwarf galaxies, decaying substructure, spreading debris trails) in §3.2.
### 2.3 Averaged Fokker-Planck equation
A number of authors have described transformation of the multivariate Fokker-Planck equation (Rosenbluth et al. 1957, Risken 1989). The approach is the familiar one: write the equation in terms of scalars, covariant and contravariant vectors and tensors and covariant derivatives only. In the first case, the authors use the Jacobian of the coordinate transformation as a metric and in the second, the authors use the diffusion matrix. We will use the first case here. Denote the Jacobian of the coordinate transformation as $`J`$. Under a change of coordinates, one can show after a fair bit of algebra that the advection and diffusion terms transform as
$`D_k^{}`$ $`=`$ $`{\displaystyle \frac{I_k^{}}{I_i}}D_i+{\displaystyle \frac{^2I_k^{}}{I_iI_j}}D_{ij},`$ (9)
$`D_{kl}^{}`$ $`=`$ $`{\displaystyle \frac{I_k^{}}{I_i}}{\displaystyle \frac{I_l^{}}{I_j}}D_{ij}.`$ (10)
The phase-space distribution function transforms as $`f^{}(𝐈)=Jf(𝐈)`$ (cf. Risken 1989) and in the new variables, the equation takes the standard Fokker-Planck form:
$$\frac{f^{}(𝐈,t)}{t}=\left[\frac{}{I_k^{}}D_k^{}+\frac{^2}{I_k^{}I_l^{}}D_{kl}^{}\right]f^{}(𝐈,t).$$
(11)
Now let $`𝐈^{}=(E,\kappa ,\mathrm{cos}\beta )`$. Assuming that the distribution function $`f`$ is time-independent and non-zero, we may integrate equation (11) over $`\kappa `$ and $`\mathrm{cos}\beta `$. Since both of these variables have a bounded domain, the flux through their boundaries must vanish, leaving a single flux term:
$$\frac{f^{}}{t}=\frac{}{E}D_Ef^{}+\frac{}{I^j}\left(D_{Ej}f^{}\right)_{iso}$$
(12)
where the angle brackets denote integration over $`\kappa `$ and $`\mathrm{cos}\beta `$ and sum over $`j`$ denote the sum over all three variables. The isotropically averaged Fokker-Planck equation is then
$$\frac{\overline{f}(E)}{t}=\frac{}{E}\left[D_E_{iso}\overline{f}(E)+\frac{}{E}\left(D_{EE}_{iso}\overline{f}(E)\right)\right]$$
(13)
where $`D_E_{iso}`$ and $`D_{EE}_{iso}`$ are the isotropically averaged diffusion coefficients:
$`{\displaystyle \genfrac{}{}{0pt}{}{D_E}{D_{EE}}}`$ $`=`$ $`{\displaystyle \frac{J_{max}^2(E)}{f(E)}}{\displaystyle 𝑑\kappa d(\mathrm{cos}\beta )\kappa \left\{\genfrac{}{}{0pt}{}{D_E}{D_{EE}}\right\}f(E,\kappa ,\beta )}`$ (14)
where
$`\overline{f}(E)`$ $`=`$ $`J_{max}^2(E){\displaystyle 𝑑\kappa d(\mathrm{cos}\beta )\kappa f(E,\kappa ,\beta )}`$ (15)
and the phase-space volume is
$`P(E)`$ $``$ $`J_{max}^2(E){\displaystyle 𝑑\kappa d(\mathrm{cos}\beta )\kappa }.`$ (16)
Note that the standard notation in the globular cluster literature is $`f(E)=\overline{f}(E)/P(E)`$.
## 3 Noise models
As described in §2, we represent distortions in the structure of halo in a biorthogonal basis. Any distortion can then be summarized by a set of coefficients. Because large spatial scales are most important in understanding global evolution, we can truncate this expansion and still recover most of the power. Internal and therefore quasi-periodic distortions contribute at a discrete spectrum of frequencies. Paper 1, §2 (see eqns. 21-22 in Paper 1 for the final result) derives the response of a halo to a point perturbation at a single frequency. Similar arguments lead to an expression for a continuous spectrum of perturbation frequencies. In the latter case, one computes the response of the stellar system to each frequency in the spectrum and then sums over all frequencies. We will begin with the development common to both cases.
The goal is calculation of the coefficients defined by equations (4) and (5). We begin by determining these coefficients for action variables and transform to $`(E,\kappa ,\mathrm{cos}\beta )`$ in the end. Because orbits in the equilibrium phase space are quasi-periodic and representable as fixed actions and constantly advancing angles, any perturbed quantity can be represented as a Fourier series in angles with coefficients depending on actions. Following Paper 1, the perturbed Hamiltonian is
$`H(𝐈,𝐰)`$ $`=`$ $`H_o(𝐈)+H_1(𝐈,𝐰)`$ (17)
$`=`$ $`H_o(𝐈)+{\displaystyle \underset{𝐥}{}}H_{1𝐥}(𝐈)e^{𝐥𝐰}`$ (18)
where
$`H_{1𝐥}(𝐈)`$ $`=`$ $`{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{m=l}{\overset{l}{}}}{\displaystyle \underset{j}{}}Y_{ll_2}(\pi /2,0)r_{l_2m}^l(\beta )W_{ll_2m}^{l_1j}(𝐈)a_j^{lm}(t)`$
$`=`$ $`{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{m=l}{\overset{l}{}}}{\displaystyle \underset{j}{}}Y_{ll_2}(\pi /2,0)r_{l_2m}^l(\beta )W_{ll_2m}^{l_1j}(𝐈){\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\omega e^{i\omega t}{\displaystyle \underset{k}{}}\left[_{jk}^{lm}(\omega )+\delta _{jk}\right]b_k^{lm}(\omega )`$
where $`𝐥=l_1,l_2,l_3`$ is a triple of integers, $`r_{ij}^l(\beta )`$ and $`W_{ll_2l_3}^{l_1j}(𝐈)`$ are the rotation matrices and gravitational potential transforms defined in Paper 1. The time dependence of the coefficients describing the response, $`a_j^{lm}(t)`$, is represented as its Fourier transform. This allows each frequency to be treated separately. The response matrix $``$ describes the reaction of the galaxy to the perturbation; so the entire response is the sum of both the response and direct forcing, $`_{jk}^{lm}+\delta _{jk}`$. We may integrate the equations of motion directly to evaluate $`𝐈(\tau +t)`$. Hamilton’s equations yield
$$\dot{I}_j=\frac{H}{w_j}=i\underset{𝐥}{}l_jH_{1𝐥}(𝐈)e^{𝐥𝐰}$$
(21)
and therefore we have
$$\mathrm{\Delta }I_j(t+\tau )I_j(t+\tau )I_j(t)=_t^{t+\tau }𝑑t\dot{I}_j(t).$$
(22)
The evolution of perturbed distribution function in time follows from the linearized collisionless Boltzmann equation and the total time derivative for a Hamiltonian system:
$$\dot{f}_1\frac{f_1}{t}+[f_1,H]=\frac{f_1}{t}+\frac{H_0}{𝐈}\frac{f_1}{𝐰}=\frac{H_1}{𝐰}\frac{f_0}{𝐈}.$$
(23)
Analogous to the development above for $`H_1(t)`$, we have
$$\dot{f}_1(𝐈,𝐰,t)=\underset{𝐥}{}e^{i𝐥𝐰}i𝐥\frac{f_0}{𝐈}H_{1𝐥}(𝐈,t)$$
(24)
and therefore
$$f_1(𝐈,𝐰,t+\tau )=\underset{𝐥}{}e^{i𝐥𝐰}i𝐥\frac{f_0}{𝐈}_t^{t+\tau }𝑑t^{}H_{1𝐥}(𝐈,t^{})$$
(25)
To evaluate equation (5) we need the first and second order action moments defined by equation (6). We assume, by adopting the time-asymptotic response matrix $`^{lm}`$ in deriving that $`f_1`$ and $`H_1`$ above, that $`\tau `$ is larger than intrinsic dynamical times, consistent with the ordering of our slow and fast time scales. Previous work, including comparison to n-body simulations, suggests that this is a very good approximation for time scales longer than several crossing times. The response to the distortion induces a shift in the actions $`𝐈`$ and the overall response causes a change in the distribution function. This is represented in the matrix equation defined by equation (3) for an external perturbation described by the coefficients $`b_j^{lm}(t)`$. The fiducial stochastic variables are the coefficients $`b_j^{lm}(t)`$ themselves.
The overall conditional probability required in equation (2) and the following development for the moments therefore has two contributions. First, the response of the galaxy changes the underlying distribution and consequently the probability of obtaining a given final state. Second, the resonant coupling changes the action of an orbit at a particular point in phase space. Altogether we have
$$P(𝐈^{},t+\tau |𝐈,t)=\left(1+\frac{f_1(𝐈,t)}{f_0(𝐈)}\right)\delta \left(𝐈^{}𝐈_t^{t+\tau }𝑑t^{}\mathrm{\Delta }\dot{𝐈}(t^{})\right).$$
(26)
The first- and second-order moments are proportional to the square of the perturbation coefficients $`b^2`$ and are therefore second-order in the perturbation amplitude. With additional work, one can show that the next contributing order is proportional to $`b^4`$ and therefore relatively negligible. We will denote the ensemble average of the fluctuating coefficients which we will denote as angle brackets, e.g.: $`b_j^{lm}(t_1)b_j^{lm}(t_2)\mathrm{}b_j^{lm}(t_n)`$. Note that $`b_k^{lm}(t_1)b_k^{lm}(t_2)`$ is related to the density correlation function:
$`b_k^{lm}(t_1)b_k^{lm}(t_2)`$ $`=`$ $`{\displaystyle d^3r_1d^3r_2Y_{lm}^{}(\theta _1,\varphi _1)Y_{lm}(\theta _2,\varphi _2)u_j^{lm}(r_1)u_k^{lm}(r_2)\rho (𝐫_\mathrm{𝟏},t_1)\rho (𝐫_\mathrm{𝟐},t_2)}.`$ (27)
Assuming that the process causing fluctuations is independent of time (e.g. a stationary process) we can write $`\rho (𝐫_\mathrm{𝟏},t_1)\rho (𝐫_\mathrm{𝟐},t_2)=C(𝐫_\mathrm{𝟏},𝐫_\mathrm{𝟐},t_1t_2)`$. The quantity $`b_k^{lm}(t_1)b_k^{lm}(t_2)=b_k^{lm}(0)b_k^{lm}(t_1t_2)`$ describes the correlation of random variables $`b_j^{lm}`$ as a function of time. The limit $`t_1t_2`$ gives the mean-squared fluctuation amplitude and was explored and compared to n-body simulations in Paper 1.
We can now use equations (LABEL:eq:h1a), (21) and (25) to evaluate the moments in equation (6). After explicit substitution and averaging over angles, we have
$`\mathrm{\Delta }I_j(t+\tau )`$ $`=`$ $`\left({\displaystyle \frac{1}{2\pi }}\right)^2{\displaystyle _{\mathrm{}}^{\mathrm{}}}d\omega _1{\displaystyle _{\mathrm{}}^{\mathrm{}}}d\omega _2i𝐥{\displaystyle \frac{\mathrm{ln}f_o}{𝐈}}{\displaystyle _t^{t+\tau }}dt_1e^{i\omega _1t}{\displaystyle _t^{t+\tau }}dt_2e^{i\omega _2t}\times `$ (28)
$`{\displaystyle \underset{\mu }{}}Y_{ll_2}(\pi /2,0)r_{l_2m}^l(\beta )W_{ll_2m}^{l_1\mu }(𝐈){\displaystyle \underset{r}{}}_{\mu r}^{lm}(\omega _1)b_r^{lm}(t_1)\times `$
$`il_j{\displaystyle \underset{\nu }{}}Y_{ll_2}(\pi /2,0)r_{l_2m}^l(\beta )W_{ll_2m}^{l_1\nu }(𝐈){\displaystyle \underset{s}{}}_{\nu s}^{lm}(\omega _2)b_s^{lm}(t_2)`$
$`=`$ $`l_j𝐥{\displaystyle \frac{\mathrm{ln}f_o}{𝐈}}|Y_{ll_2}(\pi /2,0)|^2\left({\displaystyle \frac{1}{2\pi }}\right)^2{\displaystyle _{\mathrm{}}^{\mathrm{}}}d\omega _1{\displaystyle _{\mathrm{}}^{\mathrm{}}}d\omega _2\times `$
$`{\displaystyle \underset{\mu \nu }{}}{\displaystyle \underset{rs}{}}r_{l_2m}^l(\beta )r_{l_2m}^l(\beta )W_{ll_2m}^{l_1\mu }(𝐈)W_{ll_2m}^{l_1\nu }(𝐈)_{\mu r}^{lm}(\omega _1)_{\nu s}^{lm}(\omega _2)\times `$
$`{\displaystyle _t^{t+\tau }}𝑑t_1{\displaystyle _t^{t+\tau }}𝑑t_2e^{i(\omega _1t_1+\omega _2t_2)}b_r^{lm}(\omega _1)b_s^{lm}(\omega _2).`$
The expression for $`\mathrm{\Delta }I_j(t+\tau )\mathrm{\Delta }I_k(t+\tau )`$ is nearly the same, with $`𝐥d\mathrm{ln}f_o/d𝐈`$ in equation (28) replaced with $`l_k`$. Although cumbersome in appearance, all quantities in this expression are straightforwardly computed. In particular, the rotation matrices, $`r_{l_2m}^l(\beta )`$, have closed-form analytic expressions and the response matrices, $`_{\mu r}^{lm}(\omega )`$, have elements that can be evaluated by quadrature. For a specific stochastic process, all that remains is to evaluate the Fourier transform of the density correlation function. We will do this below for quasi-periodic and transient noise sources.
### 3.1 Orbiting point mass perturbers
For a halo of black holes, we may assume that the perturbers are point masses. In other words, the density $`\rho `$ is a sum of delta functions. Expanding the distribution for a single black home in an action-angle series gives
$`b_j^{lm}(t)`$ $`=`$ $`{\displaystyle \underset{𝐥}{}}Y_{ll_2}(\pi /2,0)r_{l_2m}^l(\beta )W_{ll_2m}^{l_1\nu }(𝐈)e^{i𝐥𝐰(t)}`$ (29)
where $`𝐰(t)=𝐰_o+\mathrm{\Omega }t`$ and after Fourier transforming in time, we find
$`b_j^{lm}(\omega )`$ $`=`$ $`{\displaystyle \underset{𝐥}{}}Y_{ll_2}(\pi /2,0)r_{l_2m}^l(\beta )W_{ll_2m}^{l_1\nu }(𝐈)e^{i𝐥𝐰_o}2\pi \delta \left(\omega 𝐥\mathrm{\Omega }\right)`$ (30)
As in Paper 1, we assume that orbits of individual particles are uncorrelated. The particle wakes do in fact give rise to correlations but this is of higher order in $`1/N`$ in the BBGKY expansion than the lowest-order effect we will consider here (cf. Gilbert 1969). The number density of particles at $`𝐈_1,𝐰_1`$ at time $`t_1`$ and at $`𝐈_2,𝐰_2`$ at time $`t_2`$ is
$$𝒫(𝐈_1,𝐰_1,t_1;𝐈_2,𝐰_2,t_2)=𝒫(𝐈_1,𝐰_1)\delta (𝐈_1𝐈_2)\delta (𝐰_1𝐰_2+\mathrm{\Omega }(𝐈_1)(t_2t_1))$$
(31)
where $`𝒫(𝐈,𝐰)`$ is the equilibrium particle distribution with
$$N=d^3Id^3w𝒫(𝐈,𝐰).$$
(32)
Direct substitution demonstrates that equation (31) solves the Liouville equation with the initial condition $`𝐈_2=𝐈_1`$ and $`𝐰_2=𝐰_1`$ at $`t=t_1`$. Similarly, integrating equation (31) over all coordinates gives $`N`$. The ensemble average $`b_r^{lm}(t_1)b_s^{lm}(t_2)`$ is then the average of $`b_r^{lm}(t_1)b_s^{lm}(t_2)`$ the distribution given by equation (31).
We now apply this two-particle distribution function to explicitly evaluate the ensemble average in equation (28). The ensemble average here implies an average of possible distributions of point masses consistent with some underlying distribution. Therefore
$`b_r^{lm}(\omega _1)b_s^{l^{}m^{}}(\omega _2)`$ $`=`$ $`{\displaystyle \underset{𝐥}{}}\delta _{mm^{}}Y_{ll_2}(\pi /2,0)Y_{l^{}l_2}^{}(\pi /2,0)r_{l_2m}^l(\beta )r_{l_2}^l^{}(\beta )\times `$
$`W_{ll_2m}^{l_1\nu }(𝐈)W_{l^{}l_2m}^{l_1\mu }(𝐈)(2\pi )^2\delta (\omega _1+\omega _2)\delta (\omega _1𝐥\mathrm{\Omega })\text{}.`$
and
$$b_r^{lm}(t_1)b_s^{l^{}m}(t_2)=\left(\frac{1}{2\pi }\right)^2_{\mathrm{}}^{\mathrm{}}𝑑\omega _1_{\mathrm{}}^{\mathrm{}}𝑑\omega _2e^{i(\omega _1t_1\omega _2t_2)}b_r^{l^{}m}(\omega _1)b_s^{lm}(\omega _2).$$
(34)
Integration over angles identifies $`𝐥`$ with $`𝐥^{}`$ and therefore $`m=l_3=l_3^{}=m^{}`$. Now using equations (29), (30), (31) to evaluate the average we find that the final part of expression (28) becomes
$`\left({\displaystyle \frac{1}{2\pi }}\right)^2{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\omega _1{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\omega _2e^{i(\omega _1t_1+\omega _2t_2)}b_r^{lm}(\omega _1)b_s^{lm}(\omega _2)=`$
$`(2\pi )^3{\displaystyle \underset{𝐥}{}}{\displaystyle d^3If_o(𝐈)e^{i𝐥\mathrm{\Omega }(𝐈)(t_1t_2)}|Y_{ll_2}(\pi /2,0)|^2|r_{l_2m}^l(\beta )|^2W_{ll_2m}^{l_1r}(𝐈)W_{ll_2m}^{l_1s}(𝐈)}`$
where the integration over $`\omega _1`$ and $`\omega _2`$ has become a sum over discrete frequencies denoted $`𝐥`$ and we have exploited the orthogonality of rotation matrices:
$$𝑑\beta \mathrm{sin}(\beta )r_{\mu \nu }^l(\beta )r_{\mu \nu }^l^{}(\beta )=\frac{2}{2l+1}\delta _{ll^{}}$$
(36)
(Edmonds 1960).
We now substitute this development into equation (28) and find
$`\mathrm{\Delta }I_j(t+\tau )`$ $`=`$ $`{\displaystyle \frac{1}{f_0(𝐈)}}l_j𝐥{\displaystyle \frac{f_o}{𝐈}}|Y_{ll_2}(\pi /2,0)|^2{\displaystyle \underset{\mu \nu }{}}{\displaystyle \underset{rs}{}}r_{l_2m}^l(\beta )r_{l_2m}^l(\beta )W_{ll_2m}^{l_1\mu }(𝐈)W_{ll_2m}^{l_1\nu }(𝐈)\times `$ (37)
$`\{(2\pi )^3{\displaystyle \underset{𝐥}{}}{\displaystyle }d^3If_o(𝐈)|Y_{ll_2}(\pi /2,0)|^2r_{l_2m}^l(\beta )r_{l_2m}^l(\beta )W_{ll_2m}^{l_1r}(𝐈)W_{ll_2m}^{l_1s}(𝐈)_{\mu r}^{lm}(𝐥\mathrm{\Omega }(𝐈))_{\nu s}^{lm}(𝐥\mathrm{\Omega }(𝐈))\times `$
$`{\displaystyle _t^{t+\tau }}dt_1{\displaystyle _t^{t+\tau }}dt_2e^{i𝐥\mathrm{\Omega }(𝐈)(t_1t_2)}\}.`$
As described generally above, the term in $`\{\}`$ is the temporal correlation function for the response to the point masses fluctuations and only depends on the time difference $`t_1t_2`$. The correlation is finite and lasts some order-unity number of dynamical times. Our diffusion calculation is in the regime $`\tau 1/\mathrm{\Omega }`$. We may change the double time integration from variables $`t_1,t_2`$ to $`T=(t_1+t_2)/2,\tau =t_1t_2`$. The integral over $`\tau `$ gives a delta function $`\delta (𝐥\mathrm{\Omega }(𝐈))`$ and the integral over $`T`$ gives $`\tau `$. The delta function implies that only orbits with commensurate frequencies will give rise to secular changes. Physically, the disturbance must present an asymmetric force distribution to cause a secular change in the actions of dark matter orbits. If the orbital frequencies are not commensurate, the long-term average of the perturbing force will be axisymmetric. For orbiting black holes, the ratio $`\mathrm{\Omega }_1/\mathrm{\Omega }_2`$ ranges from 1 to 2 as the halo potential varies from that for point mass to homogeneous core. Therefore for most systems, no commensurabilities are available for harmonic orders $`l=1,2`$. In other words, a resonant $`l=1`$ orbit will look like a Keplerian orbit and a resonant $`l=2`$ orbit will look like a bisymmetric oval (e.g. stationary bar orbit) and neither exist in any significant measure in most extended stellar systems.
Returning to equation (4), we can now read off our diffusion coefficients in action variables:
$`D_j^{(1)}(𝐈,t)`$ $`=`$ $`\underset{\tau 0}{lim}{\displaystyle \frac{\mathrm{\Delta }I_j(t+\tau )}{\tau }}`$
$`=`$ $`{\displaystyle \frac{1}{f_0(𝐈)}}l_j𝐥{\displaystyle \frac{f_o}{𝐈}}|Y_{ll_2}(\pi /2,0)|^2{\displaystyle \underset{\mu \nu }{}}{\displaystyle \underset{rs}{}}r_{l_2m}^l(\beta )r_{l_2m}^l(\beta )W_{ll_2m}^{l_1\mu }(𝐈)W_{ll_2m}^{l_1\nu }(𝐈)\times `$
$`\{(2\pi )^3{\displaystyle \underset{𝐥}{}}{\displaystyle }d^3If_o(𝐈)|Y_{ll_2}(\pi /2,0)|^2r_{l_2m}^l(\beta )r_{l_2m}^l(\beta )\times `$
$`W_{ll_2m}^{l_1r}(𝐈)W_{ll_2m}^{l_1s}(𝐈)_{\mu r}^{lm}(𝐥\mathrm{\Omega }(𝐈))_{\nu s}^{lm}(𝐥\mathrm{\Omega }(𝐈))2\pi \delta (𝐥\mathrm{\Omega }(𝐈))\},`$
$`D_{jk}^{(2)}(𝐈,t)`$ $`=`$ $`\underset{\tau 0}{lim}{\displaystyle \frac{\mathrm{\Delta }I_j(t+\tau )\mathrm{\Delta }I_k(t+\tau )}{2\tau }}`$
$`=`$ $`{\displaystyle \frac{l_jl_k}{2}}|Y_{ll_2}(\pi /2,0)|^2{\displaystyle \underset{\mu \nu }{}}{\displaystyle \underset{rs}{}}r_{l_2m}^l(\beta )r_{l_2m}^l(\beta )W_{ll_2m}^{l_1\mu }(𝐈)W_{ll_2m}^{l_1\nu }(𝐈)\times `$
$`\{(2\pi )^3{\displaystyle \underset{𝐥}{}}{\displaystyle }d^3If_o(𝐈)|Y_{ll_2}(\pi /2,0)|^2r_{l_2m}^l(\beta )r_{l_2m}^l(\beta )\times `$
$`W_{ll_2m}^{l_1r}(𝐈)W_{ll_2m}^{l_1s}(𝐈)_{\mu r}^{lm}(𝐥\mathrm{\Omega }(𝐈))_{\nu s}^{lm}(𝐥\mathrm{\Omega }(𝐈))2\pi \delta (𝐥\mathrm{\Omega }(𝐈))\}.`$
(39)
Note that the limit $`\tau 0`$ is taken in the sense that $`\tau `$ is small compared to the evolutionary time scale due to the fluctuations but large compared to the dynamical time. The time dependence in the diffusion coefficients reminds us that the underlying equilibrium distribution $`f_o(𝐈)`$ changes on an evolutionary time scale but, for the purposes of computation, is held fixed on a dynamical time scale. The integrals may be simplified by noting that $`d^3I=dEdJJd(\mathrm{cos}\beta )/\mathrm{\Omega }_1(E,J)`$, we can do the integral in $`\beta `$ using the orthogonality of the rotation matrices as previously described. For a given background distribution function $`f_o(𝐈)`$, the term in curly brackets can be computed once and for all since they are independent of the local value of the actions.
### 3.2 Transient processes
In the previous application for orbiting point masses, we saw that only harmonics with $`l3`$ were effective at driving evolution. Here, we describe the results of bombarding the galaxy isotropically with bits of mass. This is an idealization of interactions during the epoch of galaxy evolution. The perturbations are shots, not quasi-periodic, and therefore lead to a continuous spectrum of perturbation frequencies. This case and the previously considered point-mass case represent two extremes. For example, a decaying orbit will have a set of broad peaks and a low-frequency continuum.
The expansion coefficients for the perturbation are straightforward for the mass inside of the galaxy:
$$b_i^{lm}(t)=d^3rY_{lm}^{}(\theta ,\varphi )u_i^{lm}(r)\delta \left(𝐫𝐫(t)\right)=Y_{lm}^{}(\theta (t),\varphi (t))u_i^{lm}(r(t)).$$
(40)
If the mass is outside the galaxy, we use the multipole expansion with the density rather than the potential member of the biorthogonal pair to evaluate the coefficients:
$`b_i^{lm}(t)`$ $`=`$ $`{\displaystyle \frac{1}{4\pi }}{\displaystyle d^3rY_{lm}^{}(\theta ,\varphi )d_i^{lm}(r)\frac{1}{|𝐫𝐫(t)|}}`$ (41)
$`=`$ $`{\displaystyle \frac{1}{4\pi }}{\displaystyle d^3rY_{lm}^{}(\theta ,\varphi )d_i^{lm}(r)\underset{l}{}\underset{m=l}{\overset{l}{}}\frac{4\pi }{2l+1}\frac{r_<^l}{r_>^{l+1}}Y_{lm}^{}(\theta (t),\varphi (t))Y_{lm}^{}(\theta ,\varphi )}`$
$`=`$ $`{\displaystyle \frac{1}{2l+1}}Y_{lm}^{}(\theta (t),\varphi (t)){\displaystyle 𝑑rr^2d_i^{lm}(r)\frac{r^l}{r(t)^{l+1}}}`$
The Fourier transform needed for equation (28) is most easily done assuming the perturber is in the equatorial plane. We denote the Fourier transform of $`b_j^{lm}(t)`$ as $`\widehat{b}_j^{lm}(\omega )`$. In practice, we will perform the transform numerically by FFT. Note that equations (40) and (41) describe the time dependence in coefficients for any trajectory $`r(t),\theta (t),\varphi (t)`$. A cloud of points and more generally any phase-space distortion yielding $`b_i^{lm}(t)`$ or coefficients at discrete times from an n-body simulation are possible input to the FFT. Now, we can evaluate the final line in equation (28) by changing coordinates from $`t_1,t_2`$ to $`T=(t_1+t_2)/2,\tau =t_1t_2`$. We have
$$e^{i(\omega _1t_1+\omega _2t_2)}=e^{i(\omega _1[T+\tau /2]+\omega _2[T\tau /2])}=e^{i(\omega _1+\omega _2)T}e^{i(\omega _1\omega _2)\tau /2}.$$
Using this in the last line of equation (28) we find
$`{\displaystyle _t^{t+\tau }}𝑑t_1{\displaystyle _t^{t+\tau }}𝑑t_2e^{i(\omega _1t_1+\omega _2t_2)}b_r^{lm}(\omega _1)b_s^{lm}(\omega _2)`$ $`=`$ $`{\displaystyle _t^{t+\tau /2}}𝑑T\mathrm{\hspace{0.17em}4}\pi \delta (\omega _1\omega _2)e^{i2\omega T}b_r^{lm}(\omega _1)b_s^{lm}(\omega _2)`$
$`=`$ $`4\pi \delta (\omega _1\omega _2)e^{i\omega \tau }{\displaystyle \frac{\mathrm{sin}\omega \tau }{\omega }}\overline{b}_r^{lm}(\omega _1)\overline{b}_s^{lm}(\omega _2).`$
where $`\omega \omega _1=\omega _2`$. In deriving the second equality, we note that the bombardment must occur between $`t`$ and $`t+\tau `$ and use the shift properties of the Fourier transform to write $`b_r^{lm}(\omega _1)b_s^{lm}(\omega _2)=e^{2\omega t}\overline{b}_r^{lm}(\omega _1)\overline{b}_s^{lm}(\omega _2)`$ where the transform $`\overline{b}(\omega )`$ denotes the transform of an event centred about the temporal origin. In the limit $`\tau 0`$, this expression becomes $`4\pi \delta (\omega _1\omega _2)\tau `$. Substituting, this back into equation (28) we can perform one of the $`\omega `$ integrals straight away. After rearranging we have
$`\mathrm{\Delta }I_j(t+\tau )`$ $`=`$ $`{\displaystyle \frac{1}{f_0(𝐈)}}l_j𝐥{\displaystyle \frac{f_o}{𝐈}}|Y_{ll_2}(\pi /2,0)|^2{\displaystyle \underset{\mu \nu }{}}{\displaystyle \underset{rs}{}}r_{l_2m}^l(\beta )r_{l_2m}^l(\beta )W_{ll_2m}^{l_1\mu }(𝐈)W_{ll_2m}^{l_1\nu }(𝐈)\times `$ (43)
$`4\pi {\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\omega _{\mu r}^{lm}(\omega )_{\nu s}^{lm}(\omega )b_r^{lm}(\omega _1)b_s^{lm}(\omega _2).`$
The expression for $`\mathrm{\Delta }I_j(t+\tau )\mathrm{\Delta }I_k(t+\tau )`$ follows by analogy with the equations (28) and (39) in §3.1.
Symmetry suggests the choice of perturbing orbits on the equatorial plane should not effect the final results. This can be explicitly demonstrated explicitly using the rotational properties of the spherical harmonics. Let $`_{mm^{}}^l(\alpha ,\beta ,\gamma )`$ be the rotation matrix with the Euler angles $`\alpha ,\beta ,\gamma `$. Then, because
$$Y_{lm}(\theta ,\varphi )=\underset{m^{}=l}{\overset{l}{}}Y_{lm^{}}(\theta ^{},\varphi ^{})_{mm^{}}^l(\alpha ,\beta ,\gamma )$$
where the primed coordinates refer to the rotated coordinate system, we have
$$b_j^{lm}(t)=\underset{m^{}=l}{\overset{l}{}}b_j^{lm^{}}(t)_{mm^{}}^l(\alpha ,\beta ,\gamma )$$
(44)
For an isotropic spherical system $`_{\mu r}^{lm}(\omega )`$ is independent of $`m`$. We can exploit this and the rotational properties of $`b_r^{lm}(\omega )`$ to simplify the computation of $`b_r^{lm}(\omega _1)b_s^{lm}(\omega _2)`$. We may express $`b_r^{lm}(\omega )`$ in any convenient coordinate system and use the rotation matrices to rotate this to any orientation:
$$b_r^{lm}(\omega _1)b_s^{lm}(\omega _2)=\underset{m^{}=l}{\overset{l}{}}R_{mm^{}}^l(\alpha ,\beta ,\gamma )b_r^{lm^{}}(\omega _1)\underset{m^{\prime \prime }=l}{\overset{l}{}}R_{mm^{\prime \prime }}^l(\alpha ,\beta ,\gamma )b_s^{lm^{\prime \prime }}(\omega _2)$$
(45)
Summing over all values $`m`$ for a given $`l`$ we have
$$\underset{m=l}{\overset{l}{}}b_r^{lm}(\omega _1)b_s^{lm}(\omega _2)=\underset{m^{}=l}{\overset{l}{}}b_r^{lm^{}}(\omega _1)b_s^{lm^{}}(\omega _2)$$
(46)
having used the orthogonality of rotation matrices. Finally, we assume that the ensemble average includes random events from all directions and therefore only the same event will be correlated in the computation of $`b_r^{lm}(\omega _1)b_s^{lm}(\omega _2)`$.
## 4 Summary
This paper presents a general equation for noise-driven evolution which incorporates the full self-gravitating response of a stochastic process in the perturbation limit. The motivating question is same as that of violent relaxation: what is the long-term response of a stellar system to a fluctuating potential? By working in the perturbation limit, the linearity guarantees that the process is Markovian and therefore the response of the stellar system to repeated stochastic events is naturally treated using the matrix-method and the statistic methods for handling general stochastic differential equations.
In general, expansions of integral equations for stochastic processes yield an infinite number of terms. A general theorem demonstrates that the truncation at quadratic order is consistent for a Markov process and the resulting evolution equation takes the Fokker-Planck form<sup>1</sup><sup>1</sup>1This is also true for two-body collisions if one eliminates strong encounters.. Evaluation of the Fokker-Planck coefficients requires the specifying the response of stellar system individual events in the stochastic process.
We explicitly develop techniques for two general situations: quasi-periodic perturbers and transient perturbers. The former case is motivated by halo of super-massive black holes (e.g. Lacey & Ostriker 1985). The latter case includes almost everything else; for example, unbound dwarf encounters (fly-by), orbiting substructure decaying due to dynamical friction, disrupting dwarfs, and mixing tidal debris. There is no practical constraint on deriving the Fokker-Planck coefficients for transient noise but that the process can be represented as an expansion in some biorthogonal basis. For dwarfs on decaying or unbound orbits, this is particularly easy and can be done by quadrature. Alternatively, one may construct an n-body simulation using the expansion method and let the simulation produce the time series of coefficients directly.
The companion paper (Paper 2) applies the apparatus described here to investigate the evolution of haloes during the noisy epoch of galaxy formation. We find that both unbound encounters and decaying substructure drives the halo profile to a self-similar form similar to those found recently in cosmological simulations. There are a number of interesting applications. For example, the distribution of binary semi-major axes $`a`$ in the field star populations is proportional to $`a^1`$ over several orders of magnitude. Preliminary work suggests that this may be explained as the result of fluctuations from the noisy environment found in proto-stellar clusters and molecular clouds. Other possible applications include effect of transient bar formation and spiral structure on overall galactic structure in the 10 billion years since formation.
## Acknowledgments
I thank Enrico Vesperini for comments and discussion. This work was support in part by NSF AST-9529328.
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# Overcritical Rotation of a Trapped Bose-Einstein Condensate
## a Normal branch:
this branch starts at $`\mathrm{\Omega }=0`$. The linear dependence at small $`\mathrm{\Omega }`$ is given by $`\alpha =\mathrm{\Omega }ϵ`$. By increasing $`\mathrm{\Omega }`$ the square radius $`y^2`$ increases and eventually diverges at $`\mathrm{\Omega }=\omega _y`$ where $`\stackrel{~}{\omega }_y0`$ and the branch has its end . Also the angular momentum and the release energy diverge at $`\mathrm{\Omega }=\omega _y`$. Notice that when $`\stackrel{~}{\omega }_y0`$ the moment of inertia takes the rigid value since $`x^2y^2`$ (see Eq.(11)).
## b Overcritical branch:
this branch starts at $`\mathrm{\Omega }=+\mathrm{}`$ where $`\alpha `$ behaves like $`\alpha =(\omega _x^2\omega _y^2)/4\mathrm{\Omega }`$. It is worth noticing that in this limit both $`\stackrel{~}{\omega }_x^2`$ and $`\stackrel{~}{\omega }_y^2`$ approach the value $`(\omega _x^2+\omega _y^2)/2`$ and therefore the shape of the density profile becomes symmetric despite the asymmetry of the confining trap. In the same limit the angular momentum tends to zero while the release energy approaches the finite value $`E_{\mathrm{rel}}=(2/7)\stackrel{~}{\mu }_{\mathrm{}}`$, where $`\stackrel{~}{\mu }_{\mathrm{}}`$ is the chemical potential (8) with $`\stackrel{~}{\omega }_x^2=\stackrel{~}{\omega }_y^2=(\omega _x^2+\omega _y^2)/2`$. In the overcritical branch the deformation of the cloud takes a sign opposite to the one of the trap. This branch exhibits a back-bending at a value of $`\mathrm{\Omega }`$ which is smaller than $`\omega _x`$, but can be higher or smaller than $`\omega _y`$, depending on whether the value of $`ϵ`$ is larger or smaller than $`0.2`$ (see Figs. 2 and 3). In both cases this branch ends, after the back-bending, at the value $`\mathrm{\Omega }=\omega _x`$, where $`\stackrel{~}{\omega }_x0`$ and $`x^2`$, $`L_z`$ and $`E_{\mathrm{rel}}`$ diverge.
It is also useful to discuss the instructive case $`ϵ0`$ corresponding to symmetric trapping in the $`x`$-$`y`$ plane ($`\omega _x=\omega _y`$). In this case one finds a solution with $`\alpha =0`$ for $`\mathrm{\Omega }<\omega _x/\sqrt{2}`$ . For higher frequencies three solutions appear: the first one still corresponds to a non-rotating configuration ($`\alpha =0`$), while two solutions, given respectively by $`\alpha =\pm \sqrt{2\mathrm{\Omega }^2\omega _x^2}`$, correspond to rotating deformed configurations. The existence of these solutions, which break the original symmetry of the Hamiltonian, is the analog of the bifurcation from the axisymmetric Maclaurin to the triaxial Jacobi ellipsoids for rotating classical fluids . It is worth pointing out that the existence of the bifurcation is the consequence of two-body interactions and is absent in the non-interacting Bose gas. When $`ϵ`$ is slightly different from zero (and positive) the two solutions are no longer degenerate, the one with $`\alpha <0`$ having the lowest energy.
The existence of stationary solutions in the rotating frame raises the important question of stability. Actually, one should distinguish between thermodynamic and dynamical instability. The former corresponds to the absence of thermodynamic equilibrium and its signature, at zero temperature, is given by the existence of excitations with negative energy . The latter is instead associated with the decay of the initial configuration due to interaction effects and is in general revealed by the appearance of excitations with complex energy. Configurations characterized by thermodynamic instability are destabilized only in the presence of dissipative processes.
Let us first discuss the stability with respect to the center of mass motion. In the presence of harmonic trapping the corresponding equations of motion, in the rotating frame, take the classical form (rotating Blackburn’s pendulum, ) and are not affected by interatomic forces. Their solutions obey the dispersion law
$$\omega ^2=\frac{1}{2}\left[\omega _x^2+\omega _y^2+2\mathrm{\Omega }^2\pm \sqrt{(\omega _x^2\omega _y^2)^2+8\mathrm{\Omega }^2(\omega _x^2+\omega _y^2)}\right],$$
(16)
and are dynamically stable ($`\omega ^2>0`$) for $`\mathrm{\Omega }<\omega _y`$ and $`\mathrm{\Omega }>\omega _x`$. So the requirement that the dipole oscillation be dynamically stable excludes the region between the dotted lines in Figs. 1, 2 and 3. Notice that this is the same region where the Schroedinger equation for the non-interacting Bose gas has no stationary solutions in the rotating frame.
We have further explored the conditions of stability by studying the quadrupole oscillations of the condensate around the equilibrium configuration in the rotating frame. The calculation is derivable by linearizing the equations of motion (1) and (2), with the choice
$`\delta \rho (𝐫)`$ $`=a_0+a_xx^2+a_yy^2+a_zz^2+a_{xy}xy`$ (17)
$`\delta 𝐯(𝐫)`$ $`=\mathbf{}(\alpha _xx^2+\alpha _yy^2+\alpha _zz^2+\alpha _{xy}xy),`$ (18)
for the fluctuations of the density and of the velocity field where the coefficients $`a_i`$ and $`\alpha _i`$ depend on time. The results of the analysis show that the window of dynamical instability for the quadrupole oscillations is different from the one of the dipole. In particular the normal branch is always dynamically stable, while the overcritical branch is stable only in its lower part since, after the back-bending, where d$`\alpha /`$d$`\mathrm{\Omega }>0`$, one of the quadrupole frequencies becomes purely imaginary ($`\omega ^2<0`$). We have investigated the stability of the quadrupole excitations also in the limiting case $`ϵ=0`$. In this case the upper and lower branches $`\alpha =\pm \sqrt{2\mathrm{\Omega }^2\omega _x^2}`$ give rise to a vanishing value for one of the quadrupole frequencies, the others being always real. This vanishing solution corresponds to the rotation of the system in the $`xy`$ plane and reflects the rotational symmetry of the Hamiltonian.
We stress again that the solutions discussed in this paper correspond, in general, to metastable configurations. For example it is well known that vortical configurations can become energetically favorable for relatively small values of $`\mathrm{\Omega }`$. Also excitations with higher multipolarities can become energetically favorable for some values of the angular velocity . In this work we assume that at very low temperature, where collisions are rare, the solutions (3)-(4) can survive, at least for useful time intervals, also in conditions of thermodynamic instability.
Let us finally briefly discuss the experimental possibility of realizing the rotations described above. The normal branch, characterized by the huge increase of the size in the $`y`$ direction and of the release energy when $`\mathrm{\Omega }\omega _y`$, could be in principle generated by an adiabatic increase of the angular velocity, starting from a cold condensate . The transition from the normal to the overcritical branch cannot be instead realized in a continuous way and the only way to realize the exotic configurations of this branch is to engineer the proper conditions of dynamical equilibrium, by building up both the proper phase of the order parameter and the shape of the density profile.
We are indebted to S. Vitale for introducing us to the problem of the overcritical rotations. Fruitful discussions with E. Cornell, J. Dalibard, L. Pitaevskii and G. Shlyapnikov are also acknowledged. This work has been supported by the Ministero della Ricerca Scientifica e Tecnologica (MURST).
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# Integral representations and Liouville theorems for solutions of periodic elliptic equations
## 1 Introduction
The topic of this paper stems from two sources. The first of them are representation theorems for certain classes of eigenfunctions of the Laplace operator in $`^n`$, or equivalently, of solutions of the Helmholtz equation
$$\mathrm{\Delta }uk^2u=0\text{ in }^n,$$
(1.1)
where $`k^{}:=\{0\}`$. Such theorems for arbitrary solutions of (1.1) were obtained in $`^2`$ and in the hyperbolic plane by S. Helgason , and in $`^n`$ by M. Hashizume et al. , M. Morimoto , and recently by S. Agmon . We remark that it should also be possible to deduce similar results from the L. Ehrenpreis’ fundamental principle. The zero set of the symbol of the operator in the left hand side of (1.1) is
$$\mathrm{\Sigma }=\left\{\xi ^n\right|\xi ^2=k^2\},$$
where $`\xi ^2=_{j=1}^n\xi _j^2`$. The L. Ehrenpreis’ “fundamental principle” in the particular case of equation (1.1) claims that any solution of (1.1) can be represented as a combination (i.e., an integral with respect to the parameter $`\xi `$) of the exponential solutions
$$e_\xi (x):=\mathrm{exp}(i\xi x),\xi \mathrm{\Sigma },$$
where $`\xi x=_{j=1}^n\xi _jx_j`$ (see the details and more precise formulation in or ). The set $`\mathrm{\Sigma }`$ is an irreducible analytic subset of $`^n`$, which is uniquely determined when $`k0`$ by its spherical subset
$$S=\left\{\xi ^n\right|\xi =k\omega ,\omega S^{n1}^n\}.$$
Here $`S^{n1}`$ denotes the unit sphere in $`^n`$. It is clear then that due to the uniqueness of analytic continuation, the exponential representation of solutions $`u(x)`$ of (1.1) should be reducible to one that involves the solutions $`e_\xi `$ with $`\xi S`$ only. Namely, consider the restriction mapping from functions analytic on the whole characteristic variety $`\mathrm{\Sigma }`$ to the sphere $`S`$. Due to the irreducibility and the uniqueness of analytic continuation, this mapping is one-to-one. Hence, there is a function space on the sphere $`S`$ which is the isomorphic image of the space of all analytic functions on $`\mathrm{\Sigma }`$. It follows that any hyperfunction (analytic functional) on $`\mathrm{\Sigma }`$ can be rewritten as a functional on $`S`$. Since the “fundamental principle” essentially expresses all solutions of (1.1) as applications of such analytic functionals to the analytic family of exponential solutions, we get our conclusion.
Now, depending on how fast the solution $`u(x)`$ grows at infinity, the corresponding representing functional on $`S`$ is actually a measure, a distribution, a hyperfunction, or a functional of a more general kind. For instance (see ), an arbitrary solution in $`^n`$ can be represented as
$$u(x)=<\varphi (\xi ),e_\xi (x)>,$$
(1.2)
where $`\varphi (\xi )`$ is a functional on $`S`$ which belongs to the dual space to the space $`:=lim_R\mathrm{}_R(S^{n1})`$. Here for every $`R>0`$ the Hilbert space $`_R(S^{n1})`$ is defined as follows:
$`_R(S^{n1}):=\left\{\psi \right|`$ $`\psi (\omega )={\displaystyle \underset{l,m}{}}a_{l,m}{\displaystyle \frac{(R/2)^l}{\mathrm{\Gamma }(l+\frac{n+1}{2})}}Y_l^m(\omega ),\text{ s.t.}`$
$`\psi __R:=({\displaystyle \underset{l,m}{}}|a_{l,m}|^2)^{1/2}<\mathrm{}\},`$
where $`Y_l^m(\omega )`$ denote spherical harmonics, and $`(S^{n1})`$ is equipped with the inductive limit topology of $`lim_R\mathrm{}_R(S^{n1})`$.
The representation (1.2) can be formally rewritten as
$$u(x)=_Se_\xi (x)𝑑\varphi (\xi ).$$
where $`\varphi (\xi )`$ is a suitable functional.
The functional $`\varphi `$ is a hyperfunction (analytic functional) on $`S`$ if and only if for arbitrary $`\epsilon >0`$ the solution $`u(x)`$ grows not faster than
$$O(\mathrm{exp}((|Imk|+\epsilon )|x|))$$
(see ). One can also describe other classes of solutions, for instance, solutions which are represented by a distribution or a measure (see and the references therein).
As we have already mentioned, these results could be probably extracted from the “fundamental principle” . The crucial factors are that $`S`$ is sufficiently massive and $`\mathrm{\Sigma }`$ is irreducible, so $`S`$ determines $`\mathrm{\Sigma }`$ uniquely. Besides, $`S`$ is a rather simple analytic manifold. These features allow more or less explicit descriptions of the needed spaces of test functions and functionals. It is easy to understand that if $`\mathrm{\Sigma }`$ were reducible, it would not be possible to obtain the representation of all solutions using only $`\xi S`$. The reason is that the solution $`e_\xi `$ with $`\xi `$ that belongs to a component not touching $`S`$ would not be representable this way. On the other hand, if one wants to deal only with solutions growing not faster than $`O(\mathrm{exp}(|Imk|+\epsilon )|x|)`$ for all $`\epsilon >0`$, then the irreducibility of $`\mathrm{\Sigma }`$ is not needed. In this case, it is only required that $`\mathrm{\Sigma }`$ is irreducible in a vicinity of $`S`$, so other components of $`\mathrm{\Sigma }`$ do not meet $`S`$.
The “fundamental principle” was extended in to solutions of certain growth (for instance, of exponential growth) of elliptic and hypoelliptic periodic equations (see also the extensions of the results of provided in ). The role of the exponential solutions is played here by the so called Floquet-Bloch solutions (see Definition 1), and an analog of the characteristic manifold $`\mathrm{\Sigma }`$ is the variety $`F`$ sometimes called the Fermi surface (see Definition 2). This raises the hope of finding representations similar to the ones discussed above for the more general case of a second order elliptic operator with periodic coefficients. This is not straightforward, however, due to several reasons. First of all, it is not that clear what should be a natural analog of $`S`$. An appropriate variety, as we explain later, is provided for second order equations by the analysis of the cone of positive solutions done by S. Agmon and by V. Lin and Y. Pinchover (see , and the references therein). The disadvantage is that the whole consideration must be done below the spectrum of the operator (more precisely, below the generalized principal eigenvalue $`\mathrm{\Lambda }_0`$, see (2.8)). Secondly, proving the irreducibility of $`F`$ happens to be a very hard nut to crack (this problem arises also in direct and inverse spectral problems, see for instance ). Fortunately enough, by appropriately restricting the growth of the solutions, one can sometimes work near a single irreducible component, and hence avoid proving the irreducibility of $`F`$. Consequently, we prove a representation theorem (Theorem 18) that characterizes all the solutions which have integral expansion into positive Bloch solutions with a hyperfunction as a “measure”.
The “fundamental principle” also suggests a point of view that is crucial for establishing representation theorems for solutions of equations with constant or periodic coefficients. Namely, it is to one’s advantage to treat solutions of the original equation in the dual sense, i.e., as functionals on appropriate spaces that are orthogonal to the range of the dual operator. We adopt this approach throughout the paper.
If one attempts now to further restrict the growth of solutions and considers the problem of the structure of all polynomially growing (or bounded) solutions, one arrives at the second topic of our study, Liouville type theorems. The classical Liouville theorem characterizes the space of all harmonic functions in $`^n`$ of polynomial growth of order $`N`$. The validity of an analog of this classical theorem has been studied in many situations (see for instance for recent results, surveys, and further references). An interesting case was considered by M. Avellaneda and F. -H. Lin , and also by J. Moser and M. Struwe . In these papers the authors dealt with polynomially growing solutions of a second order elliptic equation $`Lu=0`$ in divergence form with periodic coefficients and obtained a comprehensive answer (for related results see also and the references therein). More precisely, using the formalism of homogenization theory , it was proved that any solution $`v`$ of the equation $`Lu=0`$ in $`^n`$ of polynomial growth is representable as a finite sum of the form
$$v(x)=\underset{j=(j_1,\mathrm{},j_n)_+^n}{}x^jp_j(x),$$
(1.3)
where the functions $`p_j(x)`$ are periodic with respect to the group of periods of the equation. Moreover, the space of all solutions of the equation $`Lu=0`$ of polynomial growth of order at most $`N`$ is of dimension $`h_{n,N}`$, where
$$h_{n,N}:=\left(\begin{array}{c}n+N\\ N\end{array}\right)\left(\begin{array}{c}n+N2\\ N2\end{array}\right)$$
(1.4)
is the dimension of the space of all harmonic polynomials of degree $`N`$ in $`n`$ variables. We will also use the notation
$$q_{n,N}:=\left(\begin{array}{c}n+N\\ N\end{array}\right)$$
(1.5)
for the dimension of the space of all polynomials of degree at most $`N`$ in $`n`$ variables. Notice that $`q_{n1,N}`$ also coincides with the dimension of the space of all homogeneous polynomials of degree $`N`$ in $`n`$ variables, so in particular, $`h_{n,N}=q_{n1,N1}+q_{n1,N}`$.
We remark that the method of can be slightly modified to provide an extension of this Liouville theorem for general second order elliptic equations with periodic coefficients under the assumption that the generalized principal eigenvalue is zero (see Appendix A and also the recent paper , where a partial result of this type was independently obtained).
One can make an observation that these Liouville theorems are actually of the same nature as the representation theorems discussed above. In this case the analog of the set $`S`$ is the single point $`\xi =0`$ and the representing functional $`\varphi `$ is a distribution supported at this point.
Let us recall the following standard notion of Floquet theory (see ):
###### Definition 1
A solution $`u(x)`$ representable as a finite sum of the form
$$u(x)=e^{ikx}\left(\underset{j=(j_1,\mathrm{},j_n)_+^n}{}x^jp_j(x)\right)$$
(1.6)
with functions $`p_j(x)`$ periodic with respect to the group of periods of the equation is called a Floquet solution with a quasimomentum $`k^n`$. Here $`kx=k_lx_l`$ The maximum value of $`|j|=\underset{l=1}{\overset{n}{}}j_l`$ in the representation (1.6) is said to be the order of the Floquet solution. Floquet solutions of zero order are called Bloch solutions.
One sees that the representation (1.3) corresponds to a Floquet solution with a zero quasimomentum. A Liouville theorem of the type mentioned above implies in particular that the only real quasimomentum that can occur for the equation under consideration is $`k=0`$ (modulo the action of the lattice reciprocal to the group of periods). We show in the present paper that the finiteness of the set of real quasimomenta for a periodic elliptic equation is equivalent to the finite dimensionality of the spaces of solutions having a given polynomial growth and to their representation similar to, albeit more general than (1.3). This statement is very general and holds for any periodic elliptic equation (it is also true for hypoelliptic equations and systems, although we not address them here). If some additional information is available on the analytic behaviour of the dispersion relations, one can find the exact dimensions of these spaces (see Theorem 23). We present in Theorem 28 some classes of second order equations (including Schrödinger, magnetic Schrödinger, and general second order elliptic equations with real periodic coefficients) for which one can achieve all these sharp results. We show that the problem of calculating the dimensions of the spaces of Floquet solutions of a given polynomial growth reduces to a purely function theoretic question and can be resolved in a very general setting (Theorem 10).
The proofs of the results of this paper are largely dependent upon the techniques of the Floquet theory developed in .
The outline of the paper is as follows. The next section introduces necessary notations and preliminary results from the Floquet theory and the theory of positive solutions of periodic elliptic equations. In particular, we obtain a new general result (Theorem 10) on the dimensions of the spaces of Floquet solutions, which plays crucial role in our approach to Liouville theorems. Section 3 contains the proof of the integral representation (Theorem 18) analogous to Theorem 5.1 in . In Section 4, we discuss Liouville type theorems. In particular, Theorems 23 and 28 are established. In order to make the reading of the paper easier, we postpone the proofs of all the technical lemmas to Section 5. Some conclusions and remarks are provided in Section 6. The Appendix contains an alternative proof of a part of Theorem 28 using the homogenization technique similar to the one used in .
Results of this paper related to Liouville theorems were presented in March 2000 at the University of Toronto and at the Weizmann Institute.
When the paper was being prepared for submission, P. Li informed the authors that the statement of the third part of Theorem 28 for the special case of an operator of the form $`L=a_{ij}(x)_i_j`$ was simultaneously and independently obtained in using homogenization formalism .
## 2 Notations and preliminary results
Consider a linear (scalar) elliptic partial differential operator $`P(x,D)`$ of order $`m`$ in $`^n`$, $`n2`$ (in some parts of the paper we will restrict the class of operators further). Here we employ the standard notation $`D=\frac{1}{i}\frac{}{x}`$. The ellipticity is understood in the sense of the nonvanishing of the principal symbol $`P_m(x,\xi )`$ of the operator $`P`$ for all $`\xi ^n0`$. The dual operator (the formal adjoint) $`P^{}`$ has similar properties. Here we use the duality provided by the bilinear (rather than the sesquilinear) form
$$<g,f>=\underset{^n}{}f(x)g(x)𝑑x.$$
We assume that the coefficients of $`P`$ are smooth and periodic with respect to a lattice $`\mathrm{\Gamma }`$ in $`^n`$. The smoothness condition can be significantly reduced (see the Section 6). In fact, so far we only need that both operators $`P`$ and $`P^{}`$ define Fredholm mappings between the Sobolev space $`H^m`$ and $`L_2`$ on the torus $`𝕋^n=^n/\mathrm{\Gamma }`$.
An additional condition is required that would guarantee the discreteness of the spectrum of the “shifted” operators $`P(x,D+k)`$ on the torus $`𝕋^n`$ for all $`k\mathrm{C}^n`$. We need to exclude the possible pathological situation when the spectrum of $`P`$ on the torus coincides with the whole complex plane (like in the case of the operator $`\mathrm{exp}(i\varphi )d/d\varphi `$ on the circle). For instance, self-adjointness of $`P`$ could be such a condition. Another example is second order uniformly elliptic operators of the form (2.5). For more sufficient conditions see for example .
In what follows, the particular choice of the lattice is irrelevant and can always be reduced to the case $`\mathrm{\Gamma }=^n`$, which we will assume from now on. We will always use the word “periodic” in the meaning of “$`\mathrm{\Gamma }`$-periodic”.
We denote by $`K=[0,1]^n`$ the standard fundamental domain (the Wigner -Seitz cell) of the lattice $`\mathrm{\Gamma }=^n`$, and by $`B=[\pi ,\pi ]^n`$ the first Brillouin zone, which is a fundamental domain of the reciprocal (dual) lattice $`\mathrm{\Gamma }^{}=\left(2\pi \right)^n`$. We identify $`\mathrm{\Gamma }`$-periodic functions, in the natural way, with functions on $`𝕋^n`$.
We introduce now the set that plays the role of the characteristic variety $`\mathrm{\Sigma }`$ discussed in the introduction.
###### Definition 2
The (complex) Fermi surface $`F_P`$ of the operator $`P`$ (at the zero energy level) consists of all vectors $`k^n`$ (called quasimomenta) such that the equation $`Pu=0`$ has a nonzero Bloch solution $`u(x)=e^{ikx}p(x)`$, where $`p(x)`$ is a $`\mathrm{\Gamma }`$-periodic function.
It would be useful later on to realize that in this definition the positivity of the solution is not required, and in fact the solution is usually complex.
In many cases, it is convenient to introduce a spectral parameter $`\lambda `$. This leads to the notion of the Bloch variety:
###### Definition 3
The (complex) Bloch variety $`B_P`$ of the operator $`P`$ consists of all pairs $`(k,\lambda )^{n+1}`$ such that the equation $`Pu=\lambda u`$ has a nonzero Bloch solution $`u(x)=e^{ikx}p(x)`$ with the quasimomentum $`k`$.
It is clear that the Fermi surface is just the projection onto the $`k`$-space of the intersection of the Bloch variety with the hyperspace $`\lambda =0`$.
One can consider the Bloch variety $`B_P`$ as the graph of a (multivalued) function $`\lambda (k)`$, which is usually called the dispersion relation. Then the Fermi surfaces become the level surfaces of the dispersion relation. Since the spectra of all operators $`P(x,D+k)`$ on the torus $`𝕋^n`$ are discrete, we can single out continuous branches $`\lambda _j`$ of this multivalued dispersion relation. These branches are usually called the band functions (see ).
The following analyticity property of the Fermi and Bloch varieties is important:
###### Lemma 4
\[30, Theorems 3.1.7 and 4.4.2\] The Fermi and Bloch varieties are the sets of all zeros of entire functions of a finite order in $`^n`$ and $`^{n+1}`$, respectively.
Another property of the Bloch and Floquet varieties that we will need later is the relation between the corresponding varieties of the operators $`P`$ and $`P^{}`$.
###### Lemma 5
\[30, Theorem 3.1.5\] A quasimomentum $`k`$ belongs to $`F_P^{}`$ if and only if $`kF_P`$. Analogously, $`(k,\lambda )B_P^{}`$ if and only if $`(k,\lambda )B_P`$. In other words, the dispersion relations $`\lambda (k)`$ and $`\lambda ^{}(k)`$ for the operators $`P`$ and $`P^{}`$ are related as follows:
$$\lambda ^{}(k)=\lambda (k).$$
(2.1)
We note that the Fermi surface $`F_P`$ is periodic with respect to the reciprocal lattice $`\mathrm{\Gamma }^{}=(2\pi )^n`$. Therefore, it is sometimes useful to factor out the periodicity by considering the (analytic) exponential mapping $`\rho :^n(^{})^n`$, where
$$z=\rho (k)=\rho (k_1,\mathrm{},k_n)=(\mathrm{exp}ik_1,\mathrm{},\mathrm{exp}ik_n).$$
This mapping can be identified in a natural sense with the quotient map $`^n^n/\mathrm{\Gamma }^{}`$. We also introduce the complex torus
$$T=\rho (^n)=\left\{z^n\right|\left|z_j\right|=1,j=1,2,\mathrm{},n\}.$$
(2.2)
###### Definition 6
The image $`\mathrm{\Phi }_P=\rho (F_P)`$ of the Fermi surface $`F_P`$ under the mapping $`\rho `$ is called the Floquet surface of the operator $`P`$.
The reader familiar with the Floquet theory immediately recognizes the Floquet surface as the set of all Floquet multipliers of the equation $`Pu=0`$.
The main tool in the Floquet theory is an analog of the Fourier transform (see \[30, Section 2.2\], ), which we will call the Floquet transform $`𝒰`$ (it is sometimes also called the Gelfand transform):
$$f(x)𝒰f(z,x)=\underset{\gamma \mathrm{\Gamma }}{}f(x\gamma )z^\gamma ,z(^{})^n,$$
(2.3)
where we denote $`z^\gamma =z_1^{\gamma _1}z_2^{\gamma _2}\mathrm{}z_n^{\gamma _n}`$ .
It is often convenient to use for the Floquet transform $`𝒰`$ the quasimomentum coordinate $`k`$ instead of the multiplier $`z`$, where
$$z=\rho (k)=(\mathrm{exp}ik_1,\mathrm{},\mathrm{exp}ik_n).$$
We need to recall now some definitions from . For a point $`z(^{})^n`$, we denote by $`E_{m,z}`$ the closed subspace of the Sobolev space $`H^m(K)`$ formed by the restrictions of functions $`vH_{loc}^m(^n)`$ that satisfy the Floquet condition $`v(x+\gamma )=z^\gamma v(x)`$ for any $`\gamma \mathrm{\Gamma }`$. One can show (see Theorem 2.2.1 in ) that
$$_m:=\underset{z(^{})^n}{}E_{m,z}$$
(2.4)
forms a holomorphic sub-bundle of the trivial bundle $`(^{})^n\times H^m(K)`$. As any infinite dimensional analytic Hilbert bundle over a Stein domain, it is trivializable (see Theorems 1.3.2, 1.3.3, and 1.5.23 in ). One can also notice that for $`m=0`$ the bundle $`_0`$ coincides with the whole $`(^{})^n\times L^2(K)`$.
The following standard auxiliary result for the transform $`𝒰`$ collects several statements from Theorem XIII.97 in and Theorem 2.2.2 in :
###### Lemma 7
1. For any nonnegative integer $`m`$ the operator
$$𝒰:H^m(^n)L^2(T,_m)$$
is an isometric isomorphism, where $`L^2(T,_m)`$ denotes the space of square integrable sections over the complex torus $`T`$ of the bundle $`_m`$, equipped with the natural topology of a Hilbert space.
2. Let the space
$$\mathrm{\Theta }^m=\{fH_{loc}^m(R^n)|\underset{\gamma \mathrm{\Gamma }}{sup}f_{H^m(K+\gamma )}\mathrm{exp}(b|\gamma |)<\mathrm{},b>0\}$$
be equipped with the natural Fréchet topology. Then
$$𝒰:\mathrm{\Theta }^m\mathrm{\Gamma }((^{})^n,_m)$$
is an isomorphism, where $`\mathrm{\Gamma }((^{})^n,_m)`$ is the space of all analytic sections over $`(^{})^n`$ of the bundle $`_m`$, equipped with the topology of uniform convergence on compacta.
3. Let the elliptic operator $`P`$ be of order $`m`$. Then under the transform $`𝒰`$ the operator
$$P:H^m(^n)L^2(^n)$$
becomes the operator of multiplication by a holomorphic Fredholm morphism $`P(z)`$ between the fiber bundles $`_m`$ and $`_0`$. Here $`P(z)`$ acts on the fiber of $`_m`$ over the point $`zT`$ as the restriction to this fiber of the operator $`P`$ acting between $`H^m(K)`$ and $`L^2(K)`$.
Here is another standard way of looking at the morphism $`P(z)`$. Let $`z=\mathrm{exp}ik`$, then commuting with the exponent $`\mathrm{exp}ikx`$ one can (locally) trivialize the bundle $`_m`$ reducing it to the trivial bundle with the fiber $`H^m(𝕋^n)`$, where as before $`𝕋^n=^n/\mathrm{\Gamma }`$. At the same time the operator $`P(z)`$ takes the form $`P(x,D+k)`$ between Sobolev spaces on the torus $`𝕋^n`$.
Let us discuss the structure of the Floquet solutions (see Definition 1) and of functions of Floquet type (1.6) in general. For illustration, consider the constant coefficient case, where the role of the Floquet solutions is played by the exponential polynomials
$$e^{ikx}\underset{\left|j\right|N}{}p_jx^j.$$
It is well known that, considered as distributions, all such functions are Fourier transformed into distributions supported at the point $`\left(k\right)`$. Moreover, the converse statement is also true. A simple but extremely important and relatively unnoticed observation is that under the Floquet transform, each Floquet type function of the form (1.6) corresponds, in a similar way, to a (vector valued) distribution supported at the quasimomentum $`\left(k\right)`$. We collect below this fact and some other previously known properties of Floquet solutions, as well as a new result on the dimensions of the spaces of such solutions, which will play the crucial role in establishing the Liouville type theorems.
First of all, every Floquet type function $`u`$ (see (1.6)), being of exponential growth, determines a (continuous linear) functional on the space $`\mathrm{\Theta }^0`$. If, additionally, it satisfies the equation $`Pu=0`$ for a periodic elliptic operator of order $`m`$, then as such a functional it is clearly orthogonal to the range of the dual operator $`P^{}:\mathrm{\Theta }^m\mathrm{\Theta }^0`$. According to Lemma 7, after the Floquet transform any such functional becomes a functional on $`\mathrm{\Gamma }(\left(^{}\right)^n,_0)`$, which is orthogonal to the range of the Fredholm morphism $`P^{}(z):_m_0`$ generated by the dual operator $`P^{}:\mathrm{\Theta }^m\mathrm{\Theta }^0`$. We are now ready to formulate the following auxiliary result.
###### Lemma 8
1. A continuous linear functional $`u`$ on $`\mathrm{\Theta }^0`$ is generated by a function of the Floquet form (1.6) with a quasimomentum $`k`$ if and only if after the Floquet transform it corresponds to a functional on $`\mathrm{\Gamma }(\left(^{}\right)^n,_0)`$ which is a distribution $`\varphi `$ that is supported at the point $`\nu =\mathrm{exp}(ik)`$, i.e. has the form
$$\varphi ,f=\underset{\left|j\right|N}{}q_j,\frac{^{\left|j\right|}f}{z^j}|_\nu ,f\mathrm{\Gamma }(\left(^{}\right)^n,_0),$$
where $`q_jL^2(K)`$. The orders $`N`$ of the Floquet function (1.6) and of the corresponding distribution $`\varphi `$ are the same.
2. Let $`a_k`$ be the dimension of the kernel of the operator
$$P(x,D+k):H^m(𝕋^n)L^2(𝕋^n).$$
Then the dimension of the space of Floquet solutions of the equation $`Pu=0`$ of order at most $`N`$ with a quasimomentum $`k`$ is finite and does not exceed $`a_kq_{n,N}`$.
The estimate on the dimension given in the second part of Lemma 8 is very crude and in many cases can be significantly improved. In fact, as the following theorem shows, we obtain an explicit formula for the dimension of the space of Floquet solutions with a given quasimomentum in the case of a simple eigenvalue. This theorem seems to be new and constitutes the crucial part of the Liouville theorem proved in Section 4 (Theorem 23).
In order to formulate this result, we need to prepare some notions and notations.
###### Definition 9
Let $`Q`$ be a homogeneous polynomial in $`n`$ complex variables. A polynomial $`p(x)`$ in $`^n`$ is called $`Q`$-harmonic, if it satisfies the differential equation $`Q(D)p=0`$.
Let $`𝒫`$ denote the vector space of all polynomials in $`n`$ variables, and let $`P_l`$ be the subspace of all homogeneous polynomials of degree $`l`$. Denote by $`𝒫_N=\underset{l=0}{\overset{N}{}}P_l`$ the subspace of all polynomials of degree at most $`N`$. So, $`𝒫=\underset{l=0}{\overset{\mathrm{}}{}}P_l`$. If $`Q(k)`$ is a nonzero homogeneous polynomial of degree $`s`$, then the differential operator $`Q(D):P_{l+s}P_l`$ is surjective for any $`l`$ (this simple statement will also follow from the proof of the theorem below). Hence, the mapping $`Q(D):𝒫𝒫`$ has a (nonuniquely defined) linear right inverse $`R`$ that preserves the homogeneity of polynomials.
###### Theorem 10
Assume that zero is an eigenvalue of algebraic multiplicity $`1`$ of the operator $`P(x,D+k_0):H^m(𝕋^n)L^2(𝕋^n)`$ on the torus $`𝕋^n`$. Let $`\lambda (k)`$ be an analytic function in a neighborhood of $`k_0`$ such that $`\lambda (k)`$ is a simple eigenvalue of the operator $`P(x,D+k):H^m(𝕋^n)L^2(𝕋^n)`$ and $`\lambda (k_0)=0`$. Let
$$\lambda (k)=\underset{l=l_0}{\overset{\mathrm{}}{}}\lambda _l(kk_0)$$
be the Taylor expansion of $`\lambda (k)`$ around the point $`k_0`$ into homogeneous polynomials such that $`\lambda _{l_0}`$ is the first nonzero term of this expansion. Then for any $`N`$ the dimension of the space of Floquet solutions of the equation $`Pu=0`$ in $`^n`$ of order at most $`N`$ and with the quasimomentum $`k_0`$ is equal to the dimension of the space of all $`\lambda _{l_0}`$-harmonic polynomials of degree of at most $`N`$. Moreover, given a linear right inverse $`R`$ of the mapping $`\lambda _{l_0}(D):𝒫𝒫`$ that preserves homogeneity, one can construct an explicit isomorphism between these spaces.
Proof. It is sufficient to consider the case $`k_0=0`$, since the general case reduces to this by a change of variables. Consider the operator family
$$A(k)=P^{}(x,Dk)\lambda (k):H^m(𝕋^n)L^2(𝕋^n)$$
which is analytic in a neighborhood of $`0`$. At each point $`k`$ of this neighborhood $`A(k)`$ has by the construction a one-dimensional kernel. Then, according to Theorem 1.6.13 in , there exists an analytic non-vanishing vector $`\psi (k)KerA(k)`$. In other words, $`P^{}(x,Dk)\psi (x,k)=\lambda (k)\psi (x,k)`$. Let us choose a closed complementary subspace $`M`$ to $`KerA(0)`$ in $`H^m(𝕋^n)`$. Then it is complementary to $`KerA(k)`$ in a neighborhood of $`0`$. Since $`P^{}(x,D)`$ has zero kernel on $`M`$ and is Fredholm, we conclude that $`P^{}(x,Dk)`$ has zero kernel on $`M`$ for all $`k`$ in a neighborhood of $`0`$. We denote by $`\mathrm{\Pi }(k)`$ the closed subspace in $`L^2(𝕋^n)`$ defined as $`\mathrm{\Pi }(k)=P^{}(x,Dk)(M)`$. Applying Theorem 1.6.13 of again, we conclude that $`\mathrm{\Pi }(k)`$ depends holomorphically on $`k`$ in a neighborhood of $`0`$ (i.e., forms a Banach bundle) and hence it is complementary to $`KerA(k)`$. Representing now the operator $`P^{}(x,Dk)`$ in the block form according to the decompositions
$$H^m(𝕋^n)=MKerA(k)$$
and
$$L^2(𝕋^n)=\mathrm{\Pi }(k)KerA(k),$$
we get
$$P^{}(x,Dk)=\left(\begin{array}{cc}B(k)& 0\\ 0& \lambda (k)\end{array}\right),$$
where $`B(k)`$ is an analytic invertible operator-function between $`M`$ and $`\mathrm{\Pi }(k)`$. If we now have a functional $`\varphi `$ on $`\mathrm{\Gamma }(\left(^{}\right)^n,_0)`$ supported at $`v=\mathrm{exp}(0)`$, such that it is orthogonal to the range of the operator of multiplication by $`P^{}(k)`$, then it must be equal to zero on all sections of the bundle $`\mathrm{\Pi }(k)`$. This means that the restriction of such functionals to the sections of the one-dimensional bundle $`KerA(k)`$ is an one-to-one mapping. This reduces the problem to the following scalar one: find the dimension of the space of all distributions of order $`N`$ supported at the origin such that they are orthogonal to the ideal generated by $`\lambda (k)`$ in the ring of germs of analytic functions. One can change variables to eliminate the minus sign in front of $`k`$. Due to the finiteness of the order of the distribution, the problem further reduces to the following: find the dimension of the cokernel of the mapping
$$\mathrm{\Lambda }_N:𝒫_N𝒫_N,$$
where $`\mathrm{\Lambda }_N(p)`$ for $`p𝒫_N`$ is the Taylor polynomial of order $`N`$ at $`0`$ of the product $`\lambda (k)p(k)`$. Let us write the block matrix $`\mathrm{\Lambda }_{ij}`$ of the operator $`\mathrm{\Lambda }_N`$ that corresponds to the decomposition $`𝒫_N=\underset{l=0}{\overset{N}{}}P_l`$. It is obvious that $`\mathrm{\Lambda }_{ij}=0`$ for $`ij<l_0`$. For $`ijl_0`$ the entry $`\mathrm{\Lambda }_{ij}`$ is the operator of multiplication by $`\lambda _{ij}`$ acting from $`P_j`$ into $`P_i`$. Since $`\lambda _{l_0}0`$, for $`ij=l_0`$ the operator $`\mathrm{\Lambda }_{ij}`$ of multiplication by $`\lambda _{l_0}`$ has zero kernel. Being interested in the cokernel of $`\mathrm{\Lambda }_N`$, we need to find the kernel of the adjoint matrix $`\mathrm{\Lambda }_N^{}`$. The adjoint matrix acts in the space $`\underset{l=0}{\overset{N}{}}P_l^{}`$, where $`P_l^{}`$ can be naturally identified with the space of linear combinations of the derivatives of order $`l`$ of the Dirac’s delta-function at the origin. Here we have $`\mathrm{\Lambda }_{ij}^{}=0`$ for $`ji<l_0`$, and for $`jil_0`$ the entry $`\mathrm{\Lambda }_{ij}^{}`$ is the dual to the operator of multiplication by $`\lambda _{ji}`$ acting from $`P_i`$ into $`P_j`$. In particular, since for $`ji=l_0`$ the latter operator is injective, we conclude that the operators $`\mathrm{\Lambda }_{ij}^{}`$ are surjective. This enables one to find the dimension of the kernel of the matrix $`\mathrm{\Lambda }_N^{}`$ and even to describe its structure. Namely, let
$$\psi =(\psi _0,\mathrm{},\psi _N)\underset{l=0}{\overset{N}{}}P_l^{}$$
be such that $`\mathrm{\Lambda }^{}\psi =0`$. Due to the triangular structure of $`\mathrm{\Lambda }_N^{}`$, it is easy to solve this system. Indeed, it can be written as follows:
$$\underset{ji+l_0}{}\mathrm{\Lambda }_{ij}^{}\psi _j=0,i=0,\mathrm{},Nl_0.$$
Taking the Fourier transform, we can rewrite this system in the form
$$\underset{ji+l_0}{}\lambda _{ji}(D)\widehat{\psi _j}=0,i=0,\mathrm{},Nl_0,$$
where $`\widehat{\psi }`$ denotes the Fourier transform of $`\psi `$. Therefore, $`\widehat{\psi _j}`$ is a homogeneous polynomial of degree $`j`$ in $`^n`$. For $`i=Nl_0`$ we have
$$\lambda _{l_0}(D)\widehat{\psi _N}=0.$$
This equality means that $`\widehat{\psi _N}`$ can be chosen as an arbitrary $`\lambda _{l_0}`$-harmonic homogeneous polynomial of order $`N`$. Moving to the previous equation, we analogously obtain
$$\lambda _{l_0}(D)\widehat{\psi _{N1}}+\lambda _{l_0+1}(D)\widehat{\psi _N}=0,$$
or
$$\lambda _{l_0}(D)\widehat{\psi _{N1}}=\lambda _{l_0+1}(D)\widehat{\psi _N}.$$
The right hand side is already determined, and the nonhomogeneous equation, as we concluded before, always has a solution, for instance
$$R\left(\lambda _{l_0+1}(D)\widehat{\psi _N}\right).$$
This means that
$$\widehat{\psi _{N1}}+R\left(\lambda _{l_0+1}(D)\widehat{\psi _N}\right)$$
is a $`\lambda _{l_0}`$-harmonic homogeneous polynomial of order $`N1`$. We see that the solution $`\widehat{\psi _{N1}}`$ exists and is determined up to an addition of any homogeneous $`\lambda _{l_0}`$-harmonic polynomial of degree $`N1`$. Continuing this process until we reach $`\widehat{\psi _0}`$, we conclude that the mapping
$$\psi =(\psi _0,\mathrm{},\psi _N)\varphi =(\varphi _0,\mathrm{},\varphi _N),$$
where
$$\varphi _j=\widehat{\psi _j}+R\underset{i>j}{}\lambda _{ij+l_0}(D)\widehat{\psi _i}$$
establishes an isomorphism between the cokernel of the mapping $`\mathrm{\Lambda }_N`$ and the space of $`\lambda _{l_0}`$-harmonic polynomials of degree at most $`N`$. This proves the theorem.
In the cases of the simplest structures of the Taylor series, the theorem implies the following:
###### Corollary 11
Under the hypotheses of Theorem 10 the following hold:
1. If $`k_0`$ is a noncritical point of the band function $`\lambda (k)`$, then the dimension of the space of Floquet solutions of the equation $`Pu=0`$ in $`^n`$ of order at most $`N`$ with a quasimomentum $`k_0`$ is equal to the dimension $`q_{n1,N}`$ of the space of all polynomials of degree at most $`N`$ in $`^{n1}`$.
2. If the Taylor expansion of the band function $`\lambda (k)`$ at a point $`k_0`$ starts with a nondegenerate quadratic form, then the dimension of the space of Floquet solutions of the equation $`Pu=0`$ in $`^n`$ of order at most $`N`$ with a quasimomentum $`k_0`$ is equal to the dimension $`h_{n,N}`$ of the space of harmonic (in the standard sense) polynomials of degree at most $`N`$ in $`^n`$. In particular, this condition is satisfied at nondegenerate extrema.
In both cases an isomorphism can be provided explicitly as in the previous theorem.
Proof: 1. By our assumptions, the Taylor expansion of $`\lambda (k)`$ starts with a nonzero linear term $`\lambda _1(k)=_{j=1}^na_jk_j,a_j`$. The corresponding differential operator is
$$\lambda _1(D)=i\underset{j=1}{\overset{n}{}}a_j\frac{}{x_j}=\underset{j=1}{\overset{n}{}}\alpha _j\frac{}{x_j}+i\underset{j=1}{\overset{n}{}}\beta _j\frac{}{x_j},$$
where $`\alpha _j`$ and $`\beta _j`$ are real. Consider first the case when the vectors $`\alpha =(\alpha _j)`$ and $`\beta =(\beta _j)`$ are collinear. Then $`\lambda _1(D)`$ becomes $`\gamma _0\gamma _j\frac{}{x_j}`$, where $`\gamma _00`$ is a complex number and $`\gamma =(\gamma _j)`$ is a nonzero real vector. A linear change of coordinate system brings $`\lambda _1(D)`$ to the operator $`\frac{}{x_1}`$ (up to an irrelevant constant factor). Thus, the $`\lambda _1`$-harmonic polynomials are exactly those independent on $`x_1`$. Invoking Theorem 10, we get our conclusion in this case. Consider now the situation when $`\alpha `$ and $`\beta `$ are linearly independent. Then a linear change of variables brings $`\lambda _1`$ to the form $`/\overline{z}`$, where $`z=x_1+ix_2`$. Since any polynomial in variables $`(x_1,\mathrm{},x_n)`$ is a polynomial of the same degree in $`(z,\overline{z},x_3,\mathrm{},x_n)`$, the $`\lambda _1`$-harmonic polynomials are the ones depending on $`(z,x_3,\mathrm{},x_n)`$ only (i.e. the ones analytic in $`z`$). This again reduces the number of variables to $`n1`$.
2. By our assumptions, the first nonzero homogeneous term is a nondegenerate quadratic form $`\lambda _2(kk_0)`$, which is reducible to the sum of squares of coordinates by a linear change of variables. Therefore, in the new coordinates $`\lambda _2(D)=\mathrm{\Delta }`$. Using Theorem 10, we obtain the desired result.
In the remaining part of this section we restrict further the form of the operator. Namely, we consider now second order operators with real periodic coefficients of the form
$$L=\underset{i,j=1}{\overset{n}{}}a_{ij}(x)_i_j+\underset{i=1}{\overset{n}{}}b_i(x)_i+c(x).$$
(2.5)
It is assumed that the uniform ellipticity condition
$$\underset{i,j=1}{\overset{n}{}}a_{ij}(x)\zeta _i\zeta _ja\underset{i=1}{\overset{n}{}}\zeta _i^2$$
is satisfied for all $`x,\zeta ^n`$, where $`a`$ is a positive constant.
For such operators, we introduce the function that will play the crucial role in our considerations. Its properties were studied in detail in , , and . Consider the function $`\mathrm{\Lambda }(\xi ):^n`$ defined by the condition that the equation
$$Lu=\mathrm{\Lambda }(\xi )u$$
has a positive Bloch solution of the form
$$u_\xi (x)=e^{\xi x}p_\xi (x),$$
(2.6)
where $`p_\xi (x)`$ is $`\mathrm{\Gamma }`$-periodic.
###### Lemma 12
1. The value $`\mathrm{\Lambda }(\xi )`$ is uniquely determined for any $`\xi ^n`$.
2. The function $`\mathrm{\Lambda }(\xi )`$ is bounded from above, strictly concave, analytic, and has a nonzero gradient at all points except at its maximum point.
3. Consider the operator
$$L(\xi )=e^{\xi x}Le^{\xi x}=L(x,Di\xi )$$
on the torus $`𝕋^n`$. Then $`\mathrm{\Lambda }(\xi )`$ is the principal eigenvalue of $`L(\xi )`$ with a positive eigenfunction $`p_\xi `$. Moreover, $`\mathrm{\Lambda }(\xi )`$ is algebraically simple.
4. The Hessian of $`\mathrm{\Lambda }(\xi )`$ is nondegenerate at all points.
One should note that since the function $`\mathrm{\Lambda }(\xi )`$ is analytic, it is actually defined in a neighborhood of $`^n`$ in $`^n`$. This remark will be used in what follows.
Let us denote
$$\mathrm{\Lambda }_0=\underset{\xi ^n}{\mathrm{max}}\mathrm{\Lambda }(\xi ).$$
(2.7)
It follows from that an alternative definition of $`\mathrm{\Lambda }_0`$ is
$$\mathrm{\Lambda }_0=sup\{\lambda |u>0\text{ such that }(L\lambda )u=0\text{ in }^n\},$$
(2.8)
and that in the self-adjoint case $`\mathrm{\Lambda }_0`$ coincides with the bottom of the spectrum of the operator $`L`$. The common name for $`\mathrm{\Lambda }_0`$ is the generalized principal eigenvalue of the operator $`L`$ in $`^n`$.
We will often need to assume that $`\mathrm{\Lambda }_0`$ is either nonnegative or strictly positive. In the self-adjoint case such an assumption has a clear spectral interpretation. In the next lemma, we provide some known sufficient conditions for the nonnegativity or positivity of $`\mathrm{\Lambda }_0`$ for operators of the form (2.5).
###### Lemma 13
Consider an operator $`L`$ of the form (2.5)
1. $`\mathrm{\Lambda }_00`$ if and only if the operator $`L`$ admits a positive (super)solution. This condition is satisfied in particular when $`c(x)0`$.
2. $`\mathrm{\Lambda }_00`$ if and only if the operator $`L`$ admits a positive solution of the form (2.6).
3. $`\mathrm{\Lambda }_0=0`$ if and only if the equation $`Lu=0`$ admits exactly one normalized positive solution in $`^n`$.
4. If $`c(x)=0`$, then $`\mathrm{\Lambda }_0=0`$ if and only if $`\underset{𝕋^n}{}b(x)\psi (x)𝑑x=0`$, where $`\psi `$ is the principal eigenfunction of $`L^{}`$ on $`𝕋^n`$. In particular, divergence form operators satisfy this condition.
5. Let $`\xi ^n`$, and assume that $`u_\xi (x)=e^{\xi x}p_\xi (x)`$ and $`u_\xi ^{}`$ are positive Bloch solutions of the equations $`Lu=0`$ and $`L^{}u=0`$, respectively. Denote by $`\psi `$ the periodic function $`u_\xi u_\xi ^{}`$. Consider the function
$$\stackrel{~}{b}_i(x)=b_i(x)2\underset{j=1}{\overset{n}{}}a_{ij}(x)\{\xi _j+(p_\xi (x))^1_jp_\xi (x)\},$$
and denote
$$\gamma =(\gamma _1,\mathrm{},\gamma _n):=(\underset{𝕋^n}{}\stackrel{~}{b}_1(x)\psi (x)𝑑x,\mathrm{},\underset{𝕋^n}{}\stackrel{~}{b}_n(x)\psi (x)𝑑x).$$
Then $`\mathrm{\Lambda }_0=0`$ if and only if $`\gamma =0`$.
Let us discuss also some additional properties that will play an important role in the sequel. Assume that $`\mathrm{\Lambda }_00`$. Then Lemma 12 implies that the zero level set
$$\mathrm{\Xi }=\left\{\xi ^n\right|\mathrm{\Lambda }(\xi )=0\}$$
(2.9)
is either a strictly convex compact analytic surface in $`^n`$ of dimension $`n1`$ (this is the case if and only if $`\mathrm{\Lambda }_0>0`$), or a singleton (this is the case if and only if $`\mathrm{\Lambda }_0=0`$). The manifold $`\mathrm{\Xi }`$ consists of all $`\xi ^n`$ such that the equation $`Lu=0`$ admits a positive Bloch solution $`u_\xi (x)=e^{\xi x}p_\xi (x)`$. Moreover, the set of all such positive Bloch solutions is the set of all minimal positive solutions of the equation $`Lu=0`$ in $`^n`$ . It is also established that a function $`u`$ is a positive solution of the equation $`Lu=0`$ in $`^n`$ if and only if there exists a positive finite measure $`\mu `$ on $`\mathrm{\Xi }`$ such that
$$u(x)=_\mathrm{\Xi }u_\xi (x)𝑑\mu (\xi ).$$
We denote by $`G`$ the convex hull of $`\mathrm{\Xi }`$, and by $`\stackrel{}{G}`$ its interior. Note that if $`\mathrm{\Lambda }_00`$ then $`\mathrm{\Lambda }_0=0`$ if and only if $`\mathrm{\Xi }=G`$ and hence $`\stackrel{}{G}=\mathrm{}`$.
###### Lemma 14
Suppose that $`\mathrm{\Lambda }_0>0`$. There exists a neighborhood $`W`$ of $`G`$ in $`^n`$ and an analytic function
$$W\xi p_\xi ()H^2(𝕋^n)$$
such that for any $`\xi W`$ the function of $`x`$
$$u_\xi (x)=\mathrm{exp}(\xi x)p_\xi (x)$$
is a nonzero Bloch solution of the equation $`Lu=\mathrm{\Lambda }(\xi )u`$ with a quasimomentum $`i\xi `$. Moreover, one can choose the function $`p`$ in such a way that it is positive for all $`\xi \mathrm{\Xi }`$.
Comparing the definitions of $`\mathrm{\Xi }`$ and of the Fermi surface $`F_L`$, it follows that
$$i\mathrm{\Xi }F_L.$$
The next lemma specifies further the relation between these two varieties:
###### Lemma 15
Let $`\mathrm{\Lambda }_00`$. Then
1. The intersection of the complex Fermi surface $`F_L`$ with the tube
$$𝒯=\left\{k^n\right|Imk=(Imk_1,\mathrm{},Imk_n)G\}$$
(2.10)
coincides with the union of the surface $`i\mathrm{\Xi }`$ with its translations by the vectors of the reciprocal lattice $`\mathrm{\Gamma }^{}`$, i.e. consists of vectors $`k=i\xi +\gamma `$ where $`\xi \mathrm{\Xi }`$ and $`\gamma \mathrm{\Gamma }^{}`$. Moreover, up to a multiplicative constant, any nonzero Bloch solution with a quasimomentum in the above intersection is a positive Bloch solution.
2. If $`\mathrm{\Lambda }_0>0`$, then the intersection of $`F_L`$ with a sufficiently small neighborhood of $`i\mathrm{\Xi }`$ is a (smooth) analytic manifold that coincides with the set of zeros of the function $`\mathrm{\Lambda }(ik)`$.
Analogously to the definition of the Floquet surface $`\mathrm{\Phi }=\mathrm{\Phi }_L`$, we define the surface
$$\mathrm{\Psi }=\rho (i\mathrm{\Xi })=\left\{z\right|z=(\mathrm{exp}\xi _1,\mathrm{},\mathrm{exp}\xi _n),\xi \mathrm{\Xi }\},$$
(2.11)
and the tubular domain
$$V=\rho (𝒯),$$
(2.12)
where $`𝒯`$ was defined in (2.10). The results of Lemmas 14 and 15 can be rephrased in terms of these objects:
###### Lemma 16
Let $`\mathrm{\Lambda }_00`$. Then
1. $`\mathrm{\Phi }V=\mathrm{\Psi }`$.
If $`\mathrm{\Lambda }_0>0`$, then
2. The intersection of $`\mathrm{\Phi }`$ with a sufficiently small neighborhood of $`\mathrm{\Psi }`$ is a (smooth) connected analytic manifold.
3. The intersections of $`\mathrm{\Phi }`$ with neighborhoods of the tube $`V`$ form a basis of neighborhoods of $`\mathrm{\Psi }`$ in $`\mathrm{\Phi }`$.
4. For a sufficiently small neighborhood $`\mathrm{\Phi }_\epsilon `$ of $`\mathrm{\Psi }`$ in $`\mathrm{\Phi }`$ there exists an analytic function $`p:\mathrm{\Phi }_\epsilon H^2(𝕋^n)`$ such that for any $`z\mathrm{\Phi }_\epsilon `$ the function of $`x`$
$$u_z(x)=z^xp(z,x)$$
is a nonzero Bloch solution of the equation $`Lu=0`$.
We will also employ the following lemma:
###### Lemma 17
Consider an operator $`L`$ of the form (2.5)
1. Assume that $`c(x)0`$. Then the only solutions of the equation $`Lu=0`$ of the type $`\mathrm{exp}(ikx)p(x)`$, where $`k^n`$ and $`p`$ is a $`\mathrm{\Gamma }`$-periodic function are the constants. If such a nontrivial solution exists, then $`c(x)=0`$, and $`\mathrm{\Lambda }(0)=0`$ (i.e. $`0\mathrm{\Xi }`$).
2. Suppose that the operator $`L`$ admits a positive periodic supersolution $`\psi C^{2,\alpha }(^n)`$. Assume that $`v(x)=\mathrm{exp}(ikx)p(x)`$ is a nontrivial solution of the equation $`Lu=0`$, where $`k^n`$ and $`p`$ is a periodic function. Then there exists $`C`$ such that $`v=C\psi `$, the function $`\psi `$ is a positive periodic solution, and $`\mathrm{\Lambda }(0)=0`$.
## 3 Representation of solutions by hyperfunctions
The main result of this section (Theorem 18) is analogous to Theorem 5.1 in , which characterizes the class of solutions of the Helmholtz equation that can be represented by means of hyperfunctions on $`S`$ (see also the introduction of our paper).
In order to state it, we need to introduce a new object. Let us denote by $`h(\omega )`$, $`\omega S^{n1}`$ the indicator function of the convex set $`G`$. Namely,
$$h(\omega )=\underset{\xi G}{sup}(\omega \xi ),$$
(3.1)
where $`\omega \xi =_{j=1}^n\omega _j\xi _j`$ is the inner product in $`^n`$. The next Theorem will be stated in terms of this function.
###### Theorem 18
Suppose that $`\mathrm{\Lambda }_0>0`$. Let $`u`$ be a solution of the equation $`Lu=0`$ in $`^n`$ satisfying for any $`\epsilon >0`$ the estimate
$$\left|u(x)\right|C_\epsilon \mathrm{exp}\left(\left(h(x/\left|x\right|)+\epsilon \right)\left|x\right|\right),$$
(3.2)
where $`C_\epsilon `$ is a constant depending only on $`\epsilon `$ and $`u`$. Then $`u`$ can be represented as
$$u(x)=<\mu (\xi ),u_\xi (x)>,$$
(3.3)
where $`u_\xi `$ is the analytic positive Bloch solution corresponding to $`\xi \mathrm{\Xi }`$ (see Lemma 14), and $`\mu (\xi )`$ is a hyperfunction (analytic functional) on $`\mathrm{\Xi }`$. The converse statement is also true: for any hyperfunction $`\mu `$ on $`\mathrm{\Xi }`$, the function $`u(x)`$ in (3.3) is a solution of the equation $`Lu=0`$ in $`^n`$ which satisfies the growth condition (3.2).
###### Remark 19
Using a standard elliptic argument it follows that the pointwise growth condition (3.2) is equivalent to the growth condition
$$u(x)\mathrm{exp}\left(\left(h(x/\left|x\right|)+\epsilon \right)\left|x\right|\right)L^2(^n).$$
(3.4)
Proof of Theorem 18: Assume first that a solution $`u`$ has the representation (3.3). We need to prove that $`u`$ satisfies the growth condition (3.2). Due to the real analyticity of $`u_\xi `$ with respect to $`\xi `$ and according to lemmas 14 and 15, $`u_\xi `$ can be extended to an analytic vector function $`u_\xi (x)=\mathrm{exp}(\xi x)p_\xi (x)`$ on an $`\epsilon `$-neighborhood $`U_\epsilon `$ of $`\mathrm{\Xi }`$ in $`iF_L`$. Since $`\mu `$ is a hyperfunction (analytic functional) on $`\mathrm{\Xi }`$, we have an estimate
$$\left|u(x)\right|C_\epsilon \underset{\xi U_\epsilon }{\mathrm{max}}\left|u_\xi (x)\right|.$$
Hence we have
$$\left|u(x)\right|C_\epsilon \underset{\xi U_\epsilon }{\mathrm{max}}\left|e^{\xi x}\right|=C_\epsilon e^{\left|x\right|(h(x/\left|x\right|)+\delta (\epsilon ))},$$
where $`lim_{\epsilon 0}\delta (\epsilon )=0`$, which gives (3.2).
Suppose now that $`u`$ satisfies (3.2). We need to prove that $`u`$ can be represented as in (3.3). Let $`G_\epsilon `$ be the $`\epsilon `$-neighborhood of $`G`$ and $`h_\epsilon =h+\epsilon `$ be the indicator function of $`\overline{G_\epsilon }`$. Consider the following Fréchet spaces of test functions:
$$W_{m,\epsilon }=\left\{\varphi H_{loc}^m(^n)\right|<\varphi >_{m,\delta }<\mathrm{},\mathrm{\hspace{0.33em}0}<\delta <\epsilon \},$$
where
$$<\varphi >_{m,\delta }:=\underset{\gamma \mathrm{\Gamma }}{sup}\left\{\varphi _{H^m(K+\gamma )}e^{(h_\delta (\gamma /\left|\gamma \right|)\left|\gamma \right|)}\right\}.$$
It is obvious that the operator $`L^{}`$ maps continuously $`W_{2,\epsilon }`$ into $`W_{0,\epsilon }`$. Consequently, the linear functional
$$<u,\varphi >:=_^nu(x)\varphi (x)dx$$
is continuous on the space $`W_{0,\epsilon }`$ for any $`\epsilon >0`$. Since $`Lu=0`$, Schauder elliptic estimates together with the periodicity of the operator show that estimates similar to (3.2) hold also for the derivatives of $`u`$. One observes that $`u`$ is a continuous functional on $`W_{0,\epsilon }`$ which annihilates the range of the operator $`L^{}:W_{2,\epsilon }W_{0,\epsilon }`$. Now Floquet theory arguments analogous to the ones used in \[30, Section 3.2\] can be applied to yield (3.3). Let us make this part more precise.
Our first goal is to obtain a Paley-Wiener type theorem for the Floquet transform in the spaces $`W_{m,\epsilon }`$. Let us denote by $`V_\epsilon `$ the domain in $`(^{})^n`$
$`V_\epsilon =\{z=(z_1,\mathrm{},z_n)(^{})^n|`$ $`z_j=\mathrm{exp}ik_j\text{ such that}`$
$`Imk=(Imk_1,\mathrm{},Imk_n)(G_\epsilon )\}.`$
and let
$`V_\epsilon ^{}=\left\{z=(z_1,\mathrm{},z_n)(^{})^n\right|z^1=(z_1^1,\mathrm{},z_n^1)V_\epsilon \},`$
The domains $`V_\epsilon `$ form a basis of neighborhoods of the tube $`V`$, where $`V`$ is defined by (2.12). The following statement is a Paley-Wiener type theorem for the transform $`𝒰`$ which is suitable for our purpose.
###### Lemma 20
1. The operator
$$𝒰:W_{m,\epsilon }\mathrm{\Gamma }(V_\epsilon ^{},_m)$$
is an isomorphism, where $`\mathrm{\Gamma }(V_\epsilon ^{},_m)`$ is the space of holomorphic sections over $`V_\epsilon ^{}`$ of the bundle $`_m`$, equipped with the topology of uniform convergence on compacta.
2. Under the transform $`𝒰`$, the operator
$$L^{}:W_{2,\epsilon }W_{0,\epsilon }$$
becomes the operator $`(z)`$ of multiplication by a holomorphic Fredholm morphism between the fiber bundles $`_2`$ and $`_0`$:
$$\mathrm{\Gamma }(V_\epsilon ^{},_2)\stackrel{(z)}{}\mathrm{\Gamma }(V_\epsilon ^{},_0).$$
Here $`(z)`$ acts on each fiber of $`_2`$ as the restriction to this fiber of the operator $`L^{}`$ acting between $`H^2(K)`$ and $`L^2(K)`$.
Let us choose a value $`\epsilon _0>0`$ such that the intersection of $`\mathrm{\Phi }`$ with $`V_\epsilon `$ is smooth and connected. This is possible according to Lemma 16. From now on, we will only consider the values $`0<\epsilon <\epsilon _0`$.
Since the image $`𝒰u`$ of the solution $`u`$ under the Floquet transform $`𝒰`$ is a continuous linear functional on $`\mathrm{\Gamma }(V_\epsilon ^{},_0)`$ which is in the cokernel of the operator
$$\mathrm{\Gamma }(V_\epsilon ^{},_2)\stackrel{(z)}{}\mathrm{\Gamma }(V_\epsilon ^{},_0),$$
our task is to describe all such functionals. Several theorems of this kind were proven in . In our current situation such a representation can be obtained rather easily, due to the simplicity of the structure of the Floquet variety inside $`V_\epsilon `$. Namely, let $`u_z()=z^xp(z,)`$ be the Bloch solution of the equation $`Lu=0`$ introduced in Lemma 16. Let also $`(\mathrm{\Phi }_\epsilon )`$ be the space of holomorphic functions on $`\mathrm{\Phi }_\epsilon =\mathrm{\Phi }V_\epsilon `$ equipped with the topology of uniform convergence on compacta. We introduce the mapping
$$t:\mathrm{\Gamma }(V_\epsilon ^{},_0)(\mathrm{\Phi }_\epsilon )$$
which for any section $`f(z,x)`$ of the bundle $`_0`$ produces
$$t_f(z)=<f(z^1,),u_z>=\underset{𝕋^n}{}f(z^1,x)u_z(x)dx.$$
Here $`z^1=(z_1^1,\mathrm{},z_n^1)`$.
###### Lemma 21
Let $`0<\epsilon <\epsilon _0`$, where $`\epsilon _0`$ is the value defined above. Then the mapping $`t`$ is a topological homomorphism and the following sequence is exact:
$$\mathrm{\Gamma }(V_\epsilon ^{},_2)\stackrel{(z)}{}\mathrm{\Gamma }(V_\epsilon ^{},_0)\stackrel{t}{}(\mathrm{\Phi }_\epsilon )0.$$
This lemma practically finishes the proof of the theorem. Namely, the solution $`u`$ after the Floquet transform leads to a continuous linear functional on $`\mathrm{\Gamma }(V_\epsilon ^{},_0)`$ that annihilates the range of the operator of multiplication by $`(z)`$. Lemma 21 implies that such a functional can be pushed down to the space $`(\mathrm{\Phi }_\epsilon )`$. Since this functional, due to the estimate (3.2), is continuously extendable to $`(\mathrm{\Phi }_\epsilon )`$ for arbitrarily small values of $`\epsilon `$, it is in fact a hyperfunction (analytic functional) $`\mu `$ on $`\mathrm{\Phi }=\underset{\epsilon >0}{}\mathrm{\Phi }_\epsilon `$. Hence, the action $`<u,\varphi >`$ of the functional $`u`$ on a function $`\varphi W_{0,\epsilon }`$ can be obtained as
$$<u,\varphi >=<\mu (z),t(z)(𝒰\varphi )>.$$
Applying now the explicit formulas for the transforms $`𝒰`$ and $`t`$, one arrives to the representation (3.3). Indeed,
$`t_{(𝒰\varphi )}(z)={\displaystyle \underset{K}{}}𝒰\varphi (z^1,x)u_z(x)𝑑x`$ (3.5)
$`={\displaystyle \underset{\gamma \mathrm{\Gamma }}{}}{\displaystyle \underset{K\gamma }{}}\varphi (x)z^\gamma u_z(x+\gamma )𝑑x`$
$`={\displaystyle \underset{^n}{}}\varphi (x)u_z(x)𝑑x.`$
In this calculation we used the property of the Bloch solutions
$$u_z(x+\gamma )=z^\gamma u_z(x).$$
Therefore,
$$<u,\varphi >=<<\mu (z),u_z>,\varphi >,$$
which concludes the proof of the theorem.
## 4 Liouville-type theorem
In this section we discuss Liouville theorems for periodic equations. We will consider at the moment an arbitrary linear elliptic operator $`P(x,D)`$ with smooth $`\mathrm{\Gamma }`$-periodic coefficients which satisfies the assumptions made in Section 2 (as above, without loss of generality we can reduce the consideration to the case $`\mathrm{\Gamma }=^n`$).
###### Definition 22
We say that the Liouville theorem holds true for the operator $`P`$, if for any $`N`$ the space $`V_N(P)`$ of solutions of the equation $`Pu=0`$ in $`^n`$ that can be estimated as
$$u_{L^2(K+\gamma )}C(1+\left|\gamma \right|)^N\text{ for all }\gamma \mathrm{\Gamma }$$
is finite dimensional.
In the case when the Liouville theorem holds, we will be also interested in the dimensions $`d_N`$ of the spaces $`V_N(P)`$ and in representations of their elements analogous to (1.3).
The result below explains under what conditions on the operator $`P`$ a Liouville-type theorem holds. These conditions will then be verified for some specific classes of operators.
As was mentioned in the introduction, solutions representable as (1.3) are just Floquet solutions with zero quasimomentum. So, the Liouville theorem of cited in the introduction states that any polynomially growing solution is a Floquet solution with a zero quasimomentum. Let us also mention that any Bloch solution $`e^{ikx}p(x)`$ with a real quasimomentum $`k`$ is automatically bounded. This means that the validity of the Liouville theorem for an operator $`P`$ implies that the number of the real quasimomenta of solutions of the equation $`Pu=0`$ must be finite (modulo the action of the reciprocal lattice). In other words, the Fermi surface for $`P`$ intersects the real space at a finite number of points (modulo the reciprocal lattice). In terms of the Floquet variety it means that the set $`𝒵:=\mathrm{\Phi }_PT`$ is finite. We denote the cardinality of a set $`A`$ by $`\mathrm{\#}A`$. As the second statement of the next theorem shows, the finiteness of $`𝒵`$ is in fact the only claim of the Liouville theorem.
###### Theorem 23
1. The equation $`Pu=0`$ has a nonzero polynomially growing solution if and only if it has a nonzero bounded Bloch solution, i.e. if and only if the intersection $`F_L^n`$ of the Fermi surface for $`P`$ with the real space is not empty (or equivalently, $`𝒵=\mathrm{\Phi }_PT\mathrm{}`$).
2. The Liouville theorem holds for the operator $`P`$ if and only if the intersection $`F_P^n`$ is a finite set modulo the reciprocal lattice (or equivalently, $`\mathrm{\#}𝒵<\mathrm{}`$). Moreover, if $`\mathrm{\#}𝒵=\mathrm{}`$ then the Liouville theorem does not hold even for bounded solutions, i.e., $`d_0=dim(V_0)=\mathrm{}`$.
3. If the Liouville theorem holds, then each solution $`uV_N(P)`$ can be represented as a finite sum of Floquet solutions:
$$u(x)=\underset{qF_P^n}{}e^{iqx}\underset{|j|N}{}x^jp_{j,q}(x).$$
(4.1)
4. If the Liouville theorem holds, then for all $`N0`$ we have
$$d_Nd_0q_{n,N}<\mathrm{},$$
where $`q_{n,N}`$ is the dimension of the space of all polynomials of degree at most $`N`$ in $`n`$ variables.
5. Assume that the Liouville theorem holds and that for each real quasimomentum $`q`$ (i.e., for each $`qF_P^n`$) the conditions of Theorem 10 are satisfied. Then for each $`N0`$ the dimension $`d_N`$ of the space $`V_N(P)`$ is equal to the sum over $`q(F_P^n)/\mathrm{\Gamma }^{}`$ of the dimensions of the spaces of $`\lambda _q`$-harmonic polynomials (see Definition 9), where $`\lambda _q`$ is the first nonzero homogeneous term in the Taylor expansion at the point $`q`$ of the dispersion relation (band function) $`\lambda (k)`$.
Proof: Statements $`4`$ and $`5`$ follow from $`3`$ together with Lemma 8 and Theorem 10. So, we first prove statements $`2`$ and $`3`$ and conclude with the proof of the first statement.
In order to prove $`2`$ let us notice that if $`\mathrm{\#}𝒵=\mathrm{}`$ then each point $`z=\mathrm{exp}ik𝒵`$ provides a bounded Bloch solution with the quasimomentum $`k`$, and these solutions are linearly independent. This means that the Liouville theorem cannot hold in this case.
Assume now that $`\mathrm{\#}𝒵<\mathrm{}`$. We need to prove that the Liouville theorem and representation (4.1) hold true. Obviously, if $`u`$ has a representation of the form (4.1), then $`u`$ is of a polynomial growth. The proof that any polynomially growing solution is of the form (4.1) follows the same simple strategy as in the proofs of Theorem 18 and as is in the proof of the main Floquet representation \[30, Theorem 3.2.1\] (which, in turn, comes from the approach of and ). As in the case with the “fundamental principle” (see and ), it is more convenient to deal with a dual formulation, as it is done in . Namely, any polynomially growing solution $`u(x)`$ can be interpreted in the dual way, as a functional on an appropriate functional space, which belongs to the cokernel of the dual operator $`P^{}`$. Consequently, a representation theorem for all such functionals must be obtained. In order to make this idea precise, we need to introduce appropriate test functions spaces.
Consider the Fréchet spaces
$$C_m=\left\{\varphi H_{loc}^m(^n)\right|\underset{\gamma \mathrm{\Gamma }}{sup}\varphi _{H^m(K+\gamma )}(1+\left|\gamma \right|)^N<\mathrm{},N\}.$$
Let the order of the operator $`P`$ be $`m`$, then it is clear that $`P^{}`$ maps continuously $`C_m`$ into $`C_0`$. Due to the polynomial growth of $`u(x)`$, the linear functional
$$<u,\varphi >=_^nu(x)\varphi (x)𝑑x$$
is continuous on $`C_0`$. Since $`Pu=0`$, one easily observes that $`u`$ annihilates the range of the operator $`P^{}:C_mC_0`$. We need now a Paley-Wiener type theorem for the spaces $`C_m`$ with respect to the Floquet transform.
###### Lemma 24
1. The operator
$$𝒰:C_mC^{\mathrm{}}(T,_m)$$
is an isomorphism, where $`C^{\mathrm{}}(T,_m)`$ is the space of $`C^{\mathrm{}}`$ sections of the bundle $`_m`$ over the complex torus $`T`$, equipped with the standard topology.
2. Under the transform $`𝒰`$, the operator
$$P^{}:C_mC_0$$
becomes the operator $`𝒫(z)`$ of multiplication by a holomorphic Fredholm morphism between the fiber bundles $`_m`$ and $`_0`$:
$$C^{\mathrm{}}(T,_m)\stackrel{𝒫(z)}{}C^{\mathrm{}}(T,_0).$$
Here $`𝒫(z)`$ acts on each fiber of $`_m`$ as the restriction to this fiber of the operator $`P^{}`$ acting between $`H^m(K)`$ and $`L^2(K)`$.
3. The operator $`𝒫(z)`$ is invertible for a point $`zT`$ if and only if $`z^1\mathrm{\Phi }`$.
The next lemma is an analog of the classical theorem on the structure of distributions supported at a single point. Together with the previous lemma it essentially leads to the statement of the theorem.
###### Lemma 25
Let $`T`$ be a $`C^{\mathrm{}}`$-manifold and $`𝒫:TL(B_1,B_2)`$ be a $`C^{\mathrm{}}`$-function with values in the space $`L(B_1,B_2)`$ of bounded linear operators between Banach spaces $`B_1`$ and $`B_2`$. Assume that for each $`zT`$ the operator $`𝒫(z)`$ is a Fredholm operator. Then
1. If $`𝒫(z)`$ is surjective for all points $`z`$ in $`T`$, then the multiplication operator
$$C^{\mathrm{}}(T,B_1)\stackrel{𝒫(z)}{}C^{\mathrm{}}(T,B_2)$$
is surjective.
2. If $`𝒫(z)`$ is surjective for all points $`z`$ except a finite subset $`𝒵T`$, then any continuous linear functional $`g`$ on the space of smooth vector functions $`C^{\mathrm{}}(T,B_2)`$ that annihilates the range of the multiplication operator
$$C^{\mathrm{}}(T,B_1)\stackrel{𝒫(z)}{}C^{\mathrm{}}(T,B_2)$$
has the form
$$<g,\varphi >=\underset{z𝒵}{}\left[\underset{jN}{}D_{j,z}(<g_{j,z},\varphi >)\right]_z.$$
(4.2)
Here $`g_{j.z}`$ are continuous linear functionals on $`B_2`$, $`<g_{j,z},\varphi >`$ denotes the duality between $`B_2^{}`$ and $`B_2`$, $`D_{j,z}`$ are linear differential operators on $`T`$, and $`N`$.
We are ready now to finish the proof of the nontrivial part of the third statement of Theorem 23.
If $`u`$ is a solution of polynomial growth, it belongs, as it has been mentioned already, to the cokernel of the operator $`P^{}:C_mC_0`$. After the Floquet transform we are dealing with the cokernel of the operator
$$C^{\mathrm{}}(T,_2)\stackrel{𝒫(z)}{}C^{\mathrm{}}(T,_0).$$
By Lemma 24, the only points $`zT`$ where $`𝒫(z)`$ is not invertible are those points where $`z^1`$ belongs to the Floquet variety. Since by our assumption the set $`𝒵=T\mathrm{\Phi }`$ is finite, it follows that the operator function $`𝒫(z)`$ satisfies all the assumptions of Lemma 25. The fact that we are dealing with Banach bundles instead of fixed Banach spaces is irrelevant, since these bundles are trivial. This means that we have the representation (4.2) with $`g_jL^2(𝕋^n)`$. According to Lemma 8, functionals of the form (4.2) correspond under the inverse Floquet transform exactly to functions of the form (1.3).
It remains to prove the first statement of the theorem. Let $`u`$ be a polynomially growing solution. Assume that $`𝒵=\mathrm{}`$, i.e., the intersection of the Floquet variety $`\mathrm{\Phi }`$ with the complex torus $`T`$ is empty. Therefore, the last statement of Lemma 24 implies the invertibility of $`P(z)`$ for all $`zT`$. Now, the first statement of Lemma 25 guarantees the surjectivity of the mapping
$$C^{\mathrm{}}(T,_m)\stackrel{𝒫(z)}{}C^{\mathrm{}}(T,_0)$$
and hence the absence of any nontrivial functionals on $`C^{\mathrm{}}(T,_0)`$ that annihilate the image of this mapping. Since under the Floquet transform $`𝒰`$, a polynomially growing solution $`u(x)`$ is mapped to such a functional, we conclude that $`u=0`$.
###### Remark 26
The first statement of Theorem 23 is a part of the analog of the Bloch theorem provided in Theorem 4.3.1 of . Namely, the existence of a sub-exponentially (in particular, polynomially) growing solution implies the existence of a Bloch solution with a real quasimomentum, and hence the nonemptiness of the real Fermi variety. For completeness, we gave above an independent proof of this statement.
One realizes now that the cases when a Liouville-type theorem holds in a nonvacuous way are extremely rare. Namely, Theorem 23 shows that this happens only when the Fermi variety touches the real subspace at a finite set of points (modulo the reciprocal lattice). This means in particular, that in the selfadjoint case, one should expect this to happen only at the edges of the spectral gaps. Although it is possible to imagine interior points of the spectrum where such a thing could occur, it is hard to believe that these cases could be anything more than accidents.
One can expect the following conjecture to be true:
###### Conjecture 27
Let $`P`$ be a “generic” self-adjoint second order elliptic operator with periodic coefficients and $`(\lambda _{},\lambda _+)`$ be a nontrivial gap in its spectrum. Then each of the gap’s endpoints is a unique (modulo the dual lattice) and nondegenerate extremum of a single band function $`\lambda _j(k)`$.
The validity of this conjecture together with Theorem 23 would imply that generically at the gap ends the dimension of the space $`V_N`$ is equal to the dimension $`h_{n,N}`$ of the space of all harmonic polynomials of order at most $`N`$ in $`n`$ variables. Unfortunately, the only known theorem of this kind is the recent result of , which states that generically a gap edge is an extremum of a single band function.
At the bottom of the spectrum, however, much more is known. The theorem below combines some results of with the statement of Theorem 23 to obtain the structure and dimension of the space of polynomially growing solutions in this case. Below the spectrum, the Liouville theorem holds vacuously, according to the first statement of Theorem 23 and Theorem 5.5.1 in .
###### Theorem 28
1. Let $`H=\mathrm{\Delta }+V(x)`$ be a Schrödinger operator with a periodic real valued potential $`VL^{r/2}(𝕋^n),r>n`$. Then the lowest band function $`\lambda _1(k)`$ has a unique nondegenerate minimum $`\mathrm{\Lambda }_0`$ at $`k=0`$. All other band functions are strictly greater than $`\mathrm{\Lambda }_0`$. Every solution $`uV_N(H\mathrm{\Lambda }_0)`$ is representable in the form (1.3). The dimension of the space $`V_N(H\mathrm{\Lambda }_0)`$ is equal to $`h_{n,N}`$.
2. Let $`V`$ be like in the previous statement, then there exists $`ϵ>0`$ such that for any periodic real valued magnetic potential $`A`$ such that
$$A_{L^r(𝕋^n)}<ϵ$$
and
$$\underset{𝕋^n}{}A(x)𝑑x=0$$
(4.3)
the following statements hold true: The lowest band function $`\lambda _1(k)`$ of the magnetic Schrödinger operator $`H=(i+A)^2+V`$ attains a unique nondegenerate minimum $`\mathrm{\Lambda }_0`$ at a point $`k_0`$. All other band functions are strictly greater than $`\mathrm{\Lambda }_0`$. Every solution $`uV_N(H\mathrm{\Lambda }_0)`$ is representable in the Floquet form
$$v(x)=e^{ik_0x}\underset{|j|N}{}x^jp_j(x)$$
with periodic functions $`p_j(x)`$. The dimension of the space $`V_N(H\mathrm{\Lambda }_0)`$ is equal to $`h_{n,N}`$.
3. Suppose that $`L`$ is a second order elliptic operator of the form (2.5) such that $`\mathrm{\Lambda }_00`$.
If $`\mathrm{\Lambda }(0)=0`$ (i.e. $`0\mathrm{\Xi }`$), then the Liouville theorem holds and every solution $`uV_N(L)`$ is representable in the form (1.3). The dimension of the space $`V_N(L)`$ is equal to $`h_{n,N}`$ in the case when $`\mathrm{\Lambda }_0=0`$, and to $`q_{n1,N}`$ when $`\mathrm{\Lambda }_0>0`$.
If $`\mathrm{\Lambda }(0)>0`$ then the equation $`Lu=0`$ does not admit a nontrivial polynomially growing solution. So, the Liouville theorem holds vacuously.
Proof: 1. The result of says that the lowest band function $`\lambda _1(k)`$ has a unique nondegenerate minimum $`\mathrm{\Lambda }_0`$ at $`k=0`$ and that all other band functions are strictly greater than $`\mathrm{\Lambda }_0`$. Now Theorem 23 implies the rest of the claims of this statement.
2. When both the electric and magnetic potentials are sufficiently small, then the result of states that the lowest band function $`\lambda _1(k)`$ of the magnetic Schrödinger operator $`H=(i+A)^2+V`$ attains a unique nondegenerate minimum $`\mathrm{\Lambda }_0`$ at a point $`k_0`$, while all other band functions are strictly greater than $`\mathrm{\Lambda }_0`$. This statement, however, can be easily extended to the case of arbitrary electric and small magnetic potential. Indeed, when the magnetic potential is equal to zero, one can refer, as in the previous case, to . At this moment one has to use analyticity of the Bloch variety. Namely, the statement of Lemma 4 (see also \[30, Theorem 4.4.2\]) can be easily extended to include analyticity with respect to the potentials (see, for instance, ). More precisely, there exists an entire function $`f(k,\lambda ,A,V)`$ of all its arguments such that $`f(k,\lambda ,A,V)=0`$ is equivalent to
$$(k,\lambda )B_{(i+A)^2+V},$$
where $`B_H`$ is the Bloch variety of the operator $`H`$. Now, the result of for $`A=0`$ together with the stated analyticity property imply the required features of the lowest band function for sufficiently small magnetic potentials. The last step is to use again Theorem 23. Note that the normalization (4.3) always can be achieved by a gauge transformation which does not affect the spectrum and the Liouville property.
3. The assumption $`\mathrm{\Lambda }(0)0`$ implies that the operator $`L`$ admits a positive periodic supersolution. It follows from Lemma 17 that the Fermi surface $`F_L`$ can touch the real space only at the origin (modulo the reciprocal lattice $`\mathrm{\Gamma }^{}`$) and in this case $`\mathrm{\Lambda }(0)=0`$. Therefore, by the first part of Theorem 23, the Liouville Theorem holds vacuously if $`\mathrm{\Lambda }(0)>0`$.
Suppose now that $`\mathrm{\Lambda }(0)=0`$. Lemma 12 implies that if $`\mathrm{\Lambda }_0>0`$ then the point $`k=0`$ is a noncritical point of the dispersion relation, and if $`\mathrm{\Lambda }_0=0`$ then $`k=0`$ is a nondegenerate extremum. Now Theorem 23, as before, completes the proof.
## 5 Proofs of the lemmas
Proof of Lemma 8: The first claim of the lemma corresponds to Theorem 3.1.3 in . In order to prove the second part of the lemma, let us fix a $`k_0^n`$, and choose a closed subspace $`MH^m(𝕋^n)`$ complementary to the kernel of the operator $`P^{}(x,Dk_0)`$. Consider the (analytically depending on $`k`$ in a neighborhood of $`k_0`$) subspace
$$\mathrm{\Pi }(k):=P^{}(x,Dk)(M)L^2(𝕋^n).$$
and
$$𝒩:=\left[\mathrm{\Pi }(k_0)\right]^{}.$$
Then $`dim(𝒩)=a_{k_0}`$, and for values of $`k`$ close to $`k_0`$ the space $`𝒩`$ remains a complementary subspace to $`\mathrm{\Pi }(k)`$. Representing the operators $`P^{}(x,Dk)`$ in the matrix form according to the decompositions
$$H^m(𝕋^n)=MKerP^{}(x,Dk_0)$$
and
$$L^2(𝕋^n)=\mathrm{\Pi }(k)𝒩,$$
we get the matrix
$$\left(\begin{array}{cc}B(k)& \\ 0& C(k)\end{array}\right),$$
where $`B(k)`$ is an invertible analytic operator function, and $`C(k)`$ is an analytic matrix function of the size $`a_{k_0}\times a_{k_0}^{}`$. Here $`a_{k_0}^{}`$ is the dimension of the kernel of the operator $`P^{}(x,Dk_0)`$. (Notice that $`a_{k_0}=a_{k_0}^{}`$ if $`indP=0`$, which is true for instance, when dealing with scalar elliptic operators, due to the Atiyah-Singer theorem.) Now, the space of all distributions orthogonal to the range of $`P^{}`$ and supported at $`\mathrm{exp}(ik_0)`$ reduces to the space of all distributions supported at $`k_0`$, acting on $`^{a_k}`$-valued vector functions, and orthogonal to the range of the operator of multiplication by $`C(k)`$. If we drop the orthogonality condition, the dimension of the space of all such distributions of order at most $`N`$ is obviously equal to $`a_kq_{n,N}`$, which proves the estimate. We point out that a direct proof of this estimate for scalar operators can be also easily derived using the Leibnitz’s rule.
Proof of Lemma 12: Statements 1 through 3 of the lemma are contained in , except the statement that the geometric rather than the algebraic multiplicity of the eigenvalue $`\mathrm{\Lambda }(\xi )`$ is equal to one. The latter follows easily from Lemma 5.2 of . Alternatively, it can be deduced from general theorems on positive operators defined on an ordered Banach space (see for instance, \[29, Theorem 2.10\]). Statement 4 is proven in \[41, Theorem 5\].
Proof of Lemma 13: Statements 1–3 follow from the results of , while statements 4–5 follow from \[41, Theorem 5\].
Proof of Lemma 14: Consider the following family of operators on the torus: $`L(x,Di\xi )\mathrm{\Lambda }(\xi )`$. It follows from Lemma 12 that this family is analytic in a complex neighborhood $`W`$ of the set $`G`$ and its values are Fredholm operators between the appropriate Sobolev spaces. The same lemma implies that the dimension of the kernel of all these operators is equal to $`1`$. Hence, these kernels form an analytic fiber bundle over $`W`$ (see Theorem 1.6.13 and the corresponding references in ). One can always assume that the domain $`W`$ is convex (in the geometric sense). Then the kernel bundle (as all vector bundles on $`W`$) is topologically trivial. Since $`W`$, being convex, is a domain of holomorphy (see for instance Corollary 2.5.6 in ), therefore, the result of (an instance of the so called Oka’s principle) implies that the bundle is also analytically trivial. This means the existence of a nowhere zero analytic section $`u_\xi `$. Positivity of $`u_\xi `$ for $`\xi \mathrm{\Xi }`$ can be achieved as follows. Let us choose any nonzero analytic solution $`u_\xi `$ as above. Then for some small neighborhood $`W_1W`$ of $`G`$, we have $`u_\xi (0)0`$. So, we may normalize $`u_\xi `$ by dividing it by $`u_\xi (0)`$. The resulting solution is clearly positive for $`\xi \mathrm{\Xi }`$.
Proof of Lemma 15: 1. Let $`u(x)=e^{ikx}p(x)`$ be a nonzero Bloch solution, where $`p(x)`$ is a $`\mathrm{\Gamma }`$-periodic function, and $`kF𝒯`$. Assume first that $`Imk\stackrel{}{G}`$, so, $`\mathrm{\Lambda }_0>0`$. We need to prove that $`Reu=Imu=0`$. We show for instance, that $`u_1:=Reu=0`$. Suppose that $`u_10`$. We may assume that $`u_1(x_1)>0`$, for some $`x_1^n`$. Consider the positive solution
$$v(x)=\underset{\mathrm{\Xi }}{}u_\xi (x)𝑑\sigma (\xi ),$$
where $`d\sigma `$ is the $`(n1)`$-dimensional surface area element on $`\mathrm{\Xi }`$. For every $`M>0`$ there exists $`R>0`$ such that $`v(x)Mu_1(x)>0`$ for all $`|x|>R`$. By the generalized maximum principle, $`v(x)>Mu_1(x)`$ in $`^n`$. Since $`M`$ is arbitrarily large and $`u_1(x_1)>0`$, we arrived at a contradiction. Note that this argument applies also to any Floquet solution with a quasimomentum $`k`$ such that $`Imk\stackrel{}{G}`$.
Suppose now that $`Imk\mathrm{\Xi }`$ and $`\mathrm{\Lambda }0`$. Clearly, it is enough to show that there exists a real constant $`C`$ and $`\xi \mathrm{\Xi }`$ such that $`u_1:=Reu=Cu_\xi `$. Let $`\xi =Imk`$. Then for a sufficiently small $`\epsilon >0`$ the function $`v_\epsilon :=\frac{u_\xi }{2}\epsilon u_1`$ is a positive solution of the equation $`Lu=0`$, which is smaller than $`u_\xi `$. Recall that $`u_\xi `$ is a minimal positive solution of the equation $`Lu=0`$. Therefore, there exists $`c>0`$ such that $`v_\epsilon =cu_\xi `$, which implies that $`u_1=Cu_\xi `$ for some $`C`$.
2. Consider the zero set $`F_1`$ of the analytic function $`\mathrm{\Lambda }(ik)`$ in a small complex neighborhood of $`i\mathrm{\Xi }`$. Since $`\mathrm{\Lambda }_0>0`$, it follows that the gradient of $`\mathrm{\Lambda }(ik)`$ is not zero on $`i\mathrm{\Xi }`$. Therefore, $`F_1`$ is a smooth analytic variety. We will show that the Fermi surface $`F`$ coincides with $`F_1`$ in a neighborhood of $`i\mathrm{\Xi }`$, which will conclude the proof of the lemma. Indeed, obviously $`F_1F`$. Consider a point $`k_0=i\xi _0i\mathrm{\Xi }`$. By Lemma 12, zero is a simple eigenvalue of the operator $`L(x,D+k_0)=L(x,Di\xi _0)`$. This means that the spectral projector that corresponds to a neighborhood of zero is one-dimensional for all complex $`k`$ close to $`k_0`$. We conclude that for all $`k`$ in a complex neighborhood of $`i\mathrm{\Xi }`$ there is exactly one eigenvalue close to zero of the operator $`L(x,D+k)`$. By Lemma 14, we know this eigenvalue, namely $`\mathrm{\Lambda }(ik)`$. Let now $`k`$ belongs to a small neighborhood of $`i\mathrm{\Xi }`$ and assume that $`kF_1`$. Then $`\mathrm{\Lambda }(ik)0`$, and hence zero cannot be the eigenvalue of $`L(x,D+k)`$. This means that $`k`$ does not belong to the Fermi surface $`F`$.
Proof of Lemma 17: 1. If $`c0`$, the assertion of the lemma follows from \[36, Theorem 4.5\]. On the other hand, if $`c=0`$, then $`0i\mathrm{\Xi }`$, and in particular, $`0(G)`$. It follows from Lemma 15 that any Bloch solution with a real quasimomentum is the constant solution.
2. This assertion follows directly from the part 1 using the operator $`\psi ^1L\psi `$.
Proof of Lemma 20: The second statement of the lemma coincides with Theorem 2.2.3 in . So, we need to prove only the first statement.
Let $`\phi W_{m,\epsilon }`$. We will show that the series (2.3) converges uniformly on compacta in $`V_\epsilon ^{}`$ as a series of functions on $`V_\epsilon ^{}`$ with values in $`H^m(K)`$. This would imply that $`𝒰\phi \mathrm{\Gamma }(V_\epsilon ^{},H^m(K))`$, and that the corresponding (one-to-one) mapping $`𝒰:W_{m,\epsilon }\mathrm{\Gamma }(V_\epsilon ^{},H^m(K))`$ is continuous. Let $`0<\delta <\delta _1<\epsilon `$. Let $`z=\mathrm{exp}ikV_\delta ^{}`$ which means that, $`ImkG_\delta `$. We have
$`𝒰\phi (z,)_{H^m(K)}{\displaystyle \underset{\gamma \mathrm{\Gamma }}{}}\phi _{H^m(K\gamma )}e^{Imk\gamma }={\displaystyle \underset{\gamma \mathrm{\Gamma }}{}}\phi _{H^m(K+\gamma )}e^{Imk\gamma }`$
$`{\displaystyle \underset{\gamma \mathrm{\Gamma }}{}}\phi _{H^m(K+\gamma )}e^{(h(\gamma /\left|\gamma \right|)+\delta )\left|\gamma \right|}C_\delta <\phi >_{m,\delta _1}<\mathrm{}.`$
We need to check now that the mapping $`𝒰`$ acts from $`W_{m,\epsilon }`$ into $`\mathrm{\Gamma }(V_\epsilon ^{},_m)`$. This amounts to showing that $`𝒰\phi `$ satisfies the appropriate Floquet boundary conditions and hence is in fact a section of the sub-bundle $`_mV_\epsilon ^{}\times H^m(K)`$. This is a straightforward calculation (see also Theorem 2.2.2 in ).
On the other hand, let us assume that $`s(z)\mathrm{\Gamma }(V_\epsilon ^{},_m)`$. If $`z=\mathrm{exp}ik`$, then $`s`$ as a function of $`k`$ is periodic with respect to the reciprocal lattice $`\mathrm{\Gamma }^{}`$. Expanding it into the Fourier series, we get
$$s(z)=\underset{\gamma \mathrm{\Gamma }}{}s_\gamma z^\gamma \text{,}$$
where $`s_\gamma H^m(K)`$. We can now define a function $`\phi `$ on $`^n`$ such that $`\phi (x\gamma )=s_\gamma (x)`$ for $`xK`$ and $`\gamma \mathrm{\Gamma }`$.
The function $`\phi `$ belongs to $`H^m`$ in the interior of each of the cubes $`K+\gamma `$. One only needs to check that it belongs to $`H_{loc}^m`$ at the boundary points of these cubes. The requirement that $`s(z)`$ is a section of the bundle $`_m`$ rather than just of the bundle $`V_\epsilon ^{}\times H^m(K)`$ does exactly this (see the discussion at the top of page 96 in ).
It remains to show that $`\phi W_{m,\epsilon }`$. We use the standard formulas for the Fourier coefficients to get
$$\phi (\gamma )=s_\gamma =\frac{1}{(2\pi )^n}\underset{B}{}s(e^{i(\beta +i\alpha )})e^{i(\beta +i\alpha )\gamma }d\beta ,\alpha G_\epsilon ,$$
where $`B`$ is the first Brillouin zone, and we write $`z=\mathrm{exp}ik=\mathrm{exp}(i(\beta +i\alpha ))`$. Note that
$$\varphi _{H^m(K+\gamma )}\underset{zV_{\delta _1}^{}}{\mathrm{max}}s(z)_{H^m(K)}e^{\alpha (\gamma )}\alpha G_\epsilon ,$$
(5.1)
and therefore,
$$\varphi _{H^m(K+\gamma )}\underset{zV_{\delta _1}^{}}{\mathrm{max}}s(z)_{H^m(K)}e^{(h(\gamma /\left|\gamma \right|)+\delta _1)}.$$
(5.2)
This implies immediately that
$`<\phi >_{m,\delta }=\underset{\gamma \mathrm{\Gamma }}{sup}\left\{\phi _{H^m(K+\gamma )}e^{(h(\gamma /\left|\gamma \right|)+\delta )\left|\gamma \right|}\right\}`$ (5.3)
$`C\underset{zV_{\delta _1}^{}}{\mathrm{max}}s(z)_{H^m(K)}\underset{\gamma \mathrm{\Gamma }}{sup}e^{(\delta _1\delta )\left|\gamma \right|}<\mathrm{},`$
if $`\delta _1>\delta `$.
Proof of Lemma 21: The statement of this lemma is established in a much more general situation at the beginning of the proof of Theorem 1.7.1 in . However, for the sake of completeness we provide here the proof for our simpler particular situation. First of all, the sequence of the lemma is a complex (i.e., the composition of any two consecutive operators in it is equal to zero). One needs to prove this only in the second term of the sequence, where it follows immediately from the equality (3.5). Indeed, since $`u_z`$ solves the equation $`Lu=0`$, (3.5) followed by integration by parts proves the statement.
Let us turn to the exactness. We need to prove it in the second and third terms of the sequence. Consider the second term. Let $`f(z,x)\mathrm{\Gamma }(V_\epsilon ^{},_0)`$ be such that $`t_f(z)=0`$. This means that for any $`z\mathrm{\Phi }_\epsilon `$ the function $`f(z^1,)`$ is orthogonal to the Bloch solution $`u_z`$ of the equation $`Lu=0`$. We need to show that $`g(z)=(z^1)^1f(z)`$ is analytic, which will mean that $`f`$ belongs to the range of $``$. The function $`g(z^1)`$ is automatically analytic outside of $`\mathrm{\Phi }_ϵ`$, so we only need to make sure that it does not develop any singularities at this subset. We will show that all the necessary and sufficient conditions for the analyticity of $`g`$ have the form of orthogonality of values of $`f`$ at certain points to certain functionals. This would resolve the issue, since all such possible orthogonality conditions are the orthogonality of $`f(z^1)`$ to the kernel of $`L`$ on Bloch functions with a quasimomentum $`z`$, and hence to the vanishing of $`t_f(z)`$. As it was shown in the proof of Theorem 3.3.1 in \[30, pages 113-114\], the inverse operator to $`(z^1)`$ is the ratio of two analytic functions:
$$(z^1)^1=B(z)/\mathrm{\Delta }(z),$$
where $`B(z)`$ is an analytic function with values in bounded operators from $`L_2(𝕋^n)`$ to $`H^2(𝕋^n)`$, and $`\mathrm{\Delta }(z)`$ is a scalar analytic function, which is a regularized determinant of $`(z^1)(z_0^1)^1`$ for some point $`z_0`$ where the operator is invertible. Such regularized determinants are determined in the standard way by the eigenvalues of the corresponding operators (see for instance Section 2 of Chapter IV in for general definitions and properties of regularized determinants, and for our particular situation the proof of Theorem 3.1.7 and related discussion in Section 1.2 in ). The simplicity of the eigenvalue $`\mathrm{\Lambda }(\xi )`$ (Lemma 12) implies that if we introduce instead of $`z`$ the coordinate $`\xi `$ such that $`z=\mathrm{exp}\xi `$, then $`\mathrm{\Delta }(\mathrm{exp}\xi )=\mathrm{\Lambda }(\xi )\mathrm{\Delta }_1(\xi )`$, where $`\mathrm{\Delta }_1(\xi )`$ is an analytic function with no zeros in the domain under our consideration. We recall now that $`\mathrm{\Lambda }`$ has simple zeros. Hence, the necessary and sufficient condition for $`f`$ to belong to the range of the operator $``$ on the space of analytic sections is that the vector-function $`B(z)f(z)`$ vanishes on the set of the zeros of $`\mathrm{\Lambda }`$. These conditions obviously have the form of the orthogonality of values of $`f`$ to some functionals. As it was explained above, this implies exactness at the second term of the sequence.
Let us turn now to proving the exactness at the third term. We need to show that arbitrary analytic function on $`\mathrm{\Phi }_ϵ`$ can be obtained as $`t_f(z)`$ for some $`f\mathrm{\Gamma }(V_ϵ^{},_\mathcal{0})`$.
Let us denote by $`\mathrm{\Phi }_ϵ^{}`$ the manifold
$$\mathrm{\Phi }_ϵ^{}=\{z|z^1\mathrm{\Phi }_ϵ\}.$$
Consider the restriction mapping
$$\mathrm{\Gamma }(V_ϵ^{},_0)\mathrm{\Gamma }(\mathrm{\Phi }_ϵ^{},_0).$$
(5.4)
Notice that $`\mathrm{\Phi }_ϵ^{}`$ is an analytic subset in $`V_ϵ^{}`$ and that $`V_ϵ`$ and $`V_ϵ^{}`$ are domains of holomorphy. The latter can be easily proven using power test functions $`z^a`$ with integer (but not necessarily nonnegative) powers $`a`$ (a similar derivation can be found in the proof of the implication $`(iii)(i)`$ of Corollary 2.5.8 in ). Then Corollary 1 of the Bishop’s theorem \[43, Theorem 3.3\] (see the original theorem in ) claims that the restriction mapping (5.4) is surjective (recall that the bundle $`_0`$ is trivial). Hence, it is sufficient to prove that the mapping
$$\stackrel{~}{t}:\mathrm{\Gamma }(\mathrm{\Phi }_ϵ,_0)(\mathrm{\Phi }_ϵ).$$
defined as
$$\stackrel{~}{t}_f(z)=<f(z,),u_z>=\underset{𝕋^n}{}f(z,x)u_z(x)dx$$
is surjective. Consider the continuous operator $`T(z):L^2(K)`$ defined as $`T(z)y=<y,u_z>=\underset{𝕋^n}{}y(x)u_z(x)dx`$. Since $`u_z`$ is not zero, this operator is surjective. It is clear that it depends analytically on $`z`$. According to Allan’s theorem (see or Theorem 4.4 in ), since $`\mathrm{\Phi }_ϵ`$ is a Stein manifold, there exists an analytic right inverse operator $`R(z)`$. Now, given $`\varphi (z)(\mathrm{\Phi }_ϵ)`$, the function $`g(z)=R(z)\varphi (z)`$ satisfies $`\stackrel{~}{t}_g=\varphi `$. This proves the surjectivity that we need.
The last statement of the lemma about the mapping $`t`$ being a topological homomorphism is just the open mapping theorem.
Proof of Lemma 24: 1. We first show that the operator $`𝒰`$ maps continuously the space $`C_m`$ into $`C^{\mathrm{}}(T,H^m(K))`$. Indeed, if $`\phi C_m`$, then $`\phi _{H^m(K+\gamma )}`$ decays faster than any power of $`|\gamma |`$. This together with (2.3) leads to the immediate conclusion that $`𝒰\phi `$ belongs to $`C^{\mathrm{}}(T,H^m(K))`$ and to the continuity of the corresponding mapping. Since $`𝒰\phi `$ is a section of the sub-bundle $`^m`$ (see the Section 2.2 in ), this gives us the needed conclusion. Conversely, let
$$s(z)C^{\mathrm{}}(T,^m)C^{\mathrm{}}(T,H^m(K)).$$
One can expand the $`H^m(K)`$-valued function $`s(z)`$ into the Fourier series:
$$s(z)=\underset{\gamma \mathrm{\Gamma }}{}s_\gamma z^\gamma ,zT.$$
Here $`s_\gamma H^m(K)`$. Standard estimates of the Fourier coefficients of smooth functions apply, which show that $`s_\gamma `$ decays faster than any power of $`|\gamma |`$. Let us define now a function $`\varphi `$ on $`^n`$ such that $`\varphi (x\gamma )=s_\gamma (x)`$ for $`xK`$ and $`\gamma \mathrm{\Gamma }`$. The additional information that $`s`$ is a section of the sub-bundle $`^m`$ leads (as in \[30, page 96\]) to the conclusion that $`\varphi H_{loc}^m(^n)`$. This implies that $`\varphi C_m`$ and finishes the proof of the first statement of the lemma.
Statements 2 and 3 are correspondingly parts of Theorem 2.2.3 and 3.1.5 of .
Proof of Lemma 25: The first statement is rather obvious. Indeed, the statement is local, and locally one can construct a smooth one-sided inverse. The second statement can be proven like the similar statement in \[30, Corollary 1.7.2\]. For completeness, we provide the scheme of the proof here. Under the conditions of the second statement of the lemma, it is easy to see that any functional annihilating the range of the operator of multiplication by $`𝒫(z)`$ must be supported at the finite set $`𝒵`$ where $`𝒫(z)`$ is not surjective. This also reduces the considerations to a neighborhood $`U`$ of a point $`z_0𝒵`$. Using the Fredholm property, one can find a closed subspace $`M`$ of finite codimension in $`B_1`$ such that the operators $`𝒫(z)`$ have zero kernel on $`M`$ for all $`zU`$ (see the corresponding lemma in , or Lemma 1.2.11 and Remark 2 below it in ). Now the problem reduces to a similar one on a finite-dimensional space, where a standard representation of distributions supported at a point implies (4.2).
## 6 Further remarks
###### Remarks 6.1
1. Throughout the paper, we have assumed for simplicity that all the coefficients of the operators $`P`$ and $`P^{}`$ are $`C^{\mathrm{}}`$-smooth. In fact, we do not need such a restrictive assumption (see the discussion in \[30, Section 3.4.D\]). For example, a sufficient (but not necessary) condition for all the statements of Section 3 to hold true is that the coefficients of $`L`$ and $`L^{}`$ are Hölder continuous. Actually, even less is needed. For instance, conditions imposed on the Schrödinger operators in Theorem 28 are sufficient. It is clear that the conditions on the coefficients could be significantly relaxed, if the operators were considered in the weak sense, or by means of their quadratic forms. This should not change the general techniques of the proofs. We did not intend, however, to find the optimal requirements on the coefficients for all our results to hold.
2. It should be possible to describe the class of solutions of the equation $`Lu=0`$ that are representable by a distribution rather than by a hyperfunction. We plan to address this problem elsewhere.
3. The Liouville theorem can probably be extended to systems of equations (for instance, to the Maxwell system). In this case one would face the problems of a possibly nonzero index of the corresponding operator and of multiple eigenvalues (the latter can also occur for scalar operators). We believe that the technique of this paper might be adjusted to handle some of these situations. The extensions of the result of to the Pauli and Maxwell operators obtained in and would provide examples where the needed information on the behavior of the dispersion relations at the bottom of the spectrum is available.
## Appendix A Appendix
In this appendix we present an alternative proof of the third statement of Theorem 28 in the case when either $`\mathrm{\Lambda }_0=0`$ and $`N0`$, or $`\mathrm{\Lambda }_0>0`$ and $`0N1`$. The proof relies on some basic notions of homogenization theory and imitates the proof of Theorem 2 in , where $`L`$ is assumed to be an operator in divergence form. Therefore, we skip some details which are essentially the same as in .
We need to recall some basic definitions from homogenization theory (see, for example, ). Suppose that $`L`$ is a second order elliptic operator of the form
$$L=\underset{i,j=1}{\overset{n}{}}a_{ij}(x)_i_j+\underset{i=1}{\overset{n}{}}b_i(x)_i,$$
(A.1)
with periodic coefficients and denote the positive matrix $`\{a_{ij}(x)\}`$ by $`𝒜(x)`$ and the periodic vector $`(b_1,\mathrm{},b_n)^T`$ by $`b`$. Let $`\psi `$ be the positive normalized periodic solution of the equation $`L^{}u=0`$. Let $`\mathrm{\Psi }(x)=(\mathrm{\Psi }_1(x)\mathrm{},\mathrm{\Psi }_n(x))^T`$ be a solution of the equation
$$L\mathrm{\Psi }=b(x)+_{𝕋^n}b(x)\psi (x)𝑑x\text{ in }𝕋^n.$$
(A.2)
Consider the matrix
$$𝒬=\{q_{ij}\}:=_{𝕋^n}(I+\mathrm{\Psi })^T𝒜(x)(I+\mathrm{\Psi })\psi (x)𝑑x,$$
(A.3)
were $`I`$ is the identity matrix. The operator $`Q:=_{i,j=1}^nq_{ij}_i_j`$ is called the homogenized operator of the operator $`L`$, and the positive matrix $`𝒬=\{q_{ij}\}`$ is called the homogenized matrix (see, \[25, Section 2.5\]).
The following lemma, which is actually a new formulation of \[41, Theorem 5\]), establishes a connection between the function $`\mathrm{\Lambda }`$ and homogenization theory.
###### Lemma A.1
Let $`L`$ be an operator of the form (2.5) and suppose that $`\xi \mathrm{\Xi }`$. Let $`u_\xi `$ and $`u_\xi ^{}`$ be the positive Bloch solutions of the equations $`Lu=0`$ and $`L^{}u=0`$, respectively. Denote by $`\psi `$ the periodic function $`u_\xi u_\xi ^{}`$. Consider the operator
$$\stackrel{~}{L}=(u_\xi (x))^1Lu_\xi (x)=\underset{i,j=1}{\overset{n}{}}a_{ij}(x)_i_j+\underset{i=1}{\overset{n}{}}\stackrel{~}{b}_i(x)_i,$$
(A.4)
let
$$Q=\underset{i,j=1}{\overset{n}{}}q_{ij}_i_j$$
(A.5)
be the homogenized operator of the operator $`\stackrel{~}{L}`$, and $`𝒬=\{q_{ij}\}`$ be the homogenized matrix. Then $`\psi `$ is the principal eigenfunction of the operator $`\stackrel{~}{L}^{}`$ on the torus $`𝕋^n`$ with an eigenvalue $`0`$. Moreover, $`\text{Hess}(\mathrm{\Lambda }(\xi ))=𝒬`$.
Proof: The first statement of the lemma can be checked easily while the second statement follows directly from the formula in \[41, Theorem 5\], and the definition of the homogenized operator.
Proof of a part of the third statement of Theorem 28: We clearly may assume that $`L\mathrm{𝟏}=0`$, so,
$$L=\underset{i,j=1}{\overset{n}{}}a_{ij}(x)_i_j+\underset{i=1}{\overset{n}{}}b_i(x)_i.$$
We denote by $`\psi `$ the normalized positive solution of the equation $`L^{}u=0`$ in $`𝕋^n`$. Let $`\mathrm{\Psi }`$ be a solution of the system (A.2), and $`Q`$ be the homogenized operator of the operator $`L`$.
Assume first that $`\mathrm{\Lambda }_00`$. The case $`N=0`$ is trivial, and follows from Theorem 23 and Lemma 17. Let $`N=1`$. Recall that according to Theorem 23, $`d_1n+1`$. Moreover, by Theorem 23 and the Leibnitz’s rule, a (real) solution of linear growth is of the form
$$u(x)=\underset{j=1}{\overset{n}{}}a_jx_j+\varphi (x),$$
where $`a_j`$ and $`\varphi `$ is periodic.
By lemma 13, $`\mathrm{\Lambda }_0=0`$ if and only if for every $`1kn`$
$$\alpha _j:=\underset{𝕋^n}{}b_j(x)\psi (x)𝑑x=0.$$
(A.6)
For $`1jn`$, we write an “Ansatz” for a solution of linear growth of the form
$$F_j(x)=x_j+\varphi _j(x),$$
(A.7)
where $`\varphi _j`$ is a periodic function. Clearly, $`F_j`$ is a solution of $`Lu=0`$ in $`^n`$ if and only if $`\varphi _j(x)`$ solves the nonhomogeneous equation $`Lu=b_j`$ in $`𝕋^n`$. By the Fredholm alternative, this equation is solvable in $`𝕋^n`$ if and only if $`\alpha _j=0`$ which holds true for all $`1jn`$, if and only if $`\mathrm{\Lambda }_0=0`$ (and in this case, $`\varphi _j=\mathrm{\Psi }_j`$, see (A.2)). Therefore, $`d_1=n+1`$ if $`\mathrm{\Lambda }_0=0`$, and $`d_1<n+1`$ if $`\mathrm{\Lambda }_0>0`$.
In order to finish the proof for $`N=1`$, we need to prove that if $`\mathrm{\Lambda }_0>0`$, then $`d_1n`$. Without loss of generality, we may assume that $`\alpha _n0`$. We construct $`(n1)`$ linearly independent solutions of linear growth of the form
$$F_j(x)=x_j\alpha _j(\alpha _n)^1x_n+\varphi _j(x),$$
where $`1jn1`$, and $`\varphi _j`$ solves the equation $`Lu=b_j+\alpha _j(\alpha _n)^1b_n`$. Note that these $`(n1)`$ equations are solvable and therefore, $`d_1n`$.
For $`N2`$, we assume that $`\mathrm{\Lambda }_0=0`$. Recall that if $`uV_N`$ then by Theorem 23 and the Leibnitz’s rule
$$u(x)=u^{(N)}(x)+\underset{|\nu |<N}{}x^\nu p_\nu (x),$$
where
$$u^{(N)}(x)=\underset{|\nu |=N}{}x^\nu p_\nu ,$$
and $`p_\nu `$ are periodic functions if $`|\nu |<N`$, and $`p_\nu `$, if $`|\nu |=N`$.
Claim: Assume that $`\mathrm{\Lambda }_0=0`$. Then for all $`N0`$
$$Qu^{(N)}=0.$$
(A.8)
In particular, $`d_Nh_{n,N}`$.
Proof of the claim: Assume first that $`N=2`$. Then $`uV_2`$ is of the form
$$u(x)=\frac{1}{2}(Cxx)+\underset{j=1}{\overset{n}{}}x_jp_j(x)+p_0(x),$$
where $`C`$ is a constant symmetric matrix, and $`p_0,p_1,\mathrm{},p_n`$ are periodic functions.
A direct calculation shows that the vector $`p=(p_1,\mathrm{},p_n)^T`$ must satisfy the equation $`Lp=Cb`$ which is solvable since $`\mathrm{\Lambda }_0=0`$. Therefore, $`p=C\mathrm{\Psi }`$ (up to a constant vector). Also, $`p_0`$ must satisfy
$$Lp_0=f:=tr(A(I+2\mathrm{\Psi }^T)C^T)bC\mathrm{\Psi }.$$
The compatibility condition for this equation is $`\underset{𝕋^n}{}f(x)\psi (x)𝑑x=0`$ which after some calculations implies that
$$tr(𝒬C^T)=0,$$
where $`𝒬`$ is the homogenized matrix of the operator $`L`$ (see (A.3)). Since $`u^{(2)}:=\frac{1}{2}(Cxx)`$ is a homogeneous polynomial of degree $`2`$, it follows that $`Qu^{(2)}=tr(𝒬C^T)`$. Therefore, $`u^{(2)}`$ solves the equation $`Qu=0`$. Thus, the case $`N=2`$ is settled.
For $`N>2`$, we proceed by induction as in . Namely, assume that the claim (A.8) has been proven for $`N1`$, and let $`uV_N`$. Let $`\mathrm{\Delta }_i`$ be the difference operator $`\mathrm{\Delta }_if(x):=f(x+e_i)f(x)`$, where $`e_i`$ is the $`i`$-th vector of the standard basis of $`^n`$, and $`1in`$. Then $`v_i:=\mathrm{\Delta }_iuV_{N1}`$ and the leading part of $`v_i`$ is given by $`(\mathrm{\Delta }_iu)^{(N1)}=_iu^{(N)}`$. By the induction hypothesis, $`Q((\mathrm{\Delta }_iu)^{(N1)})=0`$. Therefore,
$$_i(Qu^{(N)})=Q(_iu^{(N)})=Q((\mathrm{\Delta }_iu)^{(N1)})=01in.$$
Hence, $`Qu^{(N)}=\text{const.}`$, and since $`Qu^{(N)}`$ is homogeneous of degree $`N2>0`$, we obtain that $`Qu^{(N)}=0`$, and the claim is proved.
It remains to prove that $`d_Nh_{n,N}`$. So, for any homogeneous polynomial $`h`$ of degree $`N`$ which is $`Q`$-harmonic, we need to find a solution $`uV_N`$ such that $`u^{(N)}=h`$. Let $`uV_N`$ and $`\epsilon >0`$. Consider the function
$$\epsilon ^Nu(\frac{x}{\epsilon })=\underset{|\nu |N}{}\epsilon ^{N|\nu |}x^\nu p_\nu (\frac{x}{\epsilon }),$$
which tends to $`u^{(N)}`$ as $`\epsilon 0`$. We consider $`x`$ and $`y=\frac{x}{\epsilon }`$ as independent variables and write
$$U(x,y,\epsilon ):=\underset{|\nu |N}{}\epsilon ^{N|\nu |}x^\nu p_\nu (y)=U_0(x)+\epsilon U_1(x,y)+\mathrm{}+\epsilon ^NU_N(x,y).$$
Then the equation $`L(x,_x)u=0`$ implies that
$$(L_0+\epsilon L_1+\epsilon ^2L_2)U=0,$$
where
$$L_0=L(y,_y);L_1=2\underset{i,j=1}{\overset{n}{}}a_{ij}(y)_{x_i,y_j}^2+\underset{i=1}{\overset{n}{}}b_i(y)_{x_i};L_2=\underset{i,j=1}{\overset{n}{}}a_{ij}(y)_{x_i,x_j}^2.$$
We look for a formal differential operator
$$\mathrm{\Phi }=\underset{j=0}{\overset{\mathrm{}}{}}\epsilon ^k\mathrm{\Phi }_j=\underset{\nu }{}\epsilon ^{|\nu |}\varphi _\nu (y)_x^\nu ,$$
where $`\varphi _\nu (y)`$ are periodic functions and $`\varphi _0=1`$. This operator should satisfy
$$(L_0+\epsilon L_1+\epsilon ^2L_2)\mathrm{\Phi }=M+L_0(y,_y)2\epsilon \underset{i,j=1}{\overset{n}{}}a_{ij}(y)_{x_i}_{y_j},$$
(A.9)
where the formal operator
$$M=\underset{j=2}{\overset{\mathrm{}}{}}\epsilon ^jM_j=\underset{|\nu |2}{}\epsilon ^{|\nu |}m_\nu _x^\nu ,$$
has constant coefficients.
Comparing the coefficients of $`\epsilon ^s`$ in (A.9) yields the following equations (the equation for $`s=0`$ is automatically satisfied).
$`L_0\mathrm{\Phi }_1+L_1=2{\displaystyle \underset{i,j=1}{\overset{n}{}}}a_{ij}(y)_{x_i}_{y_j},`$ $`s=1,`$ (A.10)
$`L_0\mathrm{\Phi }_s+L_1\mathrm{\Phi }_{s1}+L_2\mathrm{\Phi }_{s2}=M_s,`$ $`s2.`$ (A.11)
It is easily checked that for $`s=1`$ the functions $`\varphi _j(y)`$ of Equation (A.7) are the corresponding solutions for $`\mathrm{\Phi }_1`$. Also, Equation (A.11) for $`s=2`$ is solvable if $`M_2=Q`$, where $`Q`$ is the homogenized operator of $`L`$. Similarly, the constant coefficients of the operator $`M_s,s>2`$, are determined by the compatibility condition for Equation (A.11) with $`s>2`$.
Let $`R:𝒫𝒫`$ be a linear right inverse of the homogenized operator $`Q`$ that preserves the homogeneity of polynomials. Consider the formal operator $`A`$ which is defined by the equation
$$AI=R\underset{j=1}{\overset{\mathrm{}}{}}\epsilon ^jM_{j+2},$$
and let $`A^1`$ be its unique formal inverse. Note that $`\epsilon ^2M_2A=M`$.
Let $`U_0(x)`$ be a given homogeneous polynomial of degree $`N`$ which solves the equation $`Qu=0`$, and let $`V(x):=A^1U_0(x)`$. We have
$$MV=\epsilon ^2M_2AV=\epsilon ^2M_2U_0=0.$$
Define $`U(x,y,\epsilon ):=\mathrm{\Phi }A^1U_0=\mathrm{\Phi }V`$, and denote $`u(x):=U(x,x,1)`$. By inspection, $`u`$ has a polynomial growth of order $`N`$, and $`u^{(N)}(x)=U_0(x)`$. Moreover,
$`(L_0+\epsilon L_1+\epsilon ^2L_2)U=(L_0+\epsilon L_1+\epsilon ^2L_2)\mathrm{\Phi }V`$
$`=MV+L_0(y,_y)V(x)+2\epsilon {\displaystyle \underset{i,j=1}{\overset{n}{}}}a_{ij}(y)_{x_i}_{y_j}V(x)=0,`$
and $`\mathrm{\Phi }A^1`$ is the desired mapping.
###### Remark A.2
1. Let $`F_j`$ be the solutions of linear growth defined by Equation (A.7). A. Ancona proved that the map $`F(x)=(F_1(x),\mathrm{},F_n(x))`$ is a diffeomorphism on $`^n`$ if $`n2`$, while for $`n>2`$ this map is not necessarily a diffeomorphism.
2. Assume that $`L\mathrm{𝟏}=0`$ and $`\mathrm{\Lambda }_0=0`$. Let $`\mathrm{\Lambda }(\xi )=_{|\nu |2}a_\nu \xi ^\nu `$ be the Taylor expansion of the function $`\mathrm{\Lambda }`$. We conjecture that $`a_\nu =m_\nu `$, where $`m_\nu `$ are the coefficients of the operator $`M`$.
Acknowledgments
The authors express their gratitude to Professors S. Agmon and V. Lin for useful discussions and to Professor P. Li for the information about the manuscript .
The work of P. Kuchment was partially supported by the NSF Grant DMS 9610444 and by a DEPSCoR Grant. P. Kuchment expresses his gratitude to NSF, ARO, and to the State of Kansas for this support. The content of this paper does not necessarily reflect the position or the policy of the federal government of the USA, and no official endorsement should be inferred. The work of Y. Pinchover was partially supported by the Fund for the Promotion of Research at the Technion.
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# References
## List of figures
Figure 1. Cs. Benedek, Optics Letters
Figure 2. Cs. Benedek, Optics Letters
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# Langevin equation of collective modes of Bose-Einstein condensates in traps
## I Introduction
The realization of Bose-Einstein condensates of very rarefied evaporatively cooled gases of alkali atoms in magnetic traps offers the unique possibility to test ab initio many-body theories in the laboratory . One very fertile field has been the experimental and theoretical investigation of collective modes of the condensates, both for zero and finite temperatures. For reviews of the experimental and theoretical work see and respectively. As opposed to conventional superfluids like He-II in the new systems the collision-less regime is very naturally realized. In this regime the dominant damping mechanism for collective modes is Landau-damping, whose temperature dependence in spatially homogeneous condensates in the regime $`k_BT`$ large compared to the chemical potential $`\mu `$ has first been studied by Szépfalusy and Kondor . Recent investigations , , ,, though more exact, led to similar results, differing by a prefactor close to 1 for the damping rate. For condensates in traps Landau damping of low-lying modes is more difficult to calculate, and additional approximations are needed to cope with the fact that momentum is not conserved in a trap. The damping rate of collective modes in traps has been calculated in using the local density approximation and, in addition, a classical approximation for the correlation function whose Fourier-transform determines the cross-section of Landau-scattering. For the isotropic breathing mode in isotropic traps the Landau-damping has been calculated numerically by evaluating the coupling to a great number of discrete quasi-particle modes and subsequently introducing some smoothing. The quasi-continuum coupled to the collective mode under study was displayed explicitly in this work. Theories using an extension of the approach of via the dielectric formalism and an approach via a time-dependent mean field scheme have also been given.
In the present paper a very direct approach to the dissipative equilibrium and non-equilibrium dynamics of collective modes in trapped Bose-Einstein condensates via quantum Langevin equations is put forward. Because of the discreteness of the mode-spectrum in traps the problem is formally similar to the quantum-optical problem of discrete modes in a laser, for which the formulation in terms of quantum Langevin equations has been very useful .
In the next section the microscopic description of a trapped Bose-Einstein condensed gas is briefly set up. Then we recall the basics of the quantum Langevin equation of a boson mode. The derivation of the quantum Langevin equation of a collective mode follows. The damping rates are then evaluated in the local density and the Thomas-Fermi approximation. The last section contains a discussion of our results.
## II Microscopic equations of motion
The weakly interacting Bose-gas in a trap in standard notation is described by the Hamiltonian
$$\widehat{H}=d^3x\widehat{\psi }^+\left\{\frac{\mathrm{}^2}{2m}^2+V(𝒙)\mu +\frac{U_0}{2}\widehat{\psi }^+\widehat{\psi }\right\}\widehat{\psi }.$$
(1)
The presence of a Bose-Einstein condensate means that many ($`N_01`$) particles occupy a single normalized mode of a macroscopic classical matter wave, determined as the mode of lowest energy of the classical Hamiltonian corresponding to eq.(1). It satisfies a classical wave equation, the Gross-Pitaevskii equation , which we take in an extension defined by the so-called Popov-approximation , including the interaction of the condensate with the density $`n^{}`$ of thermal atoms, but neglecting its interaction with the pair amplitude $`\widehat{\psi }\widehat{\psi }\psi _0^2`$ of thermal particles
$$(\mathrm{}^2/2m)^2\psi _0+\left(V(𝒙)+U_0(N_0|\psi _0(𝒙)|^2+2n^{}(𝒙))\right)\psi _0=\mu \psi _0.$$
(2)
For given $`N_0`$ the chemical potential $`\mu `$ follows by imposing the normalization condition on $`\psi _0`$.
The presence of the highly occupied condensate mode makes the decomposition of the Heisenberg field-operator
$$\widehat{\psi }(𝒙,t)=\left(\sqrt{N_0}\psi _0(𝒙)+\widehat{\psi }^{}(𝒙,t)\right)\mathrm{exp}(i\mu t/\mathrm{})\mathrm{exp}(i\varphi )$$
(3)
useful. $`\sqrt{N_0}\mathrm{exp}(i\varphi )`$ is the complex amplitude of the classical condensate mode in equilibrium. $`\widehat{\psi }^{}(𝒙,t)`$ is the field operator for the particles outside the condensate. We shall neglect fluctuations of the number of atoms in the condensate and also fluctuations of the phase of the condensate, which can be shown to occur on a time-scale much longer than the relaxation-time of the collective modes .
The Hamiltonian splits up according to $`\widehat{H}=H_0+\widehat{H_1}+\widehat{H_2}+\widehat{H_3}+\widehat{H_4}`$, with a c-number term $`H_0`$ which need not concern us here, and
$`\widehat{H}_1=\sqrt{N_0}{\displaystyle d^3x\left(\left(V(𝒙)\mu +U_0(N_0|\psi _0|^2+2n^{})\right)\psi _0^{}\widehat{\psi }^{}+\left(hermitianconjugate\right)\right)}`$
$`\widehat{H}_2={\displaystyle d^3x\left(\widehat{\psi }^+(\frac{\mathrm{}^2^2}{2m}+V(𝒙)+2U_0n^{}\mu )\widehat{\psi }^{}+\frac{U_0N_0}{2}(\psi _0^2\widehat{\psi }^2+\psi _0^2\widehat{\psi }^{+2}+4|\psi _0|^2\widehat{\psi }^+\widehat{\psi }^{})\right)}`$
$`\widehat{H}_3=U_0\sqrt{N_0}{\displaystyle d^3x\left(\psi _0^{}(\widehat{\psi }^+\widehat{\psi }^{}2n^{})\widehat{\psi }^{}+(hermitianconjugate)\right)}`$
$`\widehat{H}_4={\displaystyle \frac{U_0}{2}}{\displaystyle }d^3x\widehat{(}\psi ^+\widehat{\psi }^+\widehat{\psi }^{}\widehat{\psi }^{}4n^{}\widehat{\psi }^+\widehat{\psi }^{}).`$
The splitting is here done in such a way that the term $`\widehat{H_1}`$ vanishes due to eq.(2) and the part $`\widehat{H_2}`$ describes the linearized quantum excitations around the solution of (2). $`\widehat{H_2}`$ is diagonalized by introducing quasi-particle operators $`\widehat{\alpha }_\nu ,\widehat{\alpha }_\nu ^+`$ by the standard Bogoliubov transformation
$`\widehat{\psi }^{}(𝒙,t)={\displaystyle \underset{\nu }{}}\left(u_\nu (𝒙)\widehat{\alpha }_\nu (t)+v_\nu ^{}(𝒙)\widehat{\alpha }_\nu ^+(t)\right),`$
where $`u_\nu `$, $`v_\nu `$ satisfy the usual Bogoliubov-De Gennes equations
$$\left(\begin{array}{cc}\frac{\mathrm{}^2}{2m}^2+U_{\mathrm{eff}}(𝒙)\mathrm{}\omega _\nu & K(𝒙)\\ K^{}(𝒙)& \frac{\mathrm{}^2}{2m}^2+U_{\mathrm{eff}}(𝒙)+\mathrm{}\omega _\nu \end{array}\right)\left(\genfrac{}{}{0pt}{}{u_\nu (𝒙)}{v_\nu (𝒙)}\right)=0,$$
(4)
with the abbreviations
$`U_{\mathrm{eff}}(𝒙)=V(𝒙)+2U_0(N_0|\psi _0(𝒙)|^2+n^{}(𝒙))\mu `$ (5)
$`K(𝒙)=N_0U_0\psi _0(𝒙)^2.`$ (6)
Eq. (4) is consistent with the ortho-normality conditions $`d^3x(u_\nu u_\mu ^{}v_\nu v_\mu ^{})=\delta _{\nu \mu }`$ and $`d^3r(u_\nu ^{}v_\mu u_\mu ^{}v_\nu )=0`$, which guarantee the Bose commutation relations of the $`\alpha _\nu `$, $`\alpha _\mu ^+`$. The decomposition of $`\widehat{\psi }`$ and $`\widehat{\psi }^{}`$ together imply that $`N=N_0+n^{}(𝒙)d^3x`$ with $`n^{}=_\mu \left(\overline{n}_\mu (|u_\mu |^2+|v_\mu |^2)+|v_\mu |^2\right)`$.
Within Bogoliubov(-Popov) theory the terms $`\widehat{H}_3,\widehat{H}_4`$ of the total Hamiltonian are neglected and the quasi-particle operators $`\widehat{\alpha }_\nu `$ in the Heisenberg-picture obey the Heisenberg equations of motion $`\dot{\widehat{\alpha }}_\nu =i\omega _\nu \widehat{\alpha }_\nu `$. In this approximation the collective modes and the quasi-particles have infinite lifetime. In reality, however, the lifetime will be limited by the scattering of quasi-particles in any given mode $`\nu `$ with other quasi-particles from the thermal reservoir, which is described by $`\widehat{H}_3`$ and $`\widehat{H}_4`$. One way to describe this is the quantum Langevin equation.
## III quantum Langevin-equation of a harmonic oscillator
Let us recall here briefly the quantum Langevin equation, in Markoff approximation, of a harmonic oscillator as it is commonly used in quantum optics . For a detailed discussion of its derivation I refer to . Eq.(3.4.63) of that reference states the quantum Langevin equation in resonance or ’rotating wave’ approximation for a harmonic oscillator, described by the Bose operators $`\widehat{a},\widehat{a}^+`$, in interaction with a thermal reservoir at temperature $`T`$. It takes the form
$$\dot{\widehat{a}}(t)=i\mathrm{\Omega }\widehat{a}(t)\gamma \widehat{a}(t)+\widehat{\xi }(t)$$
(7)
where $`\mathrm{\Omega }`$ is the frequency of the oscillator including a frequency shift due to the oscillator’s coupling to a heat reservoir, $`\gamma `$ is the damping rate, and $`\widehat{\xi }(t)`$ is a Gaussian noise-operator. In Markoff approximation it has the correlation functions
$$\widehat{\xi }^+(t)\widehat{\xi }(t^{})=\frac{2\gamma }{\mathrm{exp}(\mathrm{}\mathrm{\Omega }/k_BT)1}\delta (tt^{})$$
(8)
ensuring the correct normally ordered expectation values in equilibrium, and
$$[\widehat{\xi }(t),\widehat{\xi }^+(t^{})]=2\gamma \delta (tt^{})$$
(9)
ensuring the correct Bose commutation relations of $`\widehat{a}(t),\widehat{a}^+(t)`$ for all times. The fluctuation-dissipation relation therefore permits us to infer the properties of $`\widehat{\xi }(t)`$ if the coefficient of the dissipative term is known. Alternatively we can infer the dissipation rate $`\gamma `$ from a microscopic expression for $`\widehat{\xi }(t)`$ either by using (8) or, alternatively, (9).
## IV quantum Langevin-equation of collective modes
We shall here confine our attention to the dynamics of the low-lying collective modes in the collision-less regime. The interaction of the collective modes with the thermal quasi-particles is described by the Hamiltonian $`\widehat{H}_3+\widehat{H}_4`$ not yet contained in the Bogoliubov(-Popov) approximation. Because it contains the large factor $`\sqrt{N_0}`$ the contribution $`\widehat{H}_3`$ dominates over $`\widehat{H}_4`$ and the latter can be neglected in the following. Inserting the Bogoliubov transformation in $`\widehat{H}_3`$ and going to the interaction picture with respect to the unperturbed Bogoliubov-Popov Hamiltonian, $`\widehat{\stackrel{~}{H}}_3`$ in interaction representation takes the form
$`\widehat{\stackrel{~}{H}}_3={\displaystyle \frac{\sqrt{N_0}}{2}}{\displaystyle \underset{\kappa \nu \mu }{}}\{`$ $`(M_{\kappa ,\nu \mu }^{(1)}+(M_{\nu \mu ,\kappa }^{(2)})^{})\widehat{\alpha }_\kappa ^+\widehat{\alpha }_\nu \widehat{\alpha }_\mu \mathrm{exp}[i(\omega _\kappa \omega _\nu \omega _\mu )t]`$ (11)
$`+(hermitianconjugate)\}+(nonresonantterms).`$
where $`\overline{n}_\mu =1/(\mathrm{exp}(\mathrm{}\omega _\mu /k_BT)1)`$ is the thermal number of quasi-particles at frequency $`\omega _\mu `$.
Nonresonant terms, in which the frequencies of the quasi-particles cannot add up to zero, have not been written out explicitly, because later-on we shall restrict ourselves to the resonance or rotating wave approximation in which they don’t contribute. The relevant matrix elements $`M^{(1)},M^{(2)}`$ are
$`M_{\kappa ,\nu \mu }^{(1)}=`$ $`2U_0{\displaystyle }d^3x\psi _0v_\nu (u_\kappa ^{}u_\mu +{\displaystyle \frac{1}{2}}v_\kappa ^{}v_\mu )+(\nu \mu )`$ (12)
$`M_{\nu \mu ,\kappa }^{(2)}=`$ $`2U_0{\displaystyle }d^3x\psi _0u_\nu ^{}(v_\mu ^{}v_\kappa +{\displaystyle \frac{1}{2}}u_\mu ^{}u_\kappa )+(\nu \mu ).`$ (13)
$`M_{\kappa ,\nu \mu }^{(1)}`$ describes a scattering process in which one atom is scattered out of the condensate by the absorption of the two quasi-particles $`\nu ,\mu `$ out of and the emission of the new quasi-particle $`\kappa `$ into the thermal bath. Likewise $`M_{\nu \mu ,\kappa }^{(2)}`$ describes a scattering process where an incoming thermalquasiparticle $`\kappa `$ is absorbed, again an atom is kicked out of the condensate, and two quasi-particles $`\nu ,\mu `$ are emitted into the thermal bath. The scattering amplitudes for both processes are linearly superposed due to the phase-coherence of the condensate on the time-scale of the relaxation process induced by the scattering process.
Taking $`\widehat{\stackrel{~}{H}}_3`$ into account the equations of motion of $`\widehat{\stackrel{~}{\alpha }}_\nu (t)`$ in the interaction picture $`\widehat{\alpha }_\nu (t)=\widehat{\stackrel{~}{\alpha }}_\nu (t)\mathrm{exp}(i\omega _\nu t)`$ become
$`\dot{\widehat{\stackrel{~}{\alpha }}}_\nu ={\displaystyle \frac{i}{\mathrm{}}}{\displaystyle \frac{\sqrt{N_0}}{2}}{\displaystyle \underset{\kappa \mu }{}}\{`$ $`\left[M_{\nu ,\kappa \mu }^{(1)}+(M_{\kappa \mu ,\nu }^{(2)})^{}\right]\widehat{\stackrel{~}{\alpha }}_\kappa \widehat{\stackrel{~}{\alpha }}_\mu \mathrm{exp}[i(\omega _\nu \omega _\mu \omega _\kappa )t]`$ (14)
$`+2`$ $`[(M_{\kappa ,\nu \mu }^{(1)})^{}+M_{\nu \mu ,\kappa }^{(2)}]\widehat{\stackrel{~}{\alpha }}_\mu ^+\widehat{\stackrel{~}{\alpha }}_\kappa \mathrm{exp}[i(\omega _\nu +\omega _\mu \omega _\kappa )t]\}.`$ (15)
If the back-action of the collective mode on the quasi-particle operators $`\widehat{\alpha }_\mu ,\widehat{\alpha }_\kappa `$ in (15) can be ignored, the new term in this equation of motion acts like an effective random force operator. In the resonance approximation the average of this force operator vanishes. In addition it is white noise, in good approximation, if the frequencies $`\omega _\kappa \omega _\mu \omega _\nu `$ and $`\omega _\kappa +\omega _\mu \omega _\nu `$ it contains form a closely spaced quasi-continuum near $`0`$ in a neighborhood which is broad compared to the resulting damping rate $`\gamma _\nu `$. For an explicit display of this quasi-continuum in a concrete example see . In as much as this condition is satisfied for large condensates the Markoff-assumption made earlier is justified. All terms in the fluctuating force term not containing frequencies near frequency $`0`$ are non-resonant and can be omitted in comparison with resonant terms.
As we recalled in the previous section, the noise term is always accompanied by a dissipative term, and, due to the Kramers-Kronig relation, also by a frequency shift. Thus the complete quantum Langevin equation in resonance approximation and in Markoff approximation must take the form of (7)
$$\dot{\widehat{\alpha }}_\nu =i(\omega _\nu +\delta _\nu )\widehat{\alpha }_\nu \gamma _\nu \widehat{\alpha }_\nu +\widehat{\xi }_\nu (t)$$
(16)
where $`\widehat{\xi }_\nu (t)`$ is given by the second term in (15).
The damping rates $`\gamma _\nu `$ will be derived below, but we can also simply use (9) and represent them in the form
$`\gamma _\nu ={\displaystyle \frac{1}{2}}{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑t[\widehat{\xi }_\nu (t),\widehat{\xi }_\nu ^+(0)].`$ (17)
Evaluating the commutator, taking the thermal expectation value, and performing the time integral in (17) we obtain
$`\gamma _\nu ={\displaystyle \frac{\pi N_0}{\mathrm{}^2}}{\displaystyle \underset{\kappa ,\mu }{}}\{`$ $`|(M_{\kappa ,\nu \mu }^{(1)})^{}+M_{\nu \mu ,\kappa }^{(2)}|^2(\overline{n}_\mu \overline{n}_\kappa )\delta (\omega _\kappa \omega _\mu \omega _\nu )`$ (19)
$`+|M_{\nu ,\kappa \mu }^{(1)}+(M_{\kappa \mu ,\nu }^{(2)})^{}|^2(\overline{n}_\kappa +{\displaystyle \frac{1}{2}})\delta (\omega _\kappa +\omega _\mu \omega _\nu )\}`$
The first term describes Landau-damping of the mode $`\nu `$ by scattering a quasi-particle from mode $`\mu `$ to mode $`\kappa `$ and is equivalent to a result derived in by the golden rule. The second term in eq.(19) describes Beliaev damping, where the mode $`\nu `$ decays into two modes $`\kappa ,\mu `$. It survives even for $`T0`$ where $`\overline{n}_\kappa 0`$ for all modes. However, for low-lying modes in traps there are only very few modes, or no modes at all, into which decay under energy conservation can occur, and this contribution to the damping is then negligible.
Let us now derive the dissipative term of the quantum Langevin equation. To this end we consider the equations of motion for $`\frac{d}{dt}(\widehat{\stackrel{~}{\alpha }}_\mu ^+\widehat{\stackrel{~}{\alpha }}_\kappa )`$ and $`\frac{d}{dt}(\widehat{\stackrel{~}{\alpha }}_\mu \widehat{\stackrel{~}{\alpha }}_\kappa )`$, keeping again only the resonant terms. Integrating these equations over time from $`\mathrm{}`$ to $`t`$ and inserting the result back into the equation of motion for $`\widehat{\stackrel{~}{\alpha }}_\nu `$ we obtain
$`{\displaystyle \frac{d}{dt}}\widehat{\stackrel{~}{\alpha }}_\nu =\widehat{\stackrel{~}{\alpha }}_\nu {\displaystyle \frac{\sqrt{N_0}}{\mathrm{}^2}}{\displaystyle \underset{\mu \kappa }{}}\left({\displaystyle \frac{(\overline{n}_\mu +1/2)|M_{\nu ,\kappa \mu }^{(1)}+(M_{\kappa \mu ,\nu }^{(2)})^{}|^2}{ϵ+i(\omega _\mu +\omega _\kappa \omega _\nu )}}+{\displaystyle \frac{(\overline{n}_\mu \overline{n}_\kappa )|(M_{\kappa ,\nu \mu }^{(1)})^{}+M_{\nu \mu ,\kappa }^{(2)}|^2}{ϵ+i(\omega _\kappa \omega _\mu \omega _\nu )}}\right)`$ (20)
$`{\displaystyle \frac{i}{\mathrm{}}}{\displaystyle \frac{\sqrt{N_0}}{2}}{\displaystyle \underset{\kappa \mu }{}}\{[M_{\nu ,\kappa \mu }^{(1)}+(M_{\kappa \mu ,\nu }^{(2)})^{}]\widehat{\stackrel{~}{\alpha }}_\kappa (\mathrm{})\widehat{\stackrel{~}{\alpha }}_\mu (\mathrm{})\mathrm{exp}[i(\omega _\nu \omega _\mu \omega _\kappa )t]`$ (21)
$`+2[(M_{\kappa ,\nu \mu }^{(1)})^{}+M_{\nu \mu ,\kappa }^{(2)}]\widehat{\stackrel{~}{\alpha }}_\mu ^+(\mathrm{})\widehat{\stackrel{~}{\alpha }}_\kappa (\mathrm{})\mathrm{exp}[i(\omega _\nu +\omega _\mu \omega _\kappa )t]\}`$ (22)
where the limit $`ϵ+0`$ is implied. The second term on the right hand side is the fluctuating force term again, now more rigorously expressed in terms of the reservoir operators at the initial time at $`\mathrm{}`$. Taking the limit with $`(ϵi\omega )^1\pi \delta (\omega )+iP/\omega `$, where $`P/\omega `$ denotes the principal part under a frequency integral, we obtain the result (19) for the damping rate and also the frequency shifts $`\delta _\nu `$ in the quantum Langevin equation. They are given by the Kramers-Kronig relation
$$\delta _\nu =\frac{1}{\pi }P𝑑\omega \frac{\gamma (\omega )}{\omega \omega _\nu }$$
(23)
where we defined $`\gamma (\omega _\nu )=\gamma _\nu `$.
## V Damping rates
In the following we shall neglect the second term in (19), because as discussed it cannot contribute for low lying modes. Our goal in this section is the evaluation of the first term in (19) in a well defined approximation, the local density and the Thomas-Fermi approximation. The local density approximation amounts to the treatment of the quasi-continuum of the spectrum of frequencies $`\omega _\kappa \omega _\mu \omega _\nu `$ as a continuum whose density is given by the semiclassical mode-densities of the frequencies $`\omega _\mu ,\omega _\kappa `$. Why these frequencies lie much higher than the collective mode frequency $`\omega _\nu `$ will become clear below. The Thomas-Fermi approximation applies to large condensates and amounts to neglecting the kinetic energy term in the Gross-Pitaevskii equation. The collective modes satisfy $`E_\nu \mu =U_0|\psi _0(0)|^2`$ and can be represented as
$`u_\nu (𝒙)=`$ $`\left(\sqrt{{\displaystyle \frac{U_0n_0(𝒙)}{2\mathrm{}\omega _\nu }}}+{\displaystyle \frac{1}{2}}\sqrt{{\displaystyle \frac{\mathrm{}\omega _\nu }{2U_0n_0(𝒙)}}}\right)\chi _\nu (𝒙)`$ (24)
$`v_\nu (𝒙)=`$ $`\left(\sqrt{{\displaystyle \frac{U_0n_0(𝒙)}{2\mathrm{}\omega _\nu }}}+{\displaystyle \frac{1}{2}}\sqrt{{\displaystyle \frac{\mathrm{}\omega _\nu }{2U_0n_0(𝒙)}}}\right)\chi _\nu (𝒙)`$ (26)
with $`d^3x|\chi _\nu (𝒙)|^2=1`$. The functions of the low-lying states $`\chi _\nu (𝒙)`$ are known in the hydrodynamic (long-wavelength) and Thomas-Fermi approximation (and neglecting the influence of the thermal cloud which sits mainly outside the condensate and therefore has little influence on its collective excitations). In spatially isotropic parabolic traps they have the form
$$\chi _\nu (𝒙)=\frac{1}{r_{TF}^{3/2}}P_\mathrm{}_\nu ^{(2n_\nu )}(x/r_{TF})(x/r_{TF})_\nu ^{\mathrm{}}Y_{\mathrm{}_\nu m_\nu }(\theta ,\phi )\mathrm{\Theta }(1x/r_{TF})$$
(27)
The normalized polynomials $`P_{\mathrm{}}^{(2n)}(x)`$ are known explicitly .
The high-lying quasiparticle modes can be represented as
$$u_\kappa (𝒙)=\frac{E_\kappa +p_\kappa ^2/2m}{\sqrt{2E_\kappa p_\kappa ^2/m}}e^{i𝒑_\kappa 𝒙/\mathrm{}},v_\kappa (𝒙)=\frac{E_\kappa p_\kappa ^2/2m}{\sqrt{2E_\kappa p_\kappa ^2/m}}e^{i𝒑_\kappa 𝒙/\mathrm{}}$$
(28)
with the local energies in Thomas-Fermi approximation
$$E_\kappa =E(p_\kappa ,𝒙)=\sqrt{(\frac{p_\kappa ^2}{2m}+|U_0n_0(𝒙)|)^2U_0^2n_0^2(𝒙)\mathrm{\Theta }(\mu V(𝒙))}$$
(29)
and $`n_0(𝒙)=N_0|\psi _0(𝒙)|^2=(\mu /U_0)(1_i(x_i/r_{TF}^{(i)})^2)`$, and $`r_{TF}^{(i)}=\sqrt{2\mu /\omega _i^2}`$ are the three main Thomas-Fermi radii.
Let us consider now the Landau-damping of a low-lying phonon mode $`\omega _\nu `$. If the modes $`\mu ,\kappa `$ involved in the scattering process were also low-lying we could use eq.(LABEL:eq:L) and would obtain, with $`E_\kappa =E_\nu +E_\mu `$
$$(M_{\kappa ,\nu \mu }^{(1)})^{}+M_{\nu \mu ,\kappa }^{(2)}=\frac{3U_0}{4\sqrt{2}}d^3x\psi _0\chi _\kappa \chi _\mu ^{}\chi _\nu ^{}\sqrt{\frac{E_\mu E_\nu E_\kappa }{U_0^3n_0^3(x)}}$$
(30)
However, in the limit of low-lying modes where $`E_\nu ,E_\mu ,E_\kappa U_0n_0`$ this matrix element becomes very small, i. e. low-lying modes cannot significantly contribute to Landau damping of other low-lying modes. Therefore the relevant modes $`\mu ,\kappa `$ are in fact not low lying, local density approximation is applicable, and we can use eq.(28) for their representation. The matrix-element for $`E_\nu <<E_\mu ,E_\kappa `$ can be expanded in $`E_\nu /\mu `$ to lowest order around $`E_\kappa =E_\mu `$ and becomes then
$$(M_{\kappa ,\nu \mu }^{(1)})^{}+M_{\nu \mu ,\kappa }^{(2)}=\sqrt{\frac{E_\nu U_0}{2N_0}}d^3x\chi _\nu ^{}(𝒙)\mathrm{exp}(i(𝒑_\kappa 𝒑_\mu )\mathbf{}𝒙))F(E_\mu ,p_\mu )$$
(31)
with
$$F(E_\mu ,p_\mu )=\frac{p_\mu ^2}{2m}\frac{3E_\mu ^2+(p_\mu ^2/2m)^2}{E_\mu (E_\mu ^2+(p_\mu ^2/2m)^2}$$
(32)
It will be very convenient later to express the product $`\chi _\nu ^{}(𝒙)\chi _\nu (𝒙^{})`$ by the associated Wigner-function $`W_\nu `$ via
$$\chi _\nu ^{}(𝒙)\chi _\nu (𝒙^{})=\frac{d^3k}{(2\pi )^3}e^{i𝒌\mathbf{}(𝒙𝒙^{})}W_\nu (\frac{1}{2}(𝒙+𝒙^{}),𝒌)$$
(33)
In the following we denote
$$(𝒙+𝒙^{})/2𝒙𝒙𝒙^{}𝒓$$
(34)
The rate for Landau-damping can then be written as
$`\gamma _\nu =CE_\nu ^2`$ $`{\displaystyle d^3x\frac{d^3k}{(2\pi )^3}W_\nu (𝒙,𝒌)\underset{\mu }{}}`$ (36)
$`{\displaystyle \underset{\kappa }{}}F^2(E_\mu ,p_\mu ){\displaystyle d^3re^{\frac{i}{\mathrm{}}(𝒑_\kappa 𝒑_\mu \mathrm{}𝒌)\mathbf{}𝒓}\frac{\delta (E_\kappa E_\mu E_\nu )}{\mathrm{sinh}^2\left(\frac{E_\mu }{2k_BT}\right)}}`$
with
$$C=\frac{\pi }{8\mathrm{}}\frac{U_0}{k_BT}$$
(37)
The sums over the states $`\mu `$ and $`\kappa `$ are only symbolic, because in the local-density approximation the discrete states have been replaced by a continuum which is normalized on the $`\delta `$-function. Concretely, under the integral over $`𝒙`$ the sums over the energy levels of the scattering quasi-particles $`\mu ,\kappa `$ are in the local density approximation replaced locally by classical phase-space averages for fixed $`E_\mu ,E_\kappa `$ and final integration over $`E_\mu `$ and $`E_\kappa `$ which takes automatically care of the normalization on the $`\delta `$-function. Thus
$$_\mu \mathrm{}_\mu ^{(𝒙)}\mathrm{}=𝑑E_\mu \frac{d^3p_\mu }{(2\pi \mathrm{})^3}\delta \left(E_\mu E(p_\mu ,𝒙)\right)\mathrm{}$$
(38)
In the following $`_\mu `$ and $`_\kappa `$ will be interpreted according to eq.(38). The spatial integration over $`𝒓`$ can be done and produces a momentum-conservation factor $`\left(2\pi \mathrm{}^3\right)\delta ^{(3)}\left(𝒑_𝜿𝒑_𝝁\mathrm{}𝒌\right)`$. Next the integrations over $`𝒑_𝜿`$ and $`E_\kappa `$ contained in $`_\kappa `$ can be performed, which just cancel the $`\delta `$-functions of overall momentum and energy conservation and replace everywhere else $`E_\kappa E_\mu +E_\nu `$ and $`𝒑_𝜿𝒑_𝝁+\mathrm{}𝒌`$. Then the expression for $`\gamma _\nu `$ is reduced to
$`\gamma _\nu =CE_\nu ^2`$ $`{\displaystyle d^3x\frac{d^3k}{(2\pi )^3}W_\nu (𝒙,𝒌)𝑑E_\mu }`$ (40)
$`{\displaystyle \frac{d^3p_\mu }{(2\pi \mathrm{})^3}\delta \left(E_\mu E(p_\mu ,𝒙)\right)\frac{F^2(E_\mu ,p_\mu )}{\mathrm{sinh}^2\left(\frac{E_\mu }{2k_BT}\right)}\delta \left(E_\mu +E_\nu E(|𝒑_\mu +\mathrm{}𝒌|,𝒙)\right)}`$
Next the integration over the directions of $`𝒑_\mu `$ relative to $`𝒌`$ is carried out by using up the second of the two $`\delta `$-functions explicitly displayed in eq. (40). This produces a factor $`2\pi `$ for the azimuthal angle, and a factor $`|E(p_\mu ,𝒙)/p_\mu ^2|^1(2p_\mu \mathrm{}k)^1`$ from the integration over $`\mathrm{cos}\theta `$ between -1 and 1, where $`\theta `$ is the angle between $`𝒌`$ and $`𝒑_\mu `$. Finally the integration over the absolute value $`p_\mu `$ is done using up the last $`\delta `$-function, which picks out the $`𝒙`$-dependent momentum-value $`p_\mu ^{(0)}=\sqrt{2m}\sqrt{\sqrt{E_\mu ^2+U_0^2n_0^2(𝒙)}U_0n_0(𝒙)}.`$ leaving us with the expression
$$\gamma _\nu =\frac{CE_\nu ^2}{4\pi ^2\mathrm{}^3}d^3x\frac{d^3k}{(2\pi )^3}\frac{W_\nu (𝒙,𝒌)}{\mathrm{}k}𝑑E_\mu \frac{F^2(E_\mu ,p_\mu ^{(0)})}{4\left(\frac{E(p_\mu ^{(0)},𝒙)}{(p_\mu ^{(0)})^2}\right)^2\mathrm{sinh}^2\left(\frac{E_\mu }{2k_BT}\right)}$$
(41)
We now have to face the difficulty to evaluate the conditional average of $`(\mathrm{}k)^1`$. The rigorous way to do this, which unfortunately leads to multiple integrals which are tedious to evaluate, is to invert (33) which yields
$$\frac{d^3k}{(2\pi )^3}\frac{W_\nu (𝒙,𝒌)}{\mathrm{}k}=d^3r\frac{\chi _\nu ^{}(𝒙\mathbf{+}𝒓/2)\chi _\nu (𝒙\mathbf{}𝒓/2)}{2\pi ^2\mathrm{}r^2}$$
(42)
A much simpler way consists in expressing the desired average by the local sound-velocity $`\sqrt{\mu /m}\overline{c}_\nu (𝒙)`$ defined by
$$\frac{d^3k}{(2\pi )^3}\frac{W_\nu (𝒙,𝒌)}{\mathrm{}k}=\sqrt{\frac{\mu }{m}}\frac{\overline{c}_\nu (𝒙)}{E_\nu }|\chi _\nu (𝒙)|^2.$$
(43)
and estimating the dimensionless sound-velocity $`\overline{c}_\nu (𝒙)`$ semi-classically as $`\overline{c}_\nu (𝒙)\sqrt{1(𝒙/r_{TF})^2}`$ with the geometrical mean Thomas-Fermi radius $`r_{TF}=\left(2\mu /m\overline{\omega }^2\right)^{\frac{1}{2}}`$. Of course the use of the semi-classical approximation for the low lying collective mode is highly questionable and cannot be quantitatively accurate. Still we may like to use it as a rough estimate in a case where an accurate evaluation is not required or too time consuming. Below we shall check this approximation in two cases, where it cannot be expected to be particularly good.
We now introduce scaled variables $`\stackrel{~}{E}_\mu =E_\mu /(N_0U_0\psi _0^2(𝒙))`$ and $`\stackrel{~}{𝒙}=𝒙/r_{TF}`$, $`\stackrel{~}{𝒙}^{}=𝒙^{}/r_{TF}`$ with dimensionless mode-functions $`\stackrel{~}{\chi }_\nu (\stackrel{~}{𝒙})=r_{TF}^{\frac{3}{2}}\chi _\nu (𝒙)`$. Altogether, using (42), we are left with the result
$$\gamma _\nu =\frac{(a^3n_0(0))^{1/2}E_\nu ^2}{2(2\pi )^{3/2}\mathrm{}^2\overline{\omega }}H_\nu (\frac{k_BT}{\mathrm{}})$$
(44)
with the dimensionless function
$`H_\nu (z)={\displaystyle d^3\stackrel{~}{x}d^3\stackrel{~}{x}^{}}`$ $`{\displaystyle \frac{\stackrel{~}{\chi }_\nu ^{}(\stackrel{~}{𝒙})\stackrel{~}{\chi }_\nu (\stackrel{~}{𝒙}^{})(1(\stackrel{~}{𝒙}+\stackrel{~}{𝒙}^{})^2/4)}{(\stackrel{~}{𝒙}\stackrel{~}{𝒙}^{})^2}}`$ (46)
$`{\displaystyle \frac{1}{z}}{\displaystyle 𝑑\stackrel{~}{E}_\mu \left(\frac{2\stackrel{~}{E}_\mu +1\sqrt{\stackrel{~}{E}_\mu ^2+1}}{(\stackrel{~}{E}_\mu ^2+1)\mathrm{sinh}\left(\frac{1}{2z}\stackrel{~}{E}_\mu (1(\stackrel{~}{𝒙}+\stackrel{~}{𝒙}^{})^2/4)\right)}\right)^2}`$
For $`z>>1`$ the functions $`H_\nu `$ become linear in $`z=k_BT/\mu `$ and reduce to
$$H_\nu (z)3\pi zd^3\stackrel{~}{x}d^3\stackrel{~}{x}^{}\frac{\stackrel{~}{\chi }_\nu ^{}(\stackrel{~}{𝒙})\stackrel{~}{\chi }_\nu (\stackrel{~}{𝒙}^{})}{(\stackrel{~}{𝒙}\stackrel{~}{𝒙}^{})^2(1(\stackrel{~}{𝒙}+\stackrel{~}{𝒙}^{})^2/4)}$$
(47)
The result for the spatially homogeneous case can be recovered from eq.(47) for $`k_BT>>\mu `$ by using the scaled homogeneous condensate density $`1\stackrel{~}{x}^21`$, the phonon energy $`E_\nu =\sqrt{\mu /m}\mathrm{}k_\nu `$, and normalized plane waves to evaluate $`d^3\stackrel{~}{x}d^3\stackrel{~}{x}^{}\stackrel{~}{\chi }_\nu ^{}(\stackrel{~}{𝒙})\stackrel{~}{\chi }_\nu (\stackrel{~}{𝒙}^{})/(\stackrel{~}{𝒙}\stackrel{~}{𝒙}^{})^2=(2\pi ^2/k_\nu )\sqrt{m\overline{\omega }^2/2\mu }`$ which, together with (44,47) yields $`\gamma _\nu =\frac{3\pi }{8}ak_\nu \frac{k_BT}{\mathrm{}}`$.
For isotropic traps the asymptotic result (47) becomes
$`H_{n_\nu \mathrm{}_\nu m_\nu }(z)`$ $`6\pi z{\displaystyle \frac{(2\mathrm{}_\nu +1)(\mathrm{}_\nu m_\nu )!}{(\mathrm{}_\nu +m_\nu )!}}{\displaystyle \underset{0}{\overset{1}{}}}𝑑\stackrel{~}{x}\stackrel{~}{x}^2P_\mathrm{}_\nu ^{(2n_\nu )}(\stackrel{~}{x})\stackrel{~}{x}^\mathrm{}_\nu {\displaystyle \underset{0}{\overset{1}{}}}𝑑\stackrel{~}{x}^{}\stackrel{~}{x}^2P_\mathrm{}_\nu ^{(2n_\nu )}(\stackrel{~}{x^{}})\stackrel{~}{x^{}}^\mathrm{}_\nu `$ (49)
$`{\displaystyle \underset{1}{\overset{1}{}}}{\displaystyle \underset{1}{\overset{1}{}}}{\displaystyle \underset{0}{\overset{2\pi }{}}}{\displaystyle \frac{d(\mathrm{cos}\theta )d(\mathrm{cos}\theta ^{})d\varphi P_\mathrm{}_\nu ^{m_\nu }(\mathrm{cos}\theta )P_\mathrm{}_\nu ^{m_\nu }(\mathrm{cos}\theta ^{})\mathrm{exp}(im_\nu \varphi )}{(2\stackrel{~}{x}\stackrel{~}{x}^{}(\mathrm{cos}\theta \mathrm{cos}\theta ^{}+\mathrm{sin}\theta \mathrm{sin}\theta ^{}\mathrm{cos}\varphi )2)^2(\stackrel{~}{x}^2+\stackrel{~}{x}^22)^2}}`$
where the functions $`P_{\mathrm{}}^m(cos\theta )`$ are the associated Legendre functions appearing in the spherical harmonics.
If instead of (42) we use (43) to evaluate the conditional average of $`(\mathrm{}k)^1`$ we obtain in place of (49)
$`H_{n_\nu \mathrm{}_\nu m_\nu }(z)`$ $`z{\displaystyle \frac{3\sqrt{2}\pi ^3\mathrm{}\overline{\omega }}{E_\nu }}{\displaystyle \underset{0}{\overset{1}{}}}𝑑\stackrel{~}{x}{\displaystyle \frac{\stackrel{~}{x}^2}{\sqrt{(1\stackrel{~}{x}^2)}}}(P_\mathrm{}_\nu ^{(2n_\nu )}(\stackrel{~}{x})\stackrel{~}{x}^\mathrm{}_\nu )^2`$ (50)
As was already emphasized, this result can only serve as a rough estimate for (49).
In the simplest case $`n_\nu =1,\mathrm{}_\nu =0`$, which is the isotropic fundamental breathing mode, we have $`P_0^0(\mathrm{cos}(\theta ))=1,P_1^{(0)}(x)=\frac{3}{2}\sqrt{7}(1\frac{5}{3}x^2)`$. In this case (49) can be reduced to the numerical evaluation of a two-dimensional integral and we obtain
$$\gamma _{0,0}26.42..\omega _0(a^3n_0(0))^{1/2}\frac{k_BT}{\mu }.$$
(51)
Eq.(50) yields via elementary integration $`\gamma _{0,0}27.27..\omega _0(a^3n_0(0))^{1/2}k_BT/\mu `$ which agrees surprisingly well with the more rigorous result (51). Can this be considered typical? The answer is negative:
The simple result (50) lends itself to further evaluation for modes with $`\mathrm{}_\nu 0`$. For the surface modes with $`n_\nu =0,\mathrm{}_\nu 0`$ we obtain the estimate
$$\gamma _{0,\mathrm{}_\nu }\omega _0(a^3n_0(0))^{1/2}\frac{k_BT}{\mu }\frac{3\pi ^2}{4}\sqrt{\mathrm{}_\nu }\frac{\mathrm{\Gamma }(\mathrm{}_\nu +5/2)}{\mathrm{\Gamma }(\mathrm{}_\nu +2)}.$$
(52)
In this case a numerical comparison with the more accurate expression (49) for the case $`\mathrm{}_\nu =2`$ shows that the latter is about 30 percent smaller, probably giving us a realistic impression of the accuracy of the approximation for $`\overline{c}_\nu (𝒙)`$. For larger values of $`\mathrm{}_\nu `$ and $`n_\nu `$ the accuracy of this estimate can be expected to improve.
## VI Discussion and conclusion
In the present paper the many-body problem of collective modes in Bose-Einstein condensates in interaction with thermal quasi-particles was addressed by a method based on the equations of motion of the quasi-particle operators. This method leads directly to a quantum Langevin equation for the creation and annihilation operators of the collective modes, containing fluctuating force terms, a dissipation term, and a frequency shift term. These quantities are related by the fluctuation-dissipation relation and the Kramers-Kronig relation. Each part of the interaction-Hamiltonian beyond the unperturbed Bogoliubov-Popov Hamiltonian in principle gives rise to separate contributions to all three types of terms. We have here considered only the most important of these, namely the part of the interaction Hamiltonian giving rise to Landau-damping.
Dissipation can arise only from energy conserving real processes, which is manifest by the appearance of the energy conserving $`\delta `$-functions in the expressions for the damping rates. This means that only resonant processes can contribute to these rates. In finite systems like the trapped condensates this causes a problem, because there the mode spectrum is discrete, the spectrum of frequency differences $`\omega _\kappa \omega _\mu \omega _\nu `$ near $`0`$ is only a quasi-continuum, and the dissipation rates in a strict sense have to vanish. In other words, in a strict sense, what is seen as dissipation can only be a ’short-time’ effect; waiting for a sufficiently long time interval on the order of the inverse spacing of the quasi-continuum, revivals would have to appear. These will not be seen, however, at least in large condensates to which the local density and Thomas-Fermi approximation can be applied, because not only the energy stored in the collective mode but also the thermal energy of the system is available to bring into play a large number of modes which will lead to an irretrievable dissipation of the energy over many degrees of freedom. Therefore it is reasonable in such cases, if not required, to eliminate all recurrence effects, replacing the quasi-continuum by a true continuum, which is what the local density approximation does. Using this device we have arrived at definite results for the temperature-dependent damping rates of any collective mode, in an isotropic trap, which can be evaluated by computing numerically a multi-dimensional definite integral, e.g. by a Monte-Carlo routine.
The different pieces of the perturbation Hamiltonian also each give rise to frequency shifts. These are generated by virtual processes which do not require energy conservation, i.e. resonance. However the effect of the non-resonant processes is suppressed by corresponding energy-denominators and small. Here we have limited our considerations only to those processes which can also become resonant. We have here not evaluated the frequency shifts further using the local density approximation as we have done for the damping rates.
Experimental results for temperature dependent damping rates and frequency shifts have been obtained for anisotropic traps only , and we therefore refrain from a comparison with our explicit results for isotropic traps. Detailed comparisons have been made for anisotropic traps in , where heavier and more powerful formalisms were brought to bear together with a stronger reliance on numerical work.
The goal here has been more modest, namely to use a minimum amount of numerical work and to apply the direct and intuitive quantum Langevin approach to the fluctuations, damping rates and frequency shifts of collective modes in spatially inhomogeneous trapped Bose-Einstein condensates.
## Acknowledgment:
This work has been supported by the Deutsche Forschungsgemeinschaft through the Sonderforschungsbereich 237 “Unordnung und große Fluktuationen”. I wish to thank Martin Fliesser for useful discussions and owe thanks to Jürgen Reidl for the numerical evaluation of some integrals.
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# (1) Inequalities among numbers
## (1) Inequalities among numbers
Lemma (1) For any two numbers $`a,c[1,1]`$ the following equivalent inequalities hold:
$$|a\pm c|1\pm ac$$
$`(1)`$
Moreover equality in (1) holds if and only if either $`a=\pm 1`$ or $`c=\pm 1`$.
Proof. The equivalence of the two inequalities (1) follows from the fact that one is obtained from the other by changing the sign of $`c`$ and $`c`$ is arbitrary in $`[1,1]`$.
Since for any $`a,c[1,1]`$, $`1\pm ac0`$, (1) is equivalent to
$$|a\pm c|^2=a^2+c^2\pm 2ac(1\pm ac)^2=1+a^2c^2\pm 2ac$$
and this is equivalent to
$$a^2(1c^2)+c^21$$
which is identically satisfied because $`1c^20`$ and therefore
$$a^2(1c^2)+c^21c^2+c^2=1$$
$`(2)`$
Notice that in (2) equality holds if and only if $`a^2=1`$ i.e. $`a=\pm 1`$. Since, exchanging $`a`$ and $`c`$ in (1) the inequality remains unchanged, the thesis follows.
Corollary (2) For any three numbers $`a,b,c[1,1]`$ the following equivalent inequalities hold:
$$|ab\pm cb|1\pm ac$$
$`(3)`$
and equality holds if and only if $`b=\pm 1`$ and either $`a=\pm 1`$ or $`c=\pm 1`$.
Proof. For $`b[1,1]`$,
$$|ab\pm cb|=|b||a\pm c||a\pm c|$$
$`(4)`$
so the thesis follows from Lemma (1). In (4) the first equality holds if and only if $`b=\pm 1`$, so also the second statement follows from Lemma (1).
Lemma (3). For any numbers $`a`$, $`a^{}`$, $`b`$, $`b^{}`$, $`c[1,1]`$, one has
$$|abbc|+|ab^{}+b^{}c|2$$
$`(5)`$
$$ab+ab^{}+a^{}b^{}a^{}b2$$
$`(6)`$
In (5) equality holds if and only if $`b,b^{},a,c=\pm 1`$.
Proof. Because of (3)
$$|abcb|1ac$$
$`(7)`$
$$|ab^{}cb^{}|1+ac$$
$`(8)`$
adding (7) and (8) one finds (5). The left hand side of (6) is $``$ than
$$|abba^{}|+|ab^{}+b^{}a^{}|$$
$`(9)`$
and replacing $`a^{}`$ $`by`$ $`c`$, (8) becomes the left hand side of (5). If $`b,b^{}=\pm 1`$ and $`a=\pm 1`$ equality holds in (7) and (8) hence in (5). Conversely, suppose that equality holds in (5) and suppose that either $`|b|<1`$ or $`|b^{}|<1`$. Then we arrive to the contradiction
$$2=|b||aa^{}|+|b^{}||a+a^{}|<|aa^{}|+|a+a^{}|(1aa^{})+(1+aa^{})=2$$
$`(10)`$
So, if equality holds in (5), we must have $`|b|=|b^{}|=1`$. In this case (5) becomes
$$|aa^{}|+|a+a^{}|=2$$
$`(11)`$
and, if either $`|a|<1`$ or $`|a^{}|<1`$, then from Lemma (1) we know that
$$|aa^{}|+|a+a^{}|<(1aa^{})+(1+aa^{})=2$$
so we must also have $`a,a^{}=\pm 1`$.
Corollary (4). If $`a,a^{},b,b^{},c\{1,1\}`$, then the inequalities (3) and (5) are equivalent and equality holds in all of them. However the inequality in (6) may be strict.
Proof. From Lemma (1) we know that the inequalities (1) and (2) are equivalent. From Lemma (3) we know that (1) implies (5). Choosing $`b^{}=a`$ in (5), since $`a=\pm 1`$, Lemma (2) implies that (5) becomes
$$|abcb|1ac$$
which is equivalent to (1). (6) is equivalent to
$$a(b+b^{})+a^{}(b^{}b)2$$
$`(12)`$
In our assumptions either $`(b+b^{})`$ or $`(b^{}b)`$ is zero, so (12) is either equivalent to
$$a(b+b^{})2$$
or to
$$a^{}(b^{}b)2$$
and in both cases we can choose $`a,b,b^{}`$ or $`a^{},b,b^{}`$ so that the product is negative and the inequality is strict.
## (2) The Bell inequality
Corollary (1) (Bell inequality) Let $`A,B,C,D`$ be random variables defined on the same probability space $`(\mathrm{\Omega },,P)`$ and with values in the interval $`[1,1]`$. Then the following inequalities hold:
$$E(|ABBC|)1E(AC)$$
$`(1)`$
$$E(|AB+BC|)1+E(AC)$$
$`(2)`$
$$E(|ABBC|)+E(|AD+DC|)2$$
$`(3)`$
where $`E`$ denotes the expectation value in the probability space of the four variables. Moreover (1) is equivalent to (2) and, if either $`A`$ or $`C`$ has values $`\pm 1`$, then the three inequalities are equivalent.
Proof. Lemma (1.1) implies the following inequalities (interpreted pointwise on $`\mathrm{\Omega }`$):
$$|ABBC|1AC$$
$$|AB+BC|1+AC$$
$$|ABBC|+|AD+DC|2$$
from which (1), (2), (3) follow by taking expectation and using the fact that $`|E(X)|E(|X|)`$. The equivalence is established by the same arguments as in Lemma (1.1).
Remark (2). Bell’s original proof, as well as the almost totality of the available proofs of Bell’s inequality, deal only with the case of random variables assuming only the values $`+1`$ and $`1`$. The present generalization is not without interest because it dispenses from the assumption that the classical random variables, used to describe quantum observables, have the same set of values of the latter ones: a hidden variable theory is required to reproduce the results of quantum theory only when the hidden parameters are averaged over.
Theorem (3). Let $`S_a^{(1)},S_c^{(1)},S_b^{(2)},S_d^{(2)}`$ be random variables defined on a probability space $`(\mathrm{\Omega },,P)`$ and with values in the interval $`[1,+1]`$. Then the following inequalities holds:
$$\left|E(S_a^{(1)}S_b^{(2)})E(S_c^{(1)}S_b^{(2)})\right|1E(S_a^{(1)}S_c^{(1)})$$
$`(4)`$
$$\left|E(S_a^{(1)}S_b^{(2)})+E(S_c^{(1)}S_b^{(2)})\right|1+E(S_a^{(1)}S_c^{(1)})$$
$`(5)`$
$$\left|E(S_a^{(1)}S_b^{(2)})E(S_c^{(1)}S_b^{(2)})\right|+\left|E(S_a^{(1)}S_d^{(2)})+E(S_c^{(1)}S_d^{(2)})\right|2$$
$`(6)`$
Proof. This is a rephrasing of Corollary (2).
## (3) Implications of the Bell’s inequalities for the singlet correlations
To apply Bell’s inequalities to the singlet correlations, considered in the EPR paradox, it is enough to observe that they imply the following
Lemma (1) In the ordinary three-dimensional euclidean space there exist sets of three, unit length, vectors $`a`$, $`b`$, $`c`$, such that it is not possible to find a probability space $`(\mathrm{\Omega },,P)`$ and six random variables $`S_x^{(j)}`$ ($`x=a,b,c`$, $`j=1,2`$) defined on $`(\mathrm{\Omega },,P)`$ and with values in the interval $`[1,+1]`$, whose correlations are given by:
$$E(S_x^{(1)}S_y^{(2)})=xy;x,y=a,b,c$$
$`(1)`$
where, if $`x=(x_1,x_2,x_3)`$, $`y=(y_1,y_2,y_3)`$ are two three-dimensional vectors, $`xy`$ denotes their euclidean scalar product, i.e. the sum $`x_1y_1+x_2y_2+x_3y_3`$.
Remark. In the usual EPR–type experiments, the random variables $`S_a^{(j)},S_b^{(j)},S_c^{(j)}`$ represent the spin (or polarization) of particle $`j`$ of a singlet pair along the three directions $`a,b,c`$ in space. The expression in the right-hand side of (1) is the singlet correlation of two spin or polarization observables, theoretically predicted by quantum theory and experimentally confirmed by the Aspect-type experiments.
Proof. Suppose that, for any choice of the unit vectors $`x=a,b,c`$ there exist random variables $`S_x^{(j)}`$ as in the statement of the Lemma. Then, using Bell’s inequality in the form (2.5) with $`A=S_a^{(1)}`$, $`B=S_b^{(2)}`$, $`C=S_c^{(1)}`$), we obtain
$$\left|E(S_a^{(1)}S_b^{(2)})+E(S_b^{(2)}S_c^{(1)})\right|1+E(S_a^{(1)}S_c^{(1)})$$
$`(2)`$
Now notice that, if $`x=y`$ is chosen in (1), we obtain
$$E(S_x^{(1)}S_x^{(2)})=xx=x^2=1;x=a,b,c$$
and, since $`\left|S_x^{(1)}S_x^{(2)}\right|=1`$ this is possible if and only if $`S_x^{(1)}=S_x^{(2)}`$ $`\left(x=a,b,c\right)`$ $`P`$–almost everywhere. Using this (2) becomes equivalent to:
$$\left|E(S_a^{(1)}S_b^{(2)})+E(S_b^{(2)}S_c^{(1)})\right|1E(S_a^{(1)}S_c^{(2)})$$
or, again using (1), to:
$$\left|ab+bc\right|1+ac$$
$`(3)`$
If the three vectors $`a`$, $`b`$, $`c`$ are chosen to be in the same plane and such that $`a`$ is perpendicular to $`c`$ and $`b`$ lies between $`a`$ and $`b`$, forming an angle $`\theta `$ with $`a`$, then the inequality (3) becomes:
$$cos\theta +\mathrm{sin}\theta 1;0<\theta <\pi /2$$
$`(4)`$
But the maximum of the function of $`\theta \mathrm{sin}\theta +\mathrm{cos}\theta `$ in the interval $`[0,\pi /2]`$ is $`\sqrt{2}`$ (obtained for $`\theta =\pi /4`$). Therefore, for $`\theta `$ close to $`\pi /4`$, the left-hand side of (4) will be close to $`\sqrt{2}`$ which is more that $`1`$. In conclusion, for such a choice of the unit vectors $`a`$, $`b`$, $`c`$, random variables $`S_a^{(1)},S_b^{(2)},S_c^{(1)},S_c^{(2)}`$ as in the statement of the Lemma cannot exist.
Definition (2) A local realistic model for the EPR (singlet) correlations is defined by:
(1) a probability space $`(\mathrm{\Omega },,P)`$
(2) for every unit vector $`x`$, in the three-dimensional euclidean space, two random variables $`S_x^{(1)},S_x^{(2)}`$ defined on $`\mathrm{\Omega }`$ and with values in the interval $`[1,+1]`$ whose correlations, for any $`x,y`$, are given by equation (1).
Corollary (3) If $`a,b,c`$ are chosen so to violate (4) then a local realistic model for the EPR correlations, in the sense of Definition (2), does not exist.
Proof. Its existence would contradict Lemma (1).
Remark. In the literature one usually distinguishes two types of local realistic models – deterministic and stochastic ones. Both are included in Definition (2): the deterministic models are defined by random variables $`S_x^{(j)}`$ with values in the set$`\{1,+1\}`$; while, in the stochastic models, the random variables take values in the interval $`[1,+1]`$. The original paper \[Be64\] was devoted to the deterministic case. Starting from \[Be71\] several papers have been introduced to justify the stochastic models. We prefer to distinguish the definition of the models from their justification.
## (4) Bell on the meaning of Bell’s inequality
In the last section of \[Be66\] (submitted before \[Be64\], but published after) Bell briefly describes Bohm hidden variable interpretation of quantum theory underlining its non local character. He then raises the question: … that there is no proof that any hidden variable account of quantum mechanics must have this extraordinary character … and, in a footnote added during the proof corrections, he claims that: … Since the completion of this paper such a proof has been found \[Be64\].
In the short Introduction to \[Be64\], Bell reaffirms the same ideas, namely that the result proven by him in this paper shows that: … any such \[hidden variable\] theory which reproduces exactly the quantum mechanical predictions must have … a grossly nonlocal structure.
The proof goes along the following scheme: Bell proves an inequality in which, according to what he says (cf. statement after formula (1) in \[Be64\]):
… The vital assumption is that the result $`B`$ for particle $`2`$ does not depend on the setting $`a`$, of the magnet for particle $`1`$, nor $`A`$ on $`b`$.
The paper , mentioned in the above statement, is nothing but the Einstein, Podolsky, Rosen paper \[EPR35\] and the locality issue is further emphasized by the fact that he reports the famous Einstein’s statement \[Ein49\]: … But on one supposition we should, in my opinion, absolutely hold fast: the real factual situation of the system $`S_2`$ is independent of what is done with the system $`S_1`$, which is spatially separated from the former.
Stated otherwise: according to Bell, Bell’s inequality is a consequence of the locality assumption.
It follows that a theory which violates the above mentioned inequality also violates … the vital assumption needed, according to Bell, for its deduction, i.e. locality.
Since the experiments prove the violation of this inequality, Bell concludes that quantum theory does not admit a local completion; in particular quantum mechanics is a nonlocal theory. To use again Bell’s words:
the statistical predictions of quantum mechanics are incompatible with separable predetermination (\[Be64\], p.199). Moreover this incompatibility has to be understood in the sense that: in a theory in which parameters are added to quantum mechanics to determine the results of individual measurements, without changing the statistical predictions, there must be a mechanism whereby the setting of one measuring device can influence the reading of another instrument, howevere remote. Moreover, the signal involved must propagate instantaneously,…
## (5) Critique of Bell’s “vital assumption”
An assumption should be considered “vital” for a theorem if, without it, the theorem cannot be proved.
To favor Bell, let us require much less. Namely let us agree to consider his assumption vital if the theorem cannot be proved by taking as its hypothesis the negation of this assumption.
If even this minimal requirement is not satisfied, then we must conclude that the given assumption has nothing to do with the theorem.
Notice that Bell expresses his locality condition by the requirement that the result $`B`$ for particle $`2`$ should not depend on the setting $`a`$, of the magnet for particle $`1`$ (cf. citation in the preceeding section). Let us denote $`_1`$ ($`_2`$) the space of all possible measurement settings on system 1 (2).
Theorem (1) For each unit vector $`x`$ in the three dimensional euclidean space ($`x\mathrm{R}^3`$, $`x=1`$) let be given two random variables $`S_x^{(1)}`$, $`S_x^{(2)}`$ (spin of particle 1 (2) in direction $`x`$), defined on a space $`\mathrm{\Omega }`$ with a probability $`P`$ and with values in the $`2`$–point set $`\{+1,1\}`$. Fix $`3`$ of these unit vectors $`a,b,c`$ and suppose that the corresponding random variables satisfy the following non locality condition \[violating Bell’s vital assumption\]: suppose that the probability space $`\mathrm{\Omega }`$ has the following structure:
$$\mathrm{\Omega }=\mathrm{\Lambda }\times _1\times _2$$
$`(1)`$
so that, for some function $`F_a^{(1)},F_a^{(2)}:\mathrm{\Lambda }\times _1\times _2[1,1]`$,
$$S_a^{(1)}(\omega )=F_a^{(1)}(\lambda ,m_1,m_2)(S_a^{(1)}dependsonm_2)$$
$`(2)`$
$$S_a^{(2)}(\omega )=F_a^{(2)}(\lambda ,m_1,m_2)(S_a^{(2)}dependsonm_1)$$
$`(3)`$
with $`m_1_1,m_2_2`$ and similarly for $`b`$ and $`c`$. \[nothing changes in the proof if we add further dependences, for example $`F_a^{(2)}`$ may depend on all the $`S_x^{(1)}(\omega )`$ and $`F_a^{(1)}`$ on all the $`S_x^{(2)}(\omega )`$\].
Then the random variables $`S_a^{(1)},S_b^{(2)},S_c^{(1)}`$ satisfy the inequality
$$S_a^{(1)}S_b^{(2)}S_b^{(2)}S_c^{(1)}1S_a^{(1)}S_c^{(1)}$$
$`(4)`$
If moreover the singlet condition
$$S_x^{(1)}S_x^{(2)}=1;x=a,b,c$$
$`(5)`$
is also satisfied, then Bell’s inequality holds in the form
$$S_a^{(1)}S_b^{(2)}S_b^{(2)}S_c^{(1)}1+S_a^{(1)}S_c^{(2)}$$
$`(6)`$
Proof. The random variables $`S_a^{(1)}`$, $`S_b^{(2)}`$, $`S_c^{(1)}`$ satisfy the assumptions of Corollary (2.3) therefore (4), holds. If also condition (5) is satisfied then, since the variables take values in the set $`\{1,+1\}`$, with probability $`1`$ one must have
$$S_x^{(1)}=S_x^{(2)}(x=a,b,c)$$
$`(7)`$
and therefore $`S_a^{(1)}S_c^{(1)})=S_a^{(1)}S_c^{(2)}`$. Using this identity, (4) becomes (6).
Summing up: Theorem (1) proves that Bell’s inequality is satisfied if one takes as hypothesis the negation of his “vital assumption”. From this we conclude that Bell’s “vital assumption” not only is not “vital” but in fact has nothing to do with Bell’s inequality.
Remark. Using Lemma (14.1) below, we can allow that the observables take values in $`[1,1]`$ also in Theorem (1).
Remark. The above discussion is not a refutation of the Bell inequality: it is a refutation of Bell’s claim that his formulation of locality is an essential assumption for its validity: since the locality assumption is irrelevant for the proof of Bell’s inequality it follows that this inequality cannot discriminate between local and non local hidden variable theories, as claimed both in the introduction and the conclusions of Bell’s paper.
In particular: Theorem (1) gives an example of situations in which:
(i) Bell’s locality condition is violated while his inequality is satisfied.
In a recent experiment with M. Regoli \[AcRe99\] we have produced examples of situations in which:
(ii) Bell’s locality condition is satisfied while his inequality is violated.
## (6) The role of the counterfactual argument in Bell’s proof
Bell uses the counterfactual argument in an essential way in his proof because it is easy to check that formula (13) in \[Bell’64\] paper is the one which allows him to reduce, in the proof of his inequality, all consideration to the $`A`$–variables ($`S_a^{(1)}`$ in our notations, while Bell’s $`B`$–variables are the $`S_a^{(2)}`$ in our notations). The pairs of chameleons (cf. section (10), as well as the experiment of \[AcRe99\] provide a counterexample precisely to this formula.
## (7) Proofs of Bell’s inequality based on counting arguments
There is a widespread illusion to exorcize the above mentioned critiques by restricting one’s considerations to results of measurements. The following considerations show why this is an illusion.
The counting arguments, usually used to prove the Bell inequality are all based on the following scheme. In the same notations used up to now, consider $`N`$ simultaneous measurements of the singlet pairs of observables $`(S_a^1,S_b^2)`$, $`(S_b^2,S_c^1)`$, $`(S_c^2,S_a^1)`$ and one denotes $`S_{x,\nu }^j`$ the results of the $`\nu `$–th measurement of $`S_x^j`$ ($`j=1,2`$, $`x=a,b,c`$, $`\nu =1,\mathrm{},N`$). With these notations one can calculate the empirical correlations on the samples, that is
$$\frac{1}{N}\underset{\nu }{}S_{a,\nu }^1S_{b,\nu }^2=S_a^1S_b^2$$
$`(1)`$
(and similarly for the other ones). In the Bell inequality, 3 such correlations are involved.
$$S_a^1S_b^2,S_b^2S_c^1,S_a^1S_c^2$$
$`(2)`$
Thus in the three experiments observer $`1`$ has to measure $`S_a^1`$ in the first and third experiment and $`S_c^1`$ in the second, while observer $`2`$ has to measure $`S_b^2`$ in the first and second experiment and $`S_c^1`$ in the third. Therefore the directions $`a`$ and $`b`$ can be chosen arbitrarily by the two observers and it is not necessary that observer $`1`$ is informed of the choice of observer $`2`$ or conversely. However the direction $`c`$ has to be chosen by both observers and therefore at least on this direction there should be a preliminary agreement among the two observers. This preliminary information can be replaced it by a procedure in which each observer chooses at will the three directions only those choices are considered for which it happens (by chance) that the second choice of observer $`1`$ coincides with the third of observer $`2`$ (cf. section (15) for further discussion of this point). Whichever procedure has been chosen, after the results of the experiments one can compute the 3 empirical correlations
$$S_a^{(1)}S_b^{(2)}=\frac{1}{N}\underset{j=1}{\overset{N}{}}S_a^{(1)}(p_j^{(1)})S_b^{(2)}(p_j^{(1)})$$
$`(3)`$
$$S_b^{(2)}S_c^{(1)}=\frac{1}{N}\underset{j=1}{\overset{N}{}}S_c^{(1)}(p_j^{(2)})S_b^{(2)}(p_j^{(2)})$$
$`(4)`$
$$S_a^{(1)}S_c^{(2)}=\frac{1}{N}\underset{j=1}{\overset{N}{}}S_a^{(1)}(p_j^{(3)})S_c^{(2)}(p_j^{(3)})$$
$`(5)`$
where $`p_j^{(3)}`$ means the $`j`$–th point of the $`3`$–d experiment etc… If we try to apply the Bell argument directly to the empirical data given by the right hand sides of (3), (4), (5), we meet the expression
$$\frac{1}{N}\underset{j=1}{\overset{N}{}}S_a^{(1)}(p_j^{(1)})S_b^{(2)}(p_j^{(1)})\frac{1}{N}\underset{j=1}{\overset{N}{}}S_c^{(1)}(p_j^{(2)})S_b^{(2)}(p_j^{(2)})$$
$`(6)`$
from which we immediately see that, if we try to apply Bell’s reasoning to the empirical data, we are stuck at the first step because we find a sum of terms of the type
$$S_a^{(1)}(p_j^{(1)})S_b^{(2)}(p_j^{(1)})S_c^{(1)}(p_j^{(2)})S_b^{(2)}(p_j^{(2)})$$
$`(7)`$
to which the inequalities among numbers, of section (1), cannot be applied because in general
$$S_b^{(2)}(p_j^{(1)})S_b^{(2)}(p_j^{(2)})$$
$`(8)`$
More explicitly: since the expression (x.) above is of the form
$$abb^{}c$$
with $`a,b,b^{},c\{\pm 1\}`$, the only possible upper bound for it is $`2`$ and not $`1ac`$.
Even supposing that we, in order to uphold Bell’s thesis, can introduce a cleaning operation \[Ac98\], (cf. \[AcRe99\]), which eliminates all the points in which (8) is not satisfied, we would arrive to the inequality
$$\left|\frac{1}{N}\underset{j=1}{\overset{N}{}}S_a^{(1)}(p_j^{(1)})S_b^{(2)}(p_j^{(1)})\frac{1}{N}\underset{j=1}{\overset{N}{}}S_c^{(1)}(p_j^{(2)})S_b^{(2)}(p_j^{(2)})\right|1\frac{1}{N}\underset{j=1}{\overset{N}{}}S_a^{(1)}(p_j^{(1)})S_c^{(1)}(p_j^{(2)})$$
$`(9)`$
and, in order to deduce from this, something comparable with the experiments we need to use the counterfactual argument, assessing that
$$S_c^{(1)}(p_j^{(2)})=S_c^{(2)}(p_j^{(2)})$$
$`(10)`$
But in the second experiment $`S_b^{(2)}`$ and not $`S_c^{(2)}`$ has been measured. Thus to postulate the validity of (10) means to postulate that: the value assumed by $`S_b^{(2)}`$ in the second experiment is the same that we would have found if $`S_c^{(2)}`$ and not $`S_b^{(2)}`$ had been measured. The chameleon effect provides a counterexample to this statement.
## (8) The quantum probabilistic analysis
Given the results of section (5), (6), (7), it is then legitimate to ask:
if Bell’s vital assumption is irrelevant for the deduction of Bell’s inequality, which is the really vital assumption which guarantees the validity of this inequality?
This natural question was first answered in \[Ac81\] and this result motivated the birth of quantum probability as something more than a mere noncommutative generalization of probability theory; in fact a necessity motivated by experimental data.
Theorem (2.3) has only two assumptions:
(i) that the random variables take values in the interval $`[1,+1]`$
(ii) that the random variables are defined on the same probability space
Since we are dealing with spin variables, assumption (i) is reasonable.
Let us consider assumption (ii). This is equivalent to the claim that the three probability measures $`P_{ab},P_{ac},P_{cb}`$, representing the distributions of the pairs $`(S_a^{(1)},S_b^{(2)})`$, $`(S_c^{(1)},S_b^{(2)})`$, $`(S_a^{(1)},S_c^{(2)})`$ respectively, can be obtained by restriction from a single probability measure $`P`$, representing the distribution of the quadruple $`S_a^{(1)},S_c^{(1)},S_b^{(2)},S_c^{(2)}`$.
This is indeed a strong assumption because, due to the incompatibility of the spin variables along non parallel directions, the three correlations
$$S_a^{(1)}S_b^{(2)},S_c^{(1)}S_b^{(2)},S_a^{(1)}S_c^{(2)}$$
$`(6)`$
can only be estimated in different, in fact mutually incompatible, series of experiments. If we label each series of experiments by the corresponding pair (i.e. $`(a,b),(b,c),(c,a)`$), then we cannot exclude the possibility that also the probability measure in each series of experiments will depend on the corresponding pair. In other words, each of the measures $`P_{a,b},P_{b,c},P_{c,a}`$ describes the joint statistics of a pair of commuting observables $`(S_a^{(1)},S_b^{(2)})`$, $`(S_c^{(1)},S_b^{(2)})`$, $`(S_a^{(1)},S_c^{(2)})`$ and there is no a priori reason to postulate that all these joint distributions for pairs can be deduced from a single distribution for the quadruple $`\{S_a^{(1)},S_c^{(1)},S_b^{(2)},S_c^{(2)}\}`$.
We have already proved in Theorem (2.3) that this strong assumption implies the validity of the Bell inequality. Now let us prove that it is the truly vital assumption for the validity of this inequality, i.e. that, if this assumption is dropped, i.e. if no single distribution for quadruples exist, then it is an easy exercise to construct counterexamples violating Bell’s inequality. To this goal one can use the following lemma:
Lemma (1) . Let be given three probability measures $`P_{ab},P_{ac},P_{cb}`$ on a given (measurable) space $`(\mathrm{\Omega },)`$ and let $`S_a^{(1)},S_c^{(1)},S_b^{(2)},S_d^{(2)}`$ be functions, defined on $`(\mathrm{\Omega },)`$ with values in the interval $`[1,+1]`$, and such that the probability measure $`P_{ab}`$ (resp. $`P_{cb},P_{ac}`$) is the distribution of the pair $`(S_a^{(1)},S_b^{(2)})`$ (resp. $`(S_c^{(1)},S_b^{(2)})`$, $`(S_a^{(1)},S_c^{(2)})`$). For each pair define the corresponding correlation
$$\kappa _{ab}:=S_a^{(1)},S_b^{(2)}:=S_a^{(1)}S_b^{(2)}𝑑P_{ab}$$
and suppose that, for $`\epsilon ,\epsilon ^{}=\pm `$, the joint probabilities for pairs
$$P_{x,y}^{\epsilon \epsilon ^{}}:=P(S_x^{(1)}=\epsilon ;S_y^{(2)}=\epsilon ^{})$$
satisfy:
$$P_{xy}^{++}=P_{xy}^{};P_{xy}^+=P_{xy}^+$$
$`(1)`$
$$P_x^+=P_x^{}=1/2$$
$`(2)`$
then the Bell inequality
$$|\kappa _{ab}\kappa _{bc}|1\kappa _{ac}$$
$`(3)`$
is equivalent to
$$|P_{ab}^{++}P_{bc}^{++}|+P_{ac}^{++}\frac{1}{2}$$
$`(3a)`$
Proof. The inequality (3) is equivalent to
$$|2P_{ab}^{++}2P_{ab}^+2P_{bc}^{++}+2P_{bc}^+|12P_{ac}^{++}+2P_{ac}^+$$
$`(4)`$
Using the identity (equivalent to (2))
$$P_{xy}^+=\frac{1}{2}P_{xy}^{++}$$
$`(5)`$
the left hand side of (4) becomes the modulus of
$$2(P_{ab}^{++}P_{ab}^+)2(P_{bc}^{++}P_{bc}^+)=2\left(P_{ab}^{++}\frac{1}{2}+P_{ab}^{++}\right)2\left(P_{bc}^{++}\frac{1}{2}+P_{bc}^{++}\right)$$
$$=4(P_{ab}^{++}P_{bc}^{++})$$
$`(6)`$
and, again using (5), the right hand side of (4) is equal to
$$12\left(P_{ac}^{++}\frac{1}{2}+P_{ac}^{++}\right)=24P_{ac}^{++}$$
$`(7)`$
Summing up, (3) is equivalent to
$$|P_{ab}^{++}P_{bc}^{++}|\frac{1}{2}P_{ac}^{++}$$
$`(8)`$
which is (3a)
Corollary (2) . There exist triples of $`P_{ab},P_{ac},P_{cb}`$ on the $`4`$–point space $`\{+1,1\}\times \{+1,1\}`$ which satisfy conditions (1), (2) of Lemma (1) and are not compatible with any probability measure $`P`$ on the $`6`$–point space $`\{+1,1\}\times \{+1,1\}\times \{+1,1\}`$.
Proof. Because of conditions (1), (2) the probability measures $`P_{ab},P_{ac},P_{cb}`$ are uniquely determined by the three numbers
$$P_{ab}^{++},P_{ac}^{++},P_{cb}^{++}[0,1]$$
$`(9)`$
Thus, if we choose these three numbers so that the inequality (3a) is not satisfied, the Bell inequality (3) cannot be satisfied because of Lemma (1).
## (9) The realism of ballot boxes and the corresponding statistics
The fact that there is no a priori reason to postulate that the joint distributions of the pairs $`(S_a^{(1)},S_b^{(2)})`$, $`(S_c^{(1)},S_b^{(2)})`$, $`(S_a^{(1)},S_c^{(2)})`$ can be deduced from a single distribution for the quadruple $`S_a^{(1)},S_c^{(1)},S_b^{(2)},S_c^{(2)}`$, does not necessarily mean that such a common joint distribution does not exist.
On the contrary, in several physically meaningful situations, we have good reasons to expect that such a joint distribution should exist even if it might not be accessible to direct experimental verification.
This is a simple consequence of the so–called hypothesis of realism which is justified whenever we are entitled to believe that the results of our measurements are pre–determined. In the words of Bell: Since we can predict in advance the result of measuring any chosen component of $`\sigma _2`$, by previously measuring the same component of $`\sigma _1`$, it follows that the result of any such measurement must actually be predetermined.
Consider for example a box containing pairs of balls. Suppose that the experiments allow to measure either the color or the weight or the material of which each ball is made of, but the rules of the game are that on each ball only one measurement at a time can be performed. Suppose moreover that the experiments show that, for each property, only two values are realized and that, whenever a simultaneous measurement of the same property on the two elements of a pair is performed, the resulting answers are always discordant. Up to a change of convenction and in appropriate units, we can always suppose that these two values are $`\pm 1`$ and we shall do so in the following.
Then the joint distributions of pairs (of properties relative to different balls) are accessible to experiment, but those of triples, or quadruples, are not.
Nevertheless, it is reasonable to postulate that, in the box, there is a well defined (although purely Platonic, in the sense of not being accessible to experiment) number of balls with each given color, weight and material. These numbers give the relative frequencies of triples of properties for each element of the pair hence, using the perfect anticorrelation, a family of joint probabilities for all the possible sextuples. More precisely, due to the perfect anticorrelation, the relative frequency of the triples of properties
$$[S_a^{(1)}=a_1],[S_b^{(1)}=b_1],[S_c^{(1)}=c_1]$$
where $`a_1,b_1,c_1=\pm 1`$ are equal to the relative frequency of the sextuples of properties
$$[S_a^{(1)}=a_1],[S_b^{(1)}=b_1],[S_c^{(1)}=c_1],[S_a^{(2)}=a_1],[S_b^{(2)}=b_1],[S_c^{(2)}=c_1]$$
and, since we are confining ourselves to the case of $`3`$ properties and $`2`$ particles, the above ones, when $`a_1,b_1,c_1`$ vary in all possible ways in the set $`\{\pm 1\}`$, are all the possible configurations in this situation, the counterfactural argument is applicable and in fact we have used it to deduce the joint distribution of sextuples from the joint distributions of triples.
## (10) The realism of chameleons and the corresponding statistics
According to the quantum probabilistic interpretation, what Einstein, Podolsky, Rosen, Bell and several other who have discussed this topic, call the hypothesis of realism should be called in a more precise way the hypothesis of the ballot box realism as opposed to hypothesis of the chameleon realism.
The point is that, according to the quantum probabilistic interpretation, the term predetermined should not be confused with the term realized a priori, which has been discussed in section (9.): it might be conditionally dediced according to the scheme: if such and such will happen, I will react so and so….
The chameleon provides a simple example of this distinction: a chameleon becomes deterministically green on a leaf and brown on a log. In this sense we can surely claim that its color on a leaf is predetermined. However this does not mean that the chameleon was green also before jumping on the leaf.
The chameleon metaphora describes a mechanism which is perfectly local, even deterministic and surely classical and macroscopic; moreover there are no doubts that the situation it describes is absolutely realistic. Yet this realism, being different from the ballot box realism, allows to render free from metaphysics statements of the orthodox interpretation such as: the act of measurement creates the value of the measured observable. To many this looks metaphysic or magic; but load how natural it sounds when you think of the color of a chameleon.
Finally, and most important for its implications relatively to the EPR argument, the chameleon realism provides a simple and natural counterexample of a situation in which the results are predetermined however the counterfactual argument is not applicable.
Imagine in fact a box in which there are many pairs of chameleons. In each pair there is exactly an healthy one, which becomes green on a leaf and brown on a log, and a mutant one, which becomes brown on a leaf and green on a log; moreover exactly one of the chameleons in each pair weights $`100`$ grams and exactly one $`200`$ grams. A measurement consists in separating the members of each pair, each one in a smaller box, and in performing one and only one measurement on each member of each pair.
The color on the leaf, color on the log, and weight are $`2`$–valued observables (because we do not know a priori if we are measuring the healthy or the mutant chameleon). Thus, with respect to the observables: color on the leaf color on the long and weight the pairs of chameleons behave exactly as EPR pairs: whenever the same observable is measured on both elements of a pair, the results are opposite. However, suppose I measure the color on the leaf, of one element of a pair and the weight of the other one and suppose the answers I find are: green and $`100`$ grams. Can I conclude that the second element of the pair is brown and weights $`100`$ grams? Clearly not because there is no reason to believe that the second member of the pair, of which the weight was measured while in a box, was also on a leaf.
From this point of view the measurement interaction enters the very definition of an observable. However also in this interpretation, which is more similar to the quantum mechanical situation, the counterfactual argument cannot be applied because it amounts to answer “brown” to the question: which is the color on the leaf, if I have measured the weight and if I know that the chameleon is the mutant one? (this because the measurement of the other one gave green on the leaf). But this answer is not correct, because it could well be that inside the box there is a leaf and the chameleon is interacting with it while I am measuring its weight, but it could also be that it is interacting with a log, also contained inside the box in which case, being a mutant, it would be green.
Therefore if we can produce an example of a 2-particle system in which the Heisenberg evolution of each particle’s observable satisfies Bell’s locality condition, but the Schroedinger evolution of the state, i.e. the expectation value $``$, depends on the pair $`(a,b)`$ of measured observables, we can claim that this counterexample abides with the same definition of locality as Bell’s theorem.
## (11) Bell’s inequalities and the chamaleon effect
Definition (1) Let $`S`$ be a physical system and $`𝒪`$ a family of observable quantities relative to this system. We say that the it chamaleon effect is realized on $`S`$ if, for any measurement $`M`$ of an observable $`A𝒪`$, the dynamical evolution of $`S`$ depends on the observable $`A`$. If $`D`$ denotes the state space of $`S`$, this means that the change of state from the beginning to the end of the experiment is described by a map (a one–parameter group or semigroup in the case of continuous time)
$$T_A:DD$$
Remark. The explicit form of the dependence of $`T_A`$ on $`A`$ depends on both the system and the measurement and many concrete examples can be constructed. An example in the quantum domain is discussed in \[Ac98\] and the experiment of \[AcRe99\] realizes an example in the classical domain.
Remark If the system $`S`$ is composed of two sub–systems $`S_1`$ and $`S_2`$, we can also consider the case in which the evolutions of the two subsystems are different in the sense that, for system $`1`$, we have one form of functional dependence, $`T_A^{(1)}`$, of the evolution associated to the observable $`A`$ and, for system $`2`$, we have another form of functional dependence, $`T_A^{(2)}`$. In the experiment of \[AcRe99\], the state space is the unit disk $`D`$ in the plane, the observables are parametrized by angles in $`[0,2\pi )`$ (or equivalently by unit vectors in the unit circle) and, for each observable $`S_\alpha ^{(1)}`$ of system $`1`$
$$T_\alpha ^{(1)}:=R_\alpha $$
and, for each observable $`S_\alpha ^{(2)}`$ of system $`2`$
$$T_\alpha ^{(2)}:=R_{\alpha +\pi }$$
where $`R_\alpha `$ denotes (counterclockwise) rotation of an angle $`\alpha `$.
Let us consider Bell’s inequalities by assuming that a chamaleon effect
$$(S_a^{(1)},S_b^{(2)})(S_a^{(1)}T_a^{(1)},S_b^{(2)}T_b^{(2)})$$
is present. Denoting $`E`$ the common initial state of the composite system $`(1,2)`$, (e.g. singlet state), the state at the end of the measurement will be
$$E(S_a^{(1)}T_a^{(1)},S_b^{(2)}T_b^{(2)})$$
Now replace $`S_x^{(j)}`$ by:
$$\stackrel{~}{S}_x^{(j)}:=S_x^{(j)}T_x^{(j)}$$
Since the $`\stackrel{~}{S}_x^{(j)}`$ take values $`\pm 1`$, we know from Theorem (2.3) that, if we postulate the existence of joint probabilities for the triple $`\stackrel{~}{S}_a^{(1)},\stackrel{~}{S}_b^{(2)},\stackrel{~}{S}_c^{(1)}`$, compatible with the two correlations $`E(\stackrel{~}{S}_a^{(1)}\stackrel{~}{S}_b^{(2)}),E(\stackrel{~}{S}_c^{(1)}\stackrel{~}{S}_b^{(2)})`$, then the inequality
$$|E(\stackrel{~}{S}_a^{(1)}\stackrel{~}{S}_b^{(2)})E(\stackrel{~}{S}_c^{(1)}\stackrel{~}{S}_b^{(2)})|1E(\stackrel{~}{S}_a^{(1)}\stackrel{~}{S}_c^{(1)})$$
holds and, if we also have the singlet condition
$$E(S_c^{(1)}(T_c^{(1)}p)S_c^{(2)}(T_c^{(2)}p))=1$$
$`(1)`$
then a.e.
$$\stackrel{~}{S}_c^{(1)}=\stackrel{~}{S}_c^{(2)}$$
and we have the Bell’s inequality. Thus, if we postulate the same probability space, even the chamaleon effect alone is not sufficient to guarantee violation of the Bell’s inequality.
Therefore the fact that the three experiments are done on different and incompatible samples must play a crucial role.
As far as the chameleon effect is concerned, let us notice that, in the above statement of the problem the fact that we use a single initial probability measure $`E`$ is equivalent to postulate that, at time $`t=0`$ the three pairs of observables
$$(S_a^{(1)},S_b^{(2)}),(S_c^{(1)},S_b^{(2)}),(S_a^{(1)},S_c^{(1)})$$
admit a common joint distribution, in fact $`E`$.
## (12) Physical implausibility of Bell’s argument
In this section we show that, combining the chameleon effect with the fact that the three experiments refer to different samples, then even in very simple situations, no cleaning conditions can lead to a proof of the Bell’s inequality.
If we try to apply Bell’s reasoning to the empirical data, we have to start from the expression
$$\left|\frac{1}{N}\underset{j}{}S_a^{(1)}(T_a^{(1)}p_j^I)S_b^{(2)}(T_b^{(2)}p_j^I)\frac{1}{N}\underset{j}{}S_c^{(1)}(T_c^{(1)}p_j^{II})S_b^{(2)}(T_b^{(2)}p_j^{II})\right|$$
$`(1)`$
which we majorize by
$$\frac{1}{N}\underset{j}{}\left|S_a^{(1)}(T_a^{(1)}p_j^I)S_b^{(2)}(T_b^{(2)}p_j^I)S_c^{(1)}(T_c^{(1)}p_j^{II})S_b^{(2)}(T_b^{(2)}p_j^{II})\right|$$
$`(2)`$
But, if we try to apply the inequality among numbers to the expression
$$\left|S_a^{(1)}(T_a^{(1)}p_j^I)S_b^{(2)}(T_b^{(2)}p_j^I)S_c^{(1)}(T_c^{(1)}p_j^{II})S_b^{(2)}(T_b^{(2)}p_j^{II})\right|$$
$`(3)`$
we see that we are not dealing with the situation covered by Corollary (1.2), i.e.
$$|abcb|1ac$$
$`(4)`$
because, since
$$S_b^{(2)}(T_b^{(2)}p_j^I)S_b^{(2)}(T_b^{(2)}p_j^{II})$$
$`(5)`$
the left hand side of (4) must be replaced by
$$|abcb^{}|$$
$`(6)`$
whose maximum, for $`a,b,c,b^{}[1,+1]`$ is $`2`$ and not $`1ac`$.
Bell’s implicit assumption of the single probability space is equivalent to the postulate that, for each $`j=1,\mathrm{},N`$
$$p_j^I=p_j^{II}$$
$`(7)`$
Physically this means that:
the hidden parameter in the first experiment is the same as the hidden parameter in the second experiment
This is surely a very implausible assumption.
Notice however that, without this assumption, Bell’s argument cannot be carried over and we cannot deduce the inequality because we must stop at equation (2).
## (13) The role of the single probability space in CHSH’s proof
Clauser, Horne, Shimony, Holt \[ClHo69\] introduced the variant (2.6) of the Bell inequality for quadruples $`(a,b)`$, $`(a,b^{})`$, $`(a^{},b)`$, $`(a^{},b^{})`$ which is based on the following inequality among numbers
$$ab+ab^{}+a^{}ba^{}b^{}2$$
$`(1)`$
Section (1) already contains a proof of (1). For $`a,b,b^{},a[1,1]`$, a direct proof follows from
$$b+b^{}+bb^{}2$$
$`(2)`$
because
$$ab+ab^{}+a^{}ba^{}b^{}=a(b+b^{})+a^{}(bb^{})ab+b^{}+a^{}bb^{}b+b^{}+bb^{}2$$
The proof of (2) is obvious because it is equivalent to
$$b+b^{}^2+bb^{}^2=b^2+b^2+2bb^{}+b^2+b^22bb^{}=2b^2+2b^24$$
which is identically satisfied (cf. also Lemma (1.1)).
Remark (1) Notice that an inequality of the form
$$a_1b_1+a_2b_2^{}+a_3^{}b_3a_4^{}b_4^{}2$$
$`(3)`$
would be obviously false. In fact, for example the choice
$$a_1=b_1=a_2=b_2^{}=a_3^{}=b_3=b_4^{}=1;a_4^{}=1$$
would give
$$a_1b_1+a_2b_2^{}+a_3^{}b_3a_4^{}b_4^{}=4$$
That is: for the validity of (1) it is absolutely essential that the number $`a`$ is the same in the first and the second term and similarly for $`a^{}`$ in the 3–d and the 4–th, $`b^{}`$ in the 2–d and the 4–th, $`b`$ in the first and the 3–d.
This inequality among numbers can be extended to pairs of random variables by introducing the following postulates:
(P1) Instead of four numbers $`a,b,b^{},a[1,1]`$, one considers four functions
$$S_a^{(1)},S_b^{(2)},S_a^{}^{(1)},S_b^{}^{(2)}$$
all defined on the same space $`\mathrm{\Lambda }`$ (whose points are called hidden parameters) and with values in $`[1,1]`$.
(P2) One postulates that there exists a probability measure $`P`$ on $`\mathrm{\Lambda }`$ which defines the joint distribution of each of the following four pairs of functions
$$(S_a^{(1)},S_b^{(2)}),(S_a^{(1)},S_b^{}^{(2)}),(S_a^{}^{(1)},S_b^{(2)}),(S_a^{}^{(1)},S_b^{}^{(2)})$$
$`(4)`$
Remark (2) Notice that $`(P2)`$ automatically implies that the joint distributions of the four pairs of functions can be deduced from a joint distribution of the whole quadruple, i.e. the existence of a single Kolmogorov model for these four pairs.
With these premises, for each $`\lambda \mathrm{\Lambda }`$ one can apply the inequality (1) to the four numbers
$$S_a^{(1)}(\lambda ),S_b^{(2)}(\lambda ),S_a^{}^{(1)}(\lambda ),S_b^{}^{(2)}(\lambda )$$
and deduce that
$$S_a^{(1)}(\lambda )S_b^{(2)}(\lambda )+S_a^{(1)}(\lambda )S_b^{}^{(2)}(\lambda )+S_a^{}^{(1)}(\lambda )S_b^{(2)}(\lambda )S_a^{}^{(1)}(\lambda )S_b^{}^{(2)}(\lambda )2$$
$`(5)`$
From this, taking $`P`$–averages, one obtains
$$S_a^{(1)}S_b^{(2)}+S_a^{(1)}S_b^{}^{(2)}+S_a^{}^{(1)}S_b^{(2)}S_a^{}^{(1)}S_b^{}^{(2)}=$$
$`(6a)`$
$$\left(S_a^{(1)}(\lambda )S_b^{(2)}(\lambda )+S_a^{(1)}(\lambda )S_b^{}^{(2)}(\lambda )+S_a^{}^{(1)}(\lambda )S_b^{(2)}(\lambda )S_a^{}^{(1)}(\lambda )S_b^{}^{(2)}(\lambda )\right)𝑑P(\lambda )$$
$`(6b)`$
$$S_a^{(1)}(\lambda )S_b^{(2)}(\lambda )+S_a^{(1)}(\lambda )S_b^{}^{(2)}(\lambda )+S_a^{}^{(1)}(\lambda )S_b^{(2)}(\lambda )S_a^{}^{(1)}(\lambda )S_b^{}^{(2)}(\lambda )𝑑P(\lambda )2$$
$`(6c)`$
Remark (3) Notice that in the step from (6a) to (6b) we have used in an essential way the existence of a joint distribution for the whole quadruple, i.e. the fact that all these random variales can be realized in the same probability space.
In EPR type experiments we are interested in the case in which the four pairs $`(a,b)`$, $`(a,b^{})`$, $`(a^{},b)`$, $`(a^{},b^{})`$ come from four mutually incompatible experiments. Let us assume that there is a hidden parameter, determining the result of each of these experiments. This means that we interpret the number $`S_a^{(1)}(\lambda )`$ as the value of the spin of particle $`1`$ in direction $`a`$, determined by the hidden parameter $`\lambda `$.
There is obviously no reason to postulate that the hidden parameter, determining the result of the first experiment is exactly the same one which determines the result of the second experiment. However, when CHSH consider the quantity (5), they are implicitly doing the much stronger assumption that the same hidden parameter $`\lambda `$ determines the results of all the four experiments. This assumption is quite unreasonable from the physical point of view and in any case it is a much stronger assumption than simply postulating the existence of hidden parameters. The latter assumption would allow CHSH only to consider the expression
$$S_a^{(1)}(\lambda _1)S_b^{(2)}(\lambda _1)+S_a^{(1)}(\lambda _2)S_b^{}^{(2)}(\lambda _2)+S_a^{}^{(1)}(\lambda _3)S_b^{(2)}(\lambda _3)S_a^{}^{(1)}(\lambda _4)S_b^{}^{(2)}(\lambda _4)$$
$`(4)`$
and, as shown in Remark (1.) above the maximum of this expression is not $`2`$ but $`4`$ and this does not allow to deduce the Bell inequality.
## (14) The role of the counterfactual argument in CHSH’s proof
Contrarily to the original Bell’s argument, the CHSH proof of the Bell inequality does not use explicitly the counterfactual argument. Since one can perform experiments also on quadruples, rather than on triples, as originally proposed by Bell, has led some authors to claim that the counterfactual argument is not essential in the deduction of the Bell inequality. However we have just seen in section (7.) that the hidden assumption as in Bell’s proof, i.e. the realizability of all the random variales involved in the same probability space, is also present in the CHSH argument. The following lemma shows that, under the singlet assumption, the conclusion of the counterfactual argument follows from the hidden assumption of Bell and of CHSH.
Lemma (1) If $`f`$ and $`g`$ are random variables defined on a probability space $`(\mathrm{\Lambda },P)`$ and with values in $`[1,1]`$, then
$$fg:=_\mathrm{\Lambda }fg𝑑P=1$$
if and only if
$$P(fg=1)=1$$
Proof. If $`P(fg>1)>0`$, then
$$_\mathrm{\Lambda }fg𝑑P=P(fg=1)_{fg>1}|fg|𝑑P>P(fg=1)P(fg>1)>1$$
Corollary (2) Suppose that all the random variales in (x.3) are realized in the same probability space. Then, if the singlet condition:
$$S_x^{(1)}S_x^{(2)}=1$$
$`(1)`$
is satisfied, then the condition
$$S_x^{(1)}=S_x^{(2)}$$
$`(2)`$
(i.e. formula (13) in Bell’s ’64 paper) is true almost everywhere.
Proof. Follows from Lemma (1) with the choice $`f=S_x^{(1)}`$, $`g=S_x^{(2)}`$.
Summing up: if you want to compare the predictions of a hidden variable theory with quantum theory in the EPR experiment (so that at least we admit the validity of the singlet law) then the hidden assumption, of realizability of all the random variables in (3) in the same probability space, (without which Bell’s inequality cannot be proved) implies the same conclusion of the counterfactual argument. Stated otherwise: the counterfactual argument is implicit when you postulate the singlet condition and the realizability on a single probability space. It does not matter if you use triples or quadruples.
## (15) Physical difference between the CHSH’s and the original Bell’s inequalities
In the CHSH scheme:
$$(a,b),(a^{},b^{}),(a,b^{}),(a^{},b^{})$$
the agreement required by the experimenters is the following:
$`1`$ will measures the same observable in experiments I and III, and the same observable in experiments II and IV;
$`2`$ will measure the same observable in experiments I and II, and the same observable in experiments III and IV.
Here there is no restriction a priori on the choice of the observables to be measured.
In the Bell scheme the experimentalists agree that:
$`1`$ measures the same observable in experiments I and III,
$`2`$ measures the same observable in experiments I and II
$`1`$ and $`2`$ choose a priori, i.e. before the experiment begins, a direction $`c`$ and agree that $`1`$ will measure spin in direction $`c`$ in experiment II and $`2`$ will measure spin in direction $`c`$ in experiment III (strong agreement)
The strong agreement can be replaced by the following (weak agreement):
$`1`$ and $`2`$ choose a priori, i.e. before the experiment begins, a finite set of directions $`c_1,\mathrm{},c_K`$ and agree that $`1`$ will measure spin in a direction choosen randomly among the directions $`c_1,\mathrm{},c_K`$ in experiment II and $`2`$ will do the same in experiment III
In this scheme there is an a priori restriction on the choice of some of the observables to be measured.
If the directions, fixed a priori in the plane, are $`K`$, then the probability of a coincidence, corresponding to a totally random (equiprobable) choice, is
$$P(x_{II}^{(1)}=x_{III}^{(2)})=\underset{\alpha =1}{\overset{K}{}}(x_{II}^{(1)}=\alpha ;x_{III}^{(2)}=\alpha )=\underset{\alpha =1}{\overset{K}{}}\frac{1}{K^2}=\frac{1}{K}$$
This shows that, contrarily than in the CHSH scheme, the choice has to be restricted to a finite number of possibilities otherwise the probability of coincidence will be zero.
From this point of view we can claim that the Clauser, Horne, Shimony, Holt formulation of Bell’s inequalities realize an improvement with respect to the original Bell’s formulation.
## Bibliography
\[Ac81\] Luigi Accardi: “Topics in quantum probability”, Phys. Rep. 77 (1981) 169-192
\[Ac97\] Luigi Accardi: Urne e camaleonti. Dialogo sulla realtà, le leggi del caso e la teoria quantistica. Il Saggiatore (1997). Japanese translation, Maruzen (2000), russian translation, ed. by Igor Volovich, PHASIS Publishing House (2000), english translation by Daniele Tartaglia, to appear
\[Ac99\] Luigi Accardi: On the EPR paradox and the Bell inequality Volterra Preprint (1998) N. 350.
\[AcRe99a\] Luigi Accardi, Massimo Regoli: Quantum probability and the interpretation of quantum mechanics: a crucial experiment, Invited talk at the workshop: “The applications of mathematics to the sciences of nature: critical moments and aspetcs”, Arcidosso June 28-July 1 (1999). To appear in the proceedings of the workshop, Preprint Volterra N. 399 (1999)
\[AcRe99b\] Luigi Accardi, Massimo Regoli: Local realistic violation of Bell’s inequality: an experiment, Conference given by the first–named author at the Dipartimento di Fisica, Università di Pavia on 24-02-2000, Preprint Volterra N. 402
\[AcRe00\] Luigi Accardi, Massimo Regoli: Non–locality and quantum theory: new experimental evidence, Invited talk given by the first–named author at the Conference: “Quantum paradoxes”, University of Nottingham, on 4-05-2000, Preprint Volterra N. 421
\[Be64\] Bell J.S: On the Einstein Podolsky Rosen Paradox Physics 1 no.3. 195-200 1964.
\[Be66\] Bell J.S.: On the Problem of Hidden Variables in Quantum Mechanics. Rev. Mod. Phys. 38 (1966) 447-452
\[ClHo69\] J.F. Clauser , M.A. Horne, A. Shimony, R. A. Holt, Phys. Rev. Letters, 49, 1804-1806 (1969); J. S. Bell, Speakable and unspeakable in quantum mechanics. (Cambridge Univ. Press, 1987).
\[ClHo74\] Clauser J.F., Horne M.A.: Experimental Consequences of Objective Local Theories. Physical Review D, vol. 10, no. 2 (1974)
\[EPR35\] Einstein A., Podolsky B., Rosen N. Can quantum mechanical description of reality be considered complete ? Phys. Rev. 47 (1935) 777-780
\[Ein49\] A. Einstein in: Albert Einstein: Philosopher Scientist. Edited by P.A. Schilpp, Library of Living Philosophers, Evanston, Illinois, p.85 (1949)
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# Isotropic Conductivity of Two-Dimensional Three-Component Symmetric Composites
### Abstract
## I Introduction
The classical duality transformation of two-dimensional (2D) heterogeneous composites, discovered by Keller and independently by Dykhne , has been applied in a restricted set of physical contexts. The dual symmetry is based upon the observation that any 2D divergence-free field, when rotated locally at each point by $`90^o`$, becomes curl-free, and vice versa. This leads to the result that static effective physical properties of 2D infinite heterogeneous composites, like electric conductivity $`\widehat{\sigma }_e`$, thermal conductivity, dielectric permittivity, as well as some other static properties, satisfy some exact relationships which follow from the similarity between dual problems. For example, in the case of a two-component composite, a universal square-root law behavior of the bulk effective transport characteristics was demonstrated ,. It is also known<sup>*</sup><sup>*</sup>*It seems to be strange but we have not found throughout the papers concerned with this subject any published proof of this statement. Such proof is so useful that we give it in Appendix. that this transformation leads to exact results in 2D infinite composites made of an arbitrary number of components. Dykhne gave sufficient conditions that, when satisfied by the component conductivities $`\sigma _i`$, lead to exact results for the isotropic bulk effective conductivity $`\sigma _e(\sigma _i)`$. We know of only one attempt to consider a 3-component 2D composite with a doubly periodic arrangement of two kinds of circular inclusions embedded into the matrix. The effective permittivity was obtained using the dipole approximation for the inclusion polarizations. But we do not know any rigorous results obtained for any 2D composite microstructure, made of 3 (or more) components, with arbitrary component conductivities. Apparently, this is not an accident but reflects the disappearance of commutativity in symmetric groups when we upgrade from the $`S_2`$ permutation group to the $`S_n`$ permutation group.
In this paper we consider the effective isotropic conductivity problem for $`3`$-component 2D infinite composites with translational order, with a microstructure that is symmetric in the 3 components. We formulate an approach based on the conjecture of the algebraicity of $`\sigma _e(\sigma _1,\sigma _2,\sigma _3)`$ and its general properties. We have found that the algebraic equation of minimal order where all these properties can be satisfied is a cubic equation, which contains 1 free parameter, with coefficients made of the independent invariants of $`S_3`$. This equation is in agreement with Dykhne’ result (see Eq. (7)). Its predictions are compared with numerical solutions for $`\sigma _e`$ in some regular 3-component microstructures. It is also extended so as to apply to other types of 3-component composites, including random microstructures.
## II TWO-DIMENSIONAL THREE-COLOR COMPOSITES
The effective dc-conductivity problem for $`n`$-component symmetric, 2D, infinite composites with translational order, can be reformulated with the help of n-color plane groups.
Color groups are generalizations of the classical crystallographic groups. Different colors may correspond, for example, to different chemical species or, more generally, to different values of a physical property which is defined as a tensor of $`k`$-rank : scalar - density $`\rho `$ , vector - magnetic spin m , tensor of 2-nd rank - conductivity $`\widehat{\sigma }_{ij}`$, tensor of 3-rank - piezoelectric modulus $`\widehat{d}_{ijn}`$, etc.
Every $`n`$-color plane group has its origin in one of the $`N_1=`$17 regular (color-blind ) plane groups. The number $`N_n`$ of n-color plane groups is a non-monotonic function of $`n`$: $`N_2=46`$ ; $`N_3=23`$; $`N_4=96`$; $`N_5=14`$; $`N_6=90`$ , . The plane groups $`n60`$ are tabulated in .
Only 10 of the 23 3-color plane groups have a 3-fold rotation axis:
5 lattice equivalent groups \- $`\mathrm{𝖯𝟥}(𝖫)|\mathrm{𝖯𝟣}(𝖫)`$, $`\mathrm{𝖯𝟨}(𝖫)|\mathrm{𝖯𝟤}(𝖫)`$, $`\mathrm{𝖯𝟥𝟣𝗆}(𝖫)|\mathrm{𝖢𝗆}(𝖫)`$, $`\mathrm{𝖯𝟥𝗆𝟣}(𝖫)|\mathrm{𝖢𝗆}(𝖫)`$, $`\mathrm{𝖯𝟨𝗆𝗆}(𝖫)|\mathrm{𝖢𝗆𝗆}(𝖫)`$ and 5 class equivalent groups - $`\mathrm{𝖯𝟥}(𝖫)|\mathrm{𝖯𝟥}(𝖫^{^{}})`$, $`\mathrm{𝖯𝟨}(𝖫)|\mathrm{𝖯𝟨}(𝖫^{^{}})`$, $`\mathrm{𝖯𝟥𝟣𝗆}(𝖫)|\mathrm{𝖯𝟥𝗆𝟣}(𝖫^{^{}})`$,$`\mathrm{𝖯𝟥𝗆𝟣}(𝖫)|\mathrm{𝖯𝟥𝟣𝗆}(𝖫^{^{}})`$, $`\mathrm{𝖯𝟨𝗆𝗆}(𝖫)|\mathrm{𝖯𝟨𝗆𝗆}(𝖫^{^{}})`$ where $`𝖫^{^{}}`$ is a possible sub-latticeWe follow the notations of Ref. for three-color plane groups $`𝖦_{\mathrm{𝖼𝗈𝗅}}=𝖦(𝖫)|𝖦^{^{}}(𝖫^{^{}})`$ which means that G is the geometrical, or Fedorov, plane group and the subgroup $`𝖦^{^{}}𝖦`$ of index 3 contains operations that keep the first color fixed. We do not specify here the relationship between the lattice L and its sub-lattice $`𝖫^{^{}}`$. of $`𝖫`$ invariant under rotations of order 3 . All these groups are compatible only with hexagonal Bravais lattices. The 3-fold rotational symmetry makes the effective conductivity in structures governed by those groups isotropic. This follows from the Hermann theorem about $`k`$-rank tensors in media with an inner symmetry which includes a rotation axis of highest order $`r,r>k`$. Despite their different geometries, all these structures have one important property in common: They are all invariant under the full permutation group $`S_3`$ which exchanges the colors, therefore they are related to the permutational crystallographic color groups .
## III EFFECTIVE CONDUCTIVITY $`\sigma _e(\sigma _1,\sigma _2,\sigma _3)`$ and ITS ALGEBRAIC PROPERTIES
A direct way to solve the dc-conductivity problem for an $`n`$-component composite begins with the local field equations
$$\times 𝐄(𝐫)=0,𝐉(𝐫)=0,𝐉(𝐫)=\sigma (𝐫)𝐄(𝐫),$$
(1)
along with appropriate boundary conditions for the electrical potential. The local conductivity $`\sigma (𝐫)`$ is a discontinuous function $`\sigma (𝐫)=\sigma _i`$, if $`𝐫\mathrm{\Delta }_i,i=1,2,3`$ where $`\mathrm{\Delta }_i`$ is a homogeneous part of the composite with constant conductivity $`\sigma _i`$. The isotropic effective conductivity $`\sigma _e`$ can be defined via Ohm’s law for the current $`𝐉_e`$ and the field $`𝐄_e`$ averaged over the system
$$𝐉_e=\sigma _e𝐄_e,𝐉_e=\frac{1}{S}𝐉(𝐫)𝑑S,𝐄_e=\frac{1}{S}𝐄(𝐫)𝑑S.$$
(2)
Except for a medium with trivial 1D inhomogeneities (a layered medium) an exact solution of this problem does not exist for any regular or random structure. At the same time the function $`\sigma _e(\sigma _i)`$ must have the following general properties which we are going to exploit:
$`\mathrm{𝟣}.\mathrm{𝖧𝗈𝗆𝗈𝗀𝖾𝗇𝖾𝗂𝗍𝗒}\mathrm{𝗈𝖿}\mathsf{\hspace{0.33em}1}\mathrm{𝗌𝗍}\mathrm{𝗈𝗋𝖽𝖾𝗋}`$
$$\sigma _e(k\sigma _1,k\sigma _2,k\sigma _3)=k\sigma _e(\sigma _1,\sigma _2,\sigma _3).$$
(3)
This follows from the linearity of the static Maxwell equations (1) and from the definitions of the average current and field (2).
$`\mathrm{𝟤}.\mathrm{𝖯𝖾𝗋𝗆𝗎𝗍𝖺𝗍𝗂𝗈𝗇}\mathrm{𝗂𝗇𝗏𝖺𝗋𝗂𝖺𝗇𝖼𝖾}`$
$$\sigma _e(\widehat{𝒫}_l\{\sigma _1,\sigma _2,\sigma _3\})=\sigma _e(\sigma _1,\sigma _2,\sigma _3).$$
(4)
where $`\widehat{𝒫}_l,l=1,\mathrm{},6`$ is a permutation operator of the indices $`\{1,2,3\}`$ (or of 3 colors). The six operators $`\widehat{𝒫}_l`$ form the non-commutative group $`S_3`$. The existence of permutation invariance presumes that the 3 components are distributed with equal volume fractions $`p_1=p_2=p_3=1/3`$.
$`\mathrm{𝟥}.\mathrm{𝖣𝗎𝖺𝗅𝗂𝗍𝗒}`$
$$\sigma _e(\sigma _1,\sigma _2,\sigma _3)\times \sigma _e(\frac{1}{\sigma _1},\frac{1}{\sigma _2},\frac{1}{\sigma _3})=1.$$
(5)
See Eq. (A16).
$`\mathrm{𝟦}.\mathrm{𝖢𝗈𝗆𝗉𝖺𝗍𝗂𝖻𝗂𝗅𝗂𝗍𝗒}`$
$$\sigma _e(\sigma ,\sigma ,\sigma )=\sigma .$$
(6)
The last formula does not follow from the previous ones but reflects a natural requirement.
Dykhne proved a theorem for a symmetric composite of 3 components, namely
$$\sigma _e(\sigma _1,\sigma _2,\sigma _3)=\sqrt{\sigma _1\sigma _2},\text{if}\sigma _3^2=\sigma _1\sigma _2.$$
(7)
The last formula represents the only known rigorous result for 3-color 2D isotropic composites. It can easily be shown to follow from (4) and (5). Let us mention one more conclusion which follows from (7)
$$\sigma _e(\sigma ,\mathrm{\hspace{0.33em}0},\mathrm{\hspace{0.33em}0})=0,$$
(8)
which reflects the percolation property of such a composite. In fact, it is easy to construct 3-color 2D isotropic composites with the more specialized property
$$\sigma _e(\sigma _1,\sigma _2,\mathrm{\hspace{0.33em}0})=0.$$
(9)
Fig.1d illustrates such a structure with traps, i.e. a network of closed, simply connected loops, which enclose plaquettes of the two other colors. If any color denotes an insulator, then the effective conductivity of the composite is clearly zero. The 2-color 2D square checkerboard has a similar property. It differs from the case of Fig.1d by the value of the percolation threshold. The previous example (9) shows a non-universal behavior of $`\sigma _e(\sigma _i)`$ even for symmetric isotropic microstructures, with essential dependence on micro-structural details, in contrast to the 2-color 2D symmetric isotropic composites, where $`\sigma _e=\sqrt{\sigma _1\sigma _2}`$ always .
We will look for $`\sigma _e(\sigma _1,\sigma _2,\sigma _3)`$ satisfying the requirements (3 \- 6) among algebraic functions. This choice is inspired by the fact that the symmetric 2-color 2D isotropic composite generates a quadratic equation for $`\sigma _e(\sigma _1,\sigma _2)`$. Another motivation is the situation which exists vis-a-vis some discrete 2D models in statistical mechanics, where a duality transformation exists that is similar to the one invoked here. Such a transformation exists for the Ising and Potts models , where the critical point equations are algebraic — quadratic equation for the Ising model on a self-dual square lattice and cubic equations for the Potts model on the mutually dual triangle and honeycomb lattices.
According to the “fundamental theorem of symmetric functions” the symmetric group $`S_3`$ has 3 algebraically independent homogeneous invariants (basic invariants)
$`I_1`$ $`=`$ $`\sigma _1+\sigma _2+\sigma _3,I_2=\sigma _1\sigma _2+\sigma _2\sigma _3+\sigma _3\sigma _1,`$ (10)
$`I_3`$ $`=`$ $`\sigma _1\sigma _2\sigma _3,`$ (11)
which satisfy the obvious restrictions
$`{\displaystyle \frac{1}{3}}I_1{\displaystyle \frac{I_2}{I_1}}\mathrm{\hspace{0.33em}3}{\displaystyle \frac{I_3}{I_2}},{\displaystyle \frac{1}{3}}I_1\sqrt[3]{I_3}\mathrm{\hspace{0.33em}3}{\displaystyle \frac{I_3}{I_2}},`$
and can be used as independent variables instead of $`\sigma _1,\sigma _2,\sigma _3`$. The difference $`I_2I_1\sqrt[3]{I_3}`$ can have either sign. The function $`\sigma _e(I_1,I_2,I_3)`$ now satisfies (4) automatically.
It can be shown that the algebraic equation of minimal order for the function $`\sigma _e(I_1,I_2,I_3)`$, which satisfies all basic requirements, is a cubic equation of the form
$$\sigma _e^3+AI_1\sigma _e^2AI_2\sigma _eI_3=0,$$
(12)
where $`A`$ is a free parameter responsible for the non-universality, and $`\sigma _e`$ is a value bounded from both above and below
$$\frac{1}{3}I_1\sigma _e3\frac{I_3}{I_2}.$$
(13)
It is easy to show that the equation (12) satisfies all basic requirements (3-6) and automatically satisfies Dykhne’s theorem (7) independently of $`A`$. Indeed, the requirements (3, 4, 6) one can check immediately. The duality property (5) follows when we notice that (12) is equivalent to the cubic equation for $`\frac{1}{\sigma _e}(\frac{1}{\sigma _1},\frac{1}{\sigma _2},\frac{1}{\sigma _3})`$:
$`{\displaystyle \frac{1}{\sigma _e^3}}+`$ $`A`$ $`({\displaystyle \frac{1}{\sigma _1}}+{\displaystyle \frac{1}{\sigma _2}}+{\displaystyle \frac{1}{\sigma _3}}){\displaystyle \frac{1}{\sigma _e^2}}`$ (14)
$`A`$ $`({\displaystyle \frac{1}{\sigma _1\sigma _2}}+{\displaystyle \frac{1}{\sigma _2\sigma _3}}+{\displaystyle \frac{1}{\sigma _3\sigma _1}}){\displaystyle \frac{1}{\sigma _e}}{\displaystyle \frac{1}{\sigma _1\sigma _2\sigma _3}}=0.`$ (15)
The property (7) can be proven by straightforward substitution $`\sigma _3=\sqrt{\sigma _1\sigma _2}`$ into (12) .
As we will see later, $`A`$ reflects not only the plane group of a color tessellation but also the shape of the elementary cell. To illuminate the meaning of $`A`$ let us put into (12) $`A=\frac{1}{3}`$. After simple algebra one obtains the following equation
$$\frac{\sigma _e\sigma _1}{\sigma _e+\sigma _1}+\frac{\sigma _e\sigma _2}{\sigma _e+\sigma _2}+\frac{\sigma _e\sigma _3}{\sigma _e+\sigma _3}=0,$$
(16)
which coincides with the Bruggeman effective medium approximation for a symmetric, 3-component, 2D composite .
From restrictions (13) and after straightforward manipulations, one can show that $`A`$ is bounded from below
$$\mathrm{}A\mathrm{\hspace{0.33em}0}.$$
(17)
It is noteworthy that the lower bound ensures that (12) has only one positive root, thus avoiding the possibility of multiple physical solutions. More accurate estimation of the lower bound for $`A`$, using the Hashin-Shtrikman exact bounds $`\sigma _{HS}^\pm `$ for isotropic conductivity $`\sigma _e`$ of 3-component 2D composite, does not change (17). Indeed, for $`\sigma _1>\sigma _2>\sigma _3`$ one has
$$\frac{1}{3}I_1\sigma _{HS}^+\sigma _e\sigma _{HS}^{}3\frac{I_3}{I_2},$$
(18)
where
$`\sigma _{HS}^+`$ $`=`$ $`\sigma _1(1+{\displaystyle \frac{4(\sigma _2\sigma _3\sigma _1^2)}{\sigma _2\sigma _3+3\sigma _1(\sigma _2+\sigma _3)+5\sigma _1^2}}),`$ (19)
$`\sigma _{HS}^{}`$ $`=`$ $`\sigma _3(1+{\displaystyle \frac{4(\sigma _1\sigma _2\sigma _3^2)}{\sigma _1\sigma _2+3\sigma _3(\sigma _1+\sigma _2)+5\sigma _3^2}}).`$ (20)
Due to the ambiguity of the general case, we consider the special case $`\sigma _1>\sigma _2=\sigma _3`$. Here we have
$`\sigma _{HS}^+=\sigma _1{\displaystyle \frac{5+r}{1+5r}},\sigma _{HS}^{}=\sigma _3{\displaystyle \frac{1+2r}{2+r}},r={\displaystyle \frac{\sigma _1}{\sigma _3}}1.`$
We can then derive the following lower bounds from (12), (13)
$`A>{\displaystyle \frac{1+5r}{1+2r}}{\displaystyle \frac{1+r}{(2+r)^2}},`$ (21)
$`A>{\displaystyle \frac{r}{1+5r}}{\displaystyle \frac{2+r}{1+r}}{\displaystyle \frac{1+18r+r^2}{1+16r+r^2}}.`$ (22)
Since the last inequalities must fit all cases including $`r\mathrm{}`$, we return to (17).
## IV EFFECTIVE CONDUCTIVITY of REGULAR STRUCTURES: NUMERICAL RESULTS
In the present section we will study numerically four different infinite 2D 3-color class equivalent regular structures of $`\mathrm{𝖯𝟨𝗆𝗆}(𝖫)|\mathrm{𝖯𝟨𝗆𝗆}(𝖫^{^{}})`$ and $`\mathrm{𝖯𝟨}(𝖫)|\mathrm{𝖯𝟨}(𝖫^{^{}})`$ types. In order to avoid cumbersome notations for these structures, we will use the following symbols: He (honeycomb) (Fig.1a) ; Fl (flower) (Fig.1b) ; Co (cogrose) (Fig.1c) ; Rh (rhombus) (Fig.1d).
Due to the homogeneity (3) of Eq. (12) one can rescale $`\sigma _e(\sigma _1,\sigma _2,\sigma _3)`$ defining a new function $`\sigma _e^{}(\sigma _2^{},\sigma _3^{})`$ where $`\sigma _i^{}=\sigma _i/\sigma _1,i=e,2,3`$. A typical shape of the surface defined by $`\sigma _e^{}(\sigma _2^{},\sigma _3^{})`$ is shown in Fig. 2.
This surface always contains the following points
$`\sigma _e^{}(0,0)=0,\sigma _e^{}(1,1)=1,\sigma _e^{}(\sigma ,\sqrt{\sigma })=\sqrt{\sigma }.`$
The last equality is related to the Dykhne theorem and generates a loop on the surface where (7) is satisfied. This loop is common for all such surfaces, i.e. they intersect each other on it.
We computed $`\sigma _e^{}(\sigma _2^{},\sigma _3^{})`$ for the structures mentioned above: He, Fl, Co, Rh, and then extracted the parameter $`A`$ corresponding to those structures.
Before discussing these results we describe briefly the numerical algorithm which was used. This algorithm deals with hexagonal Bravais lattices represented as a grid of equilateral triangles. A subdivision procedure produces a set of $`N\times N`$ small similar triangles, the center of each one of them connected with the centers of 3 nearest neighboring cells by simple resistors. The resistor connecting cells $`i`$ and $`j`$ has a resistance equal to $`(\sigma _i+\sigma _j)/(2\sqrt{3}\sigma _i\sigma _j)`$, where $`\sigma _i`$ is the conductivity of the cell $`i`$. Translational symmetry of the composite is reflected by imposing periodic boundary conditions for the currents. The algorithm solves up to 10<sup>6</sup> linear equations arising for the subdivided 3-color elementary cell. Our computational procedure always gives a sequence of effective conductivities $`\sigma _N^{}`$, which converges monotonically (as $`N`$ 1000) from below. This fact was established by treating the solvable case following from Dykhne’s theorem, where $`\sigma _e`$ is known exactly. To get a sequence of upper bounds $`\sigma _N^+`$ we then used the duality property, i.e. simulating the dual problem with conductivities $`1/\sigma _2^{},\mathrm{\hspace{0.33em}1}/\sigma _3^{}`$ and consequently obtaining a monotonic convergence from above for the sequence $`\sigma _N^+`$.
The simulations were done on the 4 characteristic curves Because of the divergence of the calculational procedure at $`\sigma _i^{}=0`$, the first of the characteristic curves was calculated using $`\sigma _2^{}=0.001`$, for which the results can still be trusted. on the surface where it intersects with the planes:
$$1)\sigma _2^{}=0.001,\mathrm{\hspace{0.33em}\hspace{0.33em}2})\sigma _2^{}=1,\mathrm{\hspace{0.33em}\hspace{0.33em}3})\sigma _2^{}=\sigma _3^{},\mathrm{\hspace{0.33em}\hspace{0.33em}4})\sigma _2^{}+\sigma _3^{}=1.$$
(23)
These curves reflect the behavior of the surface relatively well in accordance with the chosen type of composite. From the assumption of the algebraic nature (12) of those curves we were able to extract, for every type of composite, a corresponding parameter $`A`$ by the following procedure. In a wide range of $`A`$, for each calculated point ($`\sigma _i^{},i=2,3`$) on the plane, relative deviations $`ϵ`$ were calculated
$$ϵ(A,\sigma _i^{})=\{\begin{array}{ccc}\frac{\sigma _e(A,\sigma _i^{})\sigma _N^+}{\sigma _e(A,\sigma _i^{})}& \mathrm{if}& \sigma _N^+<\sigma _e(A,\sigma _i^{})\hfill \\ 0& \mathrm{if}& \sigma _N^{}<\sigma _e(A,\sigma _i^{})<\sigma _N^+\hfill \\ \frac{\sigma _N^{}\sigma _e(A,\sigma _i^{})}{\sigma _e(A,\sigma _i^{})}& \mathrm{if}& \sigma _e(A,\sigma _i^{})<\sigma _N^{}\hfill \end{array}$$
(24)
and a maximal value of those deviations $`ϵ_{max}(A)`$ was determined for every value of $`A`$ by scanning over the entire area 0$`<\sigma _{2,3}^{}<`$1 . We then determined the best value of $`A`$ by minimizing the function $`ϵ_{max}(A)`$. As shown in Fig. 3, $`A`$ for the He and Fl structures is determined by very sharp minima. The other two minima are not as sharp.
As one could expect, the Rh structure has $`A=0`$ since this is the unique value of $`A`$ for which the solution of equation (12) has the features typical of the structures with traps. The values of $`A`$ which minimize $`ϵ_{max}(A)`$ for the other structures are listed in the caption of Fig. 3.
Figs. 47 show the computed results for upper and lower bounds on the relative bulk effective conductivity $`\sigma _e^{}`$, i.e. $`\sigma _N^+(\sigma _2^{},\sigma _3^{})`$ and $`\sigma _N^{}(\sigma _2^{},\sigma _3^{})`$ at maximal $`N`$, for the 4 microstructures of Fig. 1. Note that these upper and lower bounds often appear to coincide due to insufficient resolution in the figures. The axes in the figures are labeled with S.ef for the relative bulk effective conductivity $`\sigma _e^{}`$ and s1, s2, s3 for the component relative conductivities $`\sigma _1^{}=1,\sigma _2^{},\sigma _3^{}`$ respectively.
The pairs of almost merged points shown in Fig. 7 correspond to barely separated upper and lower bounds.
## V EFFECTIVE CONDUCTIVITY OF RANDOM STRUCTURES: EXTENSION OF THE ALGEBRAICITY CONJECTURE
In this chapter we discuss briefly a possible extension of the algebraicity conjecture for 2D three-component structures, which are non-regular but macroscopically homogeneous. On length scales large compared to the inhomogeneities, we can characterize the macroscopic response of such a medium by a single number, the plane effective conductivity $`\sigma _e`$. A reexamination of the basic properties (3-6) of $`\sigma _e(\sigma _1,\sigma _2,\sigma _3)`$ shows that one of them (4) must be discarded, while the other three (homogeneity of 1-st order, duality, and compatibility) continue to be valid. It is important to note that these properties hold irrespective of whether the microstructure of the (isotropic) composite is ordered or disordered.
Instead of Eq. (4), which implies full symmetry under the group $`S_3`$ of all 3-color permutations, we first consider the case where the structure is only symmetric under the group $`C_3`$ of cyclic 3-color permutations: 1 $``$ 2 $``$ 3 $``$ 1. An example of such a microstructure is shown in Fig. 8. (Note the arrangement of “flowers” near the center of that figure.) It follows that $`\sigma _e(\sigma _1,\sigma _2,\sigma _3)`$ is a cyclic permutation invariant function. The group $`C_3`$ is a subgroup of index 2 of the full permutation group $`S_3`$, which characterizes all the regular structures He, Fl, Co, Rh discussed in the previous sections.
Being non-reflecting, the group $`C_3`$ has more basic invariants than the group $`S_3`$ \[see Eq. (11)\], namely
$`I_1(C_3)`$ $`=`$ $`\sigma _1+\sigma _2+\sigma _3,I_2(C_3)=\sigma _1\sigma _2+\sigma _2\sigma _3+\sigma _3\sigma _1,`$ (25)
$`I_3(C_3)`$ $`=`$ $`\sigma _1\sigma _2\sigma _3,I_4(C_3)=(\sigma _1\sigma _2)(\sigma _2\sigma _3)(\sigma _3\sigma _1).`$ (26)
Nevertheless, the additional cubic invariant $`I_4(C_3)`$ cannot be incorporated into a cubic equation for $`\sigma _e`$, because that would violate the duality requirement. Thus, we are lead back to the cubic equation (12), which is dictated not only by the strong requirement (4) of the full permutation invariance $`S_3`$, but even by the milder requirement of cyclic permutation invariance $`C_3`$.
In the case of a random composite, we usually characterize the microstructure by a statistical distribution function of the local conductivity, which can be either $`\sigma _1`$, $`\sigma _2`$, or $`\sigma _3`$ at any point. Such a description results in an ensemble of representative structures, each one with its own form for the function $`\sigma _e(\sigma _1,\sigma _2,\sigma _3)`$. If we assume that the distribution function has the permutation symmetry of either $`S_3`$ or $`C_3`$, then the ensemble average of $`\sigma _e(\sigma _1,\sigma _2,\sigma _3)`$ will have that symmetry too, even though individual samples may violate it. It follows that, the same considerations which led us to stipulate the form (12), for the minimal polynomial equation that $`\sigma _e`$ could satisfy for regular structures, also lead to that same equation for random structures, if the statistical model for those structures is invariant under either $`S_3`$ or $`C_3`$. Numerical tests of this conjecture remain to be performed.
## VI CONCLUSION
In the present paper we have introduced the algebraicity conjecture for the effective conductivity problem of isotropic 2D three-component regular composites. This conjecture is based on the general properties which are satisfied by the effective conductivity $`\sigma _e(\sigma _1,\sigma _2,\sigma _3)`$. The algebraic equation of minimal order for $`\sigma _e`$ is a cubic equation with 1 positive free parameter $`A`$ responsible for the non-universality. This equation satisfies Dykhne’s theorem (7) independently of $`A`$ and has only one positive root, thus avoiding the possibility of multiple physical solutions. The value $`A=1/3`$ corresponds exactly to the Bruggeman effective medium approximation for a 2D composite with 3 equally partitioned components.
We have found support for this conjecture by numerical calculations on four different infinite 2D 3-color class equivalent regular structures of $`\mathrm{𝖯𝟨𝗆𝗆}(𝖫)|\mathrm{𝖯𝟨𝗆𝗆}(𝖫^{^{}})`$ and $`\mathrm{𝖯𝟨}(𝖫)|\mathrm{𝖯𝟨}(𝖫^{^{}})`$ types: He, Fl, Co, Rh. We have established that $`\sigma _e(\sigma _1,\sigma _2,\sigma _3)`$ is a non-universal function with essential dependence on the microstructure even for totally symmetric structures:
1) The cubic equation (12) with $`A`$=11.37 governs the conductivity problem in He structure with a very high precision $`ϵ_{max}10^4`$.
2) There is good agreement ($`ϵ_{max}10^3`$) between the cubic equation with $`A`$=3.76 and numerical results for the Fl structure.
3) In the Co structure the estimated value $`A`$=.305 is near $`1/3`$, which would follow from the Bruggeman effective medium approximation. This may indicate some similarity between the conducting properties of a 3-component random microstructure and those of the ordered Co structure.
4) The Rh structure needs special attention. It belongs to the family of structures with unicolor traps, i.e. with structures where the presence of just one non-conducting ($`\sigma _1=0`$) component is enough to make the composite an insulator ($`\sigma _e=0`$). In percolation theory this corresponds to a threshold $`p_3=1/3`$, in contrast with a 2-color composite where $`p_2=1/2`$. If the cubic equation is valid for this structure with $`A`$=0, then $`\sigma _e=\sqrt[3]{\sigma _1\sigma _2\sigma _3}`$. Unfortunately, these computations are unable to resolve this question for the Rh microstructure.
5) The different microstructures which belong to a common plane symmetry group (He, Rh \- $`\mathrm{𝖯𝟨𝗆𝗆}(𝖫)|\mathrm{𝖯𝟨𝗆𝗆}(𝖫^{^{}})`$, Fl, Co \- $`\mathrm{𝖯𝟨}(𝖫)|\mathrm{𝖯𝟨}(𝖫^{^{}})`$) are characterized by distinct values of $`A`$. This means that this parameter has a topological nature and is sensitive to more than just the symmetry properties of the elementary cell or the plane group of the entire color lattice.
Finally we discussed a possible extension of the algebraicity conjecture to other types of 2D three-component structures. First we showed that even if the microstructure is only invariant under the $`C_3`$ subgroup of the $`S_3`$ permutation group, a cubic equation for $`\sigma _e`$ must still have the form (12). We then extended this conjecture also to the case of random microstructures, described by any statistical model that is invariant under either $`S_3`$ or $`C_3`$. To test this idea, it would be useful to do numerical calculations on such structures. This is left for future investigations.
###### Acknowledgements.
We would like to thank J.L. Birman, A.M. Dykhne, I.M. Khalatnikov, Y.B. Levinson and A. Voronel for helpful discussion. This research was supported in part by grants from the Tel Aviv University Research Authority, from the Gileadi Fellowship program of the Ministry of Absorption of the State of Israel (LGF), and from the Aaron Gutwirth Foundation, Allied Invest. Ltd. (VSM).
## A
We now give a simple proof of the duality relations for a 2D isotropic medium composed of an arbitrary number $`n`$ of isotropic components.
Suppose a 2D medium with a continuous distribution of conductivity $`\sigma (𝐫)`$ is subjected to an average electric field $`𝐄_e`$. The system of equations consists of Ohm’ law
$$𝐉(𝐫)=\widehat{\sigma }(𝐫)𝐄(𝐫),$$
(A1)
and the equations for local fields
$$\times 𝐄(𝐫)=0,𝐉(𝐫)=0$$
(A2)
with appropriate boundary conditions on the electrical potential.
We are interested in the relation between the current $`𝐉_e`$ averaged over the system, $`𝐉_e=S^1𝐉(𝐫)𝑑S`$ and the averaged field $`𝐄_e=S^1𝐄(𝐫)𝑑S`$. By virtue of the linearity of (A1, A2) this relation will also be linear
$$𝐉_e=\widehat{\sigma }_e𝐄_e,\widehat{\sigma }_e=\widehat{\sigma }_e\{\widehat{\sigma }(𝐫)\}.$$
(A3)
where the tensor of the bulk effective conductivity $`\widehat{\sigma }_e=\widehat{\sigma }_e\{\widehat{\sigma }(𝐫)\}`$ is actually a tensorial functional. In the case of $`n`$ homogeneous anisotropic components, the latter becomes a tensorial function $`\widehat{\sigma }_e(\widehat{\sigma }_1,\widehat{\sigma }_2,\mathrm{},\widehat{\sigma }_n)`$. Further simplification arises when all components are isotropic — $`\widehat{\sigma }_e(\sigma _1,\sigma _2,\mathrm{},\sigma _n)`$ and, finally when the entire composite is also an isotropic medium — $`\sigma _e(\sigma _1,\sigma _2,\mathrm{},\sigma _n)`$.
In order to transform to the dual problem we rotate the $`x,y`$-components of $`𝐉`$ and $`𝐄`$ by $`90^o`$ in the plane
$`𝐉^{}(𝐫)=\widehat{R}𝐄(𝐫),𝐄^{}(𝐫)=\widehat{R}𝐉(𝐫),𝐉^{}(𝐫)=\widehat{\sigma }^{}(𝐫)𝐄^{}(𝐫),`$ (A4)
$`\widehat{\sigma }^{}(𝐫)=\widehat{R}\widehat{\sigma }^1(𝐫)\widehat{R}^1,\widehat{R}=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right).`$ (A7)
Eqs. A2 are thereby transformed as follows
$`𝐉^{}(𝐫)=\times 𝐄(𝐫)=0,\times 𝐄^{}(𝐫)=𝐉(𝐫)=0,`$
while $`𝐉_e^{}`$ and $`𝐄_e^{}`$ are connected by
$$𝐉_e^{}=\widehat{\sigma }_e^{}𝐄_e^{},\widehat{\sigma }_e^{}=\widehat{R}\widehat{\sigma }_e^1\widehat{R}^1,$$
(A8)
where the bulk effective conductivity tensor of the dual problem $`\widehat{\sigma }_e^{}`$ is defined in accordance with (A3)
$$\widehat{\sigma }_e^{}=\widehat{\sigma }_e\{\widehat{\sigma }^{}(𝐫)\}=\widehat{\sigma }_e\{\widehat{R}\widehat{\sigma }^1(𝐫)\widehat{R}^1\}.$$
(A9)
The last two equations (A8,A9) give a duality relation
$$\widehat{\sigma }_e\{\widehat{R}\widehat{\sigma }^1(𝐫)\widehat{R}^1\}\widehat{R}\widehat{\sigma }_e\{\widehat{\sigma }(𝐫)\}\widehat{R}^1=\widehat{I},$$
(A10)
where $`\widehat{I}`$ is the unit matrix. If, instead of the continuous fields $`\sigma (𝐫)`$ we deal with anisotropic media composed of $`n`$ homogeneous anisotropic components, then
$`\widehat{\sigma }_e(\widehat{R}\widehat{\sigma }_1^1\widehat{R}^1,\widehat{R}\widehat{\sigma }_2^1\widehat{R}^1,\mathrm{},\widehat{R}\widehat{\sigma }_n^1\widehat{R}^1)\widehat{R}`$ (A11)
$`\widehat{\sigma }_e(\widehat{\sigma }_1,\widehat{\sigma }_2,\mathrm{},\widehat{\sigma }_n)\widehat{R}^1=\widehat{I}.`$ (A12)
In the case of an anisotropic composite with $`n`$ homogeneous isotropic components we will have
$$\widehat{\sigma }_e(\sigma _1^1,\sigma _2^1,\mathrm{},\sigma _n^1)\widehat{R}\widehat{\sigma }_e(\sigma _1,\sigma _2,\mathrm{},\sigma _n)\widehat{R}^1=\widehat{I}.$$
(A13)
The principal values of $`\widehat{\sigma }_e`$ satisfy Keller’s theorem
$`\widehat{\sigma }_e^{xx}(\sigma _1^1,\sigma _2^1,\mathrm{},\sigma _n^1)\widehat{\sigma }_e^{yy}(\sigma _1,\sigma _2,\mathrm{},\sigma _n)=1,`$ (A14)
$`\widehat{\sigma }_e^{yy}(\sigma _1^1,\sigma _2^1,\mathrm{},\sigma _n^1)\widehat{\sigma }_e^{xx}(\sigma _1,\sigma _2,\mathrm{},\sigma _n)=1.`$ (A15)
In general, the directions of the principal axes depend on the values of $`\sigma _1,\sigma _2,`$ etc. But when the symmetry of the microstructure is sufficiently high, those directions will be fixed by that symmetry. Finally, in the case of an isotropic composite, the last equations reduce to the self-duality relation
$$\sigma _e(\sigma _1^1,\sigma _2^1,\mathrm{},\sigma _n^1)\sigma _e(\sigma _1,\sigma _2,\mathrm{},\sigma _n)=1.$$
(A16)
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# Kondo tunneling through real and artificial molecules
## Abstract
When a cerocene molecule is chemisorbed on metallic substrate, or when an asymmetric double dot is hybridized with itinerant electrons, its singlet ground state crosses its lowly excited triplet state, leading to a competition between the Zhang-Rice mechanism of singlet-triplet splitting in a confined cluster and the Kondo effect (which accompanies the tunneling through quantum dot under a Coulomb blockade restriction). The rich physics of an underscreened $`S=1`$ Kondo impurity in the presence of low-lying triplet/singlet excitations is exposed. Estimates of the magnetic susceptibility and the electric conductance are presented.
1. The problem of tunneling through a sequence of resonance levels was formulated about three decades ago . At that stage, little attention was paid to the specific structure of the tunnel barrier. Nowadays, novel experimental techniques enable the fabrication of artificial objects which carry some of the salient features of complex quantum systems existing in Nature, and to include them as resonance barriers in electron tunneling devices. They manifest numerous unusual properties and might be regarded as important ingredients in future microelectronics . Examples are double quantum dot structures , atomic and molecular wires and bridges.
At the same time, new methods of tunneling microscopy make it possible to elucidate the properties of single atoms and molecules adsorbed on a surface. The combination ”nanotip – atom/molecule – substrate” is then a quantum system with exceptional resonance features and potential applications .
In the present work we expose the physics of tunneling through real and artificial molecules in which there is presumably a singlet ground state with an even number of electrons, which are spatially separated into two groups with different degree of localization. Electrons in the first group are responsible for strong correlation effects (Coulomb blockade), whereas those in the second group are coupled to a metallic reservoir. Hybridization with itinerant electrons result in transformation of the nonmagnetic (singlet) ground state into a magnetic one.
Possible real molecules are lanthanocene molecules Ln(C<sub>8</sub>H<sub>8</sub>)<sub>2</sub> with the ions Ln=Ce, Yb in a cage formed by $`\pi `$-bonded carbon atoms . In these molecules the electrons in a strongly correlated $`f`$-shell are coupled with loosely bound $`\pi `$ electrons. In an analogy with Zhang-Rice (ZR) singlet in Cu-O planes of high-T<sub>c</sub> perovskites , the ground state of this molecule is a spin singlet combination $`{}_{}{}^{1}A_{1g}^{}(f\pi ^3)`$ of an f-electron and $`\pi `$-orbitals, and the energy of the first excited triplet state $`{}_{}{}^{3}E_{2g}^{}`$ is rather small ($`0.5`$ eV). In the ytterbocene (hole counterpart of cerocene) the ground state with one $`f`$-hole is a triplet, and the gap for a singlet excitation is tiny, $`0.1`$ eV. The fullerene-like molecules doped with Ce or Yb form another family with apparently similar properties. In all these systems there is no direct overlap between the strongly correlated $`f`$-electrons and the metallic reservoir. However, these electrons can influence the tunnel properties of the molecule via covalent bonding with the outer $`\pi `$-electrons which, in turn, are coupled to the metallic reservoir.
Artificial candidates are double-dot structures (say $`D_1`$ and $`D_2`$) in tunneling contact with each other, but only $`D_1`$ is coupled with the metallic leads. The respective gate voltages are such that $`V_{g1}<V_{g2}`$. Coulomb blockade then prevents double charging of $`D_2`$, so it can play the same role as $`4f`$ atom in molecular complexes described above. The dot $`D_1`$ donates the loosely bound electrons which contribute to the tunnel current .
The pertinent physics to be exposed below is that of a competition between the ZR mechanism of singlet-triplet splitting in a confined cluster and Kondo effect which accompanies the tunneling through quantum dot under Coulomb blockade confinement . Usually, tunneling through quantum dot containing an even number of electrons does not display a Kondo resonance due to its spin singlet ground state. Analysis of conditions under which the singlet ground state changes into a partially screened spin-one Kondo state due to hybridization with metallic leads is one of the goals of this study.
2. A simple model which describes this type of molecules was considered in , hereafter referred to as a ”Fulde molecule” (FM). It contains two electrons occupying a potential well which is formed by deep and shallow valleys. The Hamiltonian of an isolated FM is
$$H_d=\underset{i}{}\underset{\sigma }{}E_in_{i\sigma }+V\underset{ij}{}d_{i\sigma }^{}d_{j\sigma }+H_{corr}.$$
(1)
Here $`d_{i\sigma }^{}`$ creates a dot electron with spin $`\sigma `$ at valley $`i=f,l`$ and spin $`\sigma `$, while the coupling constant $`V=d_l|V|d_f`$ is the inter-well tunneling integral. There are two electrons in a neutral ground state, and $`H_{corr}=Qn_f(n_f1)/2`$ is the interaction term responsible for the Coulomb blockade of charged states (here $`n_f=_\sigma d_{f\sigma }^{}d_{f\sigma }`$). The energy difference $`\mathrm{\Delta }=E_lE_f`$ is postulated to exceed the overlap integral, $`\beta =V/\mathrm{\Delta }1.`$ Two-electron states $`|\mathrm{\Lambda }`$ of the FM are classified as a ground state singlet $`|S`$, low-lying triplet exciton $`|T0`$, $`|T\pm `$ and high-energy singlet charge-transfer exciton $`|L`$. To order $`\beta ^2`$ they are,
$`|S`$ $``$ $`\alpha ^2|s\sqrt{2}\beta |ex,`$ (2)
$`|T0`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}{\displaystyle \underset{\sigma }{}}d_{f\sigma }^{}d_{l\sigma }^{}|0,|T\pm =d_{l\pm }^{}d_{f\pm }^{}|0,`$ (3)
$`|L`$ $``$ $`\alpha ^2]|ex+\sqrt{2}\beta |s,`$ (4)
where $`|s=\frac{1}{\sqrt{2}}_\sigma \sigma d_{l\sigma }^{}d_{f\sigma }^{}|0,ex=d_l^{}d_l^{}|0,`$ and $`\alpha ^2=1\beta ^2`$. In this order, the energy levels $`E_\mathrm{\Lambda }`$ are :
$$E_S=ϵ_l+ϵ_f2V\beta ,E_T=ϵ_l+ϵ_f,E_L=2(ϵ_l+V\beta ).$$
(5)
The spin and charge branches of excitation spectrum of FM are characterized by rather different energy scales $`E_TE_S=\delta `$ and $`E_LE_S\mathrm{\Delta }`$, respectively. An interplay between Kondo triplet excitations (with some characteristic energy $`\mathrm{\Delta }_K`$) and ZR triplet excitations is expected when $`\delta \mathrm{\Delta }_K`$ in the regime of Kondo resonance induced by tunneling to metallic reservoir .
The tunneling problem is encoded in the Anderson Hamiltonian which incorporates $`H_d`$, together with the band Hamiltonian $`H_b=_{k\sigma }ϵ_kc_{k\sigma }^{}c_{k\sigma }`$ for the electrons in the leads, and the tunneling term $`H_t=_{ik\sigma }W_ic_{k\sigma }^{}d_{i\sigma }`$. Here $`c_{k\sigma }`$ are operators for lead electrons and $`W_{i=l,f}`$ are tunneling matrix elements (assumed to be $`k`$ independent but strongly dependent on the dot valley quantum number $`i`$. It is henceforth assumed that $`W_f=0`$). It is convenient to express the dot operators $`d_{i\sigma }`$ in terms of Hubbard operators, $`X^{\mathrm{\Lambda }\lambda }=|\mathrm{\Lambda }\lambda |`$. Here $`\mathrm{\Lambda }=S,T,L`$ stands for the neutral two-electron states (4), and the index $`\lambda =1\sigma ,3\sigma `$ is reserved for the charged one and three electron states: $`|1\sigma \alpha |f\sigma +\beta |l\sigma ,|3\sigma d_{f\sigma }^{}|ex\frac{V}{Q\mathrm{\Delta }}d_{f\sigma }^{}d_{f\overline{\sigma }}^{}d_{l\sigma }^{}|0.`$ The tunnel matrix elements in the Hubbard representation are given as $`W_\sigma ^{\mathrm{\Lambda }\lambda }=k\sigma ,\lambda |\widehat{W}|\mathrm{\Lambda }`$, where $`\widehat{W}`$ is the operator responsible for tunneling. The Anderson Hamiltonian then reads:
$`H={\displaystyle \underset{\mathrm{\Lambda }}{}}E_\mathrm{\Lambda }X^{\mathrm{\Lambda }\mathrm{\Lambda }}+{\displaystyle \underset{k\sigma }{}}ϵ_kc_{k\sigma }^{}c_{k\sigma }+`$ (6)
$`{\displaystyle \underset{\mathrm{\Lambda }\lambda }{}}\left(W_\sigma ^{\mathrm{\Lambda }\lambda }c_{k\sigma }^{}X^{\lambda \mathrm{\Lambda }}+\overline{W}_\sigma ^{\mathrm{\Lambda }\lambda }X^{\mathrm{\Lambda }\lambda }c_{k\sigma }\right).`$ (7)
Using the Wigner-Eckart theorem, one can write $`W_\sigma ^{\mathrm{\Lambda }\lambda }=C_{\sigma \lambda }^\mathrm{\Lambda }A_\lambda ,`$ where $`C_{\sigma \lambda }^\mathrm{\Lambda }`$ are Clebsh-Gordan coefficients and $`A_\lambda `$ is the reduced matrix element. In a given vector-coupling scheme the tunneling results in the following transitions: $`|S,|T0|1\sigma ,p\overline{\sigma }`$; $`|S,|T0|3\sigma ,k\overline{\sigma }`$; $`|T\pm |1\pm ,p\pm `$; $`|T\pm |3\pm ,k`$. Here $`p\sigma `$ and $`k\sigma `$ are, respectively, the states with an excess electron (and hole) above (below) the Fermi level of the lead. Let us focus on the case where the Coulomb blockade eliminates the three electron states $`|3\sigma `$ and consider the tunnel coupling involving only the states $`|1\sigma `$. The non-zero tunnel matrix elements are
$$W_\pm ^{T\pm }=W,W_{}^{T0}=\frac{1}{\sqrt{2}}W,W_{}^S=\pm \frac{\alpha ^2}{\sqrt{2}}W,$$
(8)
where $`W=\alpha W_l`$. The energy costs of these transitions are
$`E_{1p,S}`$ $`=`$ $`ϵ_pϵ_l+\beta V,E_{1p,T}=ϵ_pϵ_l,`$ (9)
$`E_{3k,S}`$ $`=`$ $`ϵ_l+4\beta V+\stackrel{~}{Q}ϵ_k,E_{3k,T}=ϵ_l+2\beta V+\stackrel{~}{Q}ϵ_k,`$ (10)
where $`\stackrel{~}{Q}Q[V^2/(Q\mathrm{\Delta })^2].`$ 3. We study the interplay between the singlet and triplet levels of the double quantum dot by the renormalization group method following the general line of ”poor man’s scaling” approach to the Anderson model . The renormalized levels $`\stackrel{~}{E}_\mathrm{\Lambda }`$ are determined by the equations
$$d\stackrel{~}{E}_\mathrm{\Lambda }/d\mathrm{ln}D=\mathrm{\Gamma }_\mathrm{\Lambda }/\pi .$$
(11)
Here $`\mathrm{\Gamma }_\mathrm{\Lambda }`$ are the tunnel coupling constants,
$$\mathrm{\Gamma }_T=\mathrm{\Gamma }\pi \rho _0W,\mathrm{\Gamma }_S=\alpha ^2\mathrm{\Gamma }_T,\rho _0D^1.$$
(12)
Integrating $`(\text{11})`$ under the conditions $`\stackrel{~}{E}_\mathrm{\Lambda }(D_0)=E_\mathrm{\Lambda },\mathrm{\Gamma }_\mathrm{\Lambda }=const`$, we find the scaling invariants $`E_\mathrm{\Lambda }^{}`$ which determine the scaling trajectories
$$E_\mathrm{\Lambda }^{}=E_\mathrm{\Lambda }\frac{\mathrm{\Gamma }_\mathrm{\Lambda }}{\pi }\mathrm{ln}\left(\frac{\pi D}{\mathrm{\Gamma }_\mathrm{\Lambda }}\right)$$
(13)
The level $`ϵ_f`$ is taken to be close to the bottom of the conduction band , so that scaling does not significantly affect it. It is then subtracted from the energies $`E_T`$ and $`E_S`$. Now we see that the energies $`E_\mathrm{\Lambda }`$ decrease together with $`D`$. Since $`\mathrm{\Gamma }_T>\mathrm{\Gamma }_S`$, the phase trajectory $`E_T(D,\mathrm{\Gamma }_T)`$ should cross that of $`E_S(D,\mathrm{\Gamma }_S)`$ at a certain point. Thus, quite remarkably, there is a crossover from singlet to triplet ground state of the FM due to tunnel contact with metallic leads. The crossing point can be estimated from eqs. (13): $`\pi \delta ^{}(\mathrm{\Gamma }_T\mathrm{\Gamma }_S)\mathrm{ln}(\pi D/\mathrm{\Gamma }),`$ or, referring to the bare parameters, the value $`\stackrel{~}{D}`$ of renormalized bandwidth corresponding to this crossing point is $`\stackrel{~}{D}=D\mathrm{exp}\left(\pi \mathrm{\Delta }/\mathrm{\Gamma }\right).`$ Another important crossing point is the energy $`\overline{D}=a\overline{ϵ}_l`$ ($`a1`$) where the scaling of $`ϵ_l`$ also stops, the charge fluctuations become irrelevant, and one reaches the Schrieffer-Wolff limit where only spin fluctuations are responsible for scaling of the model Hamiltonian . This energy is determined by the equation
$$\overline{D}=(\mathrm{\Gamma }/\pi )\mathrm{exp}\left(\pi (|\overline{ϵ_l}ϵ_l|)/\mathrm{\Gamma }\right).$$
(14)
If $`\stackrel{~}{D}>\overline{D}`$, the Schrieffer-Wolff regime is reached after crossover from singlet to triplet ground state of the FM, and a Kondo type resonant tunneling is feasible. In the opposite case there is a singlet ground state and a soft triplet exciton (see figure 1).
And yet, the $`S=1/2`$ Kondo regime is still accessible once a properly tuned external magnetic field is applied . The novel feature here is that $`ST`$ crossover can be induced by an upward shift of the dot level $`ϵ_l`$ relative to $`\epsilon _F`$ by a suitable gate voltage (Fig. 1, inset).
We focus on the physically richer case of triplet solution $`\overline{\delta }<0`$. For a two-electron FM the Schrieffer-Wolff transformation projects out the states $`|\lambda ,k\sigma `$ and maps the Hamiltonian $`H`$ onto an effective Hamiltonian $`\stackrel{~}{H}`$ acting in a subspace of two-electron configurations $`|\mathrm{\Lambda }`$ and reduced conduction band, $`\stackrel{~}{H}=\stackrel{~}{H}^S+\stackrel{~}{H}^T+\stackrel{~}{H}^{ST},`$
$`\stackrel{~}{H}^S`$ $`=`$ $`\stackrel{~}{E}_SX^{SS}+J^S{\displaystyle \underset{\sigma }{}}X^{SS}c_\sigma ^{}c_\sigma `$ (15)
$`\stackrel{~}{H}^T`$ $`=`$ $`\stackrel{~}{E}_T{\displaystyle \underset{\mu }{}}X^{\mu \mu }+J^T𝐒𝐬+{\displaystyle \frac{J_T}{2}}{\displaystyle \underset{\mu \sigma }{}}X^{\mu \mu }c_\sigma ^{}c_\sigma ,`$ (16)
$`\stackrel{~}{H}^{ST}`$ $`=`$ $`J^{ST}\left(𝐏𝐬\right).`$ (17)
$`(\mu =T0,T\pm )`$. The local electron operators are defined as usual $`c_\sigma =_kc_{k\sigma };𝐬=2^{1/2}_{kk^{}}_{\sigma \sigma ^{}}c_{k\sigma }^{}\widehat{\tau }c_{k^{}\sigma ^{}};`$ $`\widehat{\tau }`$ are the Pauli matrices. The singlet and triplet states are now intermixed, and the spin properties of FM are characterized by the vector operators $`𝐒`$ and $`𝐏`$ in accordance with the dynamical symmetry of spin rotator:
$`S^+`$ $`=`$ $`\sqrt{2}\left(X^{+0}+X^0\right),S^{}=\sqrt{2}\left(X^{0+}+X^0\right),`$ (18)
$`S^z`$ $`=`$ $`X^{++}X^{},P_z=\left(X^{0S}+X^{S0}\right),`$ (19)
$`P^+`$ $`=`$ $`\sqrt{2}\left(X^{+S}X^S\right),P^{}=\sqrt{2}\left(X^{S+}X^S\right),`$ (20)
These operators obey the moment algebra ($`i=x,y,z):`$
$$[P^i,P^j]=i\epsilon _{ijk}S^k,[P^i,S^j]=i\epsilon _{ijk}P^k,𝐒𝐏=0.$$
(21)
and the Casimir operator is $`S^2+P^2=3.`$ Surprisingly, this special representation of $`O(4)`$ played an important role in particle physics many years ago . The effective exchange integrals are
$`J^T={\displaystyle \frac{2|W_l|^2}{\mathrm{\Delta }_T}},J^S`$ $`=`$ $`{\displaystyle \frac{\alpha ^2|W_l|^2}{\mathrm{\Delta }_S}},J^{ST}={\displaystyle \frac{\alpha |W_l|^2}{\sqrt{2}\overline{\mathrm{\Delta }}}},`$ (22)
$`\overline{\mathrm{\Delta }}^1`$ $`=`$ $`\mathrm{\Delta }_T^1+\mathrm{\Delta }_S^1,`$ (23)
in which $`ϵ_k`$ is replaced by $`ϵ_F`$ in the denominators, that is, $`\mathrm{\Delta }_\mathrm{\Lambda }=ϵ_Fϵ_\mathrm{\Lambda }(\overline{D})`$, and $`ϵ_\mathrm{\Lambda }(\overline{D})`$ are the positions of the scaled level $`ϵ_l`$ on the flow diagram of Fig. 1. Thus, the pertinent physics is that of an underscreened Kondo impurity in the presence of potential scattering and low-lying triplet/singlet excitations. A similar model was considered recently in Ref. studying the physics of tunneling through a vertical quantum dot in magnetic field . In that case, the electron orbital motion in a plane perpendicular to the axis of the dot is characterized by the same quantum number both in the dot and in the leads , and two orbitals participate in the $`ST`$ transitions. The problem can then be mapped onto a special version of the two-impurity Kondo model.
Following we now apply the ”poor man scaling approach” to the Hamiltonian $`\stackrel{~}{H}`$ (17). Neglecting the irrelevant potential scattering phase shift and using the above mentioned cutoff procedure, a system of scaling equations is obtained,(cf. )
$$dj_1/d\mathrm{ln}d=\left[(j_1)^2+(j_2)^2\right],dj_2/d\mathrm{ln}d=2j_1j_2$$
(24)
(here $`j_1=\rho _0J^T,j_2=\rho _0J^{ST},d=\rho _0D`$). The corresponding RG flow diagram has the fixed point $`j_1=\mathrm{}`$, but the resulting Kondo temperature $`T_K(\overline{\delta })`$ turns out to be a sharp function of $`\overline{\delta }`$ . It is maximal when the $`T,S`$ states are quasi degenerate, $`\overline{\delta }T_K(\overline{\delta })`$. The scaling in this case is governed by the effective integral $`j_+=j_1+j_2`$, and the system (24) is reduced to a single equation
$$dj_+/d\mathrm{ln}d=(j_+)^2$$
(25)
with $`T_{K0}=\overline{D}\mathrm{exp}(1/j_+)`$. In the opposite limit $`\overline{\delta }T_K(\overline{\delta })`$ the scaling of $`J^{ST}`$ stops at $`D\overline{\delta }`$. Then $`j_{1,2}(\overline{\delta })=j_{1,2}\mathrm{ln}^1\left(\frac{\overline{\delta }}{T_{k0}}\right)`$ and $`T_K(\overline{\delta })=\overline{\delta }\mathrm{exp}\left[1/j_1(\overline{\delta })\right]T_{K0}`$. The singlet ground state $`S`$ with zero $`T_K`$ is realized when $`\overline{\delta }<0,|\overline{\delta }|>T_K(\overline{\delta })`$.
4. The salient features of FM stem from the qualitative dependence of its ground state and low-energy spectrum on the coupling constants $`V`$ and $`W_l`$. The unusual singlet-triplet crossing should show up in the magnetic properties of adsorbed molecules and tunnel transparency of asymmetric double quantum dots.
According to quantum chemical calculations of the energy spectrum of isolated cerocene molecule, the Van Vleck paramagnetic contribution of $`ST`$ excitations is too weak to overcome the Larmor diamagnetic contribution of C<sub>8</sub>H<sub>8</sub> rings . This situation can drastically change for a FM adsorbed on a metallic layer. The fixed point $`j_1=\mathrm{}`$ corresponds to the scattering phases $`\eta _\sigma (ϵ_F)=\pi /2.`$ In the case of adsorbed FM this means that the molecule has a residual spin 1/2 which interacts ferromagnetically with the conduction electrons . The temperature dependence of magnetic susceptibility $`\chi (T)`$ is predetermined by the energy parameters $`\overline{\delta }`$ and $`T_K(\overline{\delta })`$. In particular, $`\chi (T)`$ conserves its Curie-like character down to the lowest temperatures when $`\overline{\delta }<0`$, $`|\overline{\delta }|T_K`$ Then at $`TT_K`$ the underscreened FM remains paramagnetic, and its susceptibility is
$$\chi (T)=\chi _0(T)[1Z(T/T_K)]$$
(26)
Here $`\chi _0=3C/4T`$, $`C=(g\mu _b)^2`$, and $`Z(x)`$ is the invariant coupling function (solution of the Gell-Mann – Low equation, see ). The triplet spin state is restored at $`T>T_K`$ . In this regime the Kondo corrections as well as admixture of singlet state can be calculated by perturbation theory, with the result,
$`\chi (T)`$ $`=`$ $`{\displaystyle \frac{2C[3\mathrm{exp}(\overline{\delta }/T)]}{3T}}\left(11/\mathrm{ln}{\displaystyle \frac{T}{T_K(\delta )}}j_2\mathrm{ln}{\displaystyle \frac{\overline{D}}{\overline{\delta }}}\right),`$ (27)
$`\chi (T)`$ $`=`$ $`{\displaystyle \frac{2C}{T[3+\mathrm{exp}(\overline{\delta }/T)]}}\left(11/\mathrm{ln}{\displaystyle \frac{T}{T_{K0}}}\right),`$ (28)
respectively for $`\overline{\delta }T_{K0}`$ and $`\overline{\delta }T_{K0}`$. In the case of artificial FM the resonance scattering phase means perfect tunneling transparency of the quantum dot at $`T=0`$ and a logarithmic fall off at high temperatures. To calculate the tunneling transparency of FM sandwiched between two leads, one should add an index $`n=L,R`$ to the operator $`c_{nk\sigma }`$ and switch to the standing wave basis $`\sqrt{2}c_{k\sigma \pm }=c_{Lk\sigma }\pm c_{Rk\sigma }`$ (in a symmetric configuration $`W_{iL}=W_{iR}`$). Then only the wave (+) is involved in tunneling, and the zero bias anomaly in the differential conductance $`G(T)`$ (due to Kondo cotunneling) in the weak coupling regime $`T>T_K`$ is found as in (28):
$`G/G_0`$ $`=`$ $`2\mathrm{ln}^2[T/T_K(\delta )]+j_1j_2^2\mathrm{ln}(\overline{D}/\overline{\delta }),`$ (29)
$`G/G_0`$ $`=`$ $`3\mathrm{ln}^2[T/T_{K0}]`$ (30)
respectively for the two limiting cases $`\overline{\delta }TT_{K\delta }`$ and $`TT_{K0}\overline{\delta }`$. Here $`G_0=4\pi e^2/\mathrm{}`$. Again the maximum effect is achieved in a nearly degenerate case. At $`T0`$ the conductance tends to the unitarity limit.
In conclusion, the interplay between ZR-type coupling in real and/or artificial molecules and Kondo coupling between molecules and metallic reservoir may result in a crossover from a singlet spin state in a weak-coupling regime to an underscreened $`S=1`$ state at zero $`T`$. The onset of Kondo regime in double quantum dot with even occupation can be driven either by a magnetic field or by a gate voltage.
Acknowledgmet This research is supported by Israeli S.F. grants ”Nonlinear Current Response of Multilevel Quantum Systems” and ”Strongly Correlated Electron Systems in Restricted Geometries”, DIP grant “Quantum Electronics in Low Dimensions” and BSF grant “Dynamical Instabilities in Quantum Dots”. We are grateful to I. Krive for valuable discussions and to I. Kikoin for assistance in numerical calculations.
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# Local properties and Density of States in the two-dimensional p-d Model
## I Introduction
Many unexplained exotic features of transition metal oxides containing CuO<sub>2</sub> planes have inspired the developing of various theoretical methods to solve, at least approximately, the models which describe such compounds.
In the last years, a big effort has been devoted to explain experimental results in terms of Hubbard-like or t-J models and many interesting results have been successfully achieved. The Hubbard model was suggested as low energy derivative of the Emery (or p-d) three band model . The Emery model takes into account the crystalline structure of CuO<sub>2</sub> planes and initially retains the orbitals of $`p`$ and $`d`$ electrons at oxygen and copper sites, the $`d`$ orbital giving the strong Coulomb repulsion when occupied by two electrons. Due to the more detailed structure of the p-d model, one can expect that its solutions can give us a more rich physical picture of cuprates and can explain a wider range of experimental results than solutions of Hubbard-like or t-J models.
It is believed that the strong correlation of electrons in the copper $`d`$ orbitals leads both to the insulating antiferromagnetic state at low density of doping carriers and to superconductivity at higher doping. The doping dependence of the electronic state is not rigid band-like. This suggests that the doped hole changes the state of nearby localized copper $`d`$ electron and reduces the energy of the system if compared to the case of rigid band doping.
A detailed mechanism of binding $`d`$ electron to the doped hole was proposed by Zhang and Rice . As a result of binding, a local singlet forms and it can move through the lattice. Further theoretical investigations of the band dispersion of such singlet in CuO<sub>2</sub> planes gave a successful description of ARPES results .
However, to recover the doping and temperature dependence of micro- and macro-scopic observables, one needs to study not only such singlet excitations, although their energy lies near the Fermi level, but also the wider basis of excitations which are possible in the plane.
Recently, a technique called the Composite Operator Method (COM) has been proposed to describe the local and itinerant properties of strongly correlated systems. A physical model is usually described in terms of elementary particles and some interaction. However, at level of observation the identity of the original particles is lost; the macroscopic behavior of the system is described in terms of new excitation modes. When the interaction is strong the properties of these excitations will be very different from those of the original particles and hardly obtainable by a perturbation scheme. By following this scheme, in the COM a set of field operators is taken as the basis in which the theory is developed. These fields are chosen in order to describe the experimentally observed properties and are called composite fields since they are constructed from the initial set of particles. The properties of these composite fields are dynamically determined by the interaction and by the boundary conditions and must be self consistently calculated. In this process a special attention is put to preserve the symmetry properties of the model . Particularly, the conservation of the Pauli principles plays a fundamental role. As discussed in Ref. , although the method has many points in common with the projection method , the method of equation of motion and the spectral density approach , remarkable differences arise according to the different schemes of self consistency. The method has been applied to the study of several models, like Hubbard , t-J and p-d models.
The insufficiency of studying only those electronic states with energy close to the Fermi level was shown in Ref. . The reduced p-d model, where the transitions to the lower Hubbard level are neglected, was studied within a four-pole expansion for the Green’s function. The fields were taken as the bare $`p`$ electron, the upper Hubbard field, the $`p`$ electron field with spin-flips at copper sites and the $`p`$ electron operator with p-d charge-transfer.
In this work we investigate the full p-d model with a four-component field constituted by $`p`$ electrons, lower and upper Hubbard operators for $`d`$ electrons and a fourth field which describes $`p`$ electrons dressed by spin fluctuations of $`d`$ electrons. Some preliminary results were given in Ref. , where in particular we showed the necessity of introducing the fourth field. The model treated in terms of only the first three fields gives results not in agreement with the numerical simulation data and that do not give a satisfactory description of the experimental situation.
In Sec. II we present the p-d Hamiltonian and introduce the excitation modes represented in terms of composite field operators. The Green’s function for these operators is calculated in the four-pole approximation and depends on a set of parameters expressed as equal-time correlation functions. The study of these parameters is given in Sec. III, where by using the symmetry imposed by the Pauli principle we derive a set of equations which allow us to obtain a complete self-consistent solution.
The developing of numerical methods, as Quantum Monte Carlo (QMC) and exact diagonalization , offers a new challenge to theoreticians. Before looking at the physical implication of the model one should be confident with the solution by calculating some quantities and comparing them with the data of numerical simulations. According to this, in Sec. IV A we make a detailed comparison of our results with the available data obtained by QMC studies. The comparison shows a very good agreement signaling us that the choice of the composite field and the self-consistent procedure adopted are able to give a reasonable solution of the model. We then proceed in sections IV B and IV C to analyze some physical properties. Section V is devoted to concluding remarks. Details of calculations are given in Appendix.
## II The p-d Model within the framework of Composite Operator Method
In this section we consider the p-d model and we present the relevant features of the Composite Operator Method. Starting from the tight-binding model composed of $`d`$ electrons in $`3d_{x^2y^2}`$ copper orbitals and $`p`$ electrons in $`2p_{x,y}`$ oxygen orbitals in a CuO<sub>2</sub> plane and considering only bonding $`p`$ electrons , we keep the terms describing the strong intra-atomic Coulomb repulsion at Cu sites and the $`p`$ and $`d`$ electrons hopping.
Therefore, the original Hamiltonian of the p-d model can be written as
$`H`$ $`=`$ $`{\displaystyle \underset{i}{}}[(\epsilon _d\mu )d^{}(i)d(i)+(\epsilon _p\mu )p^{}(i)p(i)`$ (1)
$`+`$ $`Un_{}(i)n_{}(i)+\mathrm{\hspace{0.17em}2}t(p^\gamma ^{}(i)d(i)+d^{}(i)p^\gamma (i))]`$ (2)
where in the spinor notation $`p_\sigma ^{}(i)`$ and $`d_\sigma ^{}(i)`$ create a $`p`$ electron and a $`d`$ electron on site $`i(𝐑_i,t_i)`$ with spin projection $`\sigma \{,\}`$, respectively. $`\mu `$ is the chemical potential, the $`U`$ term describes the Coulomb repulsion at Cu sites, $`n_\sigma (i)=d_\sigma ^{}(i)d_\sigma (i)`$ is the charge density operator of $`d`$ electrons with spin $`\sigma `$. The $`t`$ term describes the p-d hopping and $`p^\gamma (i)`$ is defined by
$$p^\gamma (i)=\underset{j}{}\gamma _{ij}p(j).$$
(3)
Here $`\gamma _{ij}`$ is given by
$$\gamma _{ij}=\frac{a^2}{(2\pi )^2}_{\mathrm{\Omega }_B}d^2𝐤e^{i𝐤(𝐑_i𝐑_j)}\gamma (𝐤),$$
(4)
where $`\mathrm{\Omega }_B`$ is the volume of the first Brillouin zone, $`\gamma (𝐤)=\sqrt{1\alpha (𝐤)}`$ and where for a two-dimensional quadratic lattice with lattice constant $`a`$
$$\alpha (𝐤)=\frac{1}{2}\left[\mathrm{cos}(k_xa)+\mathrm{cos}(k_ya)\right].$$
(5)
Due to the strong correlation among electrons we introduce the composite field
$$\mathrm{\Psi }(i)\left(\begin{array}{c}p(i)\\ \xi (i)\\ \eta (i)\\ p_s(i)\end{array}\right),$$
(6)
where
$$\xi (i)=\left[1n(i)\right]d(i),\eta (i)=n(i)d(i)$$
(7)
are the Hubbard operators which describe the basic excitations $`n\left(i\right)=0n\left(i\right)=1`$ and $`n\left(i\right)=1n\left(i\right)=2`$ on the lattice site $`i`$, respectively. The fourth field is chosen as
$$p_s(i)=\sigma _kn_k(i)p^\gamma (i)\frac{3c}{I_{22}}\xi (i)\frac{3b}{I_{33}}\eta (i).$$
(8)
The parameters $`b`$ and $`c`$ and the quantities $`I_{22}`$ and $`I_{33}`$ are given in Appendix.
We have introduced the charge- $`\left(\mu =0\right)`$ and spin- $`\left(\mu =1,2,3\right)`$ density operator of $`d`$ electrons
$$n_\mu (i)=d^{}(i)\sigma _\mu d(i),$$
(9)
and we are using the following notation
$$\sigma _\mu (1,\stackrel{}{\sigma }),\sigma ^\mu (1,\stackrel{}{\sigma })$$
(10)
with $`\sigma _k`$ $`(k=1,2,3)`$ being the Pauli matrices.
The choice of the composite field (6) is dictated by the following considerations. The strong intra-atomic Coulomb interaction at Cu sites induces a splitting of the $`d`$ band into the lower and the upper Hubbard subbands. The large covalence between oxygen and copper electrons leads to a large fluctuation of the energy of $`p`$ electrons at O sites. The field $`p_s`$, describing $`p`$ electronic excitations accompanied by the nearest neighbor $`d`$ electron spin fluctuations, represents electronic excitations associated with the Cu-O bonds.
The Heisenberg equations of motion for the composite field $`\mathrm{\Psi }\left(i\right)`$ are
$`i{\displaystyle \frac{}{t}}\mathrm{\Psi }(i)=[\mathrm{\Psi }\left(i\right),H]`$ (11)
$`=\left(\begin{array}{c}(\epsilon _p\mu )p(i)+2t\left[\xi ^\gamma (i)+\eta ^\gamma (i)\right]\hfill \\ (\epsilon _d\mu )\xi (i)+2tp^\gamma (i)+2t\pi (i)\hfill \\ (\epsilon _d\mu +U)\eta (i)2t\pi (i)\hfill \\ (\epsilon _p\mu )p_s(i)+\epsilon _{pp}p^\gamma (i)+\epsilon _{p\xi }\xi (i)\hfill \\ +\epsilon _{p\eta }\eta (i)+t_p\pi (i)+2t\kappa _s(i)\hfill \end{array}\right),`$ (17)
where the coupling constants $`\epsilon _{pp}`$, $`\epsilon _{p\xi }`$, $`\epsilon _{p\eta }`$, $`t_p`$ and the higher order operators $`\pi (i)`$ and $`\kappa _s(i)`$ are defined by the following relations:
$$\epsilon _{pp}=\frac{6tc}{I_{22}},\epsilon _{p\xi }=\frac{3c(\epsilon _p\epsilon _d)}{I_{22}},$$
(18)
$$\epsilon _{p\eta }=\frac{3b(\epsilon _p\epsilon _\eta )}{I_{33}},t_p=6t\left(\frac{b}{I_{33}}\frac{c}{I_{22}}\right),$$
(19)
$$\pi (i)=\frac{1}{2}\sigma ^\mu n_\mu (i)p^\gamma (i)+\xi (i)p^\gamma ^{}(i)\eta (i),$$
(20)
$`\kappa _s(i)=`$ $`\sigma _kd^{}(i)\sigma _kp^\gamma (i)p^\gamma (i)`$ (22)
$`\sigma _kp^\gamma ^{}(i)\sigma _kd(i)p^\gamma (i)+\sigma _kn_k(i)d^{\gamma ^2}(i).`$
Let us introduce the thermal retarded Green’s function
$$S(i,j)R\left[\mathrm{\Psi }(i)\mathrm{\Psi }^{}(j)\right],$$
(23)
where $`R`$ is the usual retarded operator and the bracket $`\mathrm{}`$ denotes the thermal average on the grand canonical ensemble. Use of equation of motion (17) will generate an infinite hierarchy of coupled equations and we need to introduce some approximation. To this purpose, let us split the Heisenberg equation (17) into a linear and nonlinear part
$$i\frac{}{t}\mathrm{\Psi }(i)=\underset{j}{}\epsilon (i,j)\mathrm{\Psi }(j)+\delta J(i),$$
(24)
where the equal-time correlation matrix $`\epsilon (i,j)`$, the so-called energy matrix, is fixed by the requirement that the nonlinear term $`\delta J(i)`$ is orthogonal to the fundamental basis (6):
$$\{\delta J(i),\mathrm{\Psi }^{}(j)\}_{E.T.}=0.$$
(25)
In the framework of the pole approximation we neglect the nonlinear term $`\delta J(i)`$ in the equation of motion. Then, for a translational invariant system the Fourier transform $`S(𝐤,\omega )`$ of the retarded propagator has the following expression
$$S(𝐤,\omega )=\frac{1}{\omega \epsilon (𝐤)}I(𝐤),$$
(26)
where $`\epsilon (𝐤)`$, the Fourier transform of the energy matrix, can be expressed as
$$\epsilon (𝐤)=m(𝐤)I^1(𝐤).$$
(27)
In order to simplify the notation we introduced the so-called normalization matrix $`I(𝐤)`$ and the so-called mass matrix $`m(𝐤)`$
$$I(𝐤)F.T.\{\mathrm{\Psi }(i),\mathrm{\Psi }^{}(j)\}_{E.T.},$$
(28)
$$m(𝐤)F.T.\{i\frac{}{t}\mathrm{\Psi }(i),\mathrm{\Psi }^{}(j)\}_{E.T.},$$
(29)
where the symbol $`F.T.`$ denotes the Fourier transform. Straightforward calculations give the expressions of these matrices. The results for a paramagnetic state are given in Appendix. The Green’s function (26) can be put in the spectral form
$$S(𝐤,\omega )=\underset{n=1}{\overset{4}{}}\frac{\sigma ^{(n)}(𝐤)}{\omega E_n(𝐤)+i\delta }.$$
(30)
The energy-spectra $`E_n(𝐤)`$ are the eigenvalues of the $`\epsilon (𝐤)`$ matrix and the spectral functions $`\sigma ^{(n)}(𝐤)`$ are given by
$$\sigma _{ab}^{(n)}(𝐤)=\underset{c=1}{\overset{4}{}}\mathrm{\Lambda }_{an}(𝐤)\mathrm{\Lambda }_{nc}^1(𝐤)I_{cb}(𝐤),$$
(31)
where $`a,b=1,\mathrm{},4`$ and where the columns of the $`\mathrm{\Lambda }(𝐤)`$ matrix are the eigenvectors of the $`\epsilon (𝐤)`$ matrix.
Summarizing, in our scheme of calculations firstly we choose a fundamental basis of Heisenberg composite operators. The propagator for this basic field is evaluated in a pole approximation, where the incoherent part is neglected. The main ingredient of the calculation is the energy matrix. The knowledge of this function will allow us to calculate the excitations of the system and the spectral functions. The entire process requires a self-consistent procedure which will be discussed in the next Section.
## III Correlation functions and self-consistent equations
In this section we introduce the correlation functions and give the expressions of the self-consistent equations needed to calculate the fermionic propagator.
As shown in Appendix, the Green’s function depends on a set of parameters expressed as static correlation functions of composite operators. Some of these operators belong to the basic set (6) and their expectation values are directly connected to matrix elements of the Green’s function. Other operators are composite fields of higher order, out of the basis (6), and their correlation functions must be evaluated by other means. This aspect is a general property of the Green’s function method . These functions refer to a specific choice of the Hilbert space and one must specify the proper representation. In particular, the states must be constructed in such a way that the relations imposed by the Pauli principle must be satisfied at level of expectation values.
Let us introduce the correlation function
$$C(i,j)=\mathrm{\Psi }(i)\mathrm{\Psi }^{}(j).$$
(32)
By means of the spectral theorem and by recalling Eq. (30), this function has the expression
$`C(i,j)`$ $`=`$ $`{\displaystyle \frac{a^2}{2(2\pi )^2}}{\displaystyle \underset{n=1}{\overset{4}{}}}{\displaystyle _{\mathrm{\Omega }_B}}d^2ke^{i𝐤(𝐑_i𝐑_j)iE_n(𝐤)\left(t_it_j\right)}`$ (34)
$`\times \left[1+T_n(𝐤)\right]\sigma ^{(n)}(𝐤),`$
where
$$T_n(𝐤)=\mathrm{tanh}\left(\frac{E_n(𝐤)}{2k_BT}\right).$$
(35)
The parameters directly connected to the Green’s function are: $`n_d`$, $`a_s`$, $`b`$, $`b_s`$, $`c`$, $`D`$. By means of the definitions given in Appendix, we have the self-consistent equations
$$n_d=2(1C_{22}C_{33}),$$
(36)
$$a_s=C_{14}^\gamma +\frac{3c}{I_{22}}C_{12}^\gamma +\frac{3b}{I_{33}}C_{13}^\gamma ,$$
(37)
$$b=C_{13}^\gamma ,$$
(38)
$`b_s=`$ $`C_{24}^\alpha +C_{34}^\alpha +{\displaystyle \frac{3c}{I_{22}}}\left[C_{22}^\alpha +C_{23}^\alpha \right]`$ (40)
$`+{\displaystyle \frac{3b}{I_{33}}}\left[C_{23}^\alpha +C_{33}^\alpha \right],`$
$$c=C_{12}^\gamma ,$$
(41)
$$D=I_{33}C_{33}.$$
(42)
We are using the following notation
$$C_{\mu \nu }=C_{\mu \nu }(i,i),$$
(43)
$$C_{\mu \nu }^\gamma =\underset{j}{}\gamma _{ij}C_{\mu \nu }(j,i)_{E.T.},$$
(44)
$$C_{\mu \nu }^\alpha =\underset{j}{}\alpha _{ij}C_{\mu \nu }(j,i)_{E.T.}$$
(45)
with
$$\alpha _{ij}=\frac{a^2}{(2\pi )^2}_{\mathrm{\Omega }_B}d^2𝐤e^{i𝐤(𝐑_i𝐑_j)}\alpha (𝐤).$$
(46)
The parameters not directly connected to the Green’s function are $`\mu `$, $`d_s`$, $`f`$, $`\chi _s`$, defined in Appendix. In principle there are different ways to calculate these parameters: decoupling approximation, projection on the basis, use of linearized equations of motion. However, as discussed in Refs. , there is only one definite way to fix them: in order to have the right representation the Green’s function must satisfy the equation
$$\underset{ji^{}}{lim}S(i,j)=\mathrm{\Psi }(i)\mathrm{\Psi }^{^{}}(i),$$
(47)
where the r.h.s. is calculated by means of the Pauli principle.
Use of this equation leads to the following self-consistent equations
$$C_{23}=0,$$
(48)
$$C_{24}=3C_{12}^\gamma \frac{3C_{12}^\gamma C_{22}}{I_{22}},$$
(49)
$$C_{34}=\frac{3C_{13}^\gamma C_{33}}{I_{33}},$$
(50)
$$n_T=42(C_{11}+C_{22}+C_{33}),$$
(51)
where $`n_T=n_p+n_d`$ is the total particle number and $`n_p`$ is the number of $`p`$ electrons. Equations (36)-(42) and (48)-(51) constitute a set of coupled equations which will fix the parameters in a self-consistent way.
Summarizing, we have ten parameters and ten self-consistency equations which allow us to compute the Green’s function and the properties of the system in a fully self-consistent way.
Results of calculations are presented in the next Section. As as a comparison we also present results in the case where we do not take into account the Pauli principle and we express the parameters $`d_s,f,\chi _s`$ by means of the following decoupling equations
$$fC_{11}^{\gamma \gamma }\left[C_{12}^\gamma +C_{13}^\gamma \right],$$
(52)
$$d_s\left[C_{12}^{\gamma \alpha }+C_{13}^{\gamma \alpha }\right]\left[C_{22}^\alpha +2C_{23}^\alpha +C_{33}^\alpha \right],$$
(53)
$$\chi _s2\left[C_{22}^\alpha +2C_{23}^\alpha +C_{33}^\alpha \right]^2.$$
(54)
It should be noted that for certain values of the external parameters some instabilities appear in the iteration procedure and no solution of the self-consistent equations is found. This aspect is not present when the decoupling scheme (52)-(54) is used and may be related to the fact that more asymptotic fields are necessary in order to have a proper representation for the Green’s function .
## IV Results
This Section is organized as follows. The first part is devoted to compare the results of our calculations with the data of numerical analysis by Quantum Monte Carlo and exact diagonalization. As it will be shown in Sec. IV A, for all local quantities the agreement between numerical and COM results is excellent. Once we are confident to have a reasonable solution of the model, we go to the next step where we study some physical properties in order to verify if the p-d model is a realistic model for cuprate superconductors. This is done in the second part, Sections IV B and IV C, where the density of states and Fermi surface are analyzed. In this study the values of the model parameters have been taken according to the results suggested by ab initio calculation : $`t=1eV`$, $`U=6eV`$, $`\mathrm{\Delta }=2eV`$, where $`\mathrm{\Delta }=\epsilon _d+U\epsilon _p`$ is the charge-transfer energy.
In this work all energy are given in units of $`t`$ and measured with respect to the atomic level $`\epsilon _p=0`$.
### A Comparison with QMC data
Here we compare COM results with QMC calculations and Lanczos diagonalization results , , .
We introduce the squared local magnetic moment for $`d`$ electrons as
$$S_z^2=\frac{1}{N}\underset{i}{}\left[n_d(i)n_d(i)\right]^2.$$
(55)
This quantity can be expressed for paramagnetic case through the double occupancy and the number of $`d`$ electrons as follows: $`S_z^2=n_d2D`$.
Following Ref. we can distinguish two different regimes according to the values of the $`U`$ Hubbard repulsion and the charge-transfer energy $`\mathrm{\Delta }`$. In the region $`U>\mathrm{\Delta }`$ the insulating properties of the system are characterized by a charge-transfer gap. Instead, in the region $`U<\mathrm{\Delta }`$ the system is in a Mott-Hubbard regime. Both regions are studied in the following.
In Fig. 1 and Fig. 2 we plot the squared local magnetic moment against the parameters $`\mathrm{\Delta }`$ and $`U`$, respectively. $`S_z^2`$ is an increasing function of of these parameters. Figure 1 shows that the local magnetic moment $`S_z^2`$ takes the smallest value when $`\mathrm{\Delta }`$ approaches zero; in such case $`p`$ and $`d`$ upper Hubbard levels coincide and due to strong mixing the double occupancy takes a large value.
Instead, when $`\mathrm{\Delta }`$ becomes larger than $`U`$ the system changes from charge-transfer insulator to Mott-Hubbard insulator and the local magnetic moment $`S_z^2`$ becomes independent on $`\mathrm{\Delta }`$. In this regime we are addressing a single-band 2D Hubbard model solution.
An increasing behaviour of $`S_z^2`$ is also possible for charge-transfer insulator as it is shown in Fig. 2, where $`\mathrm{\Delta }`$ is taken equal to $`2U/3`$. In such case double occupancy of $`d`$ electrons decreases with increasing distance between bands and the local magnetic moment saturates.
The dependence of the squared local magnetic moment on the total number of particles $`n_T`$ for $`\mathrm{\Delta }=1`$ and $`\mathrm{\Delta }=4`$ and for the values of external parameters $`U=6`$, $`T=0.3`$ is shown in Fig. 3.
We also report results for the case when the parameters $`d_s,f,\chi _s`$ are expressed by means of decoupling equations (52)-(54) (dotted line).
In all cases the agreement with QMC results is very good, specially in the case when we take into account the Pauli principle.
From Fig. 3 we see that the local magnetic moment only slightly changes with doping. Such fact was also observed in neutron scattering experiments on $`La_{2x}Sr_xCuO_4`$ , where it was stated that doping destroys antiferromagnetic spin correlations but does not destroy local magnetic moments.
In a previous paper the transition between $`p`$ level and the lower Hubbard level $`\xi `$ was not taken into account, and the the number of $`d`$ electrons $`n_d`$ was calculated only approximately. As a result the agreement with QMC data was only qualitative.
In Fig. 4 the dependence of $`n_T`$ on the chemical potential is drawn for $`U=6`$, $`\mathrm{\Delta }=1`$ and $`T=0.1`$; a good agreement between COM predictions and QMC results is found. As a comparison, it is also reported the case when the Pauli principle is not taken into account; in this case the agreement with QMC simulations is not so good.
In Fig. 5 and Fig. 6 we present also a very good agreement of our results with Lanczos diagonalization and QMC calculations given in Ref. . The authors of this work write the full number of holes as $`n=n_{Cu}+2n_O`$, where the double number of oxygen holes is due to the two oxygen ions in the cell. We remind that in the p-d model there are two bonding and two nonbonding electrons per site. To make a comparison with results of Ref. we have to take into account the Fermi occupation number of these two unbonding orbitals and to rewrite equation (51) as
$`n_T`$ $`=`$ $`42(C_{11}+C_{22}+C_{33})`$ (56)
$`+`$ $`{\displaystyle \frac{2}{\mathrm{exp}((\epsilon _p\mu )/T)+1}}.`$ (57)
For the values of external parameters given in Fig. 5 and Fig. 6 the system is in the charge-transfer regime and in the vicinity of $`n_T=1`$ holes prefer to go to $`p`$ orbitals and electrons to $`d`$ orbitals.
In Fig. 7 it is shown the dependence of the charge-transfer susceptibility at $`𝐪=0`$ on charge-transfer energy $`\mathrm{\Delta }`$. This quantity is given by
$$\underset{𝐪0}{lim}\chi _{CT}(𝐪)=\frac{}{\mathrm{\Delta }}n_d2n_p.$$
(58)
The agreement between COM predictions and QMC simulations is very good.
In Fig. 8, Fig. 9 and Fig. 10 we plot the quantities $`n_d`$, $`S_z^2`$ and $`D`$ against the charge-transfer gap $`\mathrm{\Delta }`$ for the values of external parameters $`U=6`$, $`n_T=5`$ and $`T=1/8`$. The agreement between COM prediction and QMC simulation is excellent. The dependence of the calculated quantities on external parameters has the same behaviour as in Fig. 1. We want to note once more that the holes introduced in CuO<sub>2</sub> planes reside primarily on oxygen sites, as it is well-known experimentally.
In Fig. 11 and Fig. 12 the band dispersion of Zhang-Rice singlet for $`n_T=2.5`$ and $`n_T=2.75`$ is shown. The squares are QMC data from and the solid line is the COM result. The data differ significantly only for the values of dispersion energies far from the chemical potential.
Summarizing, in this section we have presented a detailed comparison of COM results with the available data by numerical simulation. It is worth noticing that the formulation is fully self-consistent and no adjustable parameter is used.
### B Density of states
In this paragraph and in the next one we present an analysis of the band structures of the model and we give a description of the relevant physics near the Fermi level.
The density of states (DOS) for $`p`$ and $`d`$ electrons is respectively given by the following expressions
$$N_p(\omega )=\frac{a^2}{(2\pi )^2}\underset{\mathrm{\Omega }_B}{}d^2𝐤\underset{n=1}{\overset{4}{}}\sigma _{11}^{(n)}(𝐤)\delta \left(\omega E_n(𝐤)\right),$$
(59)
$`N_d(\omega )=`$ $`{\displaystyle \frac{a^2}{(2\pi )^2}}{\displaystyle \underset{\mathrm{\Omega }_B}{}}d^2𝐤{\displaystyle \underset{n=1}{\overset{4}{}}}[\sigma _{22}^{(n)}(𝐤)+2\sigma _{23}^{(n)}(𝐤)`$ (61)
$`+\sigma _{33}^{(n)}(𝐤)]\delta (\omega E_n(𝐤)),`$
where the spectral functions $`\sigma _{\alpha \beta }^{(n)}(𝐤)`$ and the energy bands $`E_n(𝐤)`$ can be calculated as indicated in Sec. II.
In Figs. 13, 14 and 15 we present the DOS of $`d`$ electrons (solid line) and $`p`$ electrons (dotted lines) for $`U=6`$, $`\mathrm{\Delta }=2`$, $`T=0.001`$ and for $`n_T=2.70`$, $`2.85`$ and $`2.90`$, respectively. The solid vertical line at $`\omega =0`$ indicates the chemical potential.
As it can be seen from these pictures there are four bands. The lower Hubbard band located at $`\omega 9`$ represents the $`\xi `$ excitations; this band is filled mainly by $`d`$ electrons and it has a small fraction of $`p`$ electrons. The upper Hubbard band, located at $`\omega 3`$, comes from $`\eta `$ operator excitations; it has mainly the weight of $`d`$ electrons. The band coming from the $`\epsilon _p`$ atomic level is located around $`\omega 3`$. It is filled by $`p`$ electrons, but due to strong mixing it also contains a large fraction of $`d`$ electrons. Finally, the Zhang-Rice singlet band is situated around the chemical potential; this band has almost equal fraction of $`p`$ and $`d`$ electrons.
We want to stress the importance of studying the Zhang-Rice excitation as independent field in the initial set (6). As shown in Ref. , if one does not initially consider the $`p_s`$ excitation, the information about $`p`$ and $`d`$ electrons coupling is lost after taking averages (28), (29) and one does not get the band dispersion situated at the chemical potential when the total number of particles is $`n_T3`$.
The singlet coupling of $`p`$ and $`d`$ electrons was investigated by means of COM in the earlier work of Matsumoto et al. together with one of us . They obtained the singlet excitation band at the chemical potential and studied how the density of states depends on the variation of the correlators.
In this work we considered the full p-d model, as described by Eq. (1), and we used additional equations supplied by Pauli principle (Eqs. (48) -(50)) to calculate the values of these correlators.
### C An analysis of the experimental data
We investigated the region of total number of electrons around $`n_T=3`$, which corresponds to the initial situation where only one electron is in the copper $`d`$ orbital due to the strong Hubbard repulsion and two electrons in the bounded $`p`$ orbitals of oxygen.
In Fig. 16 the band structure for the Zhang-Rice singlet is drawn for the values of external parameters given in the pictures; we see that at the momentum $`𝐤=(\pi ,0)`$ such band has a saddle-point. This saddle-point leads to the van-Hove singularity in the density of states (see Fig. 13, Fig. 14 and Fig. 15) and it was observed in experiments .
We note that for $`U=6`$ and $`\mathrm{\Delta }=2`$ the coincidence of the chemical potential and van-Hove singularity takes place at the value of total number of particles $`n_T=2.85`$, corresponding to the hole doping $`\delta =0.15`$. This results to an enhancement of the thermodynamical properties such as specific heat and spin magnetic susceptibility, as observed in the two-dimensional Hubbard model . Upon increasing the particle number up to $`n_T=2.9`$ the gap between the upper Hubbard band and the Zhang-Rice singlet band disappears.
In Fig. 17 and Fig. 18 we plot the Fermi Surface for the values of the external parameters $`U=6`$, $`\mathrm{\Delta }=2`$ and $`T=0.001`$ and for total number of particles $`n_T=2.7`$ and $`n_T=2.9`$ respectively. The values of the temperature and of the total number of particles were chosen in order to compare the Fermi surface calculated by means of COM with the experimental one measured by Ino et al and shown by the circles in the pictures.
We see that at the total number of particles $`n_T=2.85`$, when the chemical potential crosses the van Hove singularity, the Fermi surface changes its shape from the electron-like in the overdoped regime ($`n_T=2.7`$) to the hole-like in the underdoped regime ($`n_T=2.9`$). This is in agreement with ARPES results and Hall coefficient measurements .
As a comparison, we also plot the Fermi Surface when we do not take into account the Pauli principle (dotted line); the agreement with experimental results is not good in this case.
## V Conclusions
The p-d model was studied in order to describe energetic properties of charge excitations in the CuO<sub>2</sub> planes of cuprates in the normal state. Excitations are described in terms of composite operators in the CuO<sub>2</sub> cluster. The most important excitation appears to be the Zhang-Rice singlet. The oxygen hole doped into the CuO<sub>2</sub> plane couples with Cu electron forming a singlet state. The band of these states is located between the oxygen level and the copper upper Hubbard atomic level. By hole doping the Fermi level crosses this band and at definite value of doping it coincides with a van-Hove singularity in the density of states. This singularity is formed by the saddled shape of the band dispersion at the Brillouin zone points $`(\pi ,0)`$ and $`(0,\pi )`$ as it was also obtained in experiments and in QMC calculations.
The large volume of the Fermi surface at low doping results from the small spectral weight of the field describing $`p`$ electrons dressed by spin-flips of $`d`$ electrons. Such reduction of the spectral weight is possible due to its redistribution to the other bands strongly mixed with Zhang-Rice singlet band.
By using a four-pole approximation in the framework of the Composite Operator Method various local quantities have been calculated as functions of model and physical parameters. The results of calculations have been presented in Sec. IV A and compared with the results of numerical analysis.
We want to stress the relevance of Pauli principle in treating with strongly correlated electronic systems within our method. The comparison of COM results with numerical data, concerning local properties and band structures, and experimental measurements obtained by Ino et al looks to have a good agreement when we take into account the Pauli principle in computing the fermionic propagator.
As last remark, we want to note that we can calculate the two-particle Green’s functions in the one-loop approximation by using the single-particle propagator. Calculations in this direction are now in progress.
###### Acknowledgements.
We wish to thank Dr. A. Avella for fruitful discussions and useful comments.
## The matrices $`I`$ and $`m`$ and the self-consistency parameters.
In this appendix we give the explicit form of matrices $`I`$ and $`m`$ and we also report the expression of the self-consistency parameters.
The $`I`$ matrix is diagonal and its elements are:
$$\begin{array}{ccc}I_{11}=1,& I_{22}=1n_d/2,& I_{33}=n_d/2,\end{array}$$
(62)
$`I_{44}`$ $`=`$ $`3(n_d2D\alpha (𝐤)\chi _s)+4a_s`$ (64)
$`9({\displaystyle \frac{c^2}{I_{22}}}+{\displaystyle \frac{b^2}{I_{33}}}).`$
The quantities $`n_d`$ and $`D`$ are the particle number and the double occupancy per site of $`d`$ electrons, respectively:
$$n_d=d^{}\left(i\right)d\left(i\right),$$
(65)
$$D=n_{}\left(i\right)n_{}\left(i\right).$$
(66)
The other parameters are defined by
$$\begin{array}{cc}b=p^\gamma (i)\eta ^{}(i),& c=p^\gamma (i)\xi ^{}(i),\end{array}$$
(67)
$$\begin{array}{cc}a_s=p^\gamma (i)p_s^{}(i),& \chi _s=\frac{1}{3}n_k(i)n_k^\alpha (i).\end{array}$$
(68)
The $`m`$ matrix is hermitician and its elements have the expressions:
$$\begin{array}{cc}m_{11}=\epsilon _p\mu ,\hfill & m_{12}=2tI_{22}\gamma (𝐤),\hfill \\ m_{13}=2tI_{33}\gamma (𝐤),\hfill & m_{14}=0,\hfill \\ m_{22}=(\epsilon _d\mu )I_{22}+2t(cb),\hfill & m_{23}=2(bc),\hfill \\ m_{24}=2tI_{\pi p_s},\hfill & m_{34}=m_{24},\hfill \end{array}$$
(69)
$$\begin{array}{c}m_{33}=(\epsilon _d+U\mu )I_{33}2t(bc),\hfill \\ m_{44}=(\epsilon _p\mu )I_{44}+2tI_{k_sp_s}+t_pI_{\pi p_s},\hfill \end{array}$$
(70)
where
$$I_{\pi p_s}=\frac{3}{2}(n_d2D\alpha (𝐤)\chi _s)+2a_s\frac{t_p(bc)}{2t},$$
(71)
$$I_{\kappa _s\xi }=\frac{1}{2t}(2tI_{\pi p_s}\epsilon _{p\xi }I_{22}t_p(cb)),$$
(72)
$$I_{\kappa _s\eta }=\frac{1}{2t}(2tI_{\pi p_s}+\epsilon _{p\eta }I_{33}+t_p(bc)),$$
(73)
$`I_{\kappa _sp_s}=`$ $`6(bf+d_s\alpha (𝐤))4b_s`$ (75)
$`3\left(c{\displaystyle \frac{I_{\kappa _s\xi }}{I_{22}}}+b{\displaystyle \frac{I_{\kappa _s\eta }}{I_{33}}}\right).`$
The parameters $`b_s`$, $`f`$, and $`d_s`$ are defined by
$$b_s=d^\alpha (i)p_s^{}(i),$$
(76)
$$f=p^\gamma (i)d^{}(i)p^\gamma (i)p^\gamma ^{}(i),$$
(77)
$$d_s=\frac{1}{3}n_k^\alpha (i)\sigma _kp^\gamma (i)d^{}(i).$$
(78)
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# Continuation of the Fermion-Number Operator and the Puzzle of Families11footnote 1To be presented at DPF 2000 at The Ohio State University, Columbus, OH, August 9–12, 2000. Download from the APS e-print archive at: http://publish.aps.org.eprint/. Paper number: aps2000mar19 003.
## 1.0 Introduction and Background
Given the number of flavors of quarks and leptons, and an appropriate (renormalizable) lagrangian, the so-called “accidental symmetries” of the lagrangian are known to “explain” the separate conservation of various (global) “charges” (e.g., lepton number, baryon number, strangeness, charm, truth, beauty, electron-, muon- and tau-numbers). However, there is nothing in such lagrangians, or their associated accidental symmetries, that would explain quarks and leptons, or tell us how many flavors of quarks and leptons to include. What is needed is a “spectrum-generating” mechanism for these particles.
The current “consensus” in particle physics, at least among string theorists, seems to be that fundamental-fermion flavor degrees-of-freedom arise at the Planck level —or possibly at $`TeV`$ energies in the event that the extra dimensions are “large” —in a theory of superstrings. The purpose of the present paper is to identify an apparently different, but probably complementary, (phenomenological) spectrum-generating mechanism that arises in what seems to be a most unlikely way.
We have found that an “analytic-continuation” of a Hermitian matrix $`𝐅(\mathrm{op})`$ representing the conventional fermion-number *operator*, from an *external* spacetime and Hilbert-space setting to a new *internal* (real) non-Euclidean space—$`𝐅(\mathrm{op})`$ is continued to a real, generally non-diagonal matrix $`𝐅(v)`$ involving a single *real* parameter $`v`$—“automatically” leads to a new description of fundamental fermions (quarks and leptons) in which families are *replicated* and there are just *three* families. The fact that this happens, suggests that there is some deep connection between the result of the aforementioned “analytic-continuation,” namely $`𝐅(v)`$, and Planck-level physics where flavor degrees-of-freedom, and family-replication, presumably originate.
The work presented here is based on a little-known book by the author wherein many of the consequences of using $`𝐅(v)`$ to describe fundamental fermions, are worked out in much greater detail than in the present paper. And, conversely, certain themes that were only alluded to in , such as the “analytic-continuation” $`𝐅(\mathrm{op})𝐅(v)`$, are discussed in much greater detail here. The author hopes that the following relatively brief “overview” and extension of will serve to increase awareness of these new ideas among particle physicists.
## 2.0 The Conventional Fermion-Number Operator
Consider the situation, presumably at some high energy, where we are dealing with “free” (isolated) leptons or “analytically free” quarks. Suppose we want to describe the scalar fermion-number carried by these particles. And, suppose further that the energies involved are not so high that quantum field-theory (QFT) breaks down. Under these conditions the fermion-number can be represented by a $`U(1)`$-type scalar “charge” , namely, a charge associated with the (continuous) group of unitary matrices $`U`$ of order 1 known as $`U(1)`$.
The fermion-number operator, which can be represented by a Hermitian matrix $`𝐅(\mathrm{op})`$, is said to generate these so-called “gauge” (or phase) transformations, which in turn act on fermion and antifermion quantum states in Hilbert space. That is, given that $`\alpha `$ is a *real* phase one has
$$U=e^{i\alpha 𝐅(\mathrm{op})},$$
(1)
and for infinitesimal transformations \[i.e., $`e^{i\delta \alpha 𝐅(\mathrm{op})}=1+i\delta \alpha 𝐅(\mathrm{op})`$\] acting on single-particle (free or “asymptotically free”) fermion and antifermion states $`|p`$ and $`|\overline{p}`$, respectively, one easily establishes that (the fermion-number “charges” are $`f_m=f_a=1`$ for matter and $`f_a`$ for antimatter)
$$\begin{array}{ccc}\hfill U|p& =e^{i\delta \alpha f_m}|p& \\ \hfill U|\overline{p}& =e^{i\delta \alpha f_a}|\overline{p},& \end{array}$$
(2)
since, by definition, $`𝐅(\mathrm{op})`$ obeys the eigenvalue equations
$$\begin{array}{ccc}\hfill 𝐅(\mathrm{op})|p& =f_m|p& \\ \hfill 𝐅(\mathrm{op})|\overline{p}& =f_a|\overline{p}.& \end{array}$$
(3)
Finally, the assumption that the Hamiltonian $`H`$ is invariant under $`U`$, namely
$$H=UHU^{},$$
(4)
ensures that $`H`$ and $`𝐅(\mathrm{op})`$ commute
$$[𝐅(\mathrm{op}),H]=0,$$
(5)
as can be verified by differentiating $`UHU^{}`$ with respect to $`\alpha `$. Hence, the total fermion-number (the number of *fermions* minus the number of *antifermions*) is a constant of the motion.
## 2.1 Matrix representation of the fermion-number operator
Because the matrix $`𝐅(\mathrm{op})`$ involves just *two* kinds of quantum states (3), namely $`|p`$ and $`|\overline{p}`$, it can be expressed as a $`2\times 2`$ diagonal Hermitian matrix (see below), where one of the adjustable parameters $`(\theta )`$ is a *fixed* constant (up to $`2\pi `$) and the other $`(\varphi )`$ is freely *adjustable*. In particular,
$$𝐅(\mathrm{op})=\left(\begin{array}{cc}\mathrm{cos}\theta \hfill & \mathrm{sin}\theta e^{i\varphi }\hfill \\ \mathrm{sin}\theta e^{+i\varphi }\hfill & \mathrm{cos}\theta \hfill \end{array}\right)|_{\mathrm{cos}\theta =1}=\sigma _z,$$
(6)
where
$$\sigma _z=\left(\begin{array}{cc}f_m& 0\\ 0& f_a\end{array}\right)$$
(7)
is one of the familiar Pauli matrices.
This form for $`𝐅(\mathrm{op})`$ is consistent with (3) where the normalization and orthogonality conditions, namely $`p|p=\overline{p}|\overline{p}=1`$ and $`p|\overline{p}=\overline{p}|p=0`$, respectively, directly yield
$$𝐅(\mathrm{op})=\left(\begin{array}{ccc}p|𝐅(\mathrm{op})|p& ,& p|𝐅(\mathrm{op})|\overline{p}\\ \overline{p}|𝐅(\mathrm{op})|p& ,& \overline{p}|𝐅(\mathrm{op})|\overline{p}\end{array}\right)=\sigma _z.$$
(8)
Note that owing to (7) and (8), $`\mathrm{cos}\theta <1`$ in (6) is excluded. Here it should also be noted that $`tr𝐅(\mathrm{op})=f_m+f_a=0`$, $`det𝐅(\mathrm{op})=f_mf_a=1`$, and $`𝐅^2(\mathrm{op})=𝐈_2`$ is the $`2\times 2`$ identity matrix.
## 3.0 The Continuation From F(op) to F($`v`$)
Now perform an “analytic continuation” on $`𝐅(\mathrm{op})`$, namely $`𝐅(\mathrm{op})𝐅(v)`$, which maintains $`𝐅(v)`$ *real* and $`\mathrm{cos}\theta 1`$. This can only be accomplished by continuing $`\theta `$ from a *real* to an *imaginary* number, and by maintaining $`e^{i\varphi }`$ *imaginary*. In particular, to maintain $`𝐅(v)`$ real, we must make the replacements $`\theta iv`$ and $`e^{i\varphi }i`$, where $`v`$ is a *real* number. Then
$`𝐅(v)`$ $`=`$ $`\left(\begin{array}{cc}\mathrm{cos}\theta \hfill & \mathrm{sin}\theta e^{i\varphi }\hfill \\ \mathrm{sin}\theta e^{+i\varphi }\hfill & \mathrm{cos}\theta \hfill \end{array}\right)|_{\theta =iv,e^{i\varphi }=i}`$
or
$`𝐅(v)`$ $`=`$ $`\left(\begin{array}{cc}\mathrm{cosh}v\hfill & \pm \mathrm{sinh}v\hfill \\ \mathrm{sinh}v\hfill & \mathrm{cosh}v\hfill \end{array}\right),`$ (14)
where, just as for $`𝐅(\mathrm{op})`$ the eigenvalues of $`𝐅(v)`$ are $`f_m`$ and $`f_a`$, and so we have $`tr𝐅(v)=f_m+f_a=0`$, $`det𝐅(v)=f_mf_a=1`$, and $`𝐅^2(v)=𝐈_2`$.
In \[6, p. 50 and 54\] it is shown that only the upper signs in (14) have physical significance and $`v0`$. And, just what is the physical significance of $`𝐅(v)`$?
Because the continuation “connects” $`𝐅(\mathrm{op})`$ and $`𝐅(v)`$, it is natural to assume that both $`𝐅(\mathrm{op})`$ and $`𝐅(v)`$ describe, or represent, aspects of the fermion number (i.e., the matter-antimatter “degree-of-freedom”). However, unlike $`𝐅(\mathrm{op})`$, $`𝐅(v)`$ will be shown to describe additional “degrees-of-freedom” such as the “up”-“down” and quark-lepton “degrees-of-freedom.” Moreover, it is abundantly clear from (14) that the the generally non-Hermitian (when $`v0`$) matrix $`𝐅(v)`$—unlike the Hermitian matrix $`𝐅(\mathrm{op})`$ in (8), which acts on Hilbert space—does *not* act on a Hilbert space in an *external* spacetime setting .
## 3.1 A new internal non-Euclidean space
When the matrix $`𝐅(v)`$ acts on a real column-vector $`\{a,b\}`$, it leaves the quadratic form $`a^2b^2`$, invariant. Therefore, the 2-space metric is non-Euclidean or “Lorentzian”, and can be represented by the matrix
$$𝐠=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right).$$
(15)
Given this metric, the *scalar product* of two real vectors assumes the form
$$(a,b)\{e,f\}=aebf.$$
(16)
Similarly, the *square* of a real vector is given by
$$(a,b)\{a,b\}=a^2b^2.$$
(17)
Here $`(,)`$ is a *row* vector while $`\{,\}`$ is a (conformable) *column* vector. Now let us demonstrate that these scalar-products transform like *charges* instead of *probabilities* as they would if we were still dealing with a Hilbert space.
## 3.2 Charge conjugation
The matrix $`𝐗`$, where $`𝐗=𝐗^1`$ or $`𝐗^2=𝐈_2`$ , given by
$$𝐗=\left(\begin{array}{cc}0& 1/d\\ d& 0\end{array}\right),$$
(18)
transforms $`𝐅(v)`$ to is $`𝐂`$-reversed counterpart. In particular, given that the components of $`𝐅(v)`$, in diagonal form \[i.e., $`𝐅(v)_{\mathrm{diag}}=\sigma _z`$\], transform like (global) “charges” ($`f_m`$ and $`f_a`$), the general similarity transformation must be such that
$$𝐗𝐅(v)𝐗^1=𝐅(v),$$
(19)
for *any* real $`d`$. However, since $`𝐗`$ should convert a real vector $`\{a,b\}`$ to its *orthogonal* antimatter-counterpart
$$𝐗\{a,b\}=\{b/d,ad\},$$
(20)
we are forced to require
$$(a,b)\{b/d,ad\}=0,$$
(21)
which means that
$$ab/dabd=0,$$
(22)
or $`d^2=1`$. Therefore, we have two choices for $`d`$, namely, $`d=\pm 1`$. Which one describes the physics of fundamental fermions?
Clearly, when $`d=\pm 1`$, the *square* of any vector, namely,
$$(a,b)\{a,b\}=a^2b^2,$$
(23)
changes signs (is $`𝐂`$-reversing) under $`𝐗`$. However, in general, the components of a vector $`\{a,b\}`$ do not change signs under $`𝐗`$, i.e., they are not $`𝐂`$-reversing.
For example, when $`d=+1`$ in (14), the components of real vectors such as $`\{a,b\}`$ are *not* $`𝐂`$-reversing charge-like quantities, which means that we cannot represent $`𝐂`$-reversing (Lorentz 4-scalar) charges such as electric-charge, charm, isospin or beauty as *components* of such vectors when $`d=+1`$. However, when $`d=1`$, both the square of $`\{a,b\}`$ *and* its components $`a`$ and $`b`$ change signs (are $`𝐂`$-reversing) under $`𝐗`$, and can be used to represent such charges. Therefore, from an experimental standpoint, $`d=+1`$ is excluded (see Ref. 6, p. 8).
We conclude that the matrix $`𝐗`$, which plays the role of charge-conjugation (with respect to these global 2-space charges only) in this non-Euclidean 2-space, is proportional to one of the familiar Pauli matrices, namely,
$$𝐗=\sigma _x.$$
(24)
The foregoing properties of the non-Euclidean charge-like scalars, leads to the following conjecture:
*The global (flavor-defining) charges associated with the aforementioned “accidental” symmetries (see Section 1.0), and the global charges associated with the non-Euclidean 2-space, are (essentially) one and the same charges.*
As shown in the next section if this conjecture is true it means, among other things, that simultaneous fundamental-fermion flavor eigenstates can be partially specified using an appropriate (mutually-commuting) combination of these (global) 2-space charges .
## 4.0 Representing Flavor Eigenstates and Flavor Doublets in the 2-Space
The continuation from the 2-D Hilbert space to the 2-D non-Euclidean “charge” space, turns out to mean that individual flavors of fundamental fermions can be partially represented in an unconventional way by geometric objects (in the non-Euclidean 2-space) which *differ* from a quantum state, but *from which the quantum states can be inferred or effectively constructed*. In particular, in the non-Euclidean 2-space, an object we call a “vector triad” represents “up”-“ down” type flavor doublets of fundamental fermions—the “up”-“down” type flavor dichotomy. That is, the components of the vectors associated with a given vector-triad are *observable* (Lorentz 4-scalar) “charges,” which can be used to (partially) define the two flavor-eigenstates in a flavor doublet .
## 4.1 Flavor eigenstates
As demonstrated in detail in \[6, pp. 16–18\], given the charge-like ($`𝐂`$-reversing) *observables* associated with the description involving $`𝐅(v)`$, namely, the real $`𝐂`$-reversing scalar-components of various matrices and vectors defined on the internal non-Euclidean 2-space, it is possible to write down flavor eigenstates .
What one does is to identify the mutually-commuting $`𝐂`$-reversing “charges” (call them $`C_i`$) or charge-like quantum numbers associated with a particular flavor, and then write the corresponding simultaneous flavor-eigenstate as
$$|C_1,C_2,C_3,\mathrm{},C_n.$$
(25)
Here $`C_1,C_2,C_3,\mathrm{},C_n`$, are the “good” charge-like quantum numbers (charges) associated with a particular flavor. It happens that these *observable* real-numbers can be identified with quantum numbers such as *electric charge, strangeness, charm*, the third-component of (global) *isospin, truth and beauty* (See Ref. 6, p. 72). To discover what charges describe a particular flavor we must first identify the vector-triad associated with that flavor.
Now, each *vector-triad* represents a flavor doublet, not just an individual flavor. That is, vector-triads provide information on *two* quantum states (simultaneous flavor-eigenstates) associated with flavor doublets. Therefore, vector-triads represent both individual flavors *and* flavor doublets. Here we simply summarize how it is the that non-Euclidean vector-triads represent both *individual* flavors and “up”-“down” type *flavor-doublets*.
## 4.2 Flavor doublets
Consider the eigenvector (call it $`𝐐`$) of $`𝐅(v)`$ for fundamental fermions . Since the space on which $`𝐅(v)`$ “acts” is two-dimensional, the *observable* vector $`𝐐`$ can be “resolved” into two (no more or less) *observable*, linearly-independent vectors, call them $`𝐔`$ and $`𝐕`$, as $`𝐐=𝐔+𝐕`$ . Now, because these three vectors ($`𝐐`$, $`𝐔`$, and $`𝐕`$) are *simultaneous* observables, it makes sense to speak of this “triad” of vectors as being a well defined geometric object, namely, a “vector triad.”
Recognizing that the components of $`𝐐`$, $`𝐔`$ and $`𝐕`$ are $`𝐂`$-reversing charge-like *observables* we can write these *observable* “charge” vectors as
$`𝐐`$ $`=`$ $`\{q_1,q_2\}`$ (26)
$`𝐔`$ $`=`$ $`\{u_1,u_2\}`$ (27)
$`𝐕`$ $`=`$ $`\{v_1,v_2\},`$ (28)
where $`q_1`$, $`q_2`$, $`u_1`$, $`u_2`$, $`v_1`$ and $`v_2`$ are the various *observable* “charges” (e.g., $`q_1`$ and $`q_2`$ are found to be *electric* charges). Given $`𝐐=𝐔+𝐕`$, the non-Euclidean metric (11), and Eqs. (22) through (24), we find the associated *observable* quadratic-“charges”
$`𝐐^2`$ $`=`$ $`𝐔^2+2𝐔𝐕+𝐕^2`$ (29)
$`2𝐔𝐕`$ $`=`$ $`2(u_1v_1u_2v_2)`$ (30)
$`𝐔^2`$ $`=`$ $`u_1^2u_2^2`$ (31)
$`𝐕^2`$ $`=`$ $`v_1^2v_2^2.`$ (32)
Finally, using the foregoing charges, we can express the *two* quantum states (simultaneous flavor-eigenstates) associated with a *single* vector-triad in the form of “ket”vectors as follows (Ref. 6, pp. 16–18)
$$\begin{array}{c}|q_1,u_1,v_1,𝐐^2,𝐔^2,2𝐔𝐕,𝐕^2,\hfill \\ |q_2,u_2,v_2,𝐐^2,𝐔^2,2𝐔𝐕,𝐕^2.\hfill \end{array}$$
(33)
Here, the state $`|q_1,u_1,v_1,𝐐^2,𝐔^2,2𝐔𝐕,𝐕^2`$ represents the “up”-type flavor-eigenstate, and $`|q_2,u_2,v_2,𝐐^2,𝐔^2,2𝐔𝐕,𝐕^2`$ represents the corresponding “down”- type flavor-eigenstate in a flavor doublet of fundamental fermions .
Up to this point in the discussion we have shown that the 2-space description naturally incorporates the matter-antimatter dichotomy and the “up”-“down” flavor dichotomy. Now let us demonstrate how quarks and leptons are incorporated in the new description.
## 4.3 Distinguishing quarks and leptons
Choosing the upper signs in (10), the matrix $`𝐅(v)`$ becomes
$$𝐅(v)=\left(\begin{array}{cc}\mathrm{cosh}v& \mathrm{sinh}v\\ \mathrm{sinh}v& \mathrm{cosh}v\end{array}\right),$$
(34)
where $`v`$ is a positive real number \[6, p. 50 and 54\].
As described in \[6, pp. 52–55\], the parameter $`v`$ distinguishes between *quarks* and *leptons*. In particular, the parameter $`v`$ is found to be *quantized* and obeys the “quantum condition”:
$$v=\mathrm{ln}M_c,$$
(35)
where $`M_c`$ counts both the *number* of fundamental fermions in a strongly-bound composite fermion, and the *strong-color multiplicity*. That is, $`M_c=3`$ for quarks (strong-color triplets) and $`M_c=1`$ for leptons (strong-color singlets).
### 4.3.1 Quark and lepton electric charges
It is shown in (Ref. 6, pp. 52–55, and Ref. 12) that the quark and lepton electric charges are the “up”-“down” components of the eigenvectors of the matrix $`𝐅(v)`$. In particular, the quark charges are given by $`(M_c=3)`$
$`q_1(f)`$ $`=`$ $`{\displaystyle \frac{(M_c^21)}{2M_c(M_cf)}}=+{\displaystyle \frac{2}{3}}\text{ for }f=+1\text{ and }+{\displaystyle \frac{1}{3}}\text{ for }f=1`$ (36)
$`q_2(f)`$ $`=`$ $`q_1(f)1.`$ (37)
Similarly, the lepton electric charges are given by $`(M_c=1)`$
$`q_1^{}(f)`$ $`=`$ $`{\displaystyle \frac{(M_c^21)}{2M_c(M_cf)}}=1\text{ for }f=+1\text{ and }0\text{ for }f=1`$ (38)
$`q_2^{}(f)`$ $`=`$ $`q_1^{}(f)+1.`$ (39)
In summary, $`𝐅(v)`$ is found to provide an explanation for the quark-lepton “dichotomy” of fundamental fermions in addition to the matter-antimatter, and “up”-“ down” type flavor-dichotomy.
## 4.4 Family replication and the number of families
In \[6, pp. 59–65\] it is shown that flavor doublets (hence families) are *replicated* and that there are only three families of quarks and leptons. We refer the reader to for a full and detailed account. Here we simply outline how this situation comes about.
By the definition of a linear-vector 2-space, a 2-vector such as $`𝐐`$ can always be resolved into a pair (no more, or less) of linearly-independent vectors $`𝐔`$ and $`𝐕`$ as $`𝐐=𝐔+𝐕`$ (see Sec. 4.2). And, since $`𝐐`$ *represents* a flavor doublet, so should $`𝐔`$ and $`𝐕`$ *represent* this *same* flavor doublet. But, if this is so, *different vector-resolutions of $`𝐐`$ (i.e., different vector-triads) should correspond to different flavor-doublets having the same $`𝐐`$*. In other words, flavor doublets should be *replicated*.
Since $`𝐐`$ can be resolved (mathematically) in an infinite number of ways, we might suppose that there are an *infinite* number of flavor doublets, and hence, families. But, because of various “quantum constraints,” it is possible to show that $`𝐐`$ *can be resolved in only three physically acceptable ways for $`𝐐`$-vectors associated with either quarks or leptons*. In other words, there can be only six quark flavors and six lepton flavors, which leads to the (*ex post facto*)“prediction” of three quark-lepton families.
## 5.0 Discussion
It is important to understand that the new 2-space description of fundamental fermions (quarks and leptons) provides a distinction between these particles that goes beyond differences that can be explained by mass differences alone. For example, in the standard model the only difference between an electron and a muon is that they have different masses. Otherwise, these particles experience identical electroweak interactions. Moreover, as described in Section 1.0, the separate conservation of electron- and muon-numbers can be attributed to certain unavoidable “accidental symmetries” associated with the (renormalizable) lagrangian describing the (electroweak) interactions of these particles.
Taken at face value, these accidental symmetries would seem to imply that there are no internal “wheels and gears” that would distinguish an electron from a muon, for example. But, if the string theories are correct, these particles would be associated with different “handles” on the compactified space \[see Ref. 3, Vol. 2, p. 408\], and so would be different in this *additional* sense. Likewise, in the present non-Euclidean 2-space description, a variety of (global) 2-scalars, which are only *indirectly* related to the accidental symmetries of the lagrangian, serve to provide a (further) distinction between particles such as the electron and muon (see also the conjecture in Sec. 3.2).
A probable experimental signal of such “internal” differences is to be found in the recent observations at the Super Kamiokande of bi-maximal neutrino mixing . Models which begin by positing a neutrino mass-matrix and associated mixing-parameters, such as the three-generation model proposed by Georgi and Glashow , do a good job of describing the observations. However, bi-maximal mixing may have a deeper explanation in terms of internal topological-differences (in the non-Euclidean 2-space) between $`\nu _e`$, and $`\nu _\mu `$ or $`\nu _\tau `$ neutrinos.
With respect to the internal transformation $`𝐅(v)`$, the topology of the non-Euclidean “vector triad” (see Sec. 4.2) representing the $`\nu _e`$ ($`\nu _\mu `$ or $`\nu _\tau `$), is found to be that of a cylinder (Möbius strip). And, assuming that a change in topology during neutrino mixing is suppressed by energy “barriers,” or other topological “barriers” (e.g., one cannot continuously deform a doughnut into a sphere), while neutrino mixing without topology-change is (relatively) enhanced, one can readily explain the experimental observation of (nearly) maximal $`\nu _\mu \nu _\tau `$ neutrino mixing—at least maximal $`\nu _\mu \nu _\tau `$ mixing over long distances, where the foregoing topological influences would be cumulative . If this explanation is basically correct, then it follows that the neutrino mass-matrix and associated mixing-parameters needed to explain bi-maximal neutrino mixing, would be the *result*, at least in part, of these deeper (internal) topological differences between neutrinos, and not their *cause*.
## 6.0 Conclusions
It is widely believed that the explanation for fundamental-fermion (quark and lepton) family replication is to be found in theories of quantum gravity (e.g., superstrings). And, yet, as demonstrated here and elsewhere , a simple “analytic continuation” of a Hermitian matrix representing the fermion-number operator, leads to a new, and unconventional, *internal* description of *quarks* and *leptons*, which also explains family replication. In particular, this description, unlike the conventional standard-model description, is capable of explaining, among other things, the fact that there are just *three* observed families of quarks and leptons. We take these facts to be evidence that the (phenomenological) “analytic continuation” $`𝐅(\mathrm{op})𝐅(v)`$, or at least the result of the continuation $`𝐅(v)`$, somehow reflects physics at the Planck level where flavor degrees-of-freedom presumably originate.
It seems that the best chance to show that $`𝐅(v)`$ and Planck-level physics are related, lies in an appropriate application of superstring theory. Accordingly, the author hopes that the new description of families presented here (see also Refs. 6 and 16) will encourage string theorists working on so-called realistic (free-fermionic) three-generation string models (e.g., see Refs. 17, 18), to take up the challenge of showing that these models either do, or do not, justify the new description.
## 7.0 References and Footnotes
S. Weinberg, *The Quantum Theory of Fields, Vol. I, Foundations*, Cambridge University Press, New York, NY (1995), pp. 529–531; *The Quantum Theory of Fields, Vol. II, Modern Applications*, Cambridge University Press, New York, NY (1996), p. 155.
J. M. Maldacena, “Gravity, Particle Physics and Their Unification,” \[hep–ph/0002092\].
M. B. Green, J. H. Schwarz and E. Witten, *Superstring Theory, Vol. 1 and 2*, Cambridge University Press, 1987.
N. Arkani-Hamed, S. Dinopoulos and G. Dvali, *Phys. Lett., B429*, 263 (1998) and \[hep–ph/9803315\].
Ruel V. Churchill, *Complex Variables and Applications*, McGraw-Hill Book Company, Inc., New York, NY (1960), pp. 259–268. The term “analytic continuation” usually refers to an individual analytic function. However, we are dealing here with the “analytic continuation” of a 2 by 2 *matrix* $`𝐅(\mathrm{op})`$ whose components are four different, but closely related, analytic functions ($`\pm \mathrm{cos}\theta ,\mathrm{sin}\theta e^{\pm i\varphi }`$). Because the term “analytic continuation” has a precise mathematical meaning, and because we would prefer to avoid confusion with well-established mathematical terminology, we will relax this precision in favor of the term “continuation” (without quotes) or “analytic continuation” (with quotes) whenever we refer to the continuation or “analytic continuation” of the matrix $`𝐅(\mathrm{op})`$.
Gerald L. Fitzpatrick, *The Family Problem-New Internal Algebraic and Geometric Regularities*, Nova Scientific Press, Issaquah, WA (1997). Additional information: http://physicsweb.org/TIPTOP/ or http://www.amazon.com/exec/obidos/ISBN=0965569500. In spite of the many successes of the standard model of particle physics, the observed proliferation of matter-fields, in the form of “replicated” generations or families, is a major unsolved problem. In this book I propose a new organizing principle for fundamental fermions, i.e., a minimalistic “extension” of the standard model based, in part, on the Cayley-Hamilton theorem for matrices. In particular, to introduce (internal) global degrees of freedom that are capable of distinguishing all observed flavors, I use the Cayley-Hamilton theorem to generalize the familiar standard-model concept of scalar fermion-numbers $`f`$ (i.e., $`f_m=+1`$ for all fermions and $`f_a=1`$ for all antifermions). This theorem states that *every* (*square*) *matrix satisfies its characteristic equation*. Hence, if $`f_m`$ and $`f_a`$ are taken to be the eigenvalues of some real matrix $`𝐅(v)`$—a “generalized fermion number”—it follows from this theorem that both $`f`$ and $`𝐅(v)`$ are square-roots of unity. Assuming further that the components of both $`𝐅(v)`$ and its eigenvectors are global charge-like quantum observables, and that $`𝐅(v)`$ “acts” on a (real) vector 2-space, both the form of $`𝐅(v)`$ and the 2-space metric are determined. I find that the 2-space has a non-Euclidean or “Lorentzian” metric, and that various associated 2-scalars serve as global flavor-defining “charges,” which can be identified with charges such as strangeness, charm, baryon and lepton numbers etc.. Hence, these global charges can be used to describe individual flavors (i.e., flavor eigenstates), flavor doublets and families. Moreover, because of the aforementioned non-Euclidean constraints, and certain standard-model constraints, I find that these global charges are effectively- “quantized” in such a way that families are replicated. Finally, because these same constraints dictate that there are only a limited number of values these charges can assume, I find that families always come in “threes.”
J. Bernstein, *Elementary Particles and Their Currents*, W. H. Freeman and Co., San Francisco (1968), pp. 23–25.
T. D. Lee, *Particle Physics and Introduction to Field Theory Vol. I*, Harwood Academic Publishers, New York, NY (1981), pp. 210–211.
Even though $`𝐅(v)_{\mathrm{diag}}=𝐅(\mathrm{op})=\sigma _z`$, these two matrices act on entirely different spaces, since $`𝐅(\mathrm{op})`$ is associated with the *constant* $`\mathrm{cos}\theta =1`$, and the *variable* phase-factor $`e^{i\varphi }`$, whereas $`𝐅(v)`$ is associated with the *variable* $`\mathrm{cos}\theta 1`$, and the *constant* phase-factor $`e^{+i\varphi }=+i`$.
When weak interactions are “turned off” flavor eigenstates and mass eigenstates are one and the same. For the most part, when we speak here of flavor eigenstates, we are referring to the situation where flavor- and mass-eigenstates are the same.
Strictly speaking, besides the specification of global charges, the overall quantum state of a fundamental fermion would, necessarily, involve a specification of the spin state, the energy-momentum state and so on, together with a specification of the particular mix of local color (gauge)-charges $`R`$, $`W`$, $`B`$, $`G`$ and $`Y`$ carried by each fundamental fermion. This color-mix would be determined, in turn, by a complementary, local $`SU(5)`$ color-dependent gauge description.
Any acceptable $`2\times 2`$ matrix $`𝐅(v)`$ possesses just *two*, *real* linearly-independent eigenvectors, call them $`𝐐`$ and $`𝐐^c`$, corresponding to the two *real* eigenvalues $`f_m`$ and $`f_a`$, respectively. Therefore, the matrix $`𝐅(v)`$ can be thought of as “producing” the conventional single-particle fermion numbers $`f_m`$ and $`f_a`$ via the 2-space eigenvalue equations
$$𝐅(v)𝐐=f_m𝐐$$
and
$$𝐅(v)𝐐^c=f_a𝐐^c,$$
respectively.
The 2-vector $`𝐐`$ (and its scalar components—the *electric* charges of quarks and leptons) describes *matter*, while the linearly-independent 2-vector $`𝐐^c`$ describes its antimatter counterpart. The superscript $`c`$ on $`𝐐^c`$ is merely a label signifying antimatter. It is not an exponent or a symbol for complex conjugation. As such, it signifies only that *the* 2-*vectors $`𝐐`$ and $`𝐐^c`$ are real vectors associated with* (“*carried by*,” “*representing*,” *etc.) individual fundamental-fermions or antifermions, respectively*, not state vectors in some Hilbert space.
Even though the 2-vector $`𝐐`$ $`(𝐐^c)`$ does not represent a quantum state, it is associated with the phase factor $`e^{i\alpha f_m}`$ $`(e^{i\alpha f_a})`$ associated with a quantum state describing matter (antimatter). To see this, replace $`𝐅(\mathrm{op})`$ in Eq. 1 in the main text by $`𝐅(v)`$. Then $`U`$ is replaced by $`U^{}`$ where
$$U^{}𝐐=e^{i\alpha f_m}𝐐$$
and
$$U^{}𝐐^c=e^{i\alpha f_a}𝐐^c.$$
Note that because $`e^{in(\delta \alpha )𝐅(v)}=e^{i\alpha 𝐅(v)}`$ when $`n`$ is very large and $`\delta \alpha `$ is very small, but nonzero (i.e., $`n\delta \alpha =\alpha )`$, $`\alpha `$ can be any *finite* number ranging from zero to infinity. Similar arguments apply to the matrix $`U`$ since $`U(\alpha )U(\alpha ^{})=U(\alpha +\alpha ^{})`$.
When we say that the vectors $`𝐐`$, $`𝐔`$ and $`𝐕`$ are *observables*, we mean that their associated component-“charges” are mutually-commuting simultaneous observables. Hence, all of these charge-like components can be known in principle, at the *same* time, meaning that the vectors $`𝐐`$, $`𝐔`$ and $`𝐕`$ can be known simultaneously.
T. Kajita, for the Super-Kamiokande, Kamiokande Collaboration, \[hep–ex/9810001\].
H. Georgi and S. L. Glashow, “Neutrinos on Earth and in the Heavens,” \[hep–ph/9808293\].
G. L. Fitzpatrick, “Topological Constraints on Long-Distance Neutrino Mixtures,” \[aps1999feb12 001\] available at: http://publish.aps.org/eprint/
A. E. Faraggi, “Towards the Classification of the Realistic Free Fermion Models,” \[hep-th/9708112\].
G. B. Cleaver, A. E. Faraggi, D. V. Nanopoulos and T. ter Veldhuis, “Towards String Predictions,” \[hep-ph/0002292\].
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# Effect of 𝛾-irradiation on superconductivity in polycrystalline YBa2Cu3O7-δ
## I Introduction
The influence of crystal lattice disorder on superconducting properties is one of the important problems of superconductivity. Theoretical insights into this field have been initiated by seminal Anderson’s work. In low-$`T_c`$ superconductors (LTSCs) the disorder influence on critical temperature $`T_c`$ is believed to be mainly due to dependences of density of states at the Fermi level $`N(E_F)`$, constant of electron-phonon interaction $`\lambda _{ep}`$ or Coulomb interaction on disorder (see reviews in Refs. ). Many aspects of this problem are still not understood clearly. Things get much worse in the case of high-$`T_c`$ superconductors (HTSCs) which are layered cuprates with perovskite-related structure. Really, it is hardly possible to consider the crystal-disorder effects in HTSCs quite properly in circumstances where the nature of superconducting pairing in these compounds is still controversial . Notwithstanding this handicap, the crystal-lattice disorder effects in HTSCs have become topic of a large body of experimental and theoretical research. The necessity of such type of studies is quite obvious from technological and fundamental points of view.
In HTSCs as well as in LTSCs the superconductivity involves pairs of electrons. Therefore, although the nature of superconducting paired state in the HTSCs is still obscure, it can be assumed that some features of the disorder influence should be common for both types of superconductors. This can be seen from the following consideration. The order parameter in the Ginzburg-Landau theory is written as $`\mathrm{\Psi }=\mathrm{\Delta }\mathrm{exp}(i\phi )`$ where $`\mathrm{\Delta }`$ is the amplitude and $`\phi `$ the phase of this parameter. The disorder can destroy the superconductivity either by reducing the amplitude of the order parameter or by destroying the phase coherence of superconducting electrons. The relative contribution of each of these mechanisms depends on the degree of inhomogeneity of the superconducting system. The role of inhomogeneity can be judged by comparing the two characteristic lengths $`\xi _d`$ and $`\xi _{GL}(T)`$, the first of which is characteristic scale of inhomogeneity associated with disorder and the second is the Ginzburg-Landau coherence length for dirty superconductors
$$\xi _{GL}(T)=\xi _{GL}(0)ϵ^{1/2},$$
$`(1)`$
where $`\xi _{GL}(0)=\pi D\mathrm{}/8kT_c`$ is the coherence length at $`T=0`$, $`D=v_Fl/3`$ is the electron diffusion coefficient ($`v_F`$ is Fermi velocity, $`l`$ is electron elastic scattering length), $`ϵ=\mathrm{ln}(T_c/T`$). When $`\xi _d\xi _{GL}(T)`$ (as for instance, in the case of the impurity or point defects of crystal lattice), the system behaves as homogeneous one, and disorder affects mainly the order parameter. In the opposite case, when $`\xi _d\xi _{GL}(T)`$, the system is inhomogeneous. An example of such system is granular metal consisting of metal grains separated by dielectric interlayer ( $`\xi _d`$ is grain size in this case). The reduction in the critical temperature $`T_c`$ in inhomogeneous systems is largely due to the suppression of phase coherence between weakly coupled superconducting regions. Both mechanisms of the $`T_c`$ reduction under disorder (reducing the amplitude of the order parameter and destroying the phase coherence of superconducting electrons) can operate simultaneously in real systems.
The HTSCs are usually characterized by low values of $`v_F`$ and $`D`$ as compared with that of LTSCs. Together with high $`T_c`$ values, this should lead, according Eq. (1), to exceptionally short superconducting coherence length. Indeed, it was found, for example, in YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-δ</sub> that $`\xi _{GL}(0)`$ is about $`\xi _{ab}1.4`$ nm and $`\xi _c0.2`$ nm, parallel and perpendicular to the CuO<sub>2</sub>–layer. From this it follows that the influence of structure inhomogeneity and thus the suppression of phase coherence should be in general far more important in HTSCs than in LTSCs.
Actually the HTSC films or bulk samples are always disordered and inhomogeneous to some extent. For example, in polycrystalline samples the regions of grain boundaries are highly disordered. These regions can be not only non-superconductive, but dielectric as well. In fact, polycrystalline samples are granular metals in which superconducting coherence between the grains can be established by Josephson coupling. Beside the granularity, chemical-composition inhomogeneity is quite common in HTSCs. The influence of this type of inhomogeneity on superconductivity can be rather strong since charge carrier density of HTSCs depends strongly on chemical composition. The composition disorder can contribute to destroying of superconducting phase coherence in the same fashion as the granularity. In most cases the inhomogeneous distribution of oxygen leads to this type of disorder. Two main sources for oxygen inhomogeneity can be distinguished: extrinsic and intrinsic. Extrinsic source is due to various technological factors of sample preparation. The intrinsic one is connected with phase separation of HTSCs on two phases with different oxygen and hence charge carrier concentration.
It is clear from the aforesaid that to a good approximation the two main types of disorder, which are essential for superconductivity, can be distinguished. The first of them is disorder associated with perturbations of crystal lattice on the atomic scale (impurities, vacancies, and the like). This type of disorder is often called microscopic. It can be responsible for electron localization and other phenomena which affect the superconducting order parameter. The second type of disorder is associated with structural inhomogeneity of superconductor (granular structure, phase separation, inhomogeneity due to technological reasons). The disorder scale in this case is much more than interatomic distances (at least, more than, say, 10-100 nm) and hence, this disorder can be called macroscopic. The macroscopic disorder affects mainly the superconducting phase coherence. In experimental studies it is desirable to separate the effects of microscopic and macroscopic disorder or, at least, to be aware of the possible joint action of them in the case where such separation is difficult. Ignoring this point could lead to serious errors in interpretation of experimental results or to total their misunderstanding.
In study of disorder effects in HTSCs it is important to use reliable methods of controllable disordering. Ideally, it is desirable to ”tune” disorder while keeping other parameters fixed. At present, the main method which is expected to be somewhat close to such ideal is irradiation of superconductor with fast (charged or uncharged) particles or with high-energy photons ($`\gamma `$-rays). By appropriate choice of particles and their energy the irradiation would be expected to produce only the microscopic disorder (vacancies and interstitials), although in the HTSCs a possibility of inducing of compositional disorder must be always taken into account.
In this communication we present some new results of investigation of influence of $`\gamma `$-irradiation on the critical temperature of bulk polycrystalline YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-δ</sub> (YBCO). The disordering of HTSCs with $`\gamma `$-rays was used rather often in previous experiments, especially in the first years after discovery of HTSCs. Such type of investigations could have an applied importance in the case of possible use of HTSCs in environments of nuclear reactors or space stations. The fundamental importance of experiments of this sort is also beyond question.
It is known that $`\gamma `$-rays produce ionization in solids. In metals, ionization produced by radiation is very rapidly neutralized by the conduction electrons. Beside this, $`\gamma `$-rays produce displaced atoms in solid, namely, vacancy-interstitial pairs (Frenkel defects) of small separation, randomly distributed through the lattice. Interaction of $`\gamma `$-rays with matter, which leads to atomic displacements, occurs principally by means of three mechanisms: the photoelectric effect, the Compton effect, and pair production. In all three processes electrons are ejected with energy comparable with the original $`\gamma `$-ray energy, and thus $`\gamma `$-ray irradiation inevitably causes a substance to be internally bombarded by fairly energetic electrons. At $`\gamma `$-ray energy of a few MeV, the dominant contribution to displacement production comes from Compton effect.
When compared with other irradiation sources, $`\gamma `$-rays have one unquestionable advantage. Attenuation distances of $`\gamma `$-rays with energies of few MeV are order of a few centimeters. This enables one to investigate bulk samples of HTSCs. By contrast, penetration distances for irradiation with particles or ions are small. In this case it is possible to study the disorder effects quite properly only in rather thin films of HTSCs. Otherwise, it can not be excluded the implantation and doping effects as well as inhomogeneity of damage that plagues the interpretation of the experimental results.
A disadvantage of $`\gamma `$-rays is that the available monoenergetic radiactive sources give low $`\gamma `$-ray fluxes. Taking into account the effective cross section for atomic displacement through the Compton mechanism one would not expect any significant changes in the conductivity of typical metals, like copper. Indeed, a typical dose of $`\gamma `$-ray irradiation is about $`1\times 10^{17}`$ photons/cm<sup>2</sup>. For this dose, the number of displacements per atom (dpa) for copper is about $`10^8`$. The calculated increase in resistivity of copper due to 1% of point defects (that is for 0.01 dpa) is about 1 $`\mu \mathrm{\Omega }`$ cm. From this it is apparent that the fraction of displaced atoms, which can be produced through the Compton mechanism in a convenient irradiation time, should not induce any noticeable resistance variations in good metals. However, in poor conductors with low charge density, such as semiconductors, $`\gamma `$-rays can strongly influence the conductivity through the carrier removal.
Charge-carrier density in optimally doped HTSCs (such as YBCO) is usually about $`5\times 10^{21}`$ cm<sup>-3</sup>, that is well below that of good metals. This density is not however low enough to expect some significant influence of $`\gamma `$-ray irradiation on conductivity and order parameter of HTSCs, and hence, on their $`T_c`$. In spite of this, in some known experiments a quite appreciable effect of $`\gamma `$-irradiation on $`T_c`$ and resitivity of YBCO was revealed. The maximal decrease in $`T_c`$ was found to be about 2K by high doses $`\mathrm{\Phi }1000`$ MR. In general, the published experimental results about the effect of $`\gamma `$-irradiation on YBCO are quite contradictory (all of them concern the optimally doped YBCO with $`T_c`$ above 90 K). For example, in contrast to Ref. , in some other investigations no change in $`T_c`$ was found at all up to doses about 1000 MR, or found to be much less than 1 K. The results about $`\gamma `$-ray effect in resistivity are even more conflicting. It is possible to find among them such extremal cases: a) No influence at all up to dose $`\mathrm{\Phi }1000`$ MR; b) A nearly tenfold increase in resistivity at $`\mathrm{\Phi }150`$ MR.
The above-mentioned relevant $`\gamma `$-ray papers just present experimental results without much (or any) consideration of possible mechanisms of $`\gamma `$-ray influence on conductivity and $`T_c`$ of HTSCs, or presentation of general view on such influence. In none of them the degree of radiation damage (the value of dpa) has been estimated that severely hinders the understanding and interpretation of the results. It may be inferred, therefore, that the undestanding of $`\gamma `$-ray influence on HTSCs is far from completion and that further experimental and theoretical investigations of this matter are necessary.
In this paper, we report investigation of the effect of $`\gamma `$-rays on $`T_c`$ and resistive transition to superconducting state in bulk polycrystalline sample of optimally doped YBCO. We have observed non-monotonic behavior of $`T_c`$ with increasing irradiation dose $`\mathrm{\Phi }`$ (up to about 220 MR): $`T_c`$ decreases at low doses ($`\mathrm{\Phi }50`$ MR) approximately by 2 K and then rises, forming a minimum. Quite unexpected, at highest doses ($`\mathrm{\Phi }120`$ MR) $`T_c`$ value goes clearly down again. The temperature width, $`\delta T_c`$, of resistive superconducting transition increases rather sharp with dose in the range $`\mathrm{\Phi }75`$ MR and somewhat drops at higher dose. We believe that this interesting effect is revealed for the first time. Unlike previous works, we have calculated cross sections for the displacement of different kinds of lattice atoms in YBCO by $`\gamma `$-rays due to the Compton process. To our knowledge, nobody has done it before for HTSCs. This enables us to estimate the possible dpa value and has facilitated the interpretation of results.
On the strength of the obtained results and those of previous investigations, together with the results of radiation damage calculation, we come to conclusion that for commonly used doses (up to $`\mathrm{\Phi }`$ 1000 MR) $`\gamma `$-rays should not produce any substantial influence on superconducting properties and conductivity of single-crystal HTSCs. In polycrystalline samples however this influence can be fairly strong. These samples, as was mentioned above, are in fact granular metals which consist of metallic grains separated by dielectric interlayers. In the HTSCs the regions of grain boundaries, and near environtments of them, are strongly depleted with charge carriers and thus should be very sensitive to $`\gamma `$-rays and radiation of any other kind. For this reason the conductivity and $`T_c`$ in polycrystalline HTSCs can be markedly affected by $`\gamma `$-rays even at commonly used not very high doses. This is the main conclusion of our paper. Based on this it is possible to understand, at least qualitatively, the sharp dictinction between the known published results about $`\gamma `$-ray influence on YBCO. In this paper, the possible mechanisms of impact of radiation-damage in grain boundary regions upon the conductivity and $`T_c`$ in YBCO are considered. Beside the atomic displacements, the probable effects of atom ionization in grain boundary regions by $`\gamma `$-rays are speculatively discussed.
## II Sample and experiment
For sample preparing the conventional solid-state reaction method was used. Fine powders of Y<sub>2</sub>O<sub>3</sub>, BaO and CuO were mixed in alumina mortar in the composition of YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-δ</sub>. Then the mixture was pressed to pellet at pressure about 100 MPa. This pellet was heated at $`T930`$C in air and then slowly cooled down. From the pellet a 3.1$`\times `$6.3$`\times `$22 mm<sup>3</sup> bar was cut. The oxygen content in sample was estimated to be about 6.9 ($`\delta 0.1`$). The sample was polycrystalline with rather large grain size (about 12 $`\mu `$m). Density of sample was 5.2 g/cm<sup>3</sup>, that comprises about 80% of the expected density for this compound. The 20% difference can be attributed to an occurence of small voids and pores that is quite usual for polycrystalline samples.
Irradiation was carried out using <sup>60</sup>Co $`\gamma `$-ray source at room temperature in air at different doses up to $`\mathrm{\Phi }220`$ MR. This isotope emitts the equal amounts of photons with energies 1.17 MeV and 1.33 MeV. The dose rate of used source was $`6.34\times 10^4`$ R/h. The temperature dependences of resistance $`R(T)`$ were measured by standard four-probe method with a direct measuring current 10 mA. The accuracy of temperature measurements was estimated to be about 0.1 K.
## III Radiation-damage calculation
In experimental investigations of radiation effects in solids it is imperative to realize the degree of disorder produced by irradiation in objects studied. Without such knowledge, adequate understanding and interpretation of results obtained is impossible. In many cases such ignorance can lead to wrong conclusions. For this reason, calculations of possible radiation damage, primarilly the dpa values, become nowadays an integral part of the most of radiation-effect investigations. For example, at ion irradiation studies the well-known TRIM simulation program is widely used. The surprising thing is that one cannot find even rough estimates of dpa in the known studies of $`\gamma `$-ray effect in YBCO. This type of calculations is rather cumbersome and may be for this reason they had not been done in the above-mentioned studies. The principal concepts of such calculations are however quite established and therefore the estimation of dpa or concentration of point defects induced by $`\gamma `$-rays can be done without principal difficulties.
In this work we have calculated the cross sections for atomic displacements in YBCO by $`\gamma `$-rays due to the Compton process which is a main producer of energetic electrons at photon energy about 1 MeV. Somewhat different versions of this type of calculation are presented in Ref. . In this study the version of Ref. was used. The number of atoms displaced in unit volume (cm<sup>3</sup>) per second is expressed in Ref. by
$$R_d^\beta =_0^{E_{max}}n_0\sigma _d^\beta (E)\mathrm{\Phi }_\beta (E)𝑑E,$$
$`(2)`$
where $`n_0`$ is number of atoms per unit volume; $`\sigma _d^\beta (E)`$ is cross section for an atom to be displaced by an electron of energy $`E`$; $`\mathrm{\Phi }_\beta (E)`$ is flux density of ejected Compton electrons per cm<sup>2</sup>/sec, at energy $`E`$, per unit range of energy; $`E_{max}`$ is maximal energy of Compton electron, which is given by $`E_{max}=2E_\gamma /(1+2E_\gamma )`$ with $`E_\gamma `$ being photon energy in monoenergetic $`\gamma `$-ray flux. The Eq. (2) can be written in more detail:
$$R_d^\beta =\mathrm{\Phi }_\gamma n_0^2_0^{E_{max}}𝑑E\sigma _d^\beta (E)(dE/dx)^1_E^{E_{max}}\sigma _c(E^{^{}})𝑑E^{^{}},$$
$`(3)`$
where $`\mathrm{\Phi }_\gamma `$ is $`\gamma `$-ray flux per cm<sup>2</sup>/sec, $`\sigma _c(E^{^{}})`$ is cross section per atom for the production of a Compton electron at energy $`E`$. This is given by the formula
$$\sigma _c(ϵ)=\sigma _0\left\{\frac{1}{1ϵ}+1ϵ+\frac{ϵ}{\gamma ^2(1ϵ)}\left[\frac{ϵ}{1ϵ}2\gamma \right]\right\},$$
$`(4)`$
where $`ϵ=E/E_\gamma `$, $`\gamma =E_\gamma /mc^2`$, and $`\sigma _0=\pi r_0^2Z_2mc^2/E_\gamma ^2`$. Here $`e`$ and $`m`$ are electronic charge and mass, respectively, $`c`$ is velocity of light, $`r_0=e^2/mc^2`$, and $`Z_2`$ is the atomic number. The energy loss per cm of electron path, $`dE/dx`$, is given by
$$dE/dx=a\left[\left(1+E/mc^2\right)/E\right]$$
$`(5)`$
where $`a=2\pi e^4n_0Z_2L`$, $`L10`$.
The cross section $`\sigma _d^\beta (E)`$ in Eqs. (2) and (3) for an atom to be displaced by an electron of energy $`E`$ depends essentially on the specific value of the threshold energy $`E_d`$. It is assumed that an atom is always displaced from its lattice site when it receives energy greater then $`E_d`$ and is never displaced at lower energy. Therefore, it should be put down $`\sigma _d^\beta (E)=0`$ in Eq. (3) if $`E<E_d`$. For the greater values of $`E`$ we have used the known McKinley-Feshbach formula to calculate $`\sigma _d^\beta (E)`$:
$$\sigma _d^\beta (E)=\frac{\pi }{4}b^2\left[\left(\frac{T_m}{E_d}1\right)\beta ^2\mathrm{ln}\frac{T_m}{E_d}+\pi \alpha \beta \left\{2\left[\left(\frac{T_m}{E_d}\right)^{1/2}1\right]\mathrm{ln}\frac{T_m}{E_d}\right\}\right]$$
$`(6)`$
where $`(\pi /4)b^2=\pi Z_2r_0^2(1\beta ^2)/\beta ^4`$; $`\beta =v/c`$ that is the ratio of electron and light velocities for which the relation $`\beta ^2=E(E+2)/(E+1)^2`$ is true if $`E`$ is taken in the units of $`mc^2`$; $`\alpha Z_2/137`$, and $`T_m`$ is the maximum energy which can be transferred in a collision by an electron of kinetic energy $`E`$:
$$T_m=\frac{2(E+2mc^2)}{M_2c^2}E,$$
$`(7)`$
where $`M_2`$ is the target atomic mass.
Some additional comments about the McKinley-Feshbach formula are necessary. It is assumed that it is accurate to one percent for $`Z_2`$ up to 40. Atomic numbers of elements in YBCO are 8(O), 29 (Cu), 39 (Y), and 56 (Ba). Therefore, an error more than one percent should be expected only for Ba atoms. But displacement cross section for heavy atoms is usually much less than for light ones, and thus, contribution from heavy atoms to total dpa should be rather small. For this reason we did not expect a significant error with the McKinley-Feshbach formula. This was justified by our calculations which are presented below.
If the struck atom has large enough energy it can cause secondary displacements which can also be taken into account at calculation of numbers of the displaced atoms due to $`\gamma `$-rays. It is easy to see however that at electron energy about 1 MeV the values of $`T_m`$ given by Eq. (7) are quite small and therefore this effect can be neglected. Nontheless, we have done cascade calculations following the recommendations of Refs. and have found that primary displacement cross sections differ from these of total displacement cross sections only by few percents. Although we will present below the total displacement cross sections, it should be kept in mind that they differ very slightly from those of obtained using Eqs. (4)–(7).
The calculated values of the effective cross section for atomic displacement by $`\gamma `$-rays from <sup>60</sup>Co source through the Compton effect, $`\sigma _c^\gamma =R_d^\beta /(\mathrm{\Phi }_\gamma n_0`$), at different values of threshold energy $`E_d`$ for ions in YBCO are given in Table I. The values of $`\sigma _c^\gamma `$ presented in this Table are the weighted averages for photons with energies 1.17 MeV and 1.33 MeV which are emitted by <sup>60</sup>Co source. The values of $`n_0`$ for different ions at calculation of $`\sigma _c^\gamma `$ were just partial ion densities for compound YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-δ</sub> deduced taking into account the compact YBCO mass density which is about 6.4 g/cm<sup>3</sup>. It can be seen from Table I that the $`\sigma _c^\gamma `$ values depend crucially on threshold energy $`E_d`$. It is worth noting also that the $`\sigma _c^\gamma `$ values for heavy Ba ions are the least as compared with other ions, as expected. The data of Table I will be used below for estimation of possible dpa values in studied YBCO after $`\gamma `$-irradiation.
## IV Results and discussion
The temperature dependence of sample resistivity $`\rho (T)`$ before $`\gamma `$-irradiation is presented in Fig. 1. The dependence is quite common for optimally doped YBCO: it is linear above the superconducting resistive transition up to room temperature, $`T_c`$ value is about 93 K. We have defined experimental $`T_c`$ to be the temperature at which normal resistance $`R_n`$ is halved. The dependences $`R_n(T)`$ in the range of resistive transition have been obtained by extrapolation of linear $`R(T)`$ dependence to this region. We have used the temperature $`T_{cz}`$ at which resistance goes to zero as a second characteristic of resistive transition. In fact, $`T_{cz}`$ is the end point of resistive transition to superconducting state, and is an essential characteristic which should be taken into account for granular or inhomogeneous superconductors. The difference in $`T_c`$ and $`T_{cz}`$, $`\delta T_c=T_cT_{cz}`$, is some quite definite measure of the width of resistive transition.
In early days of HTSC investigations (and quite often up to date) the experimental $`T_c`$ was defined as the temperature $`T_{cb}`$ at the onset of the superconducting transition. Although this is also quite essential characteristic of resistive transition, it can be evaluated with much less precision than $`T_c`$ or $`T_{cz}`$. For this reason only the observed changes in $`T_c`$ and $`T_{cz}`$ with $`\gamma `$-ray dose will be considered below.
The $`\rho (T)`$ curves for different $`\gamma `$-ray doses $`\mathrm{\Phi }`$ in the temperature range of superconducting transition are presented in Fig. 2. It can be seen that $`T_c`$ depends on $`\mathrm{\Phi }`$ in a non-monotonic way. More clearly this is seen in Fig. 3 where the changes in $`T_c`$ and $`T_{cz}`$ with $`\gamma `$-ray dose are shown. It follows from the Figure that $`T_c`$ decreases at low dose ($`\mathrm{\Phi }50`$ MR) by $`2`$ K and then rises, forming a minimum. At higher dose ($`\mathrm{\Phi }120`$ MR) $`T_c`$ value goes clearly down again. The zero-resistance temperature $`T_{cz}`$ has been changing with $`\gamma `$-ray dose in the nearly same way as $`T_c`$, but with greater amplitude: the initial decrease in $`T_{cz}`$ is about 4 K. The magnitude of $`\delta T_c`$ (which characterizes the width of resistive transition and, hence, the sample inhomogeneity) increases with dose in the range $`\mathrm{\Phi }75`$ MR and somewhat drops at higher doses (Fig. 4).
The observed initial $`T_c`$ decrease with $`\gamma `$-ray dose (Fig. 3) corresponds to some of previous studies, but the general picture of non-monotonic dependence of $`T_c`$ on radiation dose looks like a surprising thing. To our knowledge such behavior of $`T_c`$ in HTSC with $`\gamma `$-ray dose is found for the first time. A somewhat (or partially) similar non-monotonic behavior has been seen previously in $`\gamma `$-irradiated YBCO-related compound Y<sub>0.9</sub>Sm<sub>0.1</sub>Ba<sub>2</sub>Cu<sub>3</sub>O<sub>7-δ</sub>, but for the onset temperature $`T_{cb}`$ only. As this took place, the zero-resistance temperature $`T_{cz}`$ had not manifested any marked influence of $`\gamma `$-irradiation. The resistive-transition curves shown in Ref. are rather steep (no large difference between the $`T_c`$ and $`T_{cz}`$ values). At the same time the reported values of the onset temperature $`T_{cb}`$ (about 140 K) appear to be too high for YBCO-related compounds (compare, for example, with the data of Ref. ). In our opinion in evaluating the onset temperature $`T_{cb}`$ from $`R(T)`$ curves the significant error is possible. Since an employed procedure for evaluating of the $`T_{cb}`$ values has not been outlined in Ref. one should take the $`T_{cb}`$ values in it and the described non-monotonic behavior $`T_{cb}`$ with $`\gamma `$-ray dose with some precaution.
The effect of $`\gamma `$-rays in the resistivity of sample studied was found to be appreciable only in the temperature range of rather broad resistive superconducting transition (Fig. 2). With increasing temperature away from $`T_c`$ the $`\gamma `$-ray effect in resistivity falls off quickly. In particular, at $`T105`$ K a mere 4% increase in resistivity by $`\gamma `$-rays was found, whereas above 200 K hardly any radiation effect in resistivity can be detected. This suggests that $`\gamma `$-rays affect primarily the superconducting properties of the sample studied while electron transport in normal state remains actually not affected.
In the course of explaining of the results obtained, it should be taken into account, first of all, that the sample studied is polycrystalline or, better to say, granular. Indeed, its resistivity in normal state just above the superconducting transition (about 2.5 m$`\mathrm{\Omega }`$ cm, as can be seen in Figs. 1 and 2) is by a factor of 50 larger than that of the best quality optimally doped YBCO single-crystals (about 50 $`\mu \mathrm{\Omega }`$ cm). The increased resistivity comes from grain boundaries, since regions of grain boundaries in HTSCs can be poor conductive and even dielectric.
The conductivity of granular metals in normal state is determined by tunneling of single-particle excitations (unpaired electrons) through the boundaries. In superconducting state the superconducting coherence between grains can be established by Josephson coupling. Although the sample studied is inhomogeneous (granular) its critical temperature $`T_c=93`$ K corresponds to the highest $`T_c`$ values in the good-quality YBCO single crystals. Such situation is quite possible and understandable for inhomogeneous systems. Granular metals usually have a spatial distribution of thicknesses of the poor-conductive or dielectric grain boundaries. This spread in the boundary thickness may be very important for transpor properties, especially for dielectric grain boundaries because in this case the probability of tunneling is an exponential function of dielectric separation distance. From this follows that granular metals are actually percolating systems. In such systems conductivity can be determined by the presence of optimal “chains” of grains with maximal probability of tunneling for adjacent pairs of grains forming the chain. The same type of percolation picture is true for a system of metallic grains with Josephson coupling between them (when grains become superconducting at low enough temperature). The sample studied consists of rather large (12 $`\mu `$m) grains with critical temperature $`T_c`$ as high as that of in high-quality crystals. The global critical temperature of whole sample can be very close to this maximal value in the presence of good conducting optimal chains of grains with strong Josephson coupling (in utmost case, theoretically, even one such path will be enough to arrive to zero resistance).
Since conductivity of the granular metal is determined by both intragrain and intergrain transport properties, the influence of radiation damage in it should be considered separately for materials of grains and regions of grain boundaries. Let us consider at first a question: is the radiation damage induced by $`\gamma `$-ray in this study high enough to cause quite significant variations in resistivity and superconducting properties of material of grains. For this purpose, we have estimated the possible values of dpa in sample studied using the calculated in Sec. III effective cross sections for atomic displacement by $`\gamma `$-rays from <sup>60</sup>Co source through the Compton effect (Table I). It is easy to see that for the stardard value of $`E_d=20`$ eV the total dpa comprises very small value about $`10^7`$. There is experimental evidence that the displacement energy, $`E_d`$, for oxygen in CuO<sub>2</sub> planes in YBCO is close to 10 eV. In this case the dpa will be about $`6\times 10^7`$. From the general point of view (taking into account the rather high charge-carrier density in optimally doped YBCO) it should not be expected any significant variations in normal state resistivity and $`T_c`$ in grains at such small radiation damage. Indeed, the known experimental studies of single-crystal HTSCs irradiated with electrons or ions show that an appreciable influence of disorder on resistivity and $`T_c`$ can be detected only at dpa greater than $`10^3`$.
It follows from the above consideration that with commonly used $`\gamma `$-ray doses ($`1000`$ MR) one should not expect any detectable variations in resistivity and $`T_c`$ in fairly homogeneous single-crystal HTSC as well as in metallic grains of polycrystalline HTSC. For this reason we will not consider here the numerous theoretical models of disorder influence on superconducting properties of HTSCs (see Refs. and references therein) which were developed for homogeneous superconductors. In this study we have not seen any marked changes even in the normal state resistivity of the polycrystalline sample. But we observed the pronounced effect of $`\gamma `$-rays on $`T_c`$ and the width of resistive transition (Figs. 2-4). This is undoubtedly connected with influence of $`\gamma `$-rays on poor conductive or dielectic regions of grain boundaries and, therefore, on the superconducting phase coherence between the grains.
One of the first impressive demonstration of radiation-induced destruction of phase coherence of the superconducting wave functions between grains in polycrystalline HTSC was demonstrated in Ref. for optimally doped YBCO. They found that zero-resistance temperature $`T_{cz}`$ was very sensitive to even low doses of irradiation with 500 keV oxygen. At the same time, the onset temperature $`T_{cb}`$ declined much slowly. For doses, where $`T_{cz}`$ was already close to zero, and furthermore, even, when insulating behavior of $`R(T)`$ was evident below 40 K, $`T_{cb}`$ was close to 80 K (before irradiation $`T_{cb}`$ was 97 K, and $`T_{cz}`$ was 87 K). The authors of Ref. have assumed that $`T_{cb}`$ in this case reflects the behavior of intrinsic critical temperature of grains, which is not as sensitive to irradiation as that of the whole granular system.
In the years, following Ref. , the quality of HTSCs becomes much higher and the most of subsequent radiation studies were devoted to single-crystal (or, at least, nearly single crystal) HTSCs, which are usually far more homogeneous than polycrystalline samples, and for which, therefore, the phase coherence is not of decisive importance. It is certain, however, that not only single crystal HTSC products may be (or will be) used in advanced technology. Therefore inhomogeneity effect in superconductivity of these compound are still of fundamental and tecnological importance.
It can be said with confidence that the observed decrease in $`T_c`$ and zero-resistance temperature $`T_{cz}`$ combined with simultaneous increase in width of resistive transition $`\delta T_c`$ in sample studied at low dose ($`\mathrm{\Phi }<50`$ MR) (Figs. 3 and 4) is quite expected for percolating granular system. Optimal current paths, which have ensured the measured $`T_c`$ value about 93 K before irradiation, have sure some “weak” links. These are the grain boundaries, which are strongly enough depleted with charge carriers and, therefore, are sensitive even to such small mean radiation damage as in this study (maximum $`10^6`$ dpa). This leads to the observed decrease in the “global” $`T_c`$ and increase in $`\delta T_c`$ (Figs. 3 and 4).
The initial $`T_c`$ drop (which appears to be explicable) is followed by an increase in $`T_c`$ at higher dose $`\mathrm{\Phi }>50`$ MR (Fig. 3), and this is fairly surprising. This means that some concurrent mechanism, which causes an increase in $`T_c`$, comes into play. This mechanism shows itself only in the limited range of doses, since $`T_c`$ decreases again in the range above $`\mathrm{\Phi }>120`$ MR (Fig. 3). Somehow this mechanism should be also associated with an impact of $`\gamma `$-rays on intergrain Josephson coupling. There are known some cases of superconductivity enhancement in polycrystalline HTSCs under irradiation. For example, an increase in $`T_c`$ after low-dose irradiation with Si ions was observed in polycrystalline YBCO. Explanations of the effect in Ref. was not however connected with granular structure of sample. In Ref. it was reported that $`\gamma `$-irradiation of granular Tl<sub>2</sub>Ba<sub>2</sub>Ca<sub>2</sub>Cu<sub>3</sub>O<sub>10</sub> at room temperature causes the increase of critical current. This is believed to be an indication of strengthening of the intergrain Josephson coupling under $`\gamma `$-ray influence.
The known features of interaction of $`\gamma `$-rays with solids make it possible to suggest some mechanisms of $`T_c`$ increasing in granular HTSCs at low enough $`\gamma `$-ray doses. In doing so two main points should be considered: (i) the nature of grain boundaries in HTSCs, (ii) the ionizing influence of $`\gamma `$-rays. From the literature data (see Refs. and Refs. therein) it can be concluded that grain boundaries in HTSCs are disordered and oxygen deficient regions. A typical grain-boundary width in YBCO is about 2 nm (what encompasses about 10 atomic distances). It is known that $`T_c`$ of oxygen deficient YBCO compounds increases considerably with illumination of visible light. It was found that light leads to a change in doping (increasing in charge carrier density). This photoexcited state is persistent up to 250-270 K. The photodoping effect is negligible for optimally doped YBCO. All existing models of the effect are based on the suggestion that sample illumination generates electron-hole pairs in the CuO<sub>2</sub> planes. Electrons are transferred to the CuO chains and trapped there, while holes remain mobile in the CuO<sub>2</sub> planes.
The most important effect of $`\gamma `$-rays in solids is ionization. In primary collision of $`\gamma `$-quantum with an atom of solid the Compton electron and photon with lesser energy are produced. The energy of secondary photon is dissipated in causing further ionization. The same is true for the most of the energy of Compton electrons (only occasionally they displace atoms by elastic collision). Therefore, the number of ionized atoms at $`\gamma `$-irradiation is by several orders of magnitude larger than the number of displaced atoms. In this connection it must not be ruled out that effects like the above-mentioned photodoping by visible light take place in $`\gamma `$-irradiated bulk samples as well. Of course, there is a significant difference in the energy of $`\gamma `$-ray and visible light photons. The visible light can eject only electrons from outer orbitals. These low energy electrons are trapped in nearby sites like CuO chains. By contrast, $`\gamma `$-quantum can struck out any electron from an atom. But in this case the ejected elecrons have much higher energy and, therefore, can be brought far away from initial place. If electrons are ejected from CuO<sub>2</sub> planes in oxygen depleted grain boundary regions, this should increase the charge carrier density in the planes and, therefore, enhance the Josephson coupling between adjacent grains. If this takes place in some of “weak” links in optimal current path, this will result in increasing of measured $`T_c`$. Of course, the recombination of the electrons with excess holes should bring the charge carrier density in grain boundaries to an initial level. But some of these electrons should be sure trapped in remote defect places of crystal lattice like impurities, lattice distortions regions, including grain boundaries as well. If room temperature is not high enough for the electron-hole recombination of trapped electrons, this can lead to the observed effect of $`T_c`$ enhancement in some intermediate dose range. At higher doses, however, influence of radiation defects (displaced atoms) in grain boundaries overpowers the ionizing effect of $`\gamma `$-rays that results in $`T_c`$ decreasing. Needless to say that our explanation should not be considered as a conclusive. It needs an independent theoretical and experimental confirmation.
## V Conclusion
We have found non-monotonic behavior of $`T_c`$ in optimally doped polycrystalline YBCO with increasing $`\gamma `$-ray irradiation dose. This result can be explained primarily by the influence of $`\gamma `$-rays on intergrain Josephson coupling. The two kinds of such influence are discussed. The first is the producing of displaced atoms, and the second is the ionizing influence of $`\gamma `$-rays.
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# Polarized deformation quantization
## 1 Introduction
Let $`(M,\omega _0)`$ be a symplectic manifold, $`C_M^{\mathrm{}}`$ the sheaf of complex valued functions on $`M`$. A deformation quantization (or a star-product) on $`M`$ is a structure of associative algebra on the sheaf of formal power series $`C_M^{\mathrm{}}[[t]]`$ with the multiplication of the form
$`fg={\displaystyle \underset{i0}{}}t^i\mu _i(f,g),f,gC_M^{\mathrm{}},`$ (1)
where all $`\mu _i`$ are bidifferential operators, $`\mu _0(f,g)=fg`$ and $`\mu _1(f,g)\mu _1(g,f)=\{f,g\}`$, the Poisson bracket inverse to $`\omega _0`$.
Two deformation quantizations $`𝒜_1`$ and $`𝒜_2`$ on $`(M,\omega _0)`$ are equivalent if there exists a power series $`B=1+tB_1+\mathrm{}`$, where $`B_i`$ are differential operators, such that $`B:𝒜_1𝒜_2`$ is an isomorphism of algebras.
It is known that all equivalence classes of star-products on $`M`$ with the Poisson bracket $`\phi _0=\omega _0^1`$ can be obtained by the Fedosov method. According to this method, one constructs a flat connection, $`D`$, (called the Fedosov connection) on the bundle of Weyl algebras on $`M`$. The quantized algebra, $`𝒜`$, is realized as the subalgebra of flat sections of the Weyl algebra. The Weyl curvature of $`D`$, being a closed scalar two-form of the view $`\omega =\omega _0+t\omega _1+\mathrm{}`$ defines the Fedosov class
$`\theta (𝒜)={\displaystyle \frac{1}{t}}[\omega ]{\displaystyle \frac{1}{t}}[\omega _0]+H^2(M,[[t]]).`$ (2)
It is also known that the correspondence $`𝒜\theta (𝒜)`$ is a bijection between the set of isomorphism classes of star-products on $`(M,\omega _0)`$ and the set $`\frac{1}{t}[\omega _0]+H^2(M,[[t]])`$ modulo the group of formal symplectomorphisms of $`M`$, \[Fe\], \[NT\], \[Xu\].
Let $`P`$ be a polarization of $`(M,\omega _0)`$. That is, $`P`$ is an integrable Lagrangean subbundle of the complexified tangent bundle $`T_M^{}`$. Let $`\omega =\omega _0+t\omega _1+\mathrm{}`$ be a closed form which is a deformation of $`\omega _0`$ and $`𝒫`$ a polarization of $`\omega `$, which is a deformation of polarization $`P`$. We say that the triple $`(M,\omega ,𝒫)`$ is a deformation of the triple $`(M,\omega _0,P)`$.
###### Definition 1.1.
A polarized star-product(or quantization) on $`(M,\omega _0,P)`$ is a pair $`(𝒜,𝒪)`$ of $`[[t]]`$ algebras satisfying the conditions:
1) $`𝒜`$ is a star-product on $`(M,\omega _0)`$,
2) there exists a deformation $`(M,\omega ,𝒫)`$ of $`(M,\omega _0,P)`$ such that $`𝒪`$ is the sheaf of functions constant along $`𝒫`$,
3) the star-product in $`𝒜`$ being restricted to $`𝒪`$ coincides with the original commutative multiplication.
Two polarized star-products $`(𝒜_1,𝒪_1)`$ and $`(𝒜_2,𝒪_2)`$ are equivalent if there exists an isomorphism $`B:𝒜_1𝒜_2`$ of star-products such that $`B(𝒪_1)=𝒪_2`$.
The first result of the paper is the following theorem (cf. Theorem 2.10):
###### Theorem 1.2.
For any triple $`(M,\omega _0,P)`$, where $`P`$ is good (see Definition 2.4) and any its deformation $`(M,\omega ,𝒫)`$ there exists a polarized star-product, $`(𝒜,𝒪)`$, such that $`𝒪`$ is the sheaf of functions constant along $`𝒫`$ and $`\theta (𝒜)=\frac{1}{t}[\omega ]`$.
For example, any Kähler or real polarization is good. This result generalizes a result of N. Reshetikhin and M. Yakimov, \[RY\], who consider the case of a real polarization $`P`$ defined by a Lagrangean fiber bundle $`MB`$.
Our proof of this result uses the Fedosov method adapted for the case with polarization. The analogous method was used by M. Bordemann and S. Waldmann, \[BW\], for a construction of quantization with separation of variables on a Kähler manifold.
In fact, we prove a slightly stronger statement. Namely, the polarized quantization $`(𝒜,𝒪)`$ we constructed satisfies the property $`fg=fg`$ for $`f𝒪`$ and any $`g𝒜`$. Such a type of star-products is in the spirit of A. Karabegov (see \[Ka1\], where he constructs star-products with separation of variables on Kähler manifolds).
Let $`(𝒜,𝒪)`$ be a polarized quantization on $`(M,\omega _0,P)`$. Let $`F(𝒜)=\{f𝒜;[f,𝒪]𝒪\}`$. The sheaf of $`𝒪`$ modules $`F(𝒜)`$ has the natural structure of a sheaf of Lie algebras with $`𝒪`$ as a commutative ideal. $`T_𝒪=F(𝒜)/𝒪`$ can be considered as a sheaf of Lie algebras consisting of derivations of $`𝒪`$, and we may consider the exact sequence of sheaves of $`𝒪`$ modules and Lie algebras
$`0𝒪F(𝒜)T_𝒪0.`$ (3)
$`F(𝒜)`$ is an example of an $`𝒪`$-extension of $`T_𝒪`$, (see Definition 4.1).. We prove that if $`P`$ is a em strong polarization (see Definition 3.1), this sequence locally splits. In this case the class $`\mathrm{cl}(𝒜,𝒪)H^2(M,\frac{1}{t}[[t]])`$ of this extension is defined.
The second result of the paper is a formula that relates three classes in $`H^2(M,\frac{1}{t}[[t]])`$ associated with a polarized deformation quantization $`(𝒜,𝒪)`$: the Fedosov class $`\theta (𝒜)`$, the Chern class of the polarization $`P`$, and $`\mathrm{cl}(𝒜,𝒪)`$.
###### Theorem 1.3.
Let $`(𝒜,𝒪)`$ be a polarized quantization of $`(M,\omega _0,P)`$. Then,
$`\theta (𝒜)={\displaystyle \frac{1}{t}}\mathrm{cl}(𝒜,𝒪){\displaystyle \frac{1}{2}}c_1(P).`$ (4)
One can show that in the case of quantization on a Kähler manifold with separation of variables the class $`\mathrm{cl}(𝒜,𝒪)`$ coincides with the class defined by Karabegov in \[Ka1\]. The formula analogous to (4) relating the Karabegov and Fedosov classes was obtained in \[Ka2\].
The paper is organized as follows. In Section 2 we develop a variant of Fedosov’s method adapted to the polarized setting and prove Theorem 1.2.
In Section 3 we define strong polarizations and study their local properties. We show that if $`P`$ is a strong polarization, then any deformed polarization, $`𝒫`$, $`𝒫_0=P`$, is also strong, and any strong polarization is good. We prove that any polarized quantization of $`(M,\omega _0,P)`$ with $`P`$ a strong polarization is locally isomorphic to a standard quantization, the polarized Moyal star-product.
In Section 4 we prove Theorem 1.3. To this end we consider, for a given polarized quantization $`(𝒜,𝒪)`$, the extension (3) and study its behavior under passing to the opposite quantization, $`(𝒜^{op},𝒪)`$, and to the quantization $`(𝒜^\sigma ,𝒪)`$ obtained by the change of variable $`tt`$. Furthermore, we compare classes $`\mathrm{cl}(𝒜,𝒪)`$ with the classes introduced by Deligne in \[De\].
Acknowledgments. The authors are grateful to A. Karabegov for useful discussions and to Max-Planck-Institut für Mathematik for hospitality and very stimulating working atmosphere.
## 2 Polarized quantization
### 2.1 Setting and notation
Let $`M`$ be a smooth manifold, $`C_M^{\mathrm{}}`$ the sheaf of smooth complex valued functions on it. Let $`[[t]]`$ be the algebra of formal power series over $``$. For a sheaf $`E`$ of $`C_M^{\mathrm{}}`$ modules, $`E[[t]]`$ denotes the sheaf $`E_{}[[t]]`$ completed in the $`t`$-adic topology. $`E[[t]]`$ is a sheaf of $`C_M^{\mathrm{}}[[t]]`$ modules.
Let $``$ be a locally free sheaf of $`C_M^{\mathrm{}}[[t]]`$ modules of finite rank. Denote by $`T^k()`$ the $`k`$-th tensor power of $``$ over $`C_M^{\mathrm{}}[[t]]`$ and by $`T()`$ the corresponding tensor algebra completed in the $`\{,t\}`$-adic topology. Similarly we define the completed symmetric algebra $`S()`$. For a subsheaf $`𝒫`$ of $``$ we denote by $`sym_𝒫:S(𝒫)T()`$ the natural map of $`C_M^{\mathrm{}}[[t]]`$ modules defined by symmetrization.
Let $`\mathrm{\Lambda }()`$ be the exterior algebra of $``$ over $`C_M^{\mathrm{}}[[t]]`$. We will consider $`T\mathrm{\Lambda }=T()_{C_M^{\mathrm{}}[[t]]}\mathrm{\Lambda }()`$ as a graded super-algebra regarding a section $`xT()\mathrm{\Lambda }^k()`$ of degree $`k`$ even (odd) if $`k`$ is even (odd).
Denote by $`\delta _T=\delta _{T()}`$ the continuous $`C_M^{\mathrm{}}[[t]]`$ linear derivation of $`T\mathrm{\Lambda }`$ generated by the map $`T^1()11\mathrm{\Lambda }^1()`$, $`v11v`$, $`v`$ is a section. It is clear that $`\delta _T`$ is a derivation of degree $`1`$ and $`\delta _T^2=0`$. In the same way define a derivation $`\delta _{}=\delta _{S()}`$ in the (super-)algebra $`S\mathrm{\Lambda }=S()\mathrm{\Lambda }()`$. It is easy to see that for a subsheaf $`𝒫`$, $`\delta _T`$ on $`T()\mathrm{\Lambda }()`$ can be restricted to $`S(𝒫)\mathrm{\Lambda }(𝒫)`$ via the embedding $`sym_𝒫id_𝒫`$, and coincides with the $`\delta _{S(𝒫)}`$ defined on the algebra $`S(𝒫)\mathrm{\Lambda }(𝒫)`$.
On the algebra $`S\mathrm{\Lambda }`$ there is defined another derivation, $`\delta ^{}`$, of degree $`1`$ generated by the map $`1\mathrm{\Lambda }^1S^11`$, $`1vv1`$. It is easy to check that $`(\delta ^{})^2=0`$ and $`[\delta ,\delta ^{}]=\delta \delta ^{}+\delta ^{}\delta =deg`$, where $`deg`$ is the derivation assigning to an element $`xS^p\mathrm{\Lambda }^q`$ the element $`(p+q)x`$.
Let $``$ be presented as a direct sum of $`C_M^{\mathrm{}}[[t]]`$ submodules, $`=𝒫𝒬`$. We will identify the tensor product of algebras, $`(S(𝒫)\mathrm{\Lambda }(𝒫))(S(𝒬)\mathrm{\Lambda }(𝒬))`$, with $`(S(𝒫)S(𝒬))(\mathrm{\Lambda }(𝒫)\mathrm{\Lambda }(𝒬))=(S(𝒫)S(𝒬))\mathrm{\Lambda }()`$, where the last identification is obtained via skew-symmetrization of the last factor. Denote $`S(𝒫,𝒬)=S(𝒫)S(𝒬)`$, $`𝐒=𝐒(𝒫,𝒬)=(S(𝒫)S(𝒬))\mathrm{\Lambda }()`$, $`𝐒(𝒫)=S(𝒫)1\mathrm{\Lambda }()`$, and $`𝐒(𝒬)=1S(𝒬)\mathrm{\Lambda }()`$.
One has the embedding
$`sym_𝒫sym_𝒬id:(S(𝒫)S(𝒬))\mathrm{\Lambda }()T()\mathrm{\Lambda }().`$ (5)
It is obvious that the derivations $`\delta _𝒫`$, $`\delta _𝒫^{}`$, $`\delta _𝒬`$, $`\delta _𝒬^{}`$ induce the corresponding derivations on the algebra $`𝐒`$, and $`\delta =\delta _𝒫+\delta _𝒬`$ coincides with the restriction to $`𝐒`$ of $`\delta _T`$ on $`T()\mathrm{\Lambda }()`$ via the embedding (5).
Let us define the operator $`\delta ^1=\delta _{𝒫,𝒬}^1`$ on $`𝐒`$. There is the continuous $`C_M^{\mathrm{}}[[t]]`$ linear map that is equal to zero for $`xC_M^{\mathrm{}}[[t]]`$ and $`\delta ^1(x)=(1/(p+r+q)\delta ^{}(x)`$ for $`x(S^p(𝒫)S^r(𝒬))\mathrm{\Lambda }^q()`$, $`p+r+q>0`$.
There is the obvious relation
$`\delta \delta ^1+\delta ^1\delta =\text{ projection on }𝐒^+\text{ along }C_M^{\mathrm{}}[[t]],`$ (6)
where $`𝐒^+`$ is the closure of $`_{p+r+q>0}(S^p(𝒫)S^r(𝒬))\mathrm{\Lambda }^q()`$.
### 2.2 Fedosov algebra
Let $`\phi :C_M^{\mathrm{}}[[t]]`$ be a $`C_M^{\mathrm{}}[[t]]`$ linear skew-symmetric form and $`I`$ the closed ideal in $`T()`$ generated by relations
$`xyyx=t\phi (x,y).`$ (7)
We call $`𝒲()=T()/I`$ the Weyl algebra and $`𝐖=𝐖()=𝒲\mathrm{\Lambda }()`$ the Fedosov algebra over $``$.
Let $`=𝒫𝒬`$ be a decomposition into $`C_M^{\mathrm{}}[[t]]`$ modules.
The derivation $`\delta `$ of $`T()\mathrm{\Lambda }()`$ induces a derivation of $`𝐖`$. Indeed, $`\delta `$ applied to the both sides of (7) gives zero.
Define the Wick map, $`𝐰=𝐰_{𝒫,𝒬}`$, as the composition $`𝐒(𝒫,𝒬)T()\mathrm{\Lambda }()𝐖`$, where the first map is (5) and the second is the projection. By the PBW theorem $`𝐰`$ is an isomorphism of $`C_M^{\mathrm{}}[[t]]`$ modules.
Due to the isomorphism $`𝐰`$, all the operators $`\delta _𝒫`$, $`\delta _𝒬`$, $`\delta _{𝒫,𝒬}^1`$ carry over from $`𝐒`$ to $`𝐖`$. We retain for them the same notation. Note that while $`\delta _𝒫+\delta _𝒬`$ does not depend on the decomposition and coincides with the $`\delta `$ induced from $`T()\mathrm{\Lambda }()`$, $`\delta _{𝒫,𝒬}^1`$ is not a derivation and does depend on the decomposition. In particular, one can suppose that the decomposition is trivial, $`=0`$. In this case we denote $`\delta _{}^1=\delta _{,0}^1`$.
Note that $`\delta _𝒫`$, $`\delta _𝒬`$, and $`\delta `$ are derivations on $`𝐖`$.
###### Proposition 2.1.
One has
$$H(𝐖,\delta )=C_M^{\mathrm{}}[[t]].$$
Moreover, if $`x𝒲()\mathrm{\Lambda }^{k>0}()`$ then $`y=\delta _{𝒫,𝒬}^1x`$ is such that $`\delta y=x`$ for any decomposition $`=𝒫𝒬`$.
###### Proof.
Follows from (6). ∎
### 2.3 Lie subalgebras in $`𝒲`$
Let $`=𝒫𝒬`$ be a decomposition. We say that $`x𝐖`$ has $`𝐰_{𝒫,𝒬}`$-degree $`(p,q)`$ if $`𝐰_{𝒫,𝒬}^1(x)(S^p(𝒫S^q(𝒬))\mathrm{\Lambda }()`$. We say that $`x𝐖`$ has $`𝐰_{𝒫,𝒬}`$-degree $`k`$ if $`𝐰_{𝒫,𝒬}^1(x)_{p+q=k}(S^p(𝒫)S^q(𝒬))\mathrm{\Lambda }()`$. The $`𝐰_{}`$-degree is $`𝐰_{,0}`$-degree for the trivial decomposition $`=0`$.
Let $`\phi `$ be nondegenerate. Let $`𝔤`$ be a sheaf of Lie algebras acting on $``$. We call a $`C_M^{\mathrm{}}[[t]]`$ linear map $`\lambda :𝔤𝒲`$ a realization of $`𝔤`$ if it is a Lie algebra morphism ($`𝒲`$ is considered as a Lie algebra with respect to commutator $`\frac{1}{t}[,]`$) and for any $`x𝔤`$ and $`e`$ one has $`x(e)=\frac{1}{t}[\lambda (x),e]`$. Any two realizations differ by a Lie algebra morphism of $`𝔤`$ to the center of $`𝒲`$, so if $`𝔤`$ is a sheaf of semisimple Lie algebras there is not more than one realization of $`𝔤`$.
Denote by $`𝔰𝔭()`$ the sheaf of symplectic Lie algebras with respect to $`\phi `$. Since $`𝔰𝔭()`$ is semisimple, there is a unique realization $`\rho _{}:𝔰𝔭()𝒲`$. The image of this realization consists of elements having $`𝐰_{}`$-degree two.
Let $`=𝒫𝒬`$ be a decomposition into Lagrangean subsheaves. Denote by $`𝔰𝔭(𝒫,)`$ the subsheaf of $`𝔰𝔭()`$ preserving $`𝒫`$. It is easy to check that $`𝔰𝔭(𝒫,)`$ can be realized as the subset of elements of $`𝒲`$ having $`𝐰_{𝒫,𝒬}`$-degree $`(1,1)`$ and $`(2,0)`$. Denote this realization by $`\rho _{𝒫,}:𝔰𝔭(𝒫,)𝒲`$. On the other hand, $`𝔰𝔭(𝒫,)`$ can be realized in $`𝒲`$ by $`\rho _{}`$.
Let us define a $`[[t]]`$ linear map $`tr:𝔰𝔭(𝒫,)[[t]]`$ in the following way. Let $`a𝔰𝔭(𝒫,)`$ and $`a^{}`$ its restriction to $`𝒫`$. We put $`tr(a)=trace(a^{})`$. It is easy to see that $`\rho _{𝒫,}(a^{})`$ (here $`a^{}`$ is trivially extended to $``$) is the $`(1,1)`$ component of $`\rho _{𝒫,}(a)`$.
###### Lemma 2.2.
Let $`=𝒫𝒬`$ be a decomposition into Lagrangean subsheaves. Let $`a𝔰𝔭(𝒫,)`$. Then
$$\rho _{𝒫,}(a)t\frac{1}{2}tr(a)=\rho _{}(a).$$
###### Proof.
Straightforward. ∎
### 2.4 Filtrations on $`𝒲`$
We define two decreasing filtrations on $`𝒲`$ numbered by nonnegative integers.
The $`T`$-filtration $`F_{}^T𝒲`$ is defined as follows. We ascribe to the elements of $``$ degree 1 and to $`t`$ degree 2. Then $`F_n^T𝒲`$ consists of elements of $`𝒲`$ having the leading term of total degree $`n`$.
The $`𝒫`$-filtration, $`F_{}^𝒫𝒲`$, is firstly defined on $`S(𝒫)S(𝒬)`$ by the subsets $`F_n^𝒫=S^n(𝒫)S(𝒫)S(𝒬)`$, $`n=0,1,\mathrm{}`$, and carried over to $`𝒲`$ via the Wick isomorphism.
We extend those filtrations to $`𝐖`$ in the natural way standing, for example, $`F_n^T𝐖=F_n^T𝒲\mathrm{\Lambda }()`$. We will use the following mnemonic notation. To point out, for example, that a section $`x𝐖`$ belongs to $`F_n^T𝐖`$ we write $`F^T(x)n`$.
We call a subsheaf $`𝒫`$ a $`\phi `$-null subsheaf if $`\phi `$ when restricted to $`𝒫`$ equals zero.
###### Proposition 2.3.
Let $`=𝒫𝒬`$ be a decomposition of $``$ with $`𝒫`$ a $`\phi `$-null subsheaf. Then
a) The Wick map $`𝐰:𝐒𝐖`$ has the following property: for $`a𝐒(𝒫)`$ and arbitrary $`c𝐒`$ one has $`ac=𝐰(a)𝐰(c)`$.
The filtrations on $`𝐖`$ have the properties:
b) for $`x,y𝐖`$, if $`F^𝒫(x)k`$ then $`F^𝒫(xy)k`$;
c) $`F^𝒫(\delta _{𝒫,𝒬}^1x)F^𝒫(x)`$;
d) $`F^T(\delta _{𝒫,𝒬}^1x)>F^T(x)`$.
###### Proof.
a) follows from the fact that $`𝒫`$ is a Lagrangean subsheaf of $``$ and the definition of the Wick map. b) follows from a), c) and d) are obvious. ∎
### 2.5 Good polarizations and connections
Let $`T_M`$ be the complexified tangent bundle over $`M`$, $`𝒯=T_M[[t]]`$, and $`\omega :𝒯^2C_M^{\mathrm{}}[[t]]`$ a nondegenerate symplectic form. It means that $`\omega =\omega _0+t\omega _1+\mathrm{}`$ where $`\omega _0`$ is nondegenerate and $`d\omega =0`$. Let $`\phi =\phi _0+t\phi _1+\mathrm{}`$ be the nondegenerate Poisson bracket inverse to $`\omega `$. It may be considered as the map $`\phi :(𝒯^{})^2C_M^{\mathrm{}}[[t]]`$, where $`𝒯^{}=T_M^{}[[t]]`$.
For any function $`fC_M^{\mathrm{}}[[t]]`$ denote by $`X_f`$ the corresponding Hamiltonian vector field on $`M`$, $`X_f=\phi (df,)`$. The map $`\varphi :𝒯^{}𝒯`$ defined by $`dfX_f`$ is an isomorphism of sheaves. Its inverse map is given by $`X_f\omega (,X_f)=df`$.
###### Definition 2.4.
a) A direct subsheaf $`𝒫𝒯`$ is called Lagrangean if $`\omega (𝒫,𝒫)=0`$ and $`rank(𝒫)=\frac{1}{2}dimM=n`$.
b) $`𝒫`$ is called integrable if locally there exist functions $`a_1,\mathrm{},a_n`$ such that the Hamiltonian vector fields $`X_{a_1},\mathrm{},X_{a_k}`$ form a local basis in $`𝒫`$.
c) A Lagrangean integrable $`𝒫`$ is called a polarization of $`\omega `$. In this case all the $`X_{a_i}`$, $`i=1,\mathrm{},n`$, pairwise commute, and $`da_i`$ form a local basis in $`𝒫^{}`$.
d) A polarization $`𝒫`$ is called good if there exist additional functions $`f_1,\mathrm{},f_n`$ such that the Hamiltonian vector fields $`X_{a_i}`$, $`X_{f_i}`$, $`i=1,\mathrm{},n`$, pairwise commute and form a local basis in $`𝒯`$.
###### Definition 2.5.
Let $`𝒫𝒯`$ be a Lagrangean subsheaf. We call a connection, $``$, on $`M`$ a $`𝒫`$-connection if
a) it preserves $`\omega `$ and is torsion free, i.e. is a symplectic connection, and
b) it preserves $`𝒫`$, i.e. $`(𝒫)𝒫\mathrm{\Lambda }^2(𝒯^{})`$, and is flat on $`𝒫`$ along $`𝒫`$, i.e. for any $`X,Y𝒫`$ one has $`(_X_Y_Y_X_{[X,Y]})(𝒫)=0`$.
###### Proposition 2.6.
Let $`𝒫𝒯`$ be a good polarization of $`\omega `$. Then, there exists a $`𝒫`$-connection on $`M`$.
###### Proof.
Let $`a_1,\mathrm{},a_{2n}`$ be functions such that $`X_i=X_{a_i}`$, $`i=1,\mathrm{},n`$, form a local basis in $`𝒫`$ and all the $`X_i`$ commute and form a local basis in $`𝒯`$. Let us set $`_{X_i}X_j=[X_i,X_j]=0`$. This rule defines a local connection.
Let us prove the invariance of $`\phi `$, i.e.
$$_z(\omega (x,y))=\omega (_zx,y)+\omega (x,_zy)\text{ for any }x,y,z𝒯.$$
But for $`x=X_i`$, $`y=X_j`$, $`z=X_k`$ this equation is equivalent to the relation
$$\phi (a_k,\phi (a_i,a_j))=\phi (\phi (a_k,a_i),a_j)+\phi (a_i,\phi (a_k,a_j))=0,$$
which holds by the Jacobi identity for the Poisson bracket $`\phi `$.
Let us proof that $``$ is torsion free. Since the torsion is $`C_M^{\mathrm{}}[[t]]`$ linear, it is enough to prove that
$$_{X_i}(Y_j)_{X_j}(X_i)[X_i,X_j]=0$$
for all pairs $`X_i,X_j`$ of our basis. But this follows from the definition of $``$ and from the fact that all $`X_i`$ pairwise commute. Similarly, from the fact that all $`X_i`$ pairwise commute follows that $``$ is flat.
That $``$ preserves $`𝒫`$ is obvious.
Since $`X_f𝒫`$ is equivalent to $`df𝒫^{}`$, it is obvious that $``$ has the property: for any $`f,gC_M^{\mathrm{}}[[t]]`$ such that $`X_f,X_g𝒫`$ one has
$`_{X_f}X_g=0.`$ (8)
Now, let us prove the existence of a global connection.
Let $`\{U_\alpha \}`$ is an open covering of $`M`$ such that on each $`U_\alpha `$ there is a $`𝒫`$-connection $`_\alpha `$. Then the differences $`_\alpha _\beta `$ form on $`U_\alpha U_\beta `$ a Ĉech cocycle $`\psi _{\alpha ,\beta }Hom(𝒯𝒯,𝒯)`$, $`\psi _{\alpha ,\beta }(X,Y)=_{\alpha X}Y_{\beta X}Y`$.
$`\psi _{\alpha ,\beta }`$ satisfy the following properties.
It follows from (8) that
$`\psi _{\alpha ,\beta }(X,Y)=0\text{for }X,Y𝒫.`$ (9)
Since all $`_\alpha `$ are torsion free, $`\psi _{\alpha ,\beta }`$ are symmetric. Since all $`_\alpha `$ preserve $`𝒫`$, one has $`\psi _{\alpha ,\beta }(X,Y)𝒫`$ for $`Y𝒫`$.
In addition, $`\psi _{\alpha ,\beta }`$ considered as elements from $`Hom(𝒯,Hom(𝒯,𝒯)`$, $`X\psi _{\alpha ,\beta }(X,)`$, belong to $`Hom(𝒯,sp(𝒯))`$ where $`sp(𝒯)`$ consists of endomorphisms of $`𝒯`$ preserving $`\omega `$ invariant.
Since all the properties above are $`C_M^{\mathrm{}}[[t]]`$ linear, it is possible to find tensors $`\psi _\alpha Hom(𝒯𝒯,𝒯)`$ satisfying the same properties and such that $`\psi _\alpha \psi _\beta =\psi _{\alpha ,\beta }`$. Then $`=_\alpha \psi _\alpha =_\beta \psi _\beta `$ is a globally defined connection. Flatness of $``$ on $`𝒫`$ along $`𝒫`$ follows from the fact that for all $`\alpha `$ $`\psi _\alpha (X,Y)=0`$ for $`X,Y𝒫`$, which follows from (9). $``$ is torsion free because all $`\psi _\alpha `$ are symmetric. So, $``$ satisfies the proposition. ∎
### 2.6 Fedosov’s construction
Let $`(M,\omega _0)`$ be a symplectic manifold. It is known that all equivalent classes of star-products on $`M`$ with the Poisson bracket $`\phi _0=\omega _0^1`$ can be obtained by the Fedosov method. According to this method, one constructs a flat connection, $`D`$, (called the Fedosov connection) in the Weyl algebra defined on the cotangent bundle with the help of relations (7), where $`\phi =\phi _0`$. The quantized algebra, $`𝒜`$, is realized as the subalgebra of flat sections of the Weyl algebra. The Weyl curvature of $`D`$, being a closed scalar two-form of the view $`\stackrel{~}{\theta }=\omega _0+t\omega _1+\mathrm{}`$ defines the Fedosov class
$`\theta (𝒜)={\displaystyle \frac{1}{t}}[\stackrel{~}{\theta }]{\displaystyle \frac{1}{t}}[\omega _0]+H^2(M,[[t]]).`$ (10)
It is also known that the correspondence $`𝒜[\theta (𝒜)]`$ is a bijection between the set of isomorphism classes of star-products on $`(M,\omega _0)`$ and the set $`\frac{1}{t}[\omega _0]+H^2(M,[[t]])`$ modulo the group of formal symplectomorphisms of $`M`$, \[Fe\], \[NT\], \[Xu\].
We will adapt the Fedosov method to construct a polarized star-product corresponding to a class $`\frac{1}{t}[\stackrel{~}{\theta }]\frac{1}{t}[\omega _0]+H^2(M,[[t]])`$ with a polarization of $`\stackrel{~}{\theta }`$. We will see that in the presence of polarization by realizing the Fedosov scheme there appears also a Wick curvature which differs from the Weyl curvature by $`t`$ times a half of the first Chern class of the polarization.
Let $`𝒫^{}𝒯`$ be a good polarization of $`\omega `$. It is easy to prove that there exists a complement Lagrangean subsheaf $`𝒬^{}𝒯`$ such that $`𝒯=𝒫^{}𝒬^{}`$ on $`M`$.
In the following we set $`=𝒯^{}`$ and consider the decomposition $`=𝒫𝒬`$, there $`𝒫`$ and $`𝒬`$ correspond to $`𝒫^{}`$ and $`𝒬^{}`$ by the isomorphism $`\varphi `$ between $`𝒯`$ and $``$. Note that $`𝒫`$ and $`𝒫^{}`$ are orthogonal to each other, so that we may write $`𝒫^{}=𝒫^{}`$.
Let $``$ be a $`𝒫^{}`$-connection on $`M`$. Then the induced connection $`:\mathrm{\Lambda }^1`$ on $``$ preserves $`𝒫`$, i.e. $`(𝒫)𝒫\mathrm{\Lambda }^2()`$, and is flat on $`𝒫`$ along $`𝒫^{}`$, i.e. for any $`X,Y𝒫^{}`$ one has $`(_X_Y_Y_X_{[X,Y]})(𝒫)=0`$.
$``$ gives a derivation of the Fedosov algebra $`𝐖=𝐖()`$, which is an extension of the de Rham differential on functions. Analogously, $``$ gives such derivations of the algebras $`T()\mathrm{\Lambda }()`$, $`S()\mathrm{\Lambda }()`$, and $`(S(𝒫)S(𝒬))\mathrm{\Lambda }()`$. These derivations commutes with the maps (5) and $`𝐰`$.
For convenience, we will mark the elements of the Fedosov algebra lying in $`1`$ by letters with hat over them ($`\widehat{x}`$), and by $`dx`$ we will denote the copy of $`\widehat{x}`$ lying in $`1\mathrm{\Lambda }^1`$.
It is easy to check that for $`\stackrel{~}{\delta }=\omega _{ij}\widehat{x^i}dx^j`$ one has
$`\delta ={\displaystyle \frac{1}{t}}\mathrm{ad}(\stackrel{~}{\delta }),`$
$`\stackrel{~}{\delta }^2=t\omega .`$ (11)
From the fact that the torsion of $``$ is equal to zero follows
$`(\stackrel{~}{\delta })=0.`$ (12)
Since $`^2`$ is a $`C_M^{\mathrm{}}[[t]]`$ linear derivation of degree $`0`$ preserving $`𝒫`$, there is an element $`R\rho _{𝒫,}(𝔰𝔭)\mathrm{\Lambda }^2()`$ such that $`^2=\frac{1}{t}\mathrm{ad}(R)`$. In particular, one has to be
$`F^𝒫(R)1.`$ (13)
According to Fedosov, \[Fe\], we also define $`R^F\rho _{}(𝔰𝔭)\mathrm{\Lambda }^2()`$ satisfying $`^2=\frac{1}{t}\mathrm{ad}(R^F)`$.
From (12) follows
$`\delta (R)=\delta (R^F)=0.`$ (14)
Following to Fedosov, we will consider connections on $`𝐖`$ of the form
$`D=+{\displaystyle \frac{1}{t}}\mathrm{ad}(\gamma ),\gamma 𝒲\mathrm{\Lambda }^1().`$ (15)
We define the Wick curvature of $`D`$ as
$`\mathrm{\Omega }_D=R+(\gamma )+{\displaystyle \frac{1}{t}}\gamma ^2.`$ (16)
According to Fedosov, we also define the Weyl (or Fedosov) curvature of $`D`$ as
$`\mathrm{\Omega }_D^F=R^F+(\gamma )+{\displaystyle \frac{1}{t}}\gamma ^2=Rt{\displaystyle \frac{1}{2}}tr(R)+(\gamma )+{\displaystyle \frac{1}{t}}\gamma ^2=\mathrm{\Omega }_Dt{\displaystyle \frac{1}{2}}tr(R),`$ (17)
where the second equality is due to Lemma 2.2.
###### Proposition 2.7.
tr(R) is a closed form polarized by $`𝒫`$. The class of $`tr(R)`$ in $`H^2(M,)`$ coincides with the first Chern class of $`𝒫`$, i.e. one has
$`[tr(R)]=c_1(𝒫)=c_1(𝒫_{t=0}).`$ (18)
###### Proof.
Since $`tr(R)`$ is a $`C_M^{\mathrm{}}[[t]]`$ valued 2-form, it follows from the Bianchi identity that it is a closed 2-form. From flatness of $`𝒫`$ along $`𝒫^{}`$ follows that the $`(1,1)`$ component of $`R`$ does not contain terms of the view $`R_{i,j}dy^idy^j`$ with $`dy^i,dy^j𝒬`$. It follows that $`𝒫`$ is a polarization of $`tr(R)`$.
It is easy to see that $``$ being restricted to $`𝒫`$ defines a connection on $`𝒫`$, $`_𝒫`$, and
$$_𝒫^2=\frac{1}{t}\mathrm{ad}R_{1,1},$$
where $`R_{1,1}`$ is the $`𝐰_{𝒫,𝒬}`$-degree $`(1,1)`$ component of the $`𝔰𝔭_{𝒫,𝒬}`$ valued 2-form $`R`$. Now (18) follows from the definition of the Chern classes. Since Chern classes are integer valued elements of $`H^2(M,)`$, i.e. belong to $`H^2(M,)`$, the independence of $`c_1(𝒫)`$ on the parameter $`t`$ is obvious. ∎
One has
$`D^2={\displaystyle \frac{1}{t}}\mathrm{ad}(\mathrm{\Omega }_D)={\displaystyle \frac{1}{t}}\mathrm{ad}(\mathrm{\Omega }_D^F).`$ (19)
Let us take $`\gamma `$ in the form
$`\gamma =\stackrel{~}{\delta }+r,r𝒲\mathrm{\Lambda }^1(),F^T(r)3.`$ (20)
Then the connection $`D`$ has the form
$`+\delta +{\displaystyle \frac{1}{t}}\mathrm{ad}(r).`$ (21)
Using (2.6) and (12), we obtain that its Wick curvature is
$`\mathrm{\Omega }_D=R+(\stackrel{~}{\delta }+r)+{\displaystyle \frac{1}{t}}(\stackrel{~}{\delta }+r)^2=\omega +\delta r+R+r+{\displaystyle \frac{1}{t}}r^2.`$ (22)
###### Proposition 2.8.
There exists an element $`r𝒲()\mathrm{\Lambda }^1()`$ such that
a) $`F^T(r)3`$;
b) $`F^𝒫(r)1`$;
c) the connection $`D=+\delta +\frac{1}{t}\mathrm{ad}(r)`$ is flat, i.e. $`D^2=0`$;
d) for its Wick curvature one has
$$\mathrm{\Omega }_D=\omega ;$$
e) for the class in $`H^2(M,[[t]])`$ of its Weyl curvature one has
$`[\mathrm{\Omega }_D^F]=[\mathrm{\Omega }_D]t{\displaystyle \frac{1}{2}}c_1(𝒫).`$ (23)
###### Proof.
First of all, we apply the Fedosov method, (\[Fe\], Theorem 5.2.2), to find $`r`$ satisfying d). According to (22), $`r`$ must obey the equation
$`\delta r=(R+r+{\displaystyle \frac{1}{t}}r^2)`$ (24)
Look for $`r`$ as the limit of the sequence, $`r=limr_k`$, where $`r_k𝒲()\mathrm{\Lambda }^1()`$, $`k=3,4,\mathrm{}`$, and $`F^T(r_kr_{k1})k`$. As in Lemma 5.2.3 of \[Fe\], using Proposition (2.3) d) and the fact that $`F^T(R)2`$, such $`r_k`$ can be calculated recursively:
$`r_3=\delta _{𝒫,𝒬}^1(R)`$
$`r_{k+3}=(r_3+\delta _{𝒫,𝒬}^1(r_{k+2}+{\displaystyle \frac{1}{t}}r_{k+2}^2)).`$ (25)
So, a) and d) are proven. e) follows from (17) and Proposition 2.7.
Let us prove that $`F^𝒫(r_k)1`$ for all $`k`$. That $`F^𝒫(r_3)1`$ follows from the fact that $`F^𝒫(R)1`$ and from Proposition (2.3) c). Suppose that $`F^𝒫(r_i)1`$ for $`i<k+3`$. Then $`F^𝒫(r_{k+2})1`$ because $``$ preserves $`𝒫`$. On the other hand, $`F^𝒫(r_{k+2}^2)1`$ because of Proposition (2.3) b), therefore from (2.6) follows that $`F^𝒫(r_{k+3})1`$ as well. So, we have that $`r`$ being the limit of the convergent sequence $`r_k`$ satisfy the conditions a), b), and d) of the proposition. c) obviously follows from d) and (19). ∎
Denote by $`𝒲_D`$ the subsheaf of $`𝒲`$ consisting of flat sections $`a`$, i.e. such that $`Da=0`$. Since $`D`$ is a derivation of $`𝐖`$, it is clear that $`𝒲_D`$ is a sheaf of subalgebras. Let $`\sigma =id(\delta \delta _{𝒫,𝒬}^1+\delta _{𝒫,𝒬}^1\delta )`$. Then, as follows from (6), $`\sigma :𝒲C_M^{\mathrm{}}[[t]]`$ is a projection, where $`C_M^{\mathrm{}}[[t]]`$ considers as the center of the algebra $`𝒲`$.
###### Proposition 2.9.
a) The map $`\sigma :𝒲_DC_M^{\mathrm{}}[[t]]`$ is a bijection.
b) The inverse map $`\tau :C_M^{\mathrm{}}[[t]]𝒲_D`$ has the form $`\tau (f)=f+\widehat{f}`$, there $`F^T(\widehat{f})>F^T(f)`$.
c) If $`df𝒫`$ then $`F^𝒫(\widehat{f})1`$.
d) If $`df𝒫`$ then $`\sigma (\tau (f)\tau (g))=fg`$ for any $`gC_M^{\mathrm{}}[[t]]`$.
###### Proof.
Again, we apply the Fedosov iteration procedure. According to \[Fe\], Theorem 5.2.4, we look for $`\tau (f)`$ as a limit, $`\tau (f)=lima_k`$, there $`a_k𝒲`$ can be calculated recursively:
$`a_0=f`$
$`a_{k+1}=a_0+\delta _{𝒫,𝒬}^1(a_k+{\displaystyle \frac{1}{t}}\mathrm{ad}r(a_k)).`$ (26)
Put $`\widehat{f}=\tau (f)a_0`$. As in \[Fe\], Theorem 5.2.4, one proves that such $`\tau (f)`$ and $`\widehat{f}`$ satisfy a) and b). Now observe that $`a_1a_0=\delta _{𝒫,𝒬}^1(1df)`$ and if $`df𝒫`$ then $`F^𝒫(a_1a_0)1`$. By induction we conclude that $`F^𝒫(a_ka_0)1`$ for all $`k1`$. So $`F^𝒫(aa_0)1`$ as well, which proves c).
Let us prove d). We have $`\tau (f)\tau (g)=f\tau (g)+\widehat{f}\tau (g)`$. Since by c) $`F^𝒫(\widehat{f})1`$, $`F^𝒫(\widehat{f}\tau (g))1`$ as well. It follows that $`\sigma (\widehat{f}\tau (g))=0`$ and $`\sigma (\tau (f)\tau (g))=\sigma (f\tau (g))=fg`$, because $`\sigma `$ is a $`C_M^{\mathrm{}}[[t]]`$ linear map and $`\sigma (\tau (g))=g`$. ∎
### 2.7 Polarized star-product
Let $`T_M`$ be the complexified tangent bundle over $`M`$, $`𝒯=T_M[[t]]`$, and $`\omega :𝒯^2C_M^{\mathrm{}}[[t]]`$ a nondegenerate symplectic form, $`\omega =\omega _0+t\omega _1+\mathrm{}`$. Let $`𝒫𝒯`$ be a polarization of $`\omega `$. We say that $`\omega `$ is a deformation of the symplectic form $`\omega _0`$ and $`𝒫`$ is a deformation (or extension) of $`𝒫_0=𝒫/t𝒫`$. It is clear that $`𝒫_0`$ is a polarization of $`\omega _0`$.
Let $`\stackrel{~}{\omega }`$ be another deformation of $`\omega _0`$. We say that $`\omega `$ and $`\stackrel{~}{\omega }`$ are equivalent if they define the same class in $`\omega _0+H^2(M,t[[t]])`$.
Let $`𝒫`$ be a polarization of $`\omega `$ which is an extension of $`𝒫_0`$. It is easy to show that if $`\stackrel{~}{\omega }=\omega +d\lambda `$ and $`X_\lambda `$ is a vector field corresponding to $`\lambda `$ then $`\stackrel{~}{𝒫}=\mathrm{exp}(L_{X_\lambda })𝒫`$ is a polarization of $`\stackrel{~}{\omega }`$ extending $`𝒫_0`$, which we call an equivalent polarization. Thus, it is meaningful to speak about a polarization of a class of $`\omega _0+H^2(M,t[[t]])`$. It is the equivalence class of extensions of $`𝒫_0`$.
As we have mentioned at the beginning of Section 2.6, the classification theorem for star-products say that all star-products on $`C_M^{\mathrm{}}[[t]]`$ with the Poisson bracket inverse to $`\omega _0`$ are parameterized by elements
$`\theta ={\displaystyle \frac{1}{t}}[\stackrel{~}{\theta }]{\displaystyle \frac{1}{t}}[\omega _0]+H^2(M,[[t]])`$ (27)
where $`\stackrel{~}{\theta }`$ is the Weyl scalar curvatures of the Fedosov connections.
###### Theorem 2.10.
Let $`(M,\omega _0,)`$ be a symplectic manifold. Let $`\omega `$ be a deformation of $`\omega _0`$ and $`𝒫`$ a good polarization of $`\omega `$. Let $`𝒪`$ be the sheaf of functions from $`C_M^{\mathrm{}}[[t]]`$ constant along $`𝒫`$.
Then, there exists a star-product, $``$, on $`C_M^{\mathrm{}}[[t]]`$ having the Fedosov class $`\frac{1}{t}[\omega ]`$ and satisfying the property:
$`fg=fg\text{for all }f𝒪.`$ (28)
###### Proof.
Let $`\alpha `$ be a 2-form representing $`\frac{1}{2}c_1(𝒫)`$ and having $`𝒫`$ as its polarization (see Proposition (2.7)). Construct the star-product for $`\omega ^{}=\omega +t\alpha `$ and $`𝒫`$ as in Proposition 2.9. According to Proposition 2.8 d), e), the Wick curvature corresponding to that star-product is equal to $`\omega ^{}`$ and its Fedosov class is equal to $`\frac{1}{t}[\omega ]`$.
Property (28) follows from Proposition 2.9 d). ∎
## 3 Strong polarizations
### 3.1 Definition and deformation of a strong polarization
Let $`(M,\omega _0)`$ be a symplectic manifold of dimension $`2n`$, $`T`$ the tangent bundle on $`M`$ and $`T_{}=T^{}`$ its complexification. Let $`PT^{}`$ be a polarization of $`\omega _0`$, i.e. $`P`$ is a Lagrangean subbundle and locally there exist functions $`a_1,\mathrm{},a_n`$ such that $`da_1,\mathrm{},da_n`$ give a local basis for $`P^{}(T^{})^{}`$. Let the triple $`(M,\omega ,𝒫)`$ be a deformation of the triple $`(M,\omega _0,P)`$, i.e. $`𝒫`$ is a polarization of symplectic form $`\omega =\omega _0+t\omega _1+\mathrm{}`$, $`𝒫_0=P`$. We are going to give simple conditions on $`P`$ and $`𝒫`$ which guarantee for them to be good polarizations.
Denote by $`𝒪=𝒪_PC_M^{\mathrm{}}`$ the subsheaf of functions constant along the polarization $`P`$. In the obvious way one can form an analog of the Dolbeault complex, which we also call Dolbeault: $`(^{}P^{},\overline{d})`$, where $`\overline{d}`$ is the differential along $`P`$. By definition, $`𝒪`$ consists of functions $`f`$ such that $`\overline{d}f=0`$. Note that the Dolbeault complex is meaningful for any involutive $`P`$, i.e. when $`[P,P]P`$, and for any involutive defirmation $`𝒫`$.
###### Definition 3.1.
We call a polarization $`P`$ strong if for any point $`mM`$ there exists a neighborhood $`U`$ of $`m`$ such that the complex of sections $`(\mathrm{\Gamma }(U;^{}P^{}),\overline{d})`$ is exact and gives a resolution of $`\mathrm{\Gamma }(U;𝒪_P)`$.
We give the same definition for a deformed polarization $`𝒫`$ of $`\omega `$.
###### Proposition 3.2.
Let $`P`$ be a strong polarization of $`\omega _0`$. Let $`\omega `$ be a symplectic deformation of $`\omega _0`$ and $`𝒫𝒯=T^{}[[t]]`$ a deformation of $`P`$, $`𝒫_0=P`$, which is Lagrangean with respect to $`\omega `$ and involutive, i.e. $`[𝒫,𝒫]𝒫`$. Then $`𝒫`$ is a strong polarization of $`\omega `$.
###### Proof.
First of all, let us prove that $`𝒫`$ is a polarization, i.e. $`𝒫`$ is integrable. Since $`𝒫`$ is involutive, the Dolbeault complex $`(^{}𝒫^{},\overline{d}_t)`$ is meaningful, where $`\overline{d}_t=\overline{d}+t\overline{d}^{}`$ is a deformation of the differential $`\overline{d}`$ from the Dolbeault complex for $`P`$. Let $`U`$ be a neighborhood such that the complex $`(\mathrm{\Gamma }_U(^{}P^{}),\overline{d})`$ is exact. It is enough to prove that if $`f`$ is a function on $`U`$ satisfying $`\overline{d}f=0`$ then there is an extension $`f_t=f+tf^{}C_U^{\mathrm{}}[[t]]`$ such that $`\overline{d}_tf_t=0`$.
Suppose that we have already found a series $`f^{(n1)}=f+tf_1+\mathrm{}+t^{n1}f_{n1}`$ such that $`\overline{d}_tf^{(n1)}=0`$ mod $`t^n`$. Then $`\overline{d}_tf^{(n1)}=t^n\gamma _n`$ mod $`t^{n+1}`$. Applying $`\overline{d}_t`$ to the both sides of this equation and dividing by $`t^n`$ we obtain $`\overline{d}_t\gamma _n=0`$ mod $`t`$. Hence, $`\overline{d}\gamma _n=0`$. So, $`\gamma _n`$ is a closed 1-form in the Dolbeault complex corresponding to $`P`$, and one can find a function $`f_n`$ such that $`\overline{d}f_n=\gamma _n`$. It is clear that if we put $`f^{(n)}=f^{(n1)}+t^nf_n`$, we obtain $`\overline{d}_tf^{(n)}=0`$ mod $`t^{n+1}`$. In this way we construct step-by-step the series $`f_t`$ such that $`\overline{d}_tf_t=0`$ and prove that $`𝒫`$ is a polarization of $`\omega `$.
That $`𝒫`$ is strong is obvious because of the upper-semicontinuity argument. ∎
Proposition 3.2 shows that the deformed polarization $`𝒫`$ is strong if $`P=𝒫_0`$ is strong. There is the following proposition provides sufficient conditions for a Lagrangean subbundle $`P`$ to be a strong polarization.
###### Proposition 3.3.
Let $`(M,\omega _0)`$ be symplectic manifold with a Lagrangean subbundle $`PT_M^{}`$. Suppose $`P`$ satisfies the following conditions:
i) the subsheaf $`P\overline{P}`$ is a subbundle in $`T_M^{}`$;
ii) $`P+\overline{P}`$ is involutive.
Then $`P`$ is a strong polarization of $`\omega `$.
###### Proof.
The integrability of $`P`$ follows from the Frobenius-Nirenberg theorem. The exactness of the Dolbeault complex for $`P`$ is a theorem due to Rawnsley (\[Ra\], Thm. 2). ∎
Note that in the analytic case, i.e. when $`M`$ is an analytic manifold and $`\omega _0`$ is an analytic form, a Lagrangean involutive analytic subbundle $`P`$ is strong. This follows from the Frobenius theorem.
### 3.2 Local properties of strong polarizations
From now on we will denote by $`P`$ either a polarization of $`\omega _0`$ or a deformed polarization of $`\omega `$, and $`𝒪=𝒪_P`$ is the sheaf of functions constant along $`P`$. By $`\{,\}`$ we will denote the Poisson bracket corresponding to $`\omega _0`$ or $`\omega `$.
If $`P`$ is integrable, there exist, locally, functions $`a_i`$, $`i=1,\mathrm{}n`$, such that $`\{a_i,a_j\}=0`$ and $`X_{a_i}=\{a_i,\}`$ form a local basis in $`P`$. Denote by $`\overline{d}^i`$ the dual basis in $`P^{}`$. It is clear that $`\overline{d}^i`$ is a closed form in the Dolbeault complex and
$`\overline{d}={\displaystyle \underset{i}{}}\{a_i,\}\overline{d}^i.`$ (29)
###### Lemma 3.4.
Suppose $`P`$ is a strong polarization. Then, locally, there exist functions $`f_i`$ such that $`\{f_i,a_j\}=\delta _{ij}`$.
###### Proof.
For any $`i`$, consider the closed form $`\overline{d}^i`$. By Definition 3.1, there is a function $`f_i`$ such that $`\overline{d}(f_i)=\overline{d}^i`$. By (29), $`f_i`$ satisfies the lemma. ∎
Let $`F_𝒪=\{bC_M^{\mathrm{}};\{b,𝒪\}𝒪\}`$. It follows from the Jacobi identity that $`\{F_𝒪,F_𝒪\}F_𝒪`$. Hence, $`F_𝒪`$ is a Lie algebra acting on $`𝒪`$ by derivations. The kernel of this action is equal to $`𝒪`$ itself.
Note that $`F_𝒪`$ is a locally free $`𝒪`$ module. Locally, it is freely generated as an $`𝒪`$ module by functions $`1`$, $`f_1,\mathrm{},f_n`$, where $`f_i`$ are as in Lemma 3.4.
Let us denote by $`\mathrm{\Omega }^1=\mathrm{\Omega }_𝒪^1`$ the $`𝒪`$-submodule of $`(T^{})^{}`$ generated by the sheaf $`d𝒪`$, i.e. $`\mathrm{\Omega }^1=𝒪d𝒪`$. Let $`T_𝒪=F_𝒪/𝒪`$. It is a locally free $`𝒪`$ module and is a sheaf of Lie algebras acting on $`𝒪`$ with the local basis $`\{f_i,\}`$. It is easy to see that $`\mathrm{\Omega }^1`$ can be naturally identified with $`Hom_𝒪(T_𝒪,𝒪)`$ and $`d^i=da_i`$ form the dual basis to $`\{f_i,\}`$.
We form in the obvious way an analog of de Rham complex: $`(_{𝒪_P}^{}\mathrm{\Omega }^1,d)`$, which we call $`𝒪_P`$-de Rham complex.
###### Lemma 3.5.
Let $`P`$ be a strong polarization. Then for any point of $`M`$ there exists a neighborhood $`U`$ such that the complex $`(\mathrm{\Gamma }_U(_{𝒪_P}^{}\mathrm{\Omega }^1),d)`$ is exact and gives a resolution for $``$.
###### Proof.
Let us choose $`U`$ such that the Dolbeault complex and the usual $`C^{\mathrm{}}`$-de Rham complex $`(^{}(T^{})^{},d_{dR})`$ are exact. On the de Rham complex $`(^{}(T^{})^{},d_{dR})`$ consider the Hodge filtration generated by $`P^{}(T^{})^{}`$. Considering the “stupid” filtration on $`(_{𝒪_P}^{}\mathrm{\Omega }^1,d)`$ we have that the natural inclusion $`(_{𝒪_P}^{}\mathrm{\Omega }^1,d)(^{}(T^{})^{},d_{dR})`$ is a filtered quasiisomorphism, i.e. induces a quasiisomorphism on the associated graded complexes. This follows from the fact that, since $`P`$ is strong, the Dolbeault complex is exact. Now our lemma follows from the exactness of $`(^{}(T^{})^{},d_{dR})`$. ∎
In the following, saying that the Dolbeault or $`𝒪`$-de Rham complex of sheaves is exact we mean that for any point of $`M`$ there exists a neighborhood $`U`$ such that the corresponding complex of sections over $`U`$ is exact.
###### Lemma 3.6.
Let $`P`$ be a strong polarization. Then $`f_i`$ in Lemma 3.4 can be chosen in such a way that $`\{f_i,f_j\}=0`$.
###### Proof.
By Lemma 3.4 one can choose $`g_i`$ such that $`\{g_i,a_j\}=\delta _{ij}`$. For any functions $`b_1,\mathrm{},b_n`$ from $`𝒪`$ the functions $`f_i=g_i+b_i`$ also satisfy Lemma 3.4. Let us find $`b_i`$ such that $`\{f_i,f_j\}=0`$. We find $`b_i`$ from the equation
$$\{g_i+b_i,g_j+b_j\}=\{g_i,g_j\}+\{g_i,b_j\}\{g_j,b_i\}=0.$$
Note that $`g=_{i,j}\{g_i,g_j\}d^id^j`$ is a closed 2-form of the $`𝒪_P`$-de Rham complex and the equation can be rewritten as
$$db=g,$$
where $`b=_ib_id^i`$. By Lemma 3.5, such $`b`$ exists. ∎
Lemma 3.6, in particular, means that locally there exists a polarization, $`P^{}`$, of $`\omega `$ complement to $`P`$. This polarization is defined as annihilator of $`df_1,\mathrm{},df_n`$. It is obvious that locally the form $`\omega `$ can be written as $`\omega =_idf_ida_i`$. One can prove that $`P^{}`$ is also a strong polarization.
###### Proposition 3.7.
Let $`P`$ be a strong polarization. Then
a) $`P`$ is a good polarization and
b) the sequence
$`0𝒪_PF_𝒪\stackrel{\pi }{}T_𝒪0`$ (30)
locally splits as an exact sequence of $`𝒪`$ modules and Lie algebras.
###### Proof.
a) follows from the existence of functions $`f_i`$ satisfying Lemma 3.6.
b) One needs to construct, locally, a Lie algebra morphism $`s:T_𝒪F_𝒪`$ such that $`\pi s=id`$. Let us choose functions $`a_i`$ and $`f_i`$, $`i=1,\mathrm{},n`$, as in Lemma 3.6 and assign $`\{f_i,\}f_i`$. Since the elements $`X_{f_i}=\{f_i,\}`$ form a local basis in $`T_𝒪`$, this defines a map of $`𝒪_P`$ modules, $`s:T_𝒪F_𝒪`$. One has $`[X_{f_i},X_{f_j}]=0`$ and also $`\{f_i,f_j\}=0`$, which proves that $`s`$ is a Lie algebra morphism. ∎
###### Proposition 3.8.
Let $`(M,\omega ,𝒫)`$ and $`(M,\omega ^{},𝒫^{})`$ be two deformations of $`(M,\omega _0,P)`$. Then any point of $`M`$ has an open neighborhood, $`U`$, such that there exists a formal automorphism, $`B`$, of $`U`$ identical modulo $`t`$ that transforms $`(U,\omega ,𝒫)`$ to $`(U,\omega ^{},𝒫^{})`$.
###### Proof.
For any point of $`mM`$ there is an open neighborhood, $`V`$, such that on it $`\omega ^{}=\omega +td\gamma `$, where $`\gamma `$ is an 1-form. Let $`X_\gamma `$ be the vector field on $`V`$ corresponding to $`\gamma `$ by the isomorphism $`T_M^{}T_M`$ determined by $`\omega `$. Then, $`\mathrm{exp}(tX_\gamma )`$ is a formal automorphism of $`V`$ which transforms $`\omega `$ to $`\omega ^{}`$. Hence, one may suppose in the proposition that $`\omega =\omega ^{}`$ and $`𝒫`$ and $`𝒫^{}`$ are two polarizations of the same form $`\omega `$.
Let $`U`$ be a neighborhood of $`m`$ such that over $`U`$ the Dolbeault complexes corresponding to both $`𝒫`$ and $`𝒫^{}`$ are exact. Let functions $`a_i`$, $`a_i^{}`$ on $`U`$ be such that $`da_i`$ and $`da_i^{}`$ form local basises in $`𝒫^{}`$ and $`𝒫^{{}_{}{}^{}}`$, respectively. Suppose $`𝒫=𝒫^{}`$ mod $`t^k`$. Then, we may assume that $`a_i=a_i^{}`$ mod $`t^k`$, i.e. $`a_i^{}=a_i+t^kb_i`$ for some functions $`b_i`$ on $`U`$. Let $`\{,\}`$ denote the Poisson bracket inverse to $`\omega `$. One has $`\{a_i,a_j\}=\{a_i^{},a_j^{}\}=0`$, hence $`\{a_i,b_j\}+\{b_i,a_j\}=0`$. It follows from the exactness of Dolbeault complex that there exists a function $`g`$ on $`U`$ such that $`\{a_i,g\}=b_i`$ for all $`i`$. Let $`X_g`$ be the Hamiltonian vector field on $`U`$ corresponding to $`g`$. Then, the automorphism $`\mathrm{exp}(tX_g)`$ of $`U`$ leaves $`\omega `$ on the place and transforms $`𝒫`$ to a polarization, $`𝒫^{\prime \prime }`$, such that $`𝒫^{\prime \prime }=𝒫^{}`$ mod $`t^{k+1}`$. Proceeding by iteration proves the proposition. ∎
This proposition shows that any deformation $`(M,\omega ,𝒫)`$ of $`(M,\omega _0,P)`$ is, locally, isomorphic to the trivial deformation $`(M,\omega _0,P[[t]])`$.
### 3.3 Local properties of polarized star-products
Let $`(M,\omega _0)`$ be a symplectic manifold with a strong polarization $`P`$. Let $`\omega `$ be a symplectic deformation of $`\omega _0`$ and $`𝒫`$ a polarization of $`\omega `$ which is a deformation of $`P`$. As is proven in the previous section, $`𝒫`$ is strong and thus a good polarization. Let $`𝒪`$ be the sheaf of functions constant along $`𝒫`$ and $`(𝒜,𝒪)`$ a polarized star-product constructed in Section 2.
Our immediate goal is to prove analogs of Lemmas 3.4 and 3.6 for the commutator $`\frac{1}{t}[,]`$ in $`𝒜`$.
###### Lemma 3.9.
Let $`𝒫`$ be a strong polarization of $`\omega `$. Let $`a_1,\mathrm{},a_n𝒪`$ be functions on an open set $`UM`$ such that $`da_1,\mathrm{},da_n`$ form a local basis in $`𝒫^{}`$. Then, locally, there exist functions $`\widehat{f}_1,\mathrm{},\widehat{f}_n`$ such that $`\frac{1}{t}[\widehat{f}_i,a_j]=\delta _{ij}`$.
###### Proof.
The operators $`\frac{1}{t}[a_i,]`$ are pairwise commuting derivations of $`𝒜`$ restricted to an open set $`U`$. Let $`L`$ be the free $`C_M^{\mathrm{}}[[t]]`$ module over $`U`$ spanned on $`[a_i,]`$. Denote by $`\widehat{d}^iL^{}`$ the basis dual to $`\frac{1}{t}[a_i,]`$. Let us define in the obvious way the complex $`(_{C_M^{\mathrm{}}[[t]]}^{}L^{},\widehat{d})`$, where
$`\widehat{d}={\displaystyle \underset{i}{}}[a_i,]\widehat{d}^i.`$ (31)
Since $`\frac{1}{t}[a_i,]=\{a_i,\}+o(t)`$ and $`\mathrm{rank}L=\mathrm{rank}P`$, this complex is a deformation of the Dolbeault complex $`(_{C_M^{\mathrm{}}[[t]]}^{}P^{},\overline{d})`$ over $`U`$ with $`\overline{d}`$ defined by (29). Since $`𝒫`$ is strong, the Dolbeault complex is exact, therefore the complex $`(_{C_M^{\mathrm{}}[[t]]}^{}L^{},\widehat{d})`$ is exact, too.
Now we proceed as in the proof of Lemma 3.4. Namely, for any $`i`$, consider the closed form $`\widehat{d}^i`$. Since the complex $`(_{C_M^{\mathrm{}}[[t]]}^{}L^{},\widehat{d})`$ is exact, there is a function $`\widehat{f}_i`$ such that $`\widehat{d}(\widehat{f}_i)=\widehat{d}^i`$. By (31), $`f_i`$ satisfies the lemma. ∎
Let
$`F(𝒜)=\{bC_M^{\mathrm{}}[[t]];[b,𝒪]𝒪\}.`$ (32)
It is clear that $`F(𝒜)`$ is a sheaf of Lie algebras on $`M`$ with the bracket $`\frac{1}{t}[,]`$. It is a locally free $`𝒪`$ module with the local basis consisting of functions $`1`$, $`\widehat{f}_j`$. $`𝒪`$ sits in $`F(𝒜)`$ as a commutative Lie subalgebra The sheaf of Lie algebras $`F(𝒜)/𝒪`$ is also a locally free $`𝒪`$-module isomorphic to the sheaf $`T_𝒪=F_𝒪/𝒪`$ from the previous section: locally, the isomorphism is given by $`\frac{1}{t}[\widehat{f},]\{f_i,\}`$.
###### Lemma 3.10.
By hypothesis of Lemma 3.9, $`\widehat{f}_i`$ can be chosen in such a way that $`\frac{1}{t}[\widehat{f}_i,\widehat{f}_j]=0`$.
###### Proof.
The same as of Lemma 3.6. ∎
Lemma 3.10, in particular, means that locally there exists a “complement” commutative subalgebra in $`𝒜`$. This subalgebra is generated by functions $`\widehat{f}_i`$. But that “complement” subalgebra may not be an algebra of functions constant on a polarization.
###### Proposition 3.11.
Let $`𝒫`$ be a strong polarization of $`\omega `$ and $`(𝒜,𝒪)`$ the corresponding polarized quantization. Then the sequence
$`0𝒪F(𝒜)\stackrel{\pi }{}T_𝒪0`$ (33)
locally splits as an exact sequence of $`𝒪`$ modules and Lie algebras.
###### Proof.
The same as of Proposition 3.7. ∎
###### Proposition 3.12.
Let $`(M,\omega ,𝒫)`$, $`(M,\omega ^{},𝒫^{})`$ be two deformations of $`(M,\omega _0,P)`$ and $`(𝒜,𝒪)`$, $`(𝒜^{},𝒪^{})`$ corresponding polarized star-products. Then any point of $`M`$ has an open neighborhood, $`U`$, such that over $`U`$ these star-products are equivalent. Moreover, if $`𝒪=𝒪^{}`$ then the equivalence can be chosen to be identity on $`𝒪`$.
###### Proof.
By Proposition 3.8 one can suppose that $`\omega =\omega ^{}=\omega _0`$ and $`𝒫=𝒫^{}=P[[t]]`$. So, $`𝒪=𝒪^{}=𝒪_0[[t]]`$, where $`𝒪_0`$ is the algebra of functions constant along $`P`$.
Let $`mM`$ and $`U`$ be a sufficiently small neighborhood of $`m`$. Let $`a_i`$ be functions on $`U`$ such that $`da_i`$ form a local basis in $`P^{}`$ and $`f_i`$ be functions such that $`\{f_i,a_j\}=\delta _{ij}`$, where $`\{,\}`$ is the Poisson bracket inverse to $`\omega _0`$. Let $`\mu =_{i0}t^i\mu _i`$ and $`\mu ^{}=_{i0}t^i\mu _i^{}`$ be multiplications in $`C_M^{\mathrm{}}[[t]]`$ over $`U`$ corresponding to $`𝒜`$ and $`𝒜^{}`$ and consisting of bidifferential operators. So, $`\mu _0(a,b)=\mu _0^{}(a,b)=ab`$ and one may suppose that $`\mu _1(a,b)=\mu _1^{}(a,b)=\frac{1}{2}\{a,b\}`$.
Suppose that there exists a differential operator $`B^{(n)}=1+\mathrm{}+t^nB_n`$ which transforms $`\mu `$ to $`\mu ^{}`$ modulo $`t^{n+1}`$ and such that $`B_k`$, $`k=1,\mathrm{},n`$, take $`𝒪`$ to zero. This means, in particular, that we may suppose that $`\mu =\mu ^{}`$ modulo $`t^{n+1}`$, so that
$$\mu =\mu _0+\mathrm{}t^n\mu _n+\mu _{n+1}+\mathrm{},$$
$$\mu ^{}=\mu _0+\mathrm{}t^n\mu _n+\mu _{n+1}^{}+\mathrm{}.$$
We are going to prove that there exists a differential operator of the form $`B=1+t^nX+t^{n+1}B_{n+1}`$, where $`X`$ is a vector field on $`U`$, which transforms $`\mu ^{}`$ to $`\mu `$ modulo $`t^{n+2}`$ and such that $`X(𝒪)=B_{n+1}(𝒪)=0`$.
It is easy to check that $`c=\mu _{n+1}^{}\mu _{n+1}`$ is a Hochschild cocycle in the algebra $`C_U^{\mathrm{}}`$. Hence, $`\nu (a,b)=c(a,b)c(b,a)`$ is a bivector field and there is the representation: $`\mu _{n+1}^{}=\mu _{n+1}+\nu +\delta _H(B_{n+1}^{})`$, where $`B_{n+1}`$ is a differential operator and $`\delta _H`$ the Hochschild differential.
It is also easy to check that $`[[\mu _1,\nu ]]=0`$, where $`[[,]]`$ denotes the Schouten bracket of polyvector fields. So, there exists a vector field $`Y`$ such that $`[[\mu _1,Y]]=\nu `$. This means that for any $`f,gC_U^{\mathrm{}}`$
$$Y(\mu _1(f,g))\mu _1(Yf,g)\mu _1(f,Yg)=\nu (f,g).$$
Recall now that $`\mu _1(a,b)=\nu (a,b)=0`$ for $`a,b𝒪`$. It follows that for any $`a,b𝒪`$
$$\mu _1(Ya,b)+\mu _1(a,Yb)=0.$$
In particular, we have $`\mu _1(Ya_i,a_j)+\mu _1(a_i,Ya_j)=0`$ for any $`a_i,a_j`$ of our basis. It follows from the exactness of Dolbeault complex that there exists a function $`bC_U^{\mathrm{}}`$ such that $`Ya_i=\mu _1(b,a_i)`$ for all $`i`$. Denote $`X_b=\{b,\}`$, the Hamiltonian vector field corresponding to $`b`$. Put $`X=YX_b`$. Then, since $`[[\mu _1,X_b]]=0`$, we get $`[[\mu _1,X]]=\nu `$. Since $`Ya_i=X_ba_i`$ for all $`i`$, we get $`X(𝒪)=0`$.
If we transform the multiplication $`\mu ^{}`$ by the operator $`1+t^nX`$, we obtain the multiplication $`\mu ^{\prime \prime }`$ of the form $`\mu ^{\prime \prime }=\mu _0+\mathrm{}t^n\mu _n+\mu _{n+1}^{\prime \prime }+\mathrm{}`$, where $`\mu _{n+1}^{\prime \prime }=\mu _{n+1}+\delta _HB_{n+1}^{}`$.
Since $`\mu _{n+1}^{\prime \prime }(a,b)=\mu _{n+1}(a,b)=0`$ for all $`a,b𝒪`$, we have that $`B_{n+1}^{}`$ is a derivation from $`𝒪_U`$ to $`C_U^{\mathrm{}}`$. This derivation can be extended to a derivation of $`C_U^{\mathrm{}}`$. Indeed, put $`B_{n+1}^{}(a_i)=g_i`$. Then the operator $`D=g_i\{f_i,\}`$ is such an extension. Now put $`B_{n+1}=B_{n+1}^{}D`$. One has $`\delta _HB_{n+1}=\delta _HB_{n+1}^{}`$. It is easy to see that the operator $`1+t^{n+1}B_{n+1}`$ transforms $`\mu ^{\prime \prime }`$ to $`\mu `$ modulo $`t^{n+2}`$. It follows that the operator $`(1+t^{n+1}B_{n+1})(1+t^nX)`$ transforms $`\mu ^{}`$ to $`\mu `$ modulo $`t^{n+2}`$. ∎
Let $`(M,\omega _0,P)`$ be a polarized symplectic manifold with $`P`$ a strong polarization. Let $`U`$ be a sufficiently small neighborhood of a point of $`M`$. Let $`a_i`$ be functions on $`U`$ such that $`da_i`$ form a local basis in $`P^{}`$ and $`f_i`$ be functions such that $`\{f_i,a_j\}=\delta _{ij}`$, where $`\{,\}`$ is the Poisson bracket inverse to $`\omega _0`$. Then the bivector field
$$\pi =\underset{i}{}X_{a_i}X_{f_i}$$
represent this Poisson bracket. Since all $`X_{a_i}`$ and $`Y_{f_i}`$ pairwise commute, one can construct the Moyal star-product of the form
$`fg=\mu _0\mathrm{exp}({\displaystyle \frac{t}{2}}X_{a_i}X_{f_i})(fg),`$ (34)
where $`\mu _0`$ is the usual multiplication. This star-product is obviously polarized, and we call it the polarized Moyal star-product.
###### Corollary 3.13.
Any polarized star-product is locally equivalent to a polarized Moyal star-product.
###### Proof.
Follows from the previous proposition. ∎
## 4 A formula for polarized quantization
### 4.1 $`𝒪`$-extension of $`T_𝒪`$
The sequences (30) and (33) are examples of an $`𝒪`$-extension of $`T_𝒪`$. In the following we suppose that $`𝒪`$ is the sheaf of functions constant along a strong polarization and $`T_𝒪=Hom_𝒪(\mathrm{\Omega }_𝒪^1,𝒪)`$. As we have seen, $`T_𝒪`$ is a locally free $`𝒪`$ module of derivations of $`𝒪`$.
Following \[BK\], \[BB\], we give the following
###### Definition 4.1.
An $`𝒪`$-extension of $`T_𝒪`$ is a locally free sheaf $`\stackrel{~}{T}`$ of $`𝒪`$-modules equipped with a structure of a sheaf of Lie algebras over $`[[t]]`$ (with the bilinear operation denoted, as usual, by $`[,]`$), a section $`𝐜`$ of the center of $`\stackrel{~}{T}`$, and a surjective $`𝒪`$-linear map $`\sigma :\stackrel{~}{T}T_𝒪`$ (the anchor map) which is a Lie algebra homomorphism, whose kernel (which is a sheaf of $`𝒪`$-Lie algebras) is isomorphic to $`𝒪`$ (the latter equipped with the trivial Lie bracket). In addition the Leibnitz rule holds: for $`f𝒪`$, $`\tau _1,\tau _2\stackrel{~}{T}`$, $`[\tau _1,f\tau _2]=f[\tau _1,\tau _2]+\sigma (\tau _1)(f)\tau _2`$. Thus, there is an exact sequence (of $`𝒪`$-modules and Lie algebras)
$`0𝒪\stackrel{i}{}\stackrel{~}{T}\stackrel{\sigma }{}T_𝒪0,`$ (35)
where $`i(f)=f𝐜`$, $`f𝒪`$.
Just as in Proposition 3.7, one shows that the extension (35) admits local splittings which are $`𝒪`$-linear Lie algebra homomorphisms. Such a splitting is called a flat connection on $`\stackrel{~}{T}`$.
### 4.2 Characteristic class of an $`𝒪`$-extension of $`T_𝒪`$
To each $`\stackrel{~}{T}`$ one associates the cohomology class $`c(\stackrel{~}{T})H^1(M;\mathrm{\Omega }_𝒪^{1,cl})=H^2(M,[[t]])`$ as follows.
Let $`\{U_\alpha \}`$ be a sufficiently fine open covering of $`M`$, so that there are flat connections $`s_\alpha `$ on $`\stackrel{~}{T}|_{U_\alpha }`$. If $`U_{\alpha \beta }\stackrel{def}{=}U_\alpha U_\beta \mathrm{}`$, the difference $`s_\alpha s_\beta `$ gives rise to the section $`c_{\alpha ,\beta }`$ (defined over $`U_{\alpha \beta }`$) of $`\mathrm{\Omega }_𝒪^{1,cl}=\underset{¯}{\mathrm{Hom}}_𝒪(T_𝒪,𝒪)`$ by the formula $`i(c_{\alpha ,\beta }(\xi ))=s_\alpha (\xi )s_\beta (\xi )`$. The collection $`\{c_{\alpha ,\beta }\}`$ is a degree one Ĉech cochain with coefficients in the sheaf $`\mathrm{\Omega }_𝒪^{1,cl}`$. It is, in fact, a cocycle, whose cohomology class, denoted $`c(\stackrel{~}{T})`$ is independent of the choice of local flat connections.
The above construction recovers the Chern class of a complex line bundle. Namely, suppose that $`L`$ is a locally free sheaf of $`𝒪`$-modules of rank one. Let $`\stackrel{~}{T}_{}`$ denote the sheaf of differential operators on $`L`$ of order (at most) one. The sheaf $`\stackrel{~}{T}_{}`$ is equipped with the left $`𝒪`$-module structure, the Lie bracket given by the commutator, the central section – the identity operator (so that the map $`i`$ is simply the inclusion of operators of order zero). The principal symbol map serves as the anchor map $`\sigma `$. These data exhibit $`\stackrel{~}{T}_{}`$ as an $`𝒪`$-extension of $`T_𝒪`$.
On the other hand, the sheaf $`_𝒪C_M^{\mathrm{}}`$ is the sheaf of $`C^{\mathrm{}}`$ sections of a complex line bundle $`L`$. It is easy to see that the class $`c(\stackrel{~}{T}_{})`$ coincides with the the first Chern class of $`L`$.
For example, since $`det𝒫^{}=\mathrm{\Omega }_𝒪^n_𝒪C_M^{\mathrm{}}`$, we find that $`c_1(𝒫^{})=c(\stackrel{~}{T}_{\mathrm{\Omega }_𝒪^n})`$.
For a PDQ (polarized deformation quantization) $`(𝒜,𝒪)`$ we set $`\mathrm{cl}(𝒜,𝒪)=c(F(𝒜))`$.
### 4.3 The opposite PDQ
For a ring $`R`$ we denote by $`R^{op}`$ the opposite ring. That is, $`R^{op}`$ is a ring, there is a bijection $`()^{}:RR^{op}`$ such that $`1^{}=1`$ and $`r_1^{}r_2^{}=(r_2r_1)^{}`$. We will usually identify $`R`$ and $`R^{op}`$ using the bijection.
If $`𝒜`$ is a deformation quantization of $`(M,\omega _0)`$, then, clearly, $`𝒜^{op}`$ is a deformation quantization of $`(M,\omega _0)`$. If $`(𝒜,𝒪)`$ is a PDQ, then $`(𝒜^{op},𝒪)`$ is a PDQ (since $`()^{}`$ restricts to a ring isomorphism on $`𝒪`$). It is clear that $`F(𝒜)`$ and $`F(𝒜^{op})`$ coincide as subsheaves of $`𝒜`$. However, the respective $`𝒪`$-module structures, symbol maps, and Lie brackets are different. In fact, one has:
* $`f\tau ^{}=(f\tau t\sigma (\tau )(f))^{}=(f\tau )^{}t\sigma (\tau )(f)`$
* $`\sigma (\tau ^{})=\sigma (\tau )`$
* $`[\tau _1^{},\tau _2^{}]=[\tau _1,\tau _2]^{}`$.
###### Proposition 4.2.
Suppose that $`(𝒜,𝒪)`$ is a PDQ. Then,
$`\mathrm{cl}(𝒜^{op},𝒪)=c_1(𝒫^{})\mathrm{cl}(𝒜,𝒪)`$
###### Proof.
Choose an open covering $`\{U_\alpha \}`$ of $`M`$ such that, for each open set $`U_\alpha `$ there exist flat connections $`s_\alpha `$ on $`F(𝒜)|_{U_\alpha }`$ and $`s_\alpha ^\mathrm{\Omega }`$ on $`\stackrel{~}{T}_{\mathrm{\Omega }_𝒪^n}|_{U_\alpha }`$.
Let $`\mathrm{}`$ denote the splitting of the extension
$`0𝒪\stackrel{~}{T}_{\mathrm{\Omega }_𝒪^n}T_𝒪\stackrel{\sigma }{}0`$
given by the negative of the Lie derivative. This splitting is a Lie algebra homomorphism, but is not $`𝒪`$-linear. Instead, for $`f𝒪`$ and $`\xi T`$, $`\mathrm{}(f\xi )=f\mathrm{}(\xi )+\xi (f)`$. We denote the restriction of this splitting to $`U_\alpha `$ by the same letter.
Since both $`\mathrm{}`$ and $`s_\alpha ^\mathrm{\Omega }`$ are splittings of the same sequence, their difference $`\psi _\alpha =s_\alpha ^\mathrm{\Omega }\mathrm{}`$ is a map $`\psi _\alpha :T𝒪`$ which is not $`𝒪`$-linear, but satisfies $`\psi _\alpha (f\xi )=f\psi _\alpha (\xi )+\xi (f)`$. Note that $`\psi _\alpha \psi _\beta =s_\alpha ^\mathrm{\Omega }s_\beta ^\mathrm{\Omega }`$.
Let $`s_\alpha ^{op}(\xi )=t\psi _\alpha (\xi )s_\alpha (\xi )^{}`$. We claim that $`s_\alpha ^{op}`$ is a (locally defined) flat connection on $`F(𝒜^{op})`$. To this end we check various properties.
The calculation
$`\sigma (s_\alpha ^{op}(\xi ))=\sigma (t\psi _\alpha (\xi )s_\alpha (\xi )^{})=\sigma (s_\alpha (\xi )^{})=\sigma (s_\alpha (\xi ))=\xi `$
shows that $`s_\alpha ^{op}`$ is a (local) splitting. The calculation
$$\begin{array}{c}s_\alpha ^{op}(f\xi )=t\psi _\alpha (f\xi )s_\alpha (f\xi )^{}=\psi _\alpha (f\xi )(fs_\alpha (\xi ))^{}=tf\psi _\alpha (\xi )+t\xi (f)(fs_\alpha (\xi )^{}+t\xi (f))=\hfill \\ \hfill tf\psi _\alpha (\xi )fs_\alpha (\xi )^{}=fs_\alpha ^{op}(\xi )\end{array}$$
shows that this splitting is $`𝒪`$-linear.
One verifies other properties similarly.
Hence one can use the locally defined flat connections $`\{s_\alpha ^{op}\}`$ to calculate the characteristic class of $`F(𝒜^{op})`$. The desired formula follows from
$`s_\alpha ^{op}s_\beta ^{op}=(\psi _\alpha (\xi )s_\alpha (\xi ))(\psi _\beta (\xi )s_\beta (\xi )=(s_\alpha ^\mathrm{\Omega }s_\beta ^\mathrm{\Omega })(s_\alpha s_\beta )`$
### 4.4 $`\mathrm{cl}(𝒜,𝒪)`$ and other classes
Let $`𝒜_1`$ and $`𝒜_2`$ be two deformation quantizations on $`(M,\omega _0)`$. Deligne, \[De\], related to them cohomology classes
$`[𝒜_2:𝒜_1]H^2(M,[[t]])\text{and}`$
$`\mathrm{obs}_{DL}(𝒜_1){\displaystyle \frac{1}{t}}[\omega _0]+H^2(M,[[t]]).`$
He proved that
$`[𝒜_2:𝒜_1]=\theta (𝒜_2)\theta (𝒜_1)\text{and}`$ (36)
$`\mathrm{obs}_{DL}(𝒜_1)=t{\displaystyle \frac{d}{dt}}\theta (𝒜_1).`$ (37)
###### Lemma 4.3.
Let $`(𝒜_1,𝒪_1)`$ and $`(𝒜_2,𝒪_2)`$ be two polarized quantizations of $`(M,\omega _0,P)`$. Then $`[𝒜_2:𝒜_1]=\frac{1}{t}\mathrm{cl}(𝒜_2,𝒪_2)\frac{1}{t}\mathrm{cl}(𝒜_1,𝒪_1)`$.
###### Proof.
By Proposition 3.12, there is an open covering $`\{U_\alpha \}`$ of $`M`$ such that there exist isomorphisms of polarized quantizations
$`B_\alpha :(𝒜_1,𝒪_1)_{U_\alpha }(𝒜_2,𝒪_2)_{U_\alpha }.`$
These isomorphisms induce isomorphisms $`F(𝒜_1)F(𝒜_2)`$ over each $`U_\alpha `$. We may suppose that both $`F(𝒜_1)`$ and $`F(𝒜_2)`$ admit on each $`U_\alpha `$ flat connections $`s_{i\alpha }:T_{𝒪_i}F(𝒜_i)`$, $`i=1,2`$, such that $`B_\alpha (s_{1\alpha })=s_{2\alpha }`$. Since $`F(𝒜_1)=F(𝒜_2)`$ mod $`t`$, we may choose the connections in such a way that $`s_{1\alpha }=s_{2\alpha }`$ mod $`t`$.
Let $`f_{i\alpha \beta }𝒪_i`$ be functions on $`U_\alpha U_\beta `$ such that $`df_{i\alpha \beta }=s_{i\beta }s_{i\alpha }`$. Since $`s_{1\alpha }=s_{2\alpha }`$ mod $`t`$, we may choose $`f_{i\alpha \beta }`$ such that $`f_{1\alpha \beta }=f_{2\alpha \beta }`$ mod $`t`$. On $`U_\alpha U_\beta `$ one has $`B_\alpha ^1B_\beta =\mathrm{exp}(\frac{1}{t}ad(f_{2\alpha \beta }f_{1\alpha \beta }))`$. Since $`f_{2\alpha \beta }f_{1\alpha \beta }t𝒪_i`$, $`\mathrm{exp}(\frac{1}{t}ad(f_{2\alpha \beta }f_{1\alpha \beta }))`$ is an automorphism of $`𝒜_1`$ on $`U_\alpha U_\beta `$. According to Deligne’s definition, $`[𝒜_2:𝒜_1]`$ is equal to $`\frac{1}{t}\stackrel{ˇ}{}(f_{2\alpha \beta }f_{1\alpha \beta })`$. ∎
###### Lemma 4.4.
Let $`(𝒜,𝒪)`$ be a polarized quantization of $`(M,\omega _0,P)`$ corresponding to a deformation $`(M,\omega ,𝒫)`$. Then, $`\mathrm{obs}_{DL}(𝒜)=t\frac{d}{dt}(\frac{1}{t}\mathrm{cl}(𝒜,𝒪))`$.
###### Proof.
First, we assume that $`(𝒜,𝒪)`$ corresponds to the trivial deformation $`(M,\omega _0,P[[t]])`$.
By Proposition 3.12, there is an open covering $`\{U_\alpha \}`$ of $`M`$ such that $`(𝒜,𝒪)_{U_\alpha }`$ are isomorphic to Moyal polarized quantizations. Let $`D_\alpha `$ be $``$ linear derivations of $`𝒜_{U_\alpha }`$ such that for $`a𝒪`$ one has $`D_\alpha a=t\frac{d}{dt}a`$. Such $`D_\alpha `$ exist because they obviously exist for the Moyal star-products. Then, $`D_\alpha F(𝒜)F(𝒜)`$.
Let the covering $`\{U_\alpha \}`$ be fine enough and on each $`U_\alpha `$ there exists a flat connection $`s_\alpha :T_𝒪F(𝒜)`$. Then, on $`U_\alpha U_\beta `$ one has $`s_\beta s_\alpha =df_{\alpha \beta }`$ for $`f_{\alpha \beta }𝒪`$. Let $`a_{\alpha 1},\mathrm{},a_{\alpha n}`$ be functions not depending on $`t`$ such that $`da_{\alpha 1},\mathrm{},da_{\alpha n}`$ form a local basis in $`P^{}`$, and $`_{\alpha 1},\mathrm{},_{\alpha n}`$ be the dual basis in $`T_𝒪`$. Put $`f_{\alpha i}=s_{\alpha i}(_{\alpha i})`$.
One has $`[f_{\alpha i},a_{\alpha j}]=t\delta _{ij}`$. Hence, since $`a_{\alpha i}`$ does not depend of $`t`$, $`[D_\alpha f_{\alpha i},a_{\alpha j}]=t\delta _{ij}`$. This means that $`D_\alpha f_{\alpha i}=f_{\alpha i}+b_{\alpha i}`$ for some $`b_{\alpha i}𝒪`$. It follows from $`[f_{\alpha i},f_{\alpha j}]=0`$ that $`[f_{\alpha i},b_{\alpha j}]+[b_{\alpha i},f_{\alpha j}]=0`$. It follows from the exactness of $`𝒪`$-de Rham complex that there exist $`b_\alpha 𝒪`$ such that $`b_{\alpha i}=\frac{1}{t}[b_\alpha ,f_{\alpha i}]`$. So, $`D_\alpha f_{\alpha i}=f_{\alpha i}+\frac{1}{t}[b_\alpha ,f_{\alpha i}]`$. Taking $`D_\alpha \frac{1}{t}\mathrm{ad}(b_\alpha )`$ instead $`D_\alpha `$ we obtain that the new $`D_\alpha `$ satisfy $`D_\alpha f_{\alpha i}=f_{\alpha i}`$. This equations may be rewritten in the form
$`D_\alpha s_\alpha =s_\alpha .`$ (38)
Further, one has
$`s_\alpha =s_\beta [{\displaystyle \frac{1}{t}}f_{\alpha \beta },s_\beta ].`$ (39)
Using (38) and (39), we have
$$(D_\beta D_\alpha )s_\alpha =D_\beta (s_\beta [\frac{1}{t}f_{\alpha \beta },s_\beta ])D_\alpha s_\alpha =[D_\beta (\frac{1}{t}f_{\alpha \beta }),s_\beta ]=[t\frac{d}{dt}(\frac{1}{t}f_{\alpha \beta }),s_\beta ].$$
Since derivations of $`𝒜`$ are completely defined by their values on $`F(𝒜)`$, it follows that $`D_\beta D_\alpha =\mathrm{ad}(t\frac{d}{dt}(\frac{1}{t}f_{\alpha \beta })).`$ According to Deligne’s definition, this means that $`(t\frac{d}{dt}(\frac{1}{t}f_{\alpha \beta }))`$ represents $`\mathrm{obs}_{DL}(𝒜)`$.
Now, let $`(M,\omega ,𝒫)`$ be arbitrary. Denote by $`\mathrm{cl}`$ the class of a quantization corresponding to $`(M,\omega ,𝒫)`$ and by $`\mathrm{cl}_0`$ the class of a quantization corresponding to the trivial deformation $`(M,\omega _0,P[[t]])`$. The same meaning has the notation $`\theta `$, $`\theta _0`$ and $`\mathrm{obs}_{DL}`$, $`\mathrm{obs}_{DL}^{}{}_{0}{}^{}`$.
By Lemma 4.3 and (36) we have $`\frac{1}{t}\mathrm{cl}_0\frac{1}{t}\mathrm{cl}=\theta _0\theta `$, hence
$$\theta _0\frac{1}{t}\mathrm{cl}_0=\theta \frac{1}{t}\mathrm{cl}.$$
Applying $`t\frac{d}{dt}`$ to the left hand side of this equation, we obtain zero (by (37) and what we have just proved). Hence, applying $`t\frac{d}{dt}`$ to the right hand side and using (37), we have
$$\mathrm{obs}_{DL}=t\frac{d}{dt}\theta =t\frac{d}{dt}\mathrm{cl},$$
which proves the proposition. ∎
It follows from Lemma 4.3 that the difference $`\theta (𝒜)\frac{1}{t}\mathrm{cl}(𝒜,𝒪)`$ does not depend on the quantization $`(𝒜,𝒪)`$ of $`(M,\omega _0,P)`$, while Lemma 4.4 shows that this difference does not depend on $`t`$, i.e.
$`\theta (𝒜)={\displaystyle \frac{1}{t}}\mathrm{cl}(𝒜,𝒪)c,`$ (40)
where $`c`$.
### 4.5 Relation between $`\theta (𝒜)`$ and $`\mathrm{cl}(𝒜,𝒪)`$
In this Section we calculate the constant $`c`$ in (40).
Let $`𝒜`$ be a deformation quantization of $`(M,\omega _0)`$. Let $`𝒜^{op}`$ denote the opposite algebra and $`𝒜^\sigma `$ denote the algebra obtained by the automorphism of $`[[t]]`$ taking $`tt`$. It is clear that both $`𝒜^{op}`$ and $`𝒜^\sigma `$ are quantizations of $`(M,\omega _0)`$.
###### Lemma 4.5.
Let $`𝒜`$ be a deformation quantization on $`(M,\omega _0)`$. Then
a) $`\theta (𝒜^{op})=\theta (𝒜)`$;
b) $`\theta (𝒜^\sigma )=\theta ^\sigma (𝒜)`$.
###### Proof.
Follows from the Fedosov construction, see also \[De\]. ∎
###### Theorem 4.6.
Let $`(𝒜,𝒪)`$ be a polarized quantization of $`(M,\omega _0,P)`$. Then
$`\theta (𝒜)={\displaystyle \frac{1}{t}}\mathrm{cl}(𝒜,𝒪){\displaystyle \frac{1}{2}}c_1(P).`$ (41)
###### Proof.
By Lemma 4.3 we have the equalities
$$\theta (𝒜^{op})\theta (𝒜^\sigma )=[𝒜^{op}:𝒜^\sigma ]=\mathrm{cl}(𝒜^{op},𝒪)\mathrm{cl}((𝒜,𝒪)^\sigma ).$$
By Lemma 4.5 and Proposition 4.2 we have $`\theta (𝒜^{op})=\theta (𝒜)`$ and $`\mathrm{cl}(𝒜^{op},𝒪)=\mathrm{cl}(𝒜,𝒪)+c_1(\mathrm{\Omega }_𝒪^1)`$. Combining these with the preceeding identity we obtain
$$(\mathrm{cl}(𝒜,𝒪)\theta (𝒜))+(\mathrm{cl}((𝒜,𝒪)^\sigma )\theta (𝒜^\sigma ))=c_1(\mathrm{\Omega }_𝒪^1).$$
By (40) both differences in the left hand side of the last identity are equal to the differences of the respective constant terms, hence coincide (since the constant terms are not affected by $`\sigma `$). Thus the above expression says
$$2(\mathrm{cl}(𝒜,𝒪)\theta (𝒜))=c_1(\mathrm{\Omega }_𝒪^1).$$
The theorem follows from the fact that $`c_1(\mathrm{\Omega }_𝒪^1)=c_1(P)`$. ∎
Note that the formula (41) relates to the formula (23), which shows that $`\mathrm{cl}(𝒜,𝒪)=\frac{1}{t}[\mathrm{\Omega }_D]`$, where $`\mathrm{\Omega }_D`$ is the Wick curvature appearing in the Fedosov procedure of constructing $`(𝒜,𝒪)`$.
e-mail: bressler@ihes.fr
e-mail: donin@macs.biu.ac.il
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# Vanishing cycles and mutation
## 1. Introduction
This talk is about symplectic aspects of Picard-Lefschetz theory, and the role of Floer cohomology in that context. I have relied on two main sources for inspiration, which are the ideas of Donaldson on vanishing cycles respectively those of Kontsevich, partly in collaboration with Barannikov, on mirror symmetry for Fano varieties . On a more technical level, the basis is provided by Fukaya’s work on Floer cohomology . I have tried to be as concrete as possible: not only are all objects mentioned rigorously defined, but they can be explicitly computed, and the assertions made about them checked, in many examples. This hands-on approach has its drawbacks, one of which will be mentioned after the summary of contents.
The basic geometric notion is that of an exact Morse fibration, a symplectic analogue of a holomorphic Morse function. One way of analyzing such fibrations is through the vanishing cycles in a fibre. When the base is a disc, one conveniently makes a choice of a finite family of vanishing cycles, forming a so-called distinguished basis. Any two such bases are connected by a sequence of Hurwitz moves. These are familiar concepts except that here vanishing cycles are considered as Lagrangian submanifolds, rather than only as homology classes.
Apart from their role as geometric objects worthy of study on their own, exact Morse fibrations are also relevant to Floer theory since, when equipped with suitable Lagrangian boundary conditions, they provide homomorphisms between Floer cohomology groups. As an application of these new maps we construct a long exact sequence analogous to that of Floer in gauge theory.
After that we return to exact Morse fibrations over a disc. The goal is to associate to each such fibration a triangulated category. An additional assumption is necessary: for simplicity, let’s say that the total space of the fibration has zero first Chern class and zero first Betti number. The construction proceeds in several steps: first one chooses a distinguished basis of vanishing cycles. From that one obtains a Fukaya-type $`A_{\mathrm{}}`$-category, unique up to quasi-isomorphism; our invariant is the derived category of this. It resembles derived categories of coherent sheaves on some Fano varieties, in that it is generated by an exceptional collection. The main point is to show that Hurwitz moves of the distinguished basis correspond to mutations of the exceptional collection, which leave the derived category unchanged. The explicit nature of mutations also allows them to be used for concrete computations of Floer cohomology.
We have to admit that our triangulated categories are only well-defined in a weak sense: different choices made during the construction lead to equivalent categories, but it has not been proved that the equivalences are canonical up to isomorphism (to have a completely satisfactory theory, it would further be necessary to establish coherence relations between functor isomorphisms). Rather than trying to do these improvements, it seems better to look for an alternative definition bypassing the choice of distinguished basis. The approach envisaged by Kontsevich is of this kind, but more work would be needed to put it on a rigorous footing.
###### Acknowledgments.
Thanks go to Ivan Smith for reading the manuscript and suggesting several improvements.
## 2. Picard-Lefschetz theory
### 2.1.
Let $`(M,\omega ,\theta )`$ be an exact symplectic manifold of dimension $`2n`$. This means that $`M`$ is compact with boundary, $`\omega \mathrm{\Omega }^2(M)`$ is a symplectic form, and $`\theta \mathrm{\Omega }^1(M)`$ satisfies $`d\theta =\omega `$. We will consider only symplectic automorphisms $`\varphi `$ of $`M`$ which are equal to the identity near $`M`$. Such a $`\varphi `$ is called exact if $`[\varphi ^{}\theta \theta ]H^1(M,M;)`$ is zero. The exact symplectic automorphisms form a subgroup $`\mathrm{Symp}^e(M)\mathrm{Symp}(M)`$. Note that any isotopy within this subgroup is Hamiltonian.
Let $`S`$ be a smooth manifold with boundary. An exact symplectic fibration over $`S`$ consists of data $`(E,\pi ,\mathrm{\Omega },\mathrm{\Theta })`$ as follows. $`\pi :ES`$ is a proper differentiable fibre bundle whose fibres are $`2n`$-dimensional manifolds with boundary. This means that $`E`$ itself is a manifold with codimension two corners, with $`E=_vE_hE`$ consisting of two faces: $`_vE=\pi ^1(S)`$, while $`\pi |_hE:_hES`$ is again a differentiable fibre bundle. $`\mathrm{\Omega }\mathrm{\Omega }^2(E)`$ is closed, its vertical part $`\mathrm{\Omega }|\mathrm{ker}(D\pi )`$ nondegenerate at every point, and $`\mathrm{\Theta }\mathrm{\Omega }^1(E)`$ satisfies $`d\mathrm{\Theta }=\mathrm{\Omega }`$ . We impose a final condition of triviality near $`_hE`$. This means that there should be a neighbourhood $`WE`$ of $`_hE`$ and a diffeomorphism, for some $`zS`$, $`\mathrm{\Xi }:S\times (WE_z)W`$ lying over $`S`$, such that $`\mathrm{\Xi }^{}\mathrm{\Omega }=pr_2^{}(\mathrm{\Omega }|E_z)`$ and $`\mathrm{\Xi }^{}\mathrm{\Theta }=pr_2^{}(\mathrm{\Theta }|E_z)`$; here $`pr_2`$ is projection from $`S\times (WE_z)`$ to the second factor. Clearly each fibre $`(E_z,\omega _z=\mathrm{\Omega }|E_z,\theta _z=\mathrm{\Theta }|E_z)`$ is an exact symplectic manifold. The form $`\mathrm{\Omega }`$ defines a canonical connection on $`\pi :ES`$, with structure group $`\mathrm{Symp}^e(E_z)`$. In fact there is a bijective correspondence between fibre bundles with structure group $`\mathrm{Symp}^e(E_z)`$ in the usual sense, and cobordism classes of exact symplectic fibrations. We denote the parallel transport maps of the canonical connection by $`\rho _c:E_{c(a)}E_{c(b)}`$, for $`c:[a;b]S`$.
### 2.2.
From now on assume that $`S`$ is two-dimensional and oriented. An exact Morse fibration (this is shorthand for “exact symplectic fibration with Morse-type critical points”) over $`S`$ consists of data $`(E,\pi ,\mathrm{\Omega },\mathrm{\Theta },J_0,j_0)`$. The properties of $`E,\pi ,\mathrm{\Omega },\mathrm{\Theta }`$ are as before, except that $`\pi `$ is allowed to have finitely many critical points. Each fibre may contain at most one of these points, and there should be none at all on $`E`$. $`J_0`$ is an integrable complex structure defined in a neighbourhood of the set $`E^{\mathrm{crit}}E`$ of critical points, and $`\mathrm{\Omega }`$ must be a Kähler form for it. Similarly $`j_0`$ is a positively oriented complex structure on a neighbourhood of the set $`S^{\mathrm{crit}}S`$ of critical values. They should be such that $`\pi `$ is $`(J_0,j_0)`$-holomorphic, with nondegenerate second derivative at each critical point. We will usually denote exact Morse fibrations by $`(E,\pi )`$ only.
One thing that needs explaining is why these are supposed to be analogues of holomorphic functions. For this one considers pairs $`(j,J)`$ consisting of a positively oriented complex structure $`j`$ on $`S`$ and an almost complex structure $`J`$ on $`E`$, such that $`j=j_0`$ near $`S^{\mathrm{crit}}`$, $`J=J_0`$ near $`E^{\mathrm{crit}}`$, $`\pi `$ is $`(J,j)`$-holomorphic, and $`\mathrm{\Omega }(,J)|\mathrm{ker}(D\pi )`$ is symmetric and positive definite everywhere. In this situation we say that $`J`$ is compatible relative to $`j`$. The space of such pairs $`(j,J)`$ is always contractible, and in particular nonempty. Moreover, for a fixed pair, by adding a positive two-form from $`S`$ one can modify $`\mathrm{\Omega }`$ such that it becomes symplectic and tames $`J`$.
Restricting any exact Morse fibration to $`SS^{\mathrm{crit}}`$ yields an exact symplectic fibration. Before bringing the singular fibres into the picture, we need some more definitions. Let $`(M,\omega ,\theta )`$ be an exact symplectic manifold. A Lagrangian submanifold $`LM`$, always assumed to be disjoint from $`M`$, is called exact if $`[\theta |L]H^1(L;)`$ is zero. A framed Lagrangian sphere is a Lagrangian submanifold $`L`$ together with an equivalence class $`[f]`$ of diffeomorphisms $`f:S^nL`$. Here $`f_1,f_2`$ are equivalent if $`f_2^1f_1`$ is isotopic to some element of $`O(n+1)\mathrm{Diff}(S^n)`$. One can associate to any $`(L,[f])`$ a Dehn twist $`\tau _{(L,[f])}\mathrm{Symp}(M)`$ which is unique up to Hamiltonian isotopy. If $`L`$ is exact, so is the Dehn twist along it. In future, we will often omit the framing $`[f]`$ from the notation.
To return to our discussion, let $`(E,\pi )`$ be an exact Morse fibration. Take a path $`c:[0;1]S`$ with $`c^1(S^{\mathrm{crit}})=\{1\}`$ and $`c^{}(1)0`$. Let $`x`$ be the unique critical point in $`E_{c(1)}`$. Then the stable manifold
$$B=\{yE_{c(s)},\mathrm{\hspace{0.17em}0}s<1,\text{ with }\underset{t1}{lim}\rho _{c|[s,t]}(y)=x\}\{x\}$$
is a smoothly embedded $`(n+1)`$-dimensional ball on which $`\mathrm{\Omega }`$ vanishes identically, and therefore $`V=B=BE_{c(0)}`$ is an exact Lagrangian submanifold of $`E_{c(0)}`$ diffeomorphic to $`S^n`$. Moreover, $`V`$ has a canonical structure of a framed Lagrangian sphere, constructed by first carrying $`V`$ by parallel transport to $`BE_{c(s)}`$, for some $`s`$ close to $`1`$, then projecting orthogonally in local Kähler coordinates to $`TB_x`$, and finally projecting radially to the unit sphere in that tangent space. The composition of these maps is a diffeomorphism $`f^1:VS^n`$ whose inverse is the framing. $`(V,[f])`$ is called the vanishing cycle associated to $`c`$. A symplectic version of the Picard-Lefschetz theorem says that for $`l`$,$`c`$ as in Figure 1, the monodromy $`\rho _l\mathrm{Symp}^e(E_{c(0)})`$ is isotopic to $`\tau _{(V,[f])}`$.
### 2.3.
Suppose now that $`S=D`$ is the closed unit disc in $``$. Let $`(E,\pi )`$ be an exact Morse fibration over it with $`m`$ critical values. Let $`M`$ be the fibre at some base point $`z_0D`$, say $`z_0=i`$. An admissible choice of paths is a family $`(c_1,\mathrm{},c_m)`$ looking as in Figure 2, ordered by their tangent directions at $`z_0`$. The collection of exact framed Lagrangian spheres in $`M`$ arising from them is called a distinguished basis of vanishing cycles. Modifying the paths affects the distinguished basis in a way which can be determined using the Picard-Lefschetz theorem. The outcome is encoded in the following abstract notion:
###### Definition 2.1.
A Lagrangian configuration in an exact symplectic manifold $`M`$ is an ordered family $`\mathrm{\Gamma }=(L_1,\mathrm{},L_m)`$ of exact framed Lagrangian spheres. Two configurations are Hurwitz equivalent if they can be connected by a sequence of the following moves and their inverses:
* $`\mathrm{\Gamma }=(L_1,\mathrm{},L_m)(L_1,\mathrm{},L_{i1},\varphi (L_i),L_{i+1},\mathrm{},L_m)`$ for some $`1im`$ and some $`\varphi \mathrm{Symp}^e(M)`$ isotopic to the identity;
* $`\mathrm{\Gamma }c\mathrm{\Gamma }=(\tau _{L_1}(L_2),\tau _{L_1}(L_3),\mathrm{},\tau _{L_1}(L_m),L_1)`$;
* $`\mathrm{\Gamma }r\mathrm{\Gamma }=(L_1,\mathrm{},L_{m2},\tau _{L_{m1}}(L_m),L_{m1})`$.
Thus, the Hurwitz equivalence class of a distinguished basis is an invariant of the exact Morse fibration $`(E,\pi )`$. Conversely, from $`M`$ and that Hurwitz equivalence class one can reconstruct the fibration up to a suitable notion of deformation equivalence.
###### References.
The symplectic nature of Dehn twists in all dimensions was first noticed by Arnol’d . Global properties of these maps are discussed in , , . I am not aware of any systematic exposition of symplectic Picard-Lefschetz theory in the literature.
## 3. Floer cohomology
### 3.1.
From now on, any exact symplectic manifold $`(M,\omega ,\theta )`$ is assumed to have contact type boundary, with $`\theta |M`$ being the contact one-form (the same condition will be imposed on $`\mathrm{\Theta }|E_z`$, for $`E_z`$ any fibre of a Morse fibration). Then for any pair $`(L_1,L_2)`$ of exact Lagrangian submanifolds in $`M`$ there is a well-defined Floer cohomology group $`HF(L_1,L_2)`$, which is a finite-dimensional vector space over the field $`/2`$. We remind the reader that this is invariant under isotopies of $`L_1`$ or $`L_2`$, satisfies $`HF(\varphi L_1,\varphi L_2)HF(L_1,L_2)`$ for any $`\varphi \mathrm{Symp}^e(M)`$, and that there is a natural Poincaré duality $`HF(L_1,L_2)HF(L_2,L_1)^{}`$.
As a warm-up exercise, suppose that we have a compact oriented surface $`S`$ with boundary, and an exact Morse fibration $`(E^{2n+2},\pi )`$ over it. A Lagrangian boundary condition for $`E`$ is a closed submanifold $`Q^{n+1}_vE_hE`$ such that $`\pi |Q:QS`$ is a smooth fibration, satisfying $`\mathrm{\Omega }|Q=0`$ and $`[\mathrm{\Theta }|Q]=0H^1(Q;)`$. Then the intersection $`Q_z=QE_z`$, for any $`zS`$, is an exact Lagrangian submanifold in $`E_z`$, and parallel transport along $`S`$ takes these Lagrangian submanifolds into each other. Choose a complex structure $`j`$ on $`S`$ and an almost complex structure $`J`$ on $`E`$ which is compatible relative to $`j`$, as defined in the previous section. There is a Gromov type invariant $`\mathrm{\Phi }(E,\pi ,Q)/2`$ which counts, in the familiar sense, the number of $`(j,J)`$-holomorphic sections $`u:SE`$ with $`u(S)Q`$. The exactness assumptions imply that there can be no bubbles ($`J`$-holomorphic spheres in a fibre $`E_z`$, or $`J`$-holomorphic discs in $`E_z`$ with boundary in $`Q_z`$), so that the definition of the invariant is technically quite simple.
###### Example 3.1.
Let $`L`$ be an exact framed Lagrangian sphere in an exact symplectic manifold $`M`$. Starting from a standard local model, one can construct an exact Morse fibration $`(E,\pi )`$ over $`D`$ with $`E_{z_0}=M`$ for $`z_0=iD`$, having exactly one critical point, such that the monodromy around $`D`$ is $`\tau _L`$. Because $`\tau _L(L)=L`$, there is a unique Lagrangian boundary condition $`QE`$ with $`Q_{z_0}=L`$. $`\mathrm{\Phi }(E,\pi ,Q)`$ vanishes because the expected dimension of the space of $`(j,J)`$-holomorphic sections is always odd, hence never zero.
Now let $`S`$ be as before but with a finite set of marked points $`\mathrm{\Sigma }S`$. Suppose moreover that around each $`\zeta \mathrm{\Sigma }`$ we have preferred local coordinates, given by an oriented embedding $`\psi _\zeta :D^+S`$ of the half-disc $`D^+=D\{\mathrm{im}(z)0\}`$ with $`\psi _\zeta (0)=\zeta `$. Let $`(E,\pi )`$ be an exact Morse fibration over $`S^{}=S\mathrm{\Sigma }`$ which is trivial near the marked points. This means that we have a fixed exact symplectic manifold $`M`$ and preferred embeddings $`\mathrm{\Psi }_\zeta :(D^+\{0\})\times ME`$ lying over $`\psi _\zeta `$, satisfying some obvious conditions concerning $`\mathrm{\Omega }`$,$`\mathrm{\Theta }`$ that we do not care to write down. If $`QE`$ is a Lagrangian boundary condition, there is for each $`\zeta \mathrm{\Sigma }`$ a unique pair $`L_{\zeta ,\pm }`$ of exact Lagrangian submanifolds of $`M`$ such that $`\mathrm{\Psi }_\zeta ^1(Q)=[1;0)\times L_{\zeta ,}(0;1]\times L_{\zeta ,+}`$. After choosing a complex structure $`j`$ on $`S^{}`$ such that the $`\psi _\zeta `$ become holomorphic, and an almost complex structure $`J`$ on $`E`$ which is compatible relative to $`j`$ and satisfies some additional conditions with regard to $`\mathrm{\Psi }_\zeta `$, one can count pseudo-holomorphic sections with suitable behaviour near the marked points. The outcome is a relative invariant
(1)
$$\mathrm{\Phi }_{rel}(E,\pi ,Q)\underset{\zeta \mathrm{\Sigma }}{}HF(L_{\zeta ,+},L_{\zeta ,}).$$
These invariants satisfy the standard gluing law for a topological quantum field theory, which one can formulate in two parts as follows. First, if $`S`$ is not connected then the relative invariant decomposes into the tensor product of relative invariants associated to its connected components. Second, suppose that there are two marked points $`\zeta ,\zeta ^{}\mathrm{\Sigma }`$ with
$$L_{\zeta ,\pm }=L_{\zeta ^{},}.$$
One can define a new surface $`\overline{S}`$ by removing small half-discs around $`\zeta `$,$`\zeta ^{}`$ and gluing together the resulting half-circles, as in Figure 3. There is a natural set of marked points $`\overline{\mathrm{\Sigma }}\overline{S}`$ which is inherited from $`\mathrm{\Sigma }\{\zeta ,\zeta ^{}\}`$. A similar process applied to $`(E,\pi )`$ constructs a new exact Morse fibration over $`\overline{S}\overline{\mathrm{\Sigma }}`$ with Lagrangian boundary conditions. The gluing rule says that on the level of the invariants $`\mathrm{\Phi }_{rel}`$ this translates into contracting $`HF(L_{\zeta ,+},L_{\zeta ,})HF(L_{\zeta ^{},+},L_{\zeta ^{},})`$ by Poincaré duality.
###### Remark 3.2.
The reader is hereby warned of two possible misunderstandings: the gluing process does not take place along the boundaries $`E_z`$ of the fibres, and neither do we glue together two boundary circles of $`S`$.
### 3.2.
We next give three examples. Throughout, Poincaré duality will be used freely to write the invariants $`\mathrm{\Phi }_{rel}`$ in various equivalent ways, e.g. as multilinear maps between Floer cohomology groups.
* Let $`M`$ be an exact symplectic manifold and $`L_1,L_2,L_3M`$ exact Lagrangian submanifolds. Take $`S=D`$ and $`\mathrm{\Sigma }=\{\text{three points}\}S`$. Let $`I_1,I_2,I_3`$ be the three components of $`S^{}`$, ordered in positive sense. The trivial fibration $`E=S^{}\times M`$ with Lagrangian boundary conditions $`Q=_{\nu =1}^3I_\nu \times L_\nu `$ yields a relative invariant which can be written as a map $`HF(L_2,L_3)HF(L_1,L_2)HF(L_1,L_3)`$. This is just a reformulation of the usual “pair-of-pants” product. Reversing the orientation of $`S`$ yields another relative invariant, which is the Poincaré dual coproduct $`HF(L_1,L_3)HF(L_2,L_3)HF(L_1,L_2)`$.
* Let $`M`$ be an exact symplectic manifold, $`L_1,L_2M`$ exact Lagrangian submanifolds, and $`\varphi \mathrm{Symp}^e(M)`$ an automorphism which is isotopic to the identity. Imagine $`S=D`$ as being constructed out of a smaller disc $`D_0`$ and two other pieces $`D_1,D_2=[0;1]^2`$. The marked points are $`\mathrm{\Sigma }=\{\zeta _1,\zeta _2\}`$, where $`\zeta _1`$ is obtained by identifying $`(0,1)D_1`$ with $`(0,0)D_2`$, and $`\zeta _2`$ is similarly $`(1,1)D_1`$ or $`(1,0)D_2`$. One can construct an exact symplectic fibration $`E_0`$ over $`D_0`$, which is topologically $`D_0\times M`$ but has nontrivial forms $`\mathrm{\Omega }`$ and $`\mathrm{\Theta }`$, such that the monodromy around $`D_0`$ is $`\varphi `$. Assemble $`E_0`$ and the two trivial fibrations $`E_1=(D_1\{(0,1),(1,1)\})\times M`$, $`E_2=(D_2\{(0,0),(0,1)\})\times M`$ following the instructions in Figure 4. This gives an exact Morse fibration over $`S^{}`$; we equip it with the Lagrangian boundary condition which is the union of $`(0;1)\times \{1\}\times L_1E_1`$ and $`(0;1)\times \{0\}\times L_2E_2`$. The resulting relative invariant is a map $`HF(L_1,L_2)HF(L_1,\varphi (L_2))`$, in fact just the familiar “continuation” isomorphism.
* Let $`M,L_1,L_2`$ be as before, and $`LM`$ an exact framed Lagrangian sphere. The central piece $`E_0`$ in the previous construction can be replaced by the Morse fibration from Example 3.1. This leads to a map $`HF(L_1,L_2)HF(L_1,\tau _L(L_2))`$, which is not an isomorphism in general. In fact there seems to be no way at all of getting from our formalism a map in the inverse direction; the reason is the presence of critical points, which prevents one from reversing the orientation of the base.
###### Theorem 3.3.
Let $`(M,\omega ,\theta )`$ be an exact symplectic manifold, $`L_1,L_2M`$ exact Lagrangian submanifolds, and $`LM`$ an exact framed Lagrangian sphere. Suppose that $`2c_1(M,L)H^2(M,L)`$ is zero. Then there is a long exact sequence, with $`a`$ the pair-of-pants product and $`b`$ the map defined in the last example above,
(2)
To understand how this works, one needs to look at the Floer cochain complexes $`CF`$. To simplify, we suppress the dependence of these complexes on various additional choices, and write $`a,b`$ for the maps between them inducing the Floer cohomology maps mentioned above. The central object in the proof is a map of complexes
(3)
$$\mathrm{Cone}(a)\stackrel{(h,b)}{}CF(L_1,\tau _L(L_2)).$$
Here $`h:CF(L,L_2)CF(L_1,L)CF(L_1,\tau _L(L_2))`$ is a chain homotopy $`ba0`$ defined as follows. Consider the exact Morse fibration with $`S=D`$, $`\mathrm{\Sigma }=\{\text{three points}\}`$, which is represented schematically, together with its Lagrangian boundary condition, in Figure 5. Moving the two leftmost marked points simultaneously along $`D`$, in a way that preserves the symmetry of the picture with respect to the $`x`$-axis, yields a one-parameter family of fibrations, and $`h`$ arises from the corresponding parametrized spaces of pseudo-holomorphic sections. Once (3) has been defined, an argument using the natural filtration of the Floer complexes by the action functional shows that it is a quasi-isomorphism; the long exact sequence is an immediate consequence.
As one can see from this sketch of the argument, the definition of the third arrow in (2) uses the inverse of the quasi-isomorphism (3). Poincaré duality yields a symmetry of the exact sequence, which suggests a more direct description of that arrow. Namely, it should be the composition
$$\begin{array}{c}c:HF(L_1,\tau _L(L_2))HF(L,\tau _L(L_2))HF(L_1,L)\hfill \\ \hfill HF(L,L_2)HF(L_1,L)\end{array}$$
of the coproduct and the isomorphism $`HF(L,\tau _L(L_2))HF(\tau _L^1(L),L_2)=HF(L,L_2)`$. One can show that this is indeed the same map as that obtained by inverting (3); the proof uses invariants of the same kind as those defining $`h`$, for two-parameter families of exact Morse fibrations.
###### Remark 3.4.
The assumption $`2c_1(M,L)=0`$ is a technical one. It implies that the space of $`(j,J)`$-holomorphic sections in Example 3.1 has expected dimension $`2n1`$. From this it follows (except for $`n=1`$, which requires a separate treatment) that generically there is a positive-dimensional space of these sections $`u`$ such that $`u(z_0)`$ is a specific point in $`L`$. That enters into the description of the limiting behaviour of sections of the family in Figure 5 when both moveable points of $`\mathrm{\Sigma }`$ go towards the fixed one, and through it into the proof that $`h`$ is a homotopy $`ba0`$. Still, it may be possible to substitute some other argument at this point, and thereby remove the assumption from Theorem 3.3.
There is an important special case in which Theorem 3.3 has a graphical interpretation, namely when all Lagrangian submanifolds involved are vanishing cycles. Take an exact Morse fibration $`(E,\pi )`$ with arbitrary base $`S`$. To any path $`c:[0;1]S`$ with $`c^1(S^{\mathrm{crit}})=\{0;1\}`$ and $`c^{}(0),c^{}(1)0`$ one can associate the Floer cohomology group $`HF(V_{c,0},V_{c,1})`$ where $`V_{c,0},V_{c,1}E_{c(t)}`$, for some $`0<t<1`$, are the vanishing cycles associated to $`c|[0;t]`$ resp. $`c|[t;1]`$. This group is independent of $`t`$, so we may write temporarily $`HF(c)`$ for it. Consider four paths as in Figure 6. As three of them run together for small $`t`$, we may assume that their vanishing cycles all lie in the same fibre, called $`M`$. Then $`V_{c_{},0}=V_{c_+,0}=V_{c^{},0}`$ and by the Picard-Lefschetz theorem, $`V_{c_+,1}=\tau _{V_{c^{},1}}(V_{c_{},1})`$. Taking the vanishing cycles of $`c^{\prime \prime }`$ and moving them to $`M`$ by parallel transport shows that $`HF(V_{c^{},1},V_{c_{},1})=HF(V_{c^{},1},V_{c_+,1})=HF(V_{c^{\prime \prime },0},V_{c^{\prime \prime },1})`$. Applying Theorem 3.3 in this context yields a long exact sequence
This can be thought of as a skein rule, with $`HF(c^{\prime \prime })HF(c^{})`$ measuring the change to $`HF(c_{})`$ as the path moves over $`zS^{\mathrm{crit}}`$ and becomes $`c_+`$. For $`S=D`$, allowing oneself for a moment to believe naively that “two terms in a long exact sequence determine the third one”, one sees that all groups $`HF(c)`$ could be determined from knowing just finitely many of them, through a process of successively breaking up paths into shorter pieces.
###### References.
Our TQFT formalism for Floer cohomology is a variation of that in , which in turn generalizes earlier work of Piunikhin-Salamon-Schwarz . The pair-of-pants product and continuation map are defined in \[10, Chapter 10\] or or , respectively . The skein rule interpretation of the exact sequence is due to Donaldson.
## 4. Grading
### 4.1.
The assumption $`2c_1(M,L)=0`$ in Theorem 3.3 is very close to the condition needed to equip Floer cohomology groups with $``$-gradings, so one might just as well take advantage of it. Doing that requires some preliminaries, which we now go through.
Let $`M`$ be an arbitrary symplectic manifold, and $`𝒥_M`$ the space of all compatible almost complex structures. There is a canonical unitary line bundle $`\mathrm{\Delta }_M𝒥_M\times M`$ whose fibre at $`(J,x)`$ is $`\mathrm{\Lambda }^n(TM_x,J)^2`$. To any Lagrangian submanifold $`LM`$ one can associate a section $`det^2(TL)`$ over $`𝒥_M\times L`$ of the associated circle bundle $`S(\mathrm{\Delta }_M)`$; and for any symplectic automorphism $`\varphi `$ there is a canonical isomorphism $`det^2(D\varphi ):\mathrm{\Delta }_M\mathrm{\Delta }_M`$, covering the map $`\varphi _{}\times \varphi `$ from $`𝒥_M\times M`$ to itself. A Maslov map is a trivialization $`\delta _M:\mathrm{\Delta }_M`$. Suppose that we have chosen such a $`\delta _M`$. A grading of $`LM`$ is a lift
Any two gradings differ by a locally constant function $`L`$; we write $`\stackrel{~}{L}[\sigma ]`$ for $`\stackrel{~}{L}\sigma `$, $`\sigma `$. Similarly, a grading of $`\varphi `$ is a map $`\stackrel{~}{\varphi }:𝒥_M\times M`$ such that $`\mathrm{exp}(2\pi i\stackrel{~}{\varphi }(J,x))=\delta _M(det^2(D\varphi )(w))/\delta _M(w)`$ for any $`wS(\mathrm{\Delta }_M)`$ in the fibre over $`(J,x)`$. When dealing with connected manifolds $`M`$ with boundary and maps $`\varphi `$ trivial near $`M`$, there is often a preferred grading, characterized by being zero near $`𝒥_M\times M`$. Gradings of $`L`$ and $`\varphi `$ induce a grading of $`\varphi (L)`$:
$$\stackrel{~}{\varphi }(\stackrel{~}{L})\stackrel{\mathrm{def}}{=}(\stackrel{~}{\varphi }+\stackrel{~}{L})(\varphi _{}^1\times \varphi ^1).$$
If $`LM`$ is a framed Lagrangian sphere admitting gradings, there is a distinguished grading $`\stackrel{~}{\tau }_L`$ of $`\tau _L`$, which is zero outside $`𝒥_M\times \{\text{neighbourhood of }L\}`$ and satisfies
(4)
$$\stackrel{~}{\tau }_L(\stackrel{~}{L})=\stackrel{~}{L}[1n].$$
Graded Lagrangian submanifolds and graded symplectic automorphisms are pairs consisting of an $`L`$ resp. $`\varphi `$ together with a choice of grading. For brevity we denote such pairs by $`\stackrel{~}{L},\stackrel{~}{\varphi }`$ only. The Floer cohomology of a pair of graded Lagrangian submanifolds, whenever defined (e.g. when $`M`$ is exact with contact type boundary, and the Lagrangian submanifolds are themselves exact), has a canonical $``$-grading with the properties
$`HF^{}(\stackrel{~}{L}_1[\sigma ],\stackrel{~}{L}_2)=HF^{}(\stackrel{~}{L}_1,\stackrel{~}{L}_2[\sigma ])=HF^{+\sigma }(\stackrel{~}{L}_1,\stackrel{~}{L}_2),`$
(5) $`HF^{}(\stackrel{~}{\varphi }(\stackrel{~}{L}_1),\stackrel{~}{\varphi }(\stackrel{~}{L}_2))HF^{}(\stackrel{~}{L}_1,\stackrel{~}{L}_2),`$
$`HF^{}(\stackrel{~}{L}_2,\stackrel{~}{L}_1)HF^n(\stackrel{~}{L}_1,\stackrel{~}{L}_2)^{}.`$
### 4.2.
Now consider an exact Morse fibration $`(E,\pi )`$ over $`S`$, and the space $`𝒥_{E/S}`$ of pairs $`(j,J)`$ such that $`J`$ is compatible relative to $`j`$. Let $`\mathrm{\Delta }_{E/S}𝒥_{E/S}\times E`$ be the line bundle with fibres $`\mathrm{\Lambda }^{n+1}(TE_x,J)^2(TS_{\pi (x)},j)^2`$. A relative Maslov map is a trivialization $`\delta _{E/S}:\mathrm{\Delta }_{E/S}`$. Once such a map has been chosen, there is a canonical induced Maslov map $`\delta _M`$ on any regular fibre $`M=E_z`$. The vanishing cycles $`VM`$ associated to paths $`c:[0;1]S`$, $`c(0)=z`$ and $`c(1)S^{\mathrm{crit}}`$, admit gradings; and the monodromies $`\rho _l`$ along closed loops $`l`$ in $`SS^{\mathrm{crit}}`$, $`l(0)=z`$, even have canonical gradings, which we denote by $`\stackrel{~}{\rho }_l`$. If $`[l]H_1(S)`$ vanishes, $`\stackrel{~}{\rho }_l`$ is zero near $`𝒥_M\times M`$. This implies a graded version of the Picard-Lefschetz theorem, saying that $`\stackrel{~}{\rho }_l=\stackrel{~}{\tau }_V`$ for $`c,l`$ as in Figure 1.
### 4.3.
Any Lagrangian boundary condition $`Q`$ for $`(E,\pi )`$ comes with a canonical section of $`S(\mathrm{\Delta }_{E/S})|𝒥_{E/S}\times Q`$. We call $`\delta _{E/S}`$ adapted to $`Q`$ if there is a diagram
For compact $`S`$, the existence of an adapted relative Maslov map means that all components of the space of $`(j,J)`$-holomorphic sections taking $`S`$ to $`Q`$ have the same expected dimension $`n\chi (S)+\mathrm{deg}(\lambda )`$; we have already seen an instance of this in Remark 3.4. There is a generalization of this notion to the relative case, which ensures that the Floer groups in (1) have canonical gradings and that the invariant $`\mathrm{\Phi }_{rel}`$ is of a specific degree, not necessarily zero. Instead of explaining this in general we will just look at an example, that of the long exact sequence.
Suppose then that we have an exact symplectic manifold $`M`$ with a Maslov map $`\delta _M`$, graded Lagrangian submanifolds $`\stackrel{~}{L}_1,\stackrel{~}{L}_2`$, and a framed exact Lagrangian sphere $`L`$ which admits gradings. The map $`b:HF^{}(\stackrel{~}{L}_1,\stackrel{~}{L}_2)HF^{}(\stackrel{~}{L}_1,\stackrel{~}{\tau }_L(\stackrel{~}{L}_2))`$ has degree zero, and the same holds for the pair-of-pants product $`a:HF^{}(\stackrel{~}{L},\stackrel{~}{L}_2)HF^{}(\stackrel{~}{L}_1,\stackrel{~}{L})HF^{}(\stackrel{~}{L}_1,\stackrel{~}{L}_2)`$. The coproduct $`HF^{}(\stackrel{~}{L}_1,\stackrel{~}{\tau }_L(\stackrel{~}{L}_2))HF^{}(\stackrel{~}{L},\stackrel{~}{\tau }_L(\stackrel{~}{L}_2))HF^{}(\stackrel{~}{L}_1,\stackrel{~}{L})`$ has degree $`n`$ because of Poincaré duality, but then by (4) and (5) the isomorphism
$$HF^{}(\stackrel{~}{L},\stackrel{~}{\tau }_L(\stackrel{~}{L}_2))=HF^{}(\stackrel{~}{\tau }_L^1(\stackrel{~}{L}),\stackrel{~}{L}_2)=HF^{+1n}(\stackrel{~}{L},\stackrel{~}{L}_2)$$
has degree $`1n`$. Hence the third map $`c`$ raises degrees by one, like in any other cohomological long exact sequence.
###### References.
Floer’s discussion of grading in is essentially complete as it stands. Still, graded Lagrangian submanifolds, introduced by Kontsevich somewhat later, allow a better formulation of the results.
## 5. Mutation
### 5.1.
As in our discussion of Floer theory, we take the ground field to be $`/2`$; all categories will be linear over it. Let $`𝒞`$ be a triangulated category such that the spaces $`\mathrm{Hom}_𝒞^{}(X,Y)=_r\mathrm{Hom}_𝒞(X,Y[r])`$ are finite-dimensional for any $`X,Y`$. An exceptional collection in $`𝒞`$ is a finite family of objects $`(X^1,\mathrm{},X^m)`$ satisfying
$$\mathrm{Hom}_𝒞^{}(X^i,X^k)=\{\begin{array}{cc}/2\mathrm{id}_{X^i}\hfill & i=k,\hfill \\ 0\hfill & i>k.\hfill \end{array}$$
Such a collection is called full if the $`X^i`$ generate $`𝒞`$ in the triangulated sense. For $`X,Y\mathrm{Ob}𝒞`$ define $`T_X(Y)`$, up to isomorphism, as the object fitting into an exact triangle
$$T_X(Y)[1]\mathrm{Hom}_𝒞^{}(X,Y)X\stackrel{ev}{}YT_X(Y);$$
here the tensor product is just a finite sum of shifted copies of $`X`$, and $`ev`$ the canonical evaluation map. If $`(X^1,\mathrm{},X^m)`$ is a full exceptional collection then so are $`(Y^1,\mathrm{},Y^m)`$, $`(Z^1,\mathrm{},Z^m)`$ where
(6) $`Y^i`$ $`=\{\begin{array}{cc}T_{X^1}(X^{i+1})\hfill & i<m,\hfill \\ X^1\hfill & i=m,\hfill \end{array}`$
respectively
(7) $`Z^i`$ $`=\{\begin{array}{cc}X^i\hfill & i<m1,\hfill \\ T_{X^{m1}}(X^m)\hfill & i=m1,\hfill \\ X^{m1}\hfill & i=m.\hfill \end{array}`$
### 5.2.
There is a slightly different version of the same story for $`A_{\mathrm{}}`$-categories. Call an $`A_{\mathrm{}}`$-category $`𝒜`$ directed if it has finitely many objects numbered $`1,\mathrm{},m`$, say $`\mathrm{Ob}𝒜=\{X^1,\mathrm{},X^m\}`$, such that
$$hom_𝒜(X^i,X^k)=\{\begin{array}{cc}\text{finite-dimensional}\hfill & i<k,\hfill \\ /2\mathrm{id}_{X^i}\hfill & i=k,\hfill \\ 0\hfill & i>k.\hfill \end{array}$$
Note that $`\mu _𝒜^d=0`$ is necessarily zero for $`d>\mathrm{max}\{m1,2\}`$. We recall the definition of the (bounded) derived category $`D^b(𝒜)`$. The first step is to embed $`𝒜`$ into a larger $`A_{\mathrm{}}`$-category $`𝒜^{}`$ which has finite sums and shifts. Thus, an object of $`𝒜^{}`$ is a formal sum $`_{eE}X_e[\sigma _e]`$ with $`E`$ a finite set, $`X_e\mathrm{Ob}𝒜`$, and $`\sigma _e`$. Next, a twisted complex in $`𝒜`$ is a pair $`(C,\delta _C)`$ consisting of $`C\mathrm{Ob}𝒜^{}`$ and $`\delta _Chom_𝒜^{}^1(C,C)`$, such that the “generalized Maurer-Cartan equation”
(8)
$$\underset{d1}{}\mu _𝒜^{}^d(\delta _C,\mathrm{},\delta _C)=0$$
holds. Twisted complexes form an $`A_{\mathrm{}}`$-category $`\mathrm{Tw}𝒜`$ which again has direct sums and shifts. It also contains a cone $`\mathrm{Cone}(a)`$ for any morphism $`a`$ such that $`\mu _{\mathrm{Tw}𝒜}^1(a)=0`$, and the cohomological category $`D^b(𝒜)=H^0(\mathrm{Tw}𝒜)`$ inherits a triangulated structure from this. The objects of $`𝒜`$, seen as twisted complexes with zero differential, form a full exceptional collection in $`D^b(𝒜)`$.
###### Remark 5.1.
The definition of the derived category of a general $`A_{\mathrm{}}`$-category $`𝒜`$ uses pairs $`(C,\delta _C)`$ such that $`\delta _C`$ is strictly decreasing with respect to some finite filtration of $`C`$, which has the effect of making the sum (8) finite. The fact that this is not necessary in our case is just one of several technical simplifications which directedness brings with it.
Any $`A_{\mathrm{}}`$-functor $`F:𝒜`$ between directed $`A_{\mathrm{}}`$-categories induces another one $`\mathrm{Tw}F:\mathrm{Tw}𝒜\mathrm{Tw}`$ taking cones to cones, and therefore an exact functor $`D^b(F)=H^0(\mathrm{Tw}F):D^b(𝒜)D^b()`$. Call $`F`$ a quasi-isomorphism if $`H(F):H(𝒜)H()`$ is an isomorphism; since the objects of $`𝒜`$ generate $`D^b(𝒜`$), it follows that in this case $`D^b(F)`$ is an equivalence.
More interestingly, suppose that $`𝒜`$ is some directed $`A_{\mathrm{}}`$-category, and $`(Y^1,\mathrm{},Y^m)`$ an exceptional collection in $`D^b(𝒜)`$. One can then define a new directed $`A_{\mathrm{}}`$-category $``$, called the directed $`A_{\mathrm{}}`$-subcategory of $`\mathrm{Tw}𝒜`$ generated by the $`Y^i`$, as follows. Objects of $``$ are the $`Y^i`$, and $`hom_{}(Y^i,Y^k)=hom_{\mathrm{Tw}𝒜}(Y^i,Y^k)`$ for $`i<k`$, with the same composition maps between these groups as in $`\mathrm{Tw}𝒜`$. All other morphism groups and compositions are as dictated by the directedness condition. The embedding $`\mathrm{Tw}𝒜`$ can be extended to an $`A_{\mathrm{}}`$-functor $`\iota :\mathrm{Tw}\mathrm{Tw}𝒜`$, which induces an exact functor $`H^0(\iota ):D^b()D^b(𝒜)`$. For the same reason as before, $`H^0(\iota )`$ is full and faithful; it is even an equivalence if $`(Y^1,\mathrm{},Y^m)`$ is a full exceptional collection.
### 5.3.
Given $`X\mathrm{Ob}\mathrm{Tw}𝒜`$ and a finite-dimensional complex $`V`$ of vector spaces, one can form the tensor product $`VX\mathrm{Ob}\mathrm{Tw}𝒜`$, which is a direct sum of shifted copies of $`X`$ with a differential combining those on $`V`$ and $`X`$. Taking $`V=(hom_{\mathrm{Tw}𝒜}(X,Y),\mu _𝒜^1)`$ for some $`Y\mathrm{Ob}\mathrm{Tw}𝒜`$, one has a canonical evaluation morphism $`evhom_{\mathrm{Tw}𝒜}^0(hom_{\mathrm{Tw}𝒜}(X,Y)X,Y)`$ with $`\mu _{\mathrm{Tw}𝒜}^1(ev)=0`$. Let $`T_X(Y)\mathrm{Ob}\mathrm{Tw}𝒜`$ be the cone of $`ev`$. This is isomorphic in $`D^b(𝒜)`$ to the object of the same name introduced above, but it is now unique in a strict sense, not just up to isomorphism. Taking the exceptional collection formed by the objects of $`𝒜`$ and applying (6) or (7) yields another full exceptional collection, hence a directed $`A_{\mathrm{}}`$-subcategory $``$ of $`\mathrm{Tw}𝒜`$ with $`D^b()D^b(𝒜)`$. This process can be repeated indefinitely and leads to the following notion:
###### Definition 5.2.
Two directed $`A_{\mathrm{}}`$-categories with $`m`$ objects are mutations of each other if they can be related by a sequence of the following moves and their inverses:
* $`𝒜`$ if there is a quasi-isomorphism between them.
* It is allowed to shift each object by some degree, which means changing the grading of each group $`hom_𝒜(X^i,X^k)`$ by $`(\sigma _i\sigma _k)`$ for some $`\sigma _1,\mathrm{},\sigma _m`$, while keeping the same composition maps.
* $`𝒜c𝒜`$ where, if the objects of $`c𝒜`$ are denoted by $`\{Y^1,\mathrm{},Y^m\}`$, the nontrivial morphism spaces are
$$hom_{c𝒜}(Y^i,Y^k)=\{\begin{array}{cc}hom_𝒜(X^{i+1},X^{k+1})\hfill & i<k<m,\hfill \\ hom_𝒜(X^1,X^{i+1})^{}[1]\hfill & i<m,k=m\text{.}\hfill \end{array}$$
Here denotes the dual of a graded vector space. The compositions $`\mu _{c𝒜}^d:_{\nu =1}^dhom_{c𝒜}(Y^{i_\nu },Y^{i_{\nu +1}})hom_{c𝒜}(Y^{i_1},Y^{i_{d+1}})`$, $`i_1<\mathrm{}<i_{d+1}`$, are equal to those in $`𝒜`$ except when $`i_{d+1}=m`$, in which case one has $`\mu _{c𝒜}^d(a^d,\mathrm{},a^1),b=a^d,\mu _𝒜^d(a^{d1},\mathrm{},a^1,b)`$ with $`\mathrm{}`$ the dual pairing.
* $`𝒜r𝒜`$ with $`\mathrm{Ob}r𝒜=\{Z^1,\mathrm{},Z^m\}`$ and the following nontrivial morphism spaces $`hom_{r𝒜}(Z^i,Z^k)`$:
(9)
$$\{\begin{array}{cc}hom_𝒜(X^i,X^k)\hfill & i<km2\text{,}\hfill \\ hom_𝒜(X^i,X^{m1})\hfill & im2\text{}k=m\text{,}\hfill \\ hom_𝒜(X^{m1},X^m)^{}[1]\hfill & i=m1\text{}k=m\text{,}\hfill \\ (hom_𝒜(X^{m1},X^m)hom_𝒜(X^i,X^{m1}))[1]\hfill & \\ hom_𝒜(X^i,X^m)\hfill & im2\text{}k=m1\text{.}\hfill \end{array}$$
$`\mu _{r𝒜}^1`$ is given by $`\mu _𝒜^1`$ in the first two cases, its dual $`(\mu _𝒜^1)^{}`$ in the third, and in the final case by
$$\left(\begin{array}{cc}\mu _𝒜^1\mathrm{id}+\mathrm{id}\mu _𝒜^1& 0\\ \mu _𝒜^2& \mu _𝒜^1\end{array}\right).$$
Most of the higher order maps
(10)
$$\mu _{r𝒜}^d:_{\nu =1}^dhom_{r𝒜}(Z^{i_\nu },Z^{i_{\nu +1}})hom_{r𝒜}(Z^{i_1},Z^{i_{d+1}})$$
for $`i_1<\mathrm{}<i_{d+1}`$ are taken from those of $`𝒜`$ in a straightforward way, but there are two exceptions. One is when $`i_{d+1}=m1`$, in which case
$$\mu _{r𝒜}^d=\left(\begin{array}{cc}\mathrm{id}\mu _𝒜^d& 0\\ \mu _𝒜^{d+1}& \mu _𝒜^d\end{array}\right)$$
with respect to the obvious splittings on both sides of (10). The second exceptional case is when $`i_d=m1`$, $`i_{d+1}=m`$: then $`\mu _{r𝒜}^2`$ is
$$\begin{array}{c}hom_𝒜(X^{m1},X^m)^{}hom_𝒜(X^{m1},X^m)hom_𝒜(X^{i_1},X^{m1})\\ (hom_𝒜(X^{m1},X^m)^{}hom_𝒜(X^{i_1},X^m))[1]\\ \\ hom_𝒜(X^{i_1},X^{m1}),\\ \mu _{r𝒜}^2(a^3a^2a^1,b^2b^1)=a^3,a^2a^1,\end{array}$$
while the compositions of order $`d3`$ vanish.
Actually, while $`r𝒜`$ is precisely the directed $`A_{\mathrm{}}`$-subcategory of $`\mathrm{Tw}𝒜`$ generated by the collection (7), $`c𝒜`$ is only canonically quasi-isomorphic to that generated by (6). But it is still true that two directed $`A_{\mathrm{}}`$-categories which are mutations of each other have equivalent derived categories.
###### References.
This section makes no claim to originality. Exceptional collections in triangulated categories are discussed in the work of Bondal, Gorodentsev, Kapranov, and others. Most relevant for us is which emphasizes the role of dg categories. The short step from there to $`A_{\mathrm{}}`$-categories was made by Kontsevich , who is also responsible for introducing $`\mathrm{Tw}𝒜`$ and $`D^b(𝒜)`$ .
## 6. Fukaya categories
### 6.1.
Throughout this section $`M`$ is a fixed exact symplectic manifold, with a Maslov map $`\delta _M`$. A graded Lagrangian configuration in $`M`$ is a family $`\stackrel{~}{\mathrm{\Gamma }}=(\stackrel{~}{L}_1,\mathrm{},\stackrel{~}{L}_m)`$ of graded, exact, framed Lagrangian spheres. Hurwitz equivalence for graded configurations is defined as in the ungraded case, with the following adaptations: for the isotopy invariance one wants to take a grading $`\stackrel{~}{\varphi }`$ which is zero near $`𝒥_M\times M`$, so that $`\stackrel{~}{\varphi }`$ is isotopic to the identity in the group of graded symplectic automorphisms; the moves $`c\stackrel{~}{\mathrm{\Gamma }},r\stackrel{~}{\mathrm{\Gamma }}`$ use the canonical gradings $`\stackrel{~}{\tau }_L`$ of Dehn twists; and there is an additional shift move,
* $`\stackrel{~}{\mathrm{\Gamma }}(\stackrel{~}{L}_1[\sigma _1],\mathrm{},\stackrel{~}{L}_m[\sigma _m])`$ for any $`\sigma _1,\mathrm{},\sigma _m`$.
This is something of an anticlimax, since it cancels out the extra information contained in the grading; but in fact, the whole notion of graded configuration has been introduced only for notational convenience.
We will associate to $`\stackrel{~}{\mathrm{\Gamma }}`$ a directed Fukaya $`A_{\mathrm{}}`$-category $`\mathrm{𝐿𝑎𝑔}^{}(\stackrel{~}{\mathrm{\Gamma }})`$, unique up to quasi-isomorphism. Suppose first that the configuration is in general position, meaning that any two $`L_i`$ are transverse and there are no triple intersections. Objects of $`𝒜=\mathrm{𝐿𝑎𝑔}^{}(\stackrel{~}{\mathrm{\Gamma }})`$ are the graded Lagrangian submanifolds $`\stackrel{~}{L}_i`$, in the given order, and
$$hom_𝒜(\stackrel{~}{L}_i,\stackrel{~}{L}_k)=\{\begin{array}{cc}CF^{}(\stackrel{~}{L}_i,\stackrel{~}{L}_k)=(/2)^{L_iL_k}\hfill & i<k,\hfill \\ /2\mathrm{id}_{\stackrel{~}{L}_i}\hfill & i=k,\hfill \\ 0\hfill & i>k.\hfill \end{array}$$
Roughly speaking, $`\mu _𝒜^1`$ is the Floer boundary map, $`\mu _𝒜^2`$ the pair-of-pants product, and $`\mu _𝒜^3,\mu _𝒜^4,\mathrm{}`$ Fukaya’s generalizations of that product. Each $`\mu _𝒜^d`$ depends on the choice of a family $`𝐉^{d+1}`$ of almost complex structures on $`M`$, and these choices have to obey certain consistency conditions, which we will now outline (this is joint work of Lazzarini and the author).
In a first step one takes, for each $`1i_1<i_2m`$, a generic family of almost complex structures $`𝐉^2(i_1,i_2):[0;1]𝒥_M`$. As is well known, this causes all solutions of Floer’s equations
(11)
$$u:\times [0;1]M,\{\begin{array}{cc}& 𝐉^2(i_1,i_2,t)du(s,t)=du(s,t)j,\hfill \\ & u(\times \{1\})L_{i_1},u(\times \{0\})L_{i_2},\hfill \\ & u^{}\omega <\mathrm{},\hfill \end{array}$$
where $`j`$ is the standard complex structure on $`\times [0;1]`$, to be regular. From the one-dimensional solution spaces (zero-dimensional after dividing by translation) one builds $`\mu _𝒜^1:CF^{}(\stackrel{~}{L}_{i_1},\stackrel{~}{L}_{i_2})CF^{+1}(\stackrel{~}{L}_{i_1},\stackrel{~}{L}_{i_2})`$.
Next take $`S=D`$, three cyclically ordered marked points $`\zeta _\nu ^3S`$, $`1\nu 3`$, local coordinates $`\psi _\nu ^3:D^+S`$ around them, and set $`S^{}=S\{\zeta _1^3,\zeta _2^3,\zeta _3^3\}`$, all as in the definition of the invariants $`\mathrm{\Phi }_{rel}`$. Choose for each $`1i_1<i_2<i_3m`$ a generic family $`𝐉^3(i_1,i_2,i_3):S^{}𝒥_M`$, such that for $`s0`$ and $`t[0;1]`$,
(12)
$$𝐉^3(i_1,i_2,i_3,\psi _\nu ^3(e^{\pi (s+it)}))=\{\begin{array}{cc}𝐉^2(i_\nu ,i_{\nu +1},t)\hfill & \nu =1,2,\hfill \\ 𝐉^2(i_1,i_3,1t)\hfill & \nu =3.\hfill \end{array}$$
Denote by $`I_\nu ^3`$, $`1\nu 3`$, the connected components of $`S^{}`$, ordered cyclically such that $`I_1^3`$ lies between $`\zeta _3^3`$ and $`\zeta _1^3`$. The equation defining $`\mu _𝒜^2`$ is
$$u:S^{}M,\{\begin{array}{cc}& 𝐉^3(i_1,i_2,i_3,z)du(z)=du(z)j,\hfill \\ & u(I_\nu ^3)L_{i_\nu }\text{for }\nu =1,2,3\text{,}\hfill \\ & u^{}\omega <\mathrm{}.\hfill \end{array}$$
Condition (12) causes this to agree with a suitable equation (11) in the tubular coordinates $`\psi _\nu ^3(e^{\pi (s+it)})`$ on each end of $`S^{}`$.
The additional ingredient in the definition of the products of order $`d>2`$ are “moduli parameters” as the complex structure of the domain changes. Let $`𝒞^{d+1}(D)^{d+1}`$ be the configuration space of $`d+1`$ distinct, numbered and cyclically ordered points on $`D`$. The moduli space $`^{d+1}`$ and the universal disc bundle $`\mathrm{SS}^{d+1}`$ over it are defined as
$$\mathrm{SS}^{d+1}=𝒞^{d+1}\times _{\mathrm{Aut}(D)}D^{d+1}=𝒞^{d+1}/\mathrm{Aut}(D),$$
where $`\mathrm{Aut}(D)PSL(2,)`$ is the holomorphic automorphism group. Each fibre $`\mathrm{SS}_r^{d+1}`$ carries a canonical complex structure, and there are canonical sections $`\zeta _\nu ^{d+1}:^{d+1}\mathrm{SS}^{d+1}`$, $`1\nu d+1`$, such that the points $`\zeta _\nu ^{d+1}(r)\mathrm{SS}_r^{d+1}`$ are distinct and cyclically ordered for each $`r`$. Write
$$\mathrm{SS}^{d+1,}=\mathrm{SS}^{d+1}(_\nu \mathrm{im}\zeta _\nu ^{d+1}),\mathrm{SS}_r^{d+1,}=\mathrm{SS}^{d+1,}\mathrm{SS}_r^{d+1}.$$
For each $`1i_1<\mathrm{}i_{d+1}m`$ one has to choose a family $`𝐉^{d+1}(i_1,\mathrm{},i_{d+1}):\mathrm{SS}^{d+1,}𝒥_M`$, subject to two kinds of conditions.
* Compatibility with $`𝐉^2`$. This requires a preliminary choice of maps $`\psi _\nu ^{d+1}:^{d+1}\times D^+\mathrm{SS}^{d+1}`$, $`1\nu d+1`$, such that $`\psi _\nu ^{d+1}(r,)`$ provides local coordinates around $`\zeta _\nu ^{d+1}(r)`$ for each $`r^{d+1}`$. Then the conditions are similar to (12), requiring $`𝐉^{d+1}(i_1,\mathrm{},i_{d+1},\psi _\nu ^{d+1}(r,z))`$ to be determined by the previously chosen $`𝐉^2`$ for small $`|z|`$.
* Compatibility with $`𝐉^{e+1}`$ for $`2e<d`$. For this it is necessary to consider the Deligne-Mumford compactification of $`^{d+1}`$. Each stratum at infinity is a product of lower order spaces $`^{e+1}`$, and for a point $`r^{d+1}`$ sufficiently close to one such stratum, the fibre $`\mathrm{SS}_r^{d+1,}`$ is built by gluing together fibres of $`\mathrm{SS}^{e+1,}`$ for the various occurring $`e`$. The precise condition on $`𝐉^{d+1}(i_1,\mathrm{},i_{d+1})|\mathrm{SS}_r^{d+1,}`$ is too complicated to be written down here, but informally it says that this should be built up from the $`𝐉^{e+1}`$ in a corresponding way.
Let $`I_{r,\nu }^{d+1}`$, $`1\nu d+1`$, be the connected components of $`\mathrm{SS}_r^{d+1,}`$, ordered cyclically so that $`I_{r,1}^{d+1}`$ lies between $`\zeta _{d+1}^{d+1}(r)`$ and $`\zeta _1^{d+1}(r)`$. The consistency conditions leave enough freedom to make solutions of the equation
$$r^{d+1},u:\mathrm{SS}_r^{d+1,}M,\{\begin{array}{cc}& 𝐉^{d+1}(i_1,\mathrm{},i_{d+1},z)du(z)=du(z)j,\hfill \\ & u(I_{r,\nu }^{d+1})L_{i_\nu }\text{for }\nu =1,\mathrm{},d+1,\hfill \\ & u^{}\omega <\mathrm{}.\hfill \end{array}$$
regular for generic $`𝐉^{d+1}(i_1,\mathrm{},i_{d+1})`$. Counting such solutions defines
$$\mu _𝒜^d:CF^{}(\stackrel{~}{L}_{i_d},\stackrel{~}{L}_{i_{d+1}})\mathrm{}CF^{}(\stackrel{~}{L}_{i_1},\stackrel{~}{L}_{i_2})CF^{+2d}(\stackrel{~}{L}_{i_1},\stackrel{~}{L}_{i_{d+1}}).$$
The dependence of $`\mathrm{𝐿𝑎𝑔}^{}(\stackrel{~}{\mathrm{\Gamma }})`$ on the choice of almost complex structure can be analyzed using a one-parameter family argument. Since only finitely many moduli spaces are involved, the $`A_{\mathrm{}}`$-structure is subject to a finite number of changes in the family. At each of these exceptional times one can produce a quasi-isomorphism relating the old $`A_{\mathrm{}}`$-structure with the new one.
###### Remark 6.1.
As a technical point, note that directedness allows us to bypass some problems which plague Fukaya’s original setup, having to do with the chain complexes underlying $`HF(L,L)`$ and unit elements in them.
The next step is isotopy invariance which, as always in Floer theory, is also used to extend the definition to configurations which are not in general position.
###### Proposition 6.2.
Let $`\stackrel{~}{\mathrm{\Gamma }}=(\stackrel{~}{L}_1,\mathrm{},\stackrel{~}{L}_m)`$ be a graded Lagrangian configuration in general position. Take $`l\{1,\mathrm{},m\}`$, a symplectic automorphism $`\varphi `$ isotopic to the identity, and a grading $`\stackrel{~}{\varphi }`$ which is zero near $`𝒥_M\times M`$, such that $`\stackrel{~}{\mathrm{\Xi }}=(\stackrel{~}{L}_1,\mathrm{},\stackrel{~}{L}_{l1},\stackrel{~}{\varphi }(\stackrel{~}{L}_l),\mathrm{},\stackrel{~}{\varphi }(\stackrel{~}{L}_m))`$ is again in general position. Then there is a quasi-isomorphism $`F:\mathrm{𝐿𝑎𝑔}^{}(\stackrel{~}{\mathrm{\Gamma }})\mathrm{𝐿𝑎𝑔}^{}(\stackrel{~}{\mathrm{\Xi }})`$.
We will spend a moment discussing the structure of the proof, since it is a good example of arguments involving directed Fukaya categories. Recall that an $`A_{\mathrm{}}`$-functor $`F:𝒜`$ consists of a map $`F:\mathrm{Ob}𝒜\mathrm{Ob}`$, chain maps $`F^1:hom_𝒜(X,Y)hom_{}(FX,FY)`$ for $`X,Y\mathrm{Ob}𝒜`$, and multilinear “higher order terms” $`F^d`$, $`d2`$. In the present case, the map on objects is the obvious one, and the nontrivial chain maps
$$F^1:CF^{}(\stackrel{~}{L}_i,\stackrel{~}{L}_k)CF^{}(\stackrel{~}{L}_i,\stackrel{~}{\varphi }(\stackrel{~}{L}_k)),i<l\text{ and }kl,$$
are those underlying the continuation homomorphisms, so that $`F`$ is automatically a quasi-isomorphism. The main effort goes into defining higher order terms which satisfy the equations for an $`A_{\mathrm{}}`$-functor.
### 6.2.
We will now describe the relation between Hurwitz moves of $`\stackrel{~}{\mathrm{\Gamma }}`$ and mutations of $`D^b\mathrm{𝐿𝑎𝑔}^{}(\stackrel{~}{\mathrm{\Gamma }})`$. Proposition 6.2 says that isotopies of $`\stackrel{~}{\mathrm{\Gamma }}`$ result in a quasi-isomorphism. Shifting the gradings $`\stackrel{~}{L}_i`$ obviously corresponds to the first mutation in Definition 5.2.
###### Lemma 6.3.
$`\mathrm{𝐿𝑎𝑔}^{}(c\stackrel{~}{\mathrm{\Gamma }})`$ is quasi-isomorphic to $`c\mathrm{𝐿𝑎𝑔}^{}(\stackrel{~}{\mathrm{\Gamma }})`$.
The proof relies on the $`/(d+1)`$-action on $`^{d+1}`$ given by a cyclic shuffle of the marked points (this symmetry had not been used in the definition of directed Fukaya categories).
###### Theorem 6.4.
$`r\mathrm{𝐿𝑎𝑔}^{}(\stackrel{~}{\mathrm{\Gamma }})`$ is quasi-isomorphic to $`\mathrm{𝐿𝑎𝑔}^{}(r\stackrel{~}{\mathrm{\Gamma }})`$.
To see why this is plausible, set $`𝒜=\mathrm{𝐿𝑎𝑔}^{}(\stackrel{~}{\mathrm{\Gamma }})`$, and let $`(Z^1,\mathrm{},Z^m)`$ be the objects of $`r𝒜`$. The cohomology $`H(\mathrm{ℎ𝑜𝑚}_{r𝒜}(Z^i,Z^k),\mu _{r𝒜}^1)`$, $`i<k`$, is
$$\{\begin{array}{cc}HF^{}(\stackrel{~}{L}_i,\stackrel{~}{L}_k)\hfill & i<km2\text{,}\hfill \\ HF^{}(\stackrel{~}{L}_i,\stackrel{~}{L}_{m1})\hfill & im2\text{}k=m\text{,}\hfill \\ HF^{}(\stackrel{~}{\tau }_{L_{m1}}(\stackrel{~}{L}_m),\stackrel{~}{L}_{m1})\hfill & i=m1\text{}k=m\text{,}\hfill \\ H(\mathrm{Cone}(\mu _𝒜^2:CF^{}(\stackrel{~}{L}_{m1},\stackrel{~}{L}_m)CF^{}(\stackrel{~}{L}_i,\stackrel{~}{L}_{m1})\hfill & \\ CF^{}(\stackrel{~}{L}_i,\stackrel{~}{L}_m)))\hfill & im2\text{}k=m1\text{.}\hfill \end{array}$$
This is just (9) except that in writing down the third case we have used Poincaré duality and (4). As we saw when discussing Theorem 3.3, the cone in the last line is isomorphic to $`HF^{}(\stackrel{~}{L}_i,\stackrel{~}{\tau }_{L_{m1}}(\stackrel{~}{L}_m))`$. Therefore all cohomology groups are in fact isomorphic to the corresponding ones in $`\mathrm{𝐿𝑎𝑔}^{}(r\stackrel{~}{\mathrm{\Gamma }})`$. The remainder of the proof, as in Proposition 6.2, consists in extending this to a full-fledged $`A_{\mathrm{}}`$-functor.
By the general theory of mutations, what we have shown implies that if two graded Lagrangian configurations in $`M`$ are Hurwitz equivalent, their directed Fukaya categories have equivalent derived categories. Combining this with Picard-Lefschetz theory yields the following consequence:
###### Corollary 6.5.
Let $`(E,\pi )`$ be an exact Morse fibration over $`D`$, with a relative Maslov map $`\delta _{E/D}`$. Make an admissible choice of paths, and let $`\mathrm{\Gamma }`$ be the corresponding distinguished basis of vanishing cycles in a fibre. Choose any gradings $`\stackrel{~}{\mathrm{\Gamma }}`$ and form $`D^b\mathrm{𝐿𝑎𝑔}^{}(\stackrel{~}{\mathrm{\Gamma }})`$. This is independent of all choices up to equivalence, and hence is an invariant of $`(E,\pi )`$ and $`\delta _{E/D}`$.
### 6.3.
There is a computational aspect which Corollary 6.5 fails to convey, and which we will explain by giving an example. Let $`M`$ be an exact symplectic four-manifold with $`2c_1(M)=0`$, and $`(L_1,L_2)`$ two Lagrangian spheres in $`M`$. $`L_2^{}=\tau _{L_1}^2(L_2)`$ and $`L_2`$ are always isotopic as smooth submanifolds. There are two cases, $`L_1=L_2`$ and $`L_1L_2=\mathrm{}`$, in which $`L_2^{}`$ is also Lagrangian isotopic to $`L_2`$ for obvious reasons, but in general this is false:
###### Proposition 6.6.
Suppose that $`L_1,L_2`$ intersect transversally, with $`|L_1L_2|3`$, and that the local intersection numbers at all points are the same. Then $`L_2^{}`$ is not Lagrangian isotopic to $`L_2`$.
The proof goes as follows. Choose a Maslov map $`\delta _M`$ and gradings $`\stackrel{~}{L}_1,\stackrel{~}{L}_2`$. The directed Fukaya category $`𝒜=\mathrm{𝐿𝑎𝑔}^{}(\stackrel{~}{\mathrm{\Gamma }})`$ of the configuration $`\stackrel{~}{\mathrm{\Gamma }}=(\stackrel{~}{L}_1,\stackrel{~}{L}_1,\stackrel{~}{L}_2,\stackrel{~}{L}_2)`$ is determined up to quasi-isomorphism by the graded vector space $`R=HF^{}(\stackrel{~}{L}_1,\stackrel{~}{L}_2)`$ together with the degree two maps $`q_1,q_2\mathrm{End}(R)`$ given by the pair-of-pants product with the unique nontrivial element in $`HF^2(\stackrel{~}{L}_1,\stackrel{~}{L}_1)`$ resp. $`HF^2(\stackrel{~}{L}_2,\stackrel{~}{L}_2)`$. These satisfy $`q_1q_2=q_2q_1`$ and $`q_1^2=q_2^2=0`$. Lemma 6.3 and Theorem 6.4 give an explicit sequence of mutations which transforms $`𝒜`$ into the directed Fukaya category $``$ associated to the Hurwitz equivalent configuration
$$ccrc^1rc^1\stackrel{~}{\mathrm{\Gamma }}=(\stackrel{~}{\tau }_{L_1}^2(\stackrel{~}{L}_2),\stackrel{~}{L}_1,\stackrel{~}{L}_1,\stackrel{~}{L}_2).$$
In particular this determines $`HF^{}(\stackrel{~}{\tau }_{L_1}^2(\stackrel{~}{L}_2),\stackrel{~}{L}_2)`$, since that is a morphism space in $`H()`$. We will not write down the actual computation; the outcome is the total cohomology of the complex
$$/2/2[2]\stackrel{(\mathrm{id},q_2)}{}\mathrm{End}(R)\stackrel{𝜓}{}\mathrm{End}(R)[2],$$
where the second arrow is $`\psi (x)=q_1xxq_1`$. Since $`\psi ^2(x)=2q_1xq_1=0`$, linear algebra tells us that the dimension of $`\mathrm{coker}(\psi )`$ is $`(dimR)^2/2`$. With the assumption $`dimR3`$ this implies that $`HF(L_2^{},L_2)`$ is bigger than $`HF(L_2,L_2)`$, which completes the argument.
###### References.
The presence of $`A_{\mathrm{}}`$-structures in Floer theory was discovered by Fukaya ; at the time of writing, there are still no published proofs of the basic analytic results. Our approach to transversality is joint work with Lazzarini. The formal resemblance between Hurwitz moves and mutations was pointed out by Kontsevich. With the benefit of hindsight one can see that Theorem 6.4 is, in conjectural form, implicit in his discussion of that phenomenon . Proposition 6.6 takes up a topic discussed in and . However, the result itself is new, and can apparently not be proved by the elementary methods used in those earlier papers.
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# I The quadrupole moment of the J = 1+ state in 6Li
## I The quadrupole moment of the J = 1<sup>+</sup> state in <sup>6</sup>Li
Whereas the quadrupole moment of the deuteron is positive(Q = +2.74mb), that of the J=1<sup>+</sup> state of <sup>6</sup>Li is negative, Q=-0.818(17)mb. The magnetic moment of the deuteron is $`\mu =0.85741`$ nm while that of <sup>6</sup>Li is 0.822 nm.
There appears to be a big discrepancy between cluster model calculations and the shell model calculations. In nearly all cluster model calculations Q comes out positive. However in many shell model calculations Q comes out negative, sometimes too negative. This is an important problem that deserves further attention. See for example arguments in the literature between the cluster group and the shell model group. See also the recent compendium of A=6 by D.R. Tilley et. al..
For example in a modern shell model approach by Forest et. al. gets about -8 mb for Q, a factor of 10 too large but of the correct sign. On the other hand in a dynamical microscopic three cluster description of <sup>6</sup>Li where the clusters are $`\alpha `$, n, and p the result is Q = 2.56 mb.
In shell model calculations that we performed we started with 2 valence particles in the 0p shell (0$`\mathrm{}\omega `$). Then we allowed up to 2 $`\mathrm{}\omega `$ and then up to 4 $`\mathrm{}\omega `$ excitations. In the 0 $`\mathrm{}\omega `$ space if you do not have a tensor interaction Q comes out positive. With a ’realistic’ tensor interaction Q comes out negative but too negative Q=-3.5mb. However with a former student Zheng, who at Arizona also developed the no core approximation with Barrett et. al., we showed that when higher shell admixtures were admitted Q became smaller in magnitude and closer to experiment as shown in the following table.
The results are shown in the following table
| SPACE | Q(mb) | $`\mu `$(mm) |
| --- | --- | --- |
| 0 $`\mathrm{}\omega `$ | -3.60 | 0.866 |
| 2 $`\mathrm{}\omega `$ | -2.51 | 0.848 |
| 4 $`\mathrm{}\omega `$ | -0.085 | 0.846 |
| Experiment | -0.82 | 0.822 |
Note that the shell model calculations cannot get the magnetic moment low enough. With up to 4 $`\mathrm{}\omega `$ admixtures we actually overshoot and get a quadrupole moment that is too small but still negative. Some cluster models appear to explain the low magnetic moment.
An excellent discussion of many shell model calcultion of Q and $`\mu `$ has been given by Karataglidis et. al . The value of Q that they obtain with what they call the “Zheng” interaction in the up to 0,2,4 and 6 $`\mathrm{}\omega `$ spaces are -2.64,-2.08,-0.12, and 0.17 mb respectively. Thus they get Q to become positive at the 6 $`\mathrm{}\omega `$ level. But then they quote Zheng et. al as getting a value of -0.67mb in the same 6 $`\mathrm{}\omega `$ space. It is not clear why the two calculations give different answers. The changes in $`\mu `$ in ref are more moderate 0.869,0.848, 0.845 and 0.840 nm in the up to 0,2,4 and 6 $`\mathrm{}\omega `$ spaces.
Looking at all the calculations by all groups (including our own), the situation is certainly confusing, and the problem deserves further attention. This is certainly a basic problem, the deuteron embedded in the nuclear medium. This problem has wider implications whether or not there is T=0 pairing can depend on how higher order configurations affect the tensor interaction in the valence space.
## II Absence of Low Lying Intruders in <sup>8</sup>Be and the $`\alpha `$ particle model
The 0<sup>+</sup> bandhead for low lying intruders in <sup>16</sup>0, <sup>12</sup>C, and <sup>10</sup>Be are at 6.05 MeV, 7.65 MeV and 6.18 MeV respectively. In <sup>16</sup>0 and <sup>12</sup>C, these are predominantly 4 particle 4 hole excitations. In <sup>12</sup>C, we identify the intrinsic state as a linear chain. In the 7th edition of the ’Tables of Isotopes’ possible intruders in <sup>8</sup> Be were indicated, a J=0<sup>+</sup> state at 6 MeV and J=2<sup>+</sup> state at 9 MeV.
In shell model calculations allowing 2 particle 2 hole excitations we were able with a quadrupole-quadrupole force to get a J=0<sup>+</sup> state at 9.7 MeV in <sup>10</sup>Be, too high but in the right ballpark. But we could not get low lying intruders in <sup>8</sup>Be below 30 Mev. We used a deformed oscillator model to show why one gets intruders in <sup>12</sup>C and <sup>10</sup>Be but not <sup>8</sup> Be.
But perhaps the simplest explanation as suggested to us by E. Vogt is given by the $`\alpha `$ particle model. In <sup>12</sup>C we can rearrange the $`\alpha `$ particles from a triangle to a linear chain. In <sup>8</sup>Be we have only 2 $`\alpha `$ particles. One can get a rotational band by having the 2 $`\alpha `$’s rotate around each other but that is all.
The mere existence of these intruder states is of astrophysical importance. In the beta decay <sup>8</sup>B$`^8`$Be + e$`{}_{}{}^{+}+\nu `$ one goes from a J=1<sup>+</sup> T=1 to J=2<sup>+</sup> states. This is the famous ’Ray Davis’ neutrino. If there were a 2<sup>+</sup> state at 9 MeV then there would be more high energy $`\alpha `$’s than there would be if the decay were to the $`2_1^+`$ state at 3.04 MeV. The alpha spectrum from the decay of <sup>8</sup>Be seems to show more high energy alphas, but we would say that they are not due to low lying intruders.
## III Clustering and shell model in the $`f_{7/2}`$ region
In a previous cluster conference in Santorini (1993) a spectrum of <sup>44</sup>Ti was shown in an $`\alpha `$ cluster model. The spectrum looked reasonable except that there was a wide gap between the $`10^+`$ and $`12^+`$ states. However, these states are sufficiently close together that the $`12^+`$ state is isomeric. In a single j shell basis (j=f<sub>7/2</sub>) <sup>52</sup>Fe is the 4 hole system and it should have an identical spectrum to that of <sup>44</sup>Ti provided the same interaction is used. However in <sup>52</sup>Fe the $`12^+`$ lies below the $`10^+`$ state. It is extremely isomeric and has a lifetime of 12 minutes.
We have studied this and other topics by calculating the spectrum of <sup>44</sup>Ti (<sup>52</sup>Fe) with a variety of interactions designated as Model X. (See Tables I,II)
Model I: Use the spectrum of <sup>42</sup>Sc as input (particle-particle) Identify $`<(j^2)^JV(j^2)^J>=E(J)`$ experimental. For isospin T=0 J can be 1,3,5 and 7 while for T=1 J is even 0,2,4, and 6.
Model II: Use the spectrum of <sup>54</sup>Co as input (hole - hole). If there were no configuration mixing these two spectra would be identical. However, there are some differences eg the $`7^+`$ state is much lower in <sup>54</sup>Co than in <sup>42</sup>Sc.
Model III: Now we play games. We want to find out how important are the T=0 matrix elements for the structure of the nuclei. (e.g. Is T=0 pairing important?) Noticing that in <sup>42</sup>Sc the J=2,3 and 5 states are nearly degenerate in this model we set all the T=0 matrix elements to be the same and to all equal $`E(2^+)=1.5863`$ MeV.
In model III we then have $`V^{T=1}=V(^{42}Sc)^{T=1}`$ J =0,2,4,6 and $`V^{T=0}=`$ constant$`=E(2^+)`$ J=1,3,5,7. We can then write $`V^{T=0}`$=constant $`(1/4t_1t_2)`$. We can then write $`V^{T=0}=c(1/4t_1t_2)`$ where c is a constant. Hence $`_{i<j}V_{ij}^{T=0}=c/8(n(n1)+6)c/2T(T+1)`$. This means that the spectrum of states of a given isospin e.g. T=0 in <sup>44</sup>Ti(<sup>52</sup>Fe) is independent of what the constant is, it might as well be zero. Of course the relative splitting of T=1 and T=0 states will be affected. Model III will be the standard from which we derive Model IV.
Model IV: Relative to the degenerate case above, we now move the $`J=1^+`$ state down in energy to 0.5863 MeV. Our motivation is based on numerous discussions about the importance of T=0 S=1 “pairing” in nuclei. We hope to simulate the T=0 pairing by this lowering.
Model V: Relative to Model III we bring the J=1<sup>+</sup> and J=7<sup>+</sup> states down to an energy of 0.5863 MeV but keep the J=3<sup>+</sup> and 5<sup>+</sup> at E=E(2<sup>+</sup>)=1.5863 MeV. This spectrum is very close to that of <sup>42</sup>Sc.
## IV Discussion of Results
Let us first compare Model III (all T=0 matrix elements are degenerate) with Model I (spectra of <sup>42</sup>Sc). As already mentioned, making T=0 matrix elements degenerate is equivalent to making them zero as far as T=0 states are concerned.
The main difference is that the states with J=6,4,7, and 8 come down in energy as does J=9<sup>+</sup>. Also the 12-10 gap is a bit greater than for the <sup>42</sup>Sc spectra case, reminiscent of the $`\alpha `$ particle model. The J=9<sup>+</sup> state is below the 10<sup>+</sup> and 12<sup>+</sup> in the degenerate case. Clearly it is the high energy side of the spectrum which is most sensitive to the change from experimental spectrum to the “T=0 degenerate” case.
Despite the changes, we can say that the T=1 two body matrix elements give the dominant structure of the spectrum whilst the T=0 matrix elements provide the fine tuning.
We now compare Model IV with Model III. The only difference is that we break the T=0 degeneracy by lowering the J=$`1^+`$ state from 1.5863 MeV to 0.5863 Mev. We hope that this simulates to some extent T=0 S=1 neutron-proton pairing. The change from degenerate case is not that large. There is a tendency to go towards the spectrum of <sup>42</sup> Sc. The J=3,5,7, and 8 states are raised somewhat in energy. However it is hard to find a clear signature of this S=1 pairing.
Not shown is Model V where we bring down both the J=1<sup>+</sup> and 7<sup>+</sup> states to 0.5863 MeV keeping J=3 and 5 at E(2) = 1.5863 MeV. This input spectrum is close to that of <sup>42</sup> Sc so that it is not surprising that the <sup>44</sup>Ti spectrum is likewise close.
We lastly consider the results using the spectrum of <sup>54</sup>Co. This was done some time ago by Geesaman Note that there are significant changes, all at the high energy high angular momentum part of the spectrum. Relative to the <sup>42</sup>Sc case the 10<sup>+</sup> and 12<sup>+</sup> states are down in energy with the 12<sup>+</sup> below the 10<sup>+</sup> thus leading to a long lifetime for the 12<sup>+</sup> state. Note that the 9<sup>+</sup> state is now at a much high energy than the 10<sup>+</sup> or 12<sup>+</sup>. Recently the 10<sup>+</sup> state in <sup>52</sup>Fe, which lies above the 12<sup>+</sup> has been found in <sup>52</sup>Fe by Ur et. al. . It would also be of interest to find the 9<sup>+</sup> state.
## V Many particle, Many Hole States in <sup>40</sup>Ca
This is a topic we discussed in previous cluster meetings so we will be brief. We just want to remind the reader that there are all sorts of many particle-many hole highly deformed states in <sup>40</sup>Ca. One cannot properly describe <sup>40</sup>Ca in a cluster model consisting of <sup>36</sup>Ar plus an alpha particle. At the very least one has to start with <sup>32</sup>S plus two alpha particles.
In a Skyrme Hartree-Fock calculations (SK III) we obtain a near degeneracy of the 4p-4h and 8p - 8h intrinsic state. The respective energies are 12.1 MeV and 11.4 MeV. The 8p-8h intrinsic state energy is lower than the 4p-4h. By the time projection and pairing are included, the 4p-4h comes lower than the 8p-8h (6.85 MeV vs 8.02 MeV) in agreement with the order the J = 0 excitation energies of 3.0 and 5.1 Mev. Pairing will lower the states even more. We actually found many more deformed states of the form np-nh n=2,3,4,5,6,7 and 8. The intrinsic states are nearly degenerate in energy - we called this a deformation condensate. We also found for these states that the deformation parameter was approximately proportional to n i.e. the value of $`\beta `$ for 8p-8h is approximately twice the value of $`\beta `$ for 4p-4h.
In <sup>80</sup>Zr one of the np-nh states becomes the ground state. This is the 12p-12h state which has a calculated value of $`\beta =0.4`$ A more superdeformed 16p-16h state with $`\beta =0.6`$ is calculated to be at an excitation energy of about 8 MeV.
## VI Two different Views of the $`f_{7/2}`$ region
In March 2000 issue of the Physical Review C 61 there are two papers side by side. One is by our group and one by H. Hasegawa and K. Koneko . We both do calculations in the f<sub>7/2</sub> shell. We emphasize shell model behavior whilst the other authors the $`\alpha `$ cluster behaviors, even though their model space is limited to f<sub>7/2</sub>.
The other authors point out that we can get an excellent approximation to the ground states of n<sub>p</sub> = n<sub>n</sub>=2m nuclei (n<sub>p</sub> is the number of protons e.t.c.).
$$|(f_{7/2})^{4m}I=T=0>=\frac{1}{\sqrt{N_0}}(\alpha _0^{})^m|A_0>$$
(1)
where $`(\alpha _0^{})`$ creates a 4 nucleon cluster.
$$\alpha _0^{}=\underset{J,\tau }{}(J\tau ,J\tau :I=T=0)(A_{J\tau }^{}A_{J\tau }^{})_{I=0T=0}$$
(2)
For <sup>48</sup>Cr this approximation give -32.04 MeV for the ground state energy where as the exact value is 32.70 MeV.
We on the other hand have emphasized the shell model aspects. In the previously mentioned paper, we find an approximation for the excitation energies of single and double analog states in the f<sub>7/2</sub> region and in a an earlier paper “Fermionic Symmetries: Extension of the two to one relationship between spectra of even even and neighboring odd mass nuclei” we noted two things.
A. There is often a two to one relation between spectra of even-even and even odd nuclei, and in some cases the single j shell model predicts this.
B. Excitation energies of analog state are approximately the same if the neutron excess (or equivalently the ground state isospin) is the same.
The above results can be parametrized by the following formulae
SINGLE ANALOG EXCITATION (SA)
$$E(SA)=b(T+X)$$
(3)
DOUBLE ANALOG EXCITATION (DA)
$$E(DA)=2b(T+X+1/2)$$
(4)
This formula will give a two to one ratio for E(DA)/E(SA) for (<sup>44</sup>Ti,<sup>43</sup>Ti), (<sup>51</sup>Cr,<sup>50</sup>Cr), (<sup>47</sup>Sc,<sup>48</sup>Ti) etc.
The experimental SA and DA are shown in the following table.
| T=0 | <sup>44</sup>Ti (9.340), <sup>48</sup>Cr (8.75), <sup>52</sup>Fe (8.559) |
| --- | --- |
| T=1/2 | <sup>43</sup>Sc (4.274)<sup>a</sup>, <sup>43</sup>Ti (4.338)<sup>a</sup>, <sup>45</sup>Ti (4.176), <sup>49</sup>Cr(4.49), |
| | <sup>51</sup>Mn (4.451), <sup>53</sup>Co (4.390), <sup>53</sup>Fe (4.250) |
| T=1 | <sup>46</sup>Ti (14.153), <sup>50</sup>Cr (13.222) |
| T=3/2 | <sup>45</sup>Sc (6.752)<sup>a</sup>, <sup>47</sup>Ti (7.187), <sup>51</sup>Cr (6.611) |
| T=2 | <sup>48</sup>Ti (17.379) |
| T=5/2 | <sup>47</sup>Sc (8.487)<sup>a</sup>, <sup>49</sup>Ti (8.724) |
$`a`$ \- obtained from binding energy data.
In Table II we compare the theoretical single j shell calculations with the linear formula. We take b = 2.32 MeV X=1.30. Note that in the SU(4) limit X=2.5. The fact that SU(3) gives the linear formula is not sufficient for it to be the correct theory. For a simple monopole-monopole interaction a+bt(1)t(2) X=1.
Some of the two to one ratio’s hold rigoursly in the single j shell model. This holds for 3 particle and 4 particle systems or 3 holes and 4 holes. eg (<sup>43</sup>Ti,<sup>44</sup>Ti), (<sup>43</sup>Sc,<sup>44</sup>Ti),(<sup>53</sup>Fe,<sup>52</sup>Fe). Here not only single or double analog but all the J=j states in the odd spectrum are at half the energy of the J=0+ states in the even system.
Some of the relations hold approximately in the single j shell model eg for (<sup>45</sup>Sc,<sup>46</sup>Ti) and for the cross conjugate pair (<sup>51</sup>Cr,<sup>50</sup>Cr) we would get a 2 to 1 ratio if seniority four states could be neglected and one only had v=0 and v=2.
Miraculously the 2 to 1 ratio holds remarkably well experimentally for (<sup>51</sup>Cr,<sup>50</sup>Cr) - the values are 6.511 and 13.022 MeV respectively, this despite the fact that in the single j shell it should only hold approximately. Ironically the simplest system for which 2 to 1 should hold exactly does not work so well. That is to say for (<sup>43</sup>Ti,<sup>44</sup>Ti) the values are 4.338 and 9.340 MeV. When configuration mixing is included, agreement with the deviation is explained. This might be an example of a 4 particle clustering. For the hole system (<sup>53</sup>Fe, <sup>52</sup>Fe) on the other hand 2 to 1 works much better.
The fact that there is in general a close relation between even-even and even-odd puts to question whether in those many cases there is any $`\alpha `$ particle clustering.
The single j shell calculation does not predict exact equality for the S.A. excitation energies in <sup>43</sup>Sc and <sup>45</sup>Ti. The relative values are very close however 4.142 and 4.112 MeV respectively. This is fascinating. We take <sup>43</sup>Sc, jam a deuteron into it to form <sup>45</sup>Ti and it seems hardly to make any difference for S.A. excitations.
## VII Two views of cross conjugate relations
In the single j shell model, the spectra of cross conjugate nuclei should be identical. (for j<sup>n</sup> states.). A cross conjugate nucleus is one obtained by changing protons into neutron holes and neutrons into proton holes. The cross conjugate of <sup>46</sup>Ti is <sup>50</sup>Cr. Let us compare the spectra
| J | | <sup>46</sup>Ti | | <sup>50</sup>Cr | | Ratio |
| --- | --- | --- | --- | --- | --- | --- |
| 0 | | 0.000 | | 0.000 | | |
| 2 | | 0.889 | | 0.787 | | 0.8853 |
| 4 | | 2.010 | | 1.884 | | 0.9373 |
| 6 | | 3.297 | | 3.164 | | 0.9597 |
| 8 | | 4.896 | | 4.740 | | 0.9681 |
The fits are very good. The <sup>50</sup>Cr excitations are slightly smaller - it could be a universal A dependence. Where does the remarkable agreement leave room for $`\alpha `$ clustering?
However we can look for other things besides the spectra. In a recent experiment theory collaboration where the leading experimentalists were N. Koller and H.A. Speidel good agreement was obtained for g(2<sup>+</sup>) in <sup>50</sup>Cr but bad agreement for <sup>46</sup>Ti. The shell model predicts a high value of g(2<sup>+</sup>) and g(4<sup>+</sup>) for <sup>50</sup>Cr but low values for <sup>46</sup>Ti 0.25 mm. The high values are confirmed for <sup>50</sup>Cr but for <sup>46</sup>Ti the measured g factors are closer to 0.5 which suggest the rotational value g<sub>R</sub>=Z/A. These results suggest that there must be considerable clustering in <sup>46</sup>Ti that is not present in <sup>50</sup>Cr. In general the shell model appears to work better in the upper half of the “f<sub>7/2</sub> shell” than in the lower half. There appears to be much more going on in the lower half and probably this is due to intruder/cluster mixing in with the basic shell model states.
## VIII Closing Remarks
We have provided several examples where Cluster Models and the Shell Model confront each other usually to the mutual benefit of both models even though in the short term there might be some arguments. The two models give the opposite sign for the quadrupole moment of <sup>6</sup>Li, and this has to be resolved. The cluster model provides insight into some results of detailed shell model calculations e.g. why there are no low lying intruders in <sup>8</sup>Be. The low lying intruder states e.g. 4p-4h and 8p-8h in <sup>40</sup>Ca are essentially impossible to calculate in the shell model. However here cluster models and Skyrme-Hartree Fock go together in describing such states. In the f<sub>7/2</sub> region we raise the question (without fully answering it) of how to distinguish symmetry energy from clustering energy. Finally we point out the issue of hidden clustering.
This work was supported in part by the U.S. Department of Energy DE-FG02-95ER-40940
References
1. A. Csoto and R.G. Lovas, Phys. Rev C53, (1996) 1444
2. D.C. Zheng, J.P. Vory and B.R. Barrett, Phys. Rev C53 (1996) 1447
3. D.R. Tilley et. al., Energy Levels of Light Nuclei A=6, preliminary version 2
4. J.L. Forest et. al. Phys Rev C54, (1996) 646
5. L. Zamick, D.C. Zheng and M. Fayache, Phys Rev C51,(1995) 1251
6. D.C. Zheng, B.R. Barrett, J.P. Vary and H. Muther Phys Rev C51, (1995) 2471
7. S. Karataglidis, B.A. Brown, K. Amos and P.J. Dortmans, Phys. Rev C55 (1997) 2826
8. D.C. Zheng, B.R. Barrett, J.P. Vary, W.C. Haxton and C.-L. Song, Phys. Rev C52,2488 (1995)
9. M.S. Fayache, L. Zamick, and Y.Y. Sharon, Phys. Rev C55 (1997) 2389
10. M.S. Fayache, E. Moya de Guerra, P Sarrigueren, Y.Y. Sharon, and L. Zamick, Phys Rev C57, (1998) 2351
11. F.C. Barker, Aust. J. Phys 42, (1989)25; 41(1988)743
12. P. Hodgeson, Survey Talk, Atomic and Nuclear Clusters Proceedings of the Second International Conference at Santorini, Greece June 28th - July 2 1993 Z. Phys A 349,(1994) 197
13. D.F. Geesaman, R. Malmin, R.L. McGrath and J.W. Noe, P.R.L. 34, (1975) 326
14. D.F. Geesaman, R.L. McGrath, J.W. Noe, and R.E. Malmin, Physc. Rev 19, (1979) 1938
15. C.A. Ur et. al. Phys. Rev C 58,(1998) 3163
16. D.C. Zheng, D. Berdichovsky, and L. Zamick, Phys Rev C38 (1988) 437; J. Phys Soz Japan 158 (1989) Suppl. 649; Phys Rev C42 (1990) 1408
17. Y. Durga Devi, S. Robinson, and L. Zamick, Phys Rev C61 (2000) 037305
18. M. Hasegawa and K. Kaneko, Phys Rev C61 (2000) 037306
19. L. Zamick and Y. Durga Devi, Phys Rev C60, (1999) 054317
20. J.D. McCullen, B.F. Bayman and L. Zamick, Phys. Rev. 134B, 515 (1964); Technical Report NYO-9891.
21. R. Ernst, K.-H. Speidel, O. Kenn, U. Nachman, J. Gerber, P. Maier-Komor, N. Benczer-Koller, G. Jacob, G. Kumbartzki, L. Zamick, and F. Nowacki , Phys. Rev. Lett. 84, (2000) 416
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# 1 Introduction
## 1 Introduction
The main aim of this article is to demonstrate that equations of the relativistic quantum theory may have solutions which have no nonrelativstic analogue. Some of such solutions are well-known, e.g. the so-called abnormal solutions. Nakanishi wrote about the latters: ”The non-relativistic common sense does not necessarily remain valid in the relativistic quantum field theory” .
The main ideas of this article were published as a preprint of Tbilisi Mathematical Institute in collaboration with A. Tavkhelidze and L. Vachnadze .
At first the interest to the problem was initiated by the observed narrow peaks around $`1.5MeV`$ in the $`e^+e^{}`$ distribution in the experiments of heavy ion collisions . These resonances are known as the GSI resonances. They are situated above the $`2m_e`$ threshold and it naturally triggered interest to studies of bound states in continuum. In the papers of the Serpukhov and other groups the single-photon exchange quasipotential equations in QED were studied and, relying on the numeric methods, large number of peaks were reported.
The quasipotential equation can be obtained from the Bethe-Salpeter equation after certain approximation. So it was natural to try to cope with the same problem within the framework of the ladder approximation of the Bethe-Salpeter equation itself. It was particularly interesting because there was no such observation made on the earlier stages of studies of the Bethe-Salpeter (BS) equation (for a review see ). However, we think that nobody has searched for bound states in the continuum there. It is worth noting that when the total mass $`M>2m`$ the traditional method — Wick’s rotation is, in general, inadmissible. So the problem must be considered in the Minkowski space. Structure of the ladder BS equation for this case was studied by Günter and it was demonstrated that Wick’s rotation is inadmissible. In the recent years interest to the BS equations in the Minkowski space has increased ,. Below, employing standard methods, we will reduce the BS equation to the system of the one-dimensional equations and study unusual solutions in the continuous spectrum. Note that though further experiments have not confirmed existence of the GSI resonances the possibility of bound states embedded in the continuum is interesting on its own.
## 2 Bethe-Salpeter equation in the Ladder approximation and its reduction
We will consider the BS equation in the ladder approximation for a simple model , of two equal mass ($`m`$) scalar particles exchanging massless scalar meson, the so called Wick-Cutkosky model with the interaction Lagrangian:
$$L_{int}=g_1\varphi _1^{}\varphi _2A+g_2\varphi _2^{}\varphi _1A.$$
Standard equation has the form (in the rest frame):
$$\psi (p)=\frac{\lambda }{\pi ^2i}d^4q\frac{\psi (q)}{\left[\left(pq\right)^2+io\right]\left[\left(\frac{P}{2}q\right)^2m^2\right]\left[\left(\frac{P}{2}+q\right)^2m^2\right]}.$$
(1)
Here $`P=p_1+p_2`$ is the total 4-momentum, $`p=\frac{1}{2}(p_1p_2)`$ and $`q`$ are relative 4-momenta, $`\lambda =\frac{g_1g_2}{(4\pi )^2}`$ is the product of appropriate coupling constants and has dimension of the mass squared. In the rest frame where $`\stackrel{}{P}=0`$ and $`\stackrel{}{p}_1=\stackrel{}{p}_2=\stackrel{}{p}`$ we will use the following notation for the total energy (mass) $`P_0=E=M=2W`$.
Usually the Wick-Cutkosky model is considered as an auxiliary (fictitious) one. But in some reasonable approximation more realistic models reduce to the Wick-Cutkosky one. Indeed, in the single-photon (gluon) approximation in QED (QCD) for fermion and antifermion bound state BS amplitude we have an equation:
$$[m\gamma (\frac{P}{2}+p)]\psi (p)[m+\gamma (\frac{P}{2}p)]==\frac{\lambda _F}{\pi ^2i}d^4q\frac{\gamma _\mu \psi (q)\gamma _\mu }{(pq)^2+io}$$
where $`\lambda _F\left(\frac{g}{4\pi }\right)^2`$ and the Feynman gauge is used.
Usually a new function
$$\mathrm{\Psi }(p)=\left[m\gamma (\frac{P}{2}+p)\right]\psi (p)\left[m+\gamma (\frac{P}{2}p)\right]$$
is introduced and the decomposition into the Dirac structures is considered:
$$\mathrm{\Psi }=\mathrm{\Psi }^S\gamma _5\mathrm{\Psi }^P+\gamma ^\mu \mathrm{\Psi }_\mu ^V\gamma ^\mu \gamma _5\mathrm{\Psi }_\mu ^A\frac{1}{2}[\gamma ^\mu ,\gamma ^\nu ]\mathrm{\Psi }_{\mu \nu }^T$$
Further one can obtain a chain of coupled equations for Dirac structures which can be considerably simplified by taking binding energies small (i.e. for weakly coupled system). All the components of $`\mathrm{\Psi }`$ that survive in that limit of small binding energies ($`\mathrm{\Psi }^P,\mathrm{\Psi }_\mu ^A,\mathrm{\Psi }^V`$) satisfy the same equations:
$$\mathrm{\Psi }^i(p)==\frac{4\lambda _Fm^2}{\pi ^2i}d^4q\frac{\mathrm{\Psi }^i(q)}{\left[(pq)^2+io\right]\left[\left(\frac{P}{2}q\right)^2m^2\right]\left[\left(\frac{P}{2}+q\right)^2m^2\right]}$$
which coincides with the equations of Wick-Cutkosky model if we set $`4\lambda _Fm^2=\lambda `$ or $`\lambda \pi =m^2\alpha `$ where $`\alpha =e^2/4\pi `$. Therefore for weekly coupled systems in single-photon (gluon) exchange approximation in QED (QCD) we have in fact to deal with the equations of Wick-Cutkosky model.
It is well known , that in the case $`W<m`$ and $`\lambda >0`$ (which corresponds to attraction) with small $`\lambda `$ eq.(1) has a discrete spectrum in accordance with the Balmer series for the hydrogen atom. Now our task is to show whether this equation has any discrete solutions when $`W>m`$ i.e. $`M>2m`$.
It is convenient to transform the integral BS equation into a differential one in momentum space by using the identity:
$$\left[\frac{^2}{p_0^2}\frac{^2}{\stackrel{}{p}^2}\right]\frac{1}{(pq)^2+io}=4\pi ^2i\delta ^{(4)}(pq).$$
(2)
We obtain:
$$\left[\frac{^2}{p_0^2}\frac{^2}{\stackrel{}{p}^2}\right]\mathrm{\Psi }(p)=\frac{4\lambda _F\mathrm{\Psi }(p)}{(p_0^2\stackrel{}{p}^2+W^2m^2)^24p_0^2W^2)}$$
(3)
This equation is mathematically equivalent to eq.(1) if $`\mathrm{\Psi }(p)`$ satisfies the following boundary conditions : $`p^2\mathrm{\Psi }(p)`$ remains finite for large values of $`p`$ and $`\mathrm{\Psi }(p)`$ must be finite when $`p=0`$.
We think that like in the Coulomb problem in our problem too it is possible to separate variables due to the hidden additional symmetry of the Coulombic interaction . For earlier attempts of variable separation see ,,.
First of all let us single out the spherical angles in a standard way:
$$\mathrm{\Psi }(p)=\mathrm{\Phi }(p_0,p_s)Y_{lm}(\theta ,\varphi )$$
where the spherical coordinates are defined as follows:
$$p_1=p_s\mathrm{sin}(\theta )\mathrm{cos}(\varphi ),p_2=p_s\mathrm{sin}(\theta )\mathrm{sin}(\varphi ),p_3=p_s\mathrm{cos}(\theta ),p_s=\pm |\stackrel{}{p}|.$$
Then eq.(3) is reduced to the partial differential equation with two variables:
$$[\frac{^2}{p_0^2}\frac{^2}{p_s^2}](p_s\mathrm{\Phi }(p)=\frac{4\lambda _F(p_s\mathrm{\Phi }(p)}{(p_0^2p_s^2+W^2m^2)^24p_0^2W^2)}$$
(4)
Boundary conditions for $`\mathrm{\Phi }(p_0,p_s)`$ can be rewritten as
$$p_s\mathrm{\Phi }(p_0,p_s)|_{p_s0}0p_s\mathrm{\Phi }(p_0,p_s)|_{p_00}<\mathrm{}$$
$$\frac{p^2}{p_s}(p_s\mathrm{\Phi }(p_0,p_s))|_{p0}<\mathrm{}$$
(5)
To investigate the energy spectrum in the region$`W>m`$ let us use the method of bipolar transformations ,,. Peculiarity of our case is that the whole $`(p_0,p_s)`$-plane can be covered only by the means of several such transformations :
$`(\mathit{1})`$ $`p_0={\displaystyle \frac{k\mathrm{sinh}\alpha _1}{\mathrm{cosh}\alpha _1+\mathrm{cosh}\beta _1}},p_s={\displaystyle \frac{k\mathrm{sinh}\beta _1}{\mathrm{cosh}\alpha _1+\mathrm{cosh}\beta _1}}`$
$`(\mathit{2})`$ $`p_0={\displaystyle \frac{k\mathrm{sinh}\alpha _2}{\mathrm{cosh}\alpha _2\mathrm{cosh}\beta _2}},p_s={\displaystyle \frac{k\mathrm{sinh}\beta _2}{\mathrm{cosh}\alpha _2\mathrm{cosh}\beta _2}}`$ $`|\alpha _2||\beta _2|`$
$`(\mathit{3})`$ $`p_0={\displaystyle \frac{k\mathrm{cosh}\alpha _3}{\mathrm{sinh}\alpha _3+\mathrm{sinh}\beta _3}},p_s={\displaystyle \frac{k\mathrm{cosh}\beta _3}{\mathrm{sinh}\alpha _3+\mathrm{sinh}\beta _3}}`$ $`\alpha _3\beta _3`$
$`(\mathit{4})`$ $`p_0={\displaystyle \frac{k\mathrm{cosh}\alpha _4}{\mathrm{sinh}\alpha _4+\mathrm{sinh}\beta _4}},p_s={\displaystyle \frac{k\mathrm{cosh}\beta _4}{\mathrm{sinh}\alpha _4+\mathrm{sinh}\beta _4}}`$ $`\alpha _4=\beta _4.`$
The hyperbolic angles in all these transformations run over the whole infinite intervals:
$$\mathrm{}<\alpha _i,\beta _i<\mathrm{}$$
while the parameter $`k`$ is determined by the condition of separabilty in eq.(3) and equals to $`(W^2m^2)^{1/2}`$.
Action of these transformations is shown in Fig.1. After performing these transformations variables in eq.(4) separate:
$$p_s\mathrm{\Phi }(p_0,p_s)=f(\alpha _i)g(\beta _i)i=1,2,3,4$$
and we obtain the following equations for each of the functions $`f`$ and $`g`$:
(a) For transformations $`(1)(2)`$ (the shaded regions in Fig.1):
$$\left[\frac{d^2}{d\beta ^2}+h_1\frac{l(l+1)}{\mathrm{sinh}^2\beta }\right]g(\beta )=0,$$
(6)
$$\left[\frac{d^2}{d\alpha ^2}+h_1\frac{\lambda /m^2}{a\mathrm{cosh}^2\alpha }\right]f(\alpha )=0;$$
(7)
(b) For transformations $`(3)(4)`$ (the unshaded regions in Fig.1):
$$\left[\frac{d^2}{d\beta ^2}+h_2+\frac{l(l+1)}{\mathrm{cosh}^2\beta }\right]g(\beta )=0,$$
(8)
$$\left[\frac{d^2}{d\alpha ^2}+h_2\frac{\lambda /m^2}{a+\mathrm{sinh}^2\alpha }\right]f(\alpha )=0.$$
(9)
In these equations
$$a\frac{W^2}{m^2}>1$$
and $`h_1,h_2`$ are the separation parameters (new quantum numbers). The fact of separability is provided by the symmetry of relativistic problem in the 4-dimensional world .
Equations obtained above describe one-dimensional motions in the space of hyperbolic angles. Equations for $`g(\beta )`$ are kinematic ones defining the values of separation parameters $`h_{1,2}`$. Equations for $`f(\alpha )`$ contain dynamical quantities $`(\lambda ,a)`$ and hence must allow to find the energy spectrum.
Boundary conditions for $`g(\beta )`$ and $`f(\alpha )`$ follow immediately from (5) and are the same in different regions of the $`(p_0,p_s)`$-plane
$`g(0)=0`$ $`g(\pm \mathrm{})<\mathrm{}`$ (10)
$`f(0)<\mathrm{}`$ $`f(\pm \mathrm{})=0`$ (11)
Note that the boundaries of different regions in $`(p_0,p_s)`$-plane can be approached by taking limits $`\alpha ,\beta \pm \mathrm{}`$ in the relevant transformations $`(1)(4)`$.
## 3 Solutions of the one-dimensional equations
First consider the pair (6)-(7). Eq.(6) with the boundary conditions (10) has the following solution:
$$\genfrac{}{}{0pt}{}{g(\beta )=g_l(\beta )\{\mathrm{sinh}^{l+1}\beta \left[\frac{1}{\mathrm{sinh}\beta }\frac{d}{d\beta }\right]^{l+1}\mathrm{cosh}(\sqrt{h_1}\beta ),h_1>0}{0,h_1<0}$$
(12)
In general, the second equation of this pair can give the relation between energy $`W`$ and the parameter $`h_1`$
$$W=\rho _n(h_1),n=0,1,2\mathrm{}$$
Despite the fact that $`n`$ can take discrete values, energy is not quantized because $`h_1`$ remains to be an artbitrary positive number.
Things are essentially different in the unshaded regions $`(3)(4)`$. Eq.(8) has solutions with positive and negative signs of $`h_2`$. The positive $`h_2`$ are not resulted while the negative ones are quanttizes for $`l>0`$.
$$h_2=(lN)^2,N<l.$$
where $`N`$ takes odd integer values.
Next consider eq.(9) in more detail. Expression
$$V(\alpha )=\frac{\lambda /m^2}{a+\mathrm{sinh}^2\alpha }$$
plays the role of ”potential”.
For positive $`\lambda `$ (i.e. attraction) $`V(\alpha )`$ is the positive definite and convex function of $`\alpha `$ (Fig.2). Obviously we can have only continuous spectrum in this case witth posititve $`h_2`$.
When we take $`\lambda `$ negative (repulsion) $`V(\alpha )`$ becomes negative definitte and concave function of $`\alpha `$ (Fig.2). Therefore in this case there appears principal possibility to have discrete eigenvalues provided the depth of $`V(\alpha )`$ is high enough for a given negative value of $`h_2=(lN)^2`$.
One can examine easily that continuity on the boundaries of shaded and unshaded regions imposes no new constraints on separation parameters $`h_{1,2}`$. So there is no discrete specttrum for positive values of $`\lambda `$ when $`W>m`$.
Now we turn back to eqs.(8)-(9) with $`h_2<0`$ and $`\lambda <0`$. Eq.(8) is a Heun’s differential equation . Its solution with $`h_2=(lN)^2`$ and $`\lambda =|\lambda |<0`$ imposed by boundary and normalization conditions leads to complicated transcendential equations. To estimate qualitatively the coupling constant dependence of the energy eigenvalues one can make use of tthe following approximate expression
$$W^2\frac{|\lambda |}{\nu (\nu +1)}\nu =1,2,\mathrm{}$$
As $`W^2>m^2`$ it follows that
$$\frac{\lambda ^2}{m^2}\nu (\nu +1)=2,6,12,\mathrm{}$$
i.e. discrete levels apeear when $`|\lambda |`$ exceeds some critical values. It is worth noting that for some fixed values of $`\lambda ^2/m^2`$ there is finite number of discrete levels. For example, when
$$2<\frac{\lambda ^2}{m^2}<6$$
there is only a single level. And for
$$6<\frac{\lambda ^2}{m^2}<12$$
the number of levels equals $`2`$, etc.
If we take $`|\lambda |`$ in accordance with Balmer formula mentioned above in the following form $`\lambda \pi =m^2Z^2\alpha `$, then the above evaluation can be cast into
$$\frac{Z^2\alpha }{\pi }>\nu (\nu +1)$$
i.e. $`Z_{min}>30`$.
We want to note tthat eqs.(8)-(9) were derived earlier in but with erroneous sign of $`\lambda `$. Hence author deduced bound states for positive values of $`\lambda `$.
## 4 Normalization condition
Let us now discuss the normalization condition for the BS amplitude, which is nonvanishing only in the unshaded regions and has the form:
$$\psi (p)\frac{(\mathrm{sinh}\alpha +\mathrm{sinh}\beta )^3f(\alpha )g(\beta )}{(W^2m^2)^{3/2}\mathrm{cosh}^3\beta (W^2+m^2\mathrm{sinh}^2\alpha )}Y_{lm}(\theta ,\varphi )$$
(13)
The role of normalization is to provide the correct relation between the aplitude $`\psi `$ and four-point Green function. The exact form of the normalization condititon for the model under consideration looks like
$$id^4p\overline{\psi }(p)G_0^1(p)\psi (p)=\epsilon _B\lambda \frac{dS_B(\lambda )}{d\lambda }$$
(14)
where $`S_B(\lambda )=4W_B^2(\lambda )`$ and $`\epsilon _B`$ is the so called norm factor determining the sign of the pole contribution in the total Green function at $`SS_B(\lambda )`$
$$Gi\epsilon _B\frac{\psi _B(p)\overline{\psi }_B(p)}{SS_B}+\text{Reg}$$
(15)
Integration in the normalization condition (14) is carried over the unshaded region only.
The spectral representation for $`\psi (p)`$ and $`\overline{\psi }(p)`$ that we are going to use has the form :
$$\psi (p_0,\stackrel{}{p})=\frac{1}{2\pi i}_{w_+(\stackrel{}{p})}^{\mathrm{}}𝑑q_0\frac{f(q_0,\stackrel{}{p},W)}{p_0q_0+io}+\frac{1}{2\pi i}_{\mathrm{}}^{w_{}(\stackrel{}{p})}𝑑q_0\frac{g(q_0,\stackrel{}{p},W)}{p_0q_0io};$$
(16)
$$\overline{\psi }(p_0,\stackrel{}{p})=\frac{1}{2\pi i}_{w_+(\stackrel{}{p})}^{\mathrm{}}𝑑q_0\frac{f^{}(q_0,\stackrel{}{p},W)}{p_0q_0+io}+\frac{1}{2\pi i}_{\mathrm{}}^{w_{}(\stackrel{}{p})}𝑑q_0\frac{g^{}(q_0,\stackrel{}{p},W)}{p_0q_0io}$$
(17)
where $`w_{}(\stackrel{}{p})=W(\stackrel{}{p}^2+m^2)^{1/2},w_+(\stackrel{}{p})=w_{}(\stackrel{}{p})`$ are the enndpoints where cuts in the complex $`p_0`$-plane extend from. As long as 3-momentum $`\stackrel{}{p}`$ satisfies $`0p_s^2W^2m^2`$ these two cuts overlap (Fig.3).
The spectral reprezentations (16)-(17) imply
$$\overline{\psi }(z)=\psi ^{}(z^{}),z=\mathrm{}p_0+i\mathrm{}p_0$$
(18)
There are “dangerous” regions in the normalization integral where left-hand and right-hand cuts overlap. This overlapping happens when
$`p_0w_0,0`$ $`\text{for right-hand cut},`$
$`p_00,w_0,0`$ $`\text{for left-hand cut}.`$
In these intervals $`0p_s^2k^2=W^2m^2`$. In the $`(p_0,p_s)`$-plane “dangerous” points are located between the profiles of functions $`w_+(\stackrel{}{p})`$ and $`w_{}(\stackrel{}{p})`$ for $`p_s^2k^2`$. This region is drawn in Fig.4.
Note that the “dangerous” points completely lie in the region where our BS amplitude vanishes identically. Therefore, owing to this specific behaviour of the BS amplitude for bound states with $`W>m`$, intervals of $`p_0`$ where cuts overlap do not contribute to the normalization integral.
It seems to us that this property is very important and can show up in more interesting cases too. Specifically this fact may be employed for confining kernels, because sometimes bound states above threshold ($`W>m`$) do arise for such kernels. Example considered above indicates that in similar cases the $`(p_0,p_s)`$-plane may bear some structure — namely the BS amplitude may have support region in the $`(p_0,p_s)`$-plane. So effectively we can take that cuts in the spectral representation start from the origin ($`p_0=0`$) and, consequently, perform Wick rotation without any complications. This is the case when Wick rotation is admissible for dynamical reasons. As a result we obtain:
$$i^2d^3p𝑑p_4\overline{\psi }(ip_4,\stackrel{}{p})G_0^1(ip_4,\stackrel{}{p})\psi (ip_4,\stackrel{}{p})=\epsilon \lambda \frac{dS_B(\lambda )}{d\lambda }$$
(19)
Now, according to (18) and symmetry of eq.(3) under the relative time reflection $`p_0p_0`$ we have:
$$\overline{\psi }(ip_4,\stackrel{}{p})=\psi ^{}(ip_4,\stackrel{}{p})=\mathrm{\Pi }_B\psi ^{}(ip_4,\stackrel{}{p}),$$
(20)
where $`\mathrm{\Pi }_B`$ denotes the so-called relative-time parity (“$`p_0`$-parity”) of the bound state.
Using (20) in the normalization integral we get
$$d^3p𝑑p_4\overline{\psi }(ip_4,\stackrel{}{p})G_0^1(ip_4,\stackrel{}{p})\psi (ip_4,\stackrel{}{p})=\epsilon \lambda \frac{dS_B(\lambda )}{d\lambda }$$
As long as $`G_0^1(ip_4,\stackrel{}{p})`$ is positive definite in the unshaded regions the full integral above is positive definite too. Besides, from the explicit form of our solution we know that
$$\lambda \frac{dS_B(\lambda )}{d\lambda }=|\lambda |\frac{dS_B(\lambda )}{d|\lambda |}>0$$
and therefore with appropriate choice of overall multiplicative constant up to which the BS amplitude is determined the norm factor $`\epsilon _B`$ coinsides with “$`p_0`$-parity” of the bound state
$$\epsilon _B=\mathrm{\Pi }_B.$$
In other words the states with positive (negative) “$`p_0`$-parity” have the positive (negative) norm.
States with the negative norm can be be eliminated in from the physical sector as the correspondinng “one-time” quasipotential wave functions vanish identically:
$$_{\mathrm{}}^{\mathrm{}}𝑑p_0\psi (p_0,\stackrel{}{p})=0.$$
On the other hand the positive norm states, i.e. the states with positive “$`p_0`$-parity” do survive and produce the pole contributions to the total Green function (and S-matrix) with the correct sign.
## 5 Conclusions
We have demonstrated that in the BS equation for the equal mass scalar particles’ repulsive ($`\lambda <0`$) interaction via scalar massless particle exchange there emerges an effective attraction above thershold. This attraction can lead to existence of bound states for values of the coupling $`\lambda `$ larger then the critical one:
$$\frac{|\lambda _c|}{m^2}>2.$$
In QED with the one-photon exchange this corresponds to
$$\frac{\alpha _c}{\pi }>2.$$
It is worth mentioning that attraction is present for all values of $`\lambda `$ as a purely relativistic effect in the kinematic region where the retardation effects can not be treated by perturbative methods.
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# Spin-Polarized Transport Across an La0.7Sr0.3MnO3/YBa2Cu3O7-x Interface: Role of Andreev Bound States
## Abstract
Transport across an La<sub>0.7</sub>Sr<sub>0.3</sub>MnO<sub>3</sub>/YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-x</sub> (LSMO/YBCO) interface is studied as a function of temperature and surface morphology. For comparison, control measurements are performed in non-magnetic heterostructures of LaNiO<sub>3</sub>/YBCO and Ag/YBCO. In all cases, YBCO is used as bottom layer to eliminate the channel resistance and to minimize thermal effects. The observed differential conductance reflects the role of Andreev bound states in $`ab`$ plane, and brings out for the first time the suppression of such states by the spin-polarized transport across the interface. The theoretical analysis of the measured data reveals decay of the spin polarization near the LSMO surface with temperature, consistent with the reported photoemission data.
The pioneering work on spin-polarized transport in conventional superconductors , performed in the tunneling limit of strong interfacial scattering, has been recently extended to high transparency ferromagnet/superconductor heterostructures , where the two-particle process of Andreev reflection plays an important role. There is generally a good understanding of transport in normal metal/high temperature superconductor (HTSC) junctions and the important role of the Andreev bound states (midgap states) , which lead to the formation of the zero bias conductance peak (ZBCP). In contrast, there remain many open questions in the spin-polarized case involving heterostructures of colossal magnetoresistance materials (CMR) and HTSC . Systematic studies of these systems with high spin polarization are particularly important as they hold potential for probing ferromagnetic interfaces, unconventional superconductivity and for modifying superconducting properties, such as the critical current and T<sub>c</sub> . Understanding and control of interface properties in CMR/HTSC heterostructures could also lead to novel spin-based devices. The key factors in such studies are the electronic and magnetic quality of the interface region on the sides of the ferromagnet and the superconductor, and the thermal management of the device configuration, since these can critically influence the outcome of the measurement.
In this work we report and analyze some observations of spin-polarized transport at a high quality La<sub>0.7</sub>Sr<sub>0.3</sub>MnO<sub>3</sub>/YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-x</sub> (LSMO/YBCO) interface. In contrast to previous studies , YBCO is used as the bottom electrode. This choice of geometry (Fig. 1a) eliminates the channel resistance and minimizes heating effects, thereby unfolding some peculiar features in the differential conductance-voltage (G-V) characteristics, not reported so far. By analyzing these features we are able to show that the Andreev bound states, observed in the $`ab`$ plane HTSC tunneling experiments , have a major influence on the transport properties across the CMR/HTSC interface. These results also reveal, for the first time, the suppression of such bound states by the spin-polarized transport across the interface, as predicted theoretically . Our analysis of the interfacial spin polarization is consistent with the findings of photoemission experiments, which show that the surface spin polarization of LSMO decreases more rapidly with temperature than that of the bulk .
All films involved in the measurement were made by the pulsed laser deposition (PLD) technique. Before we address the issue of high quality film growth it is helpful to discuss our testing structure, shown in Fig. 1a. Several important considerations were involved in fabricating this seemingly simple structure to achieve high quality and reproducibility of the results. In particular, the use of YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-x</sub> as the bottom layer offers three significant advantages: a) eliminating the channel resistance contribution at temperatures below the T<sub>c</sub> of YBCO. Thus, the CMR/HTSC interface contribution dominates the resistance; b) providing an equipotential surface leading to a uniform and perpendicular current distribution; and c) avoiding heating effects from the CMR layer, which is significant in other structures. Au pads of 150 $`\times `$ 200 $`\mu `$m separated by 500 $`\mu `$m were patterned using optical lithography, and ion milled down to the YBCO layer in Argon ambient . In order to ensure absence of possible shorting, the YBCO layer was always 15$`\%`$ over-etched. 25 $`\mu `$m-thick gold wires were directly bonded to the top gold contact layer for measurements of current-voltage characteristic in a four-probe configuration. The results of such measurements were differentiated digitally to extract G-V curves. In the PLD procedure: first, a 2000 Å film of c-axis YBCO was grown on a (100) SrTiO<sub>3</sub> substrate using an energy density of 1.7 J/cm<sup>2</sup> in O<sub>2</sub> pressure of 150 mtorr at 800 C giving a T<sub>c</sub> of 90 K. A 700 Å film of LSMO was then deposited on the YBCO layer using an energy density of 2 J/cm<sup>2</sup> in 400 mtorr O<sub>2</sub> at 780 C and the structure was allowed to cool naturally in 400 torr O<sub>2</sub>. Finally, 200 Å gold was deposited in situ in a vacuum of 1$`\times 10^6`$ torr at 90 C to protect the surface, followed by another ex situ deposition of a 3000 Å Au layer by thermal evaporation. Rocking curves routinely have full width of half maximum of less than 0.3 for both YBCO and LSMO peaks, indicating the high quality of the films produced by this procedure. After all the processing steps the T<sub>c</sub> of the YBCO films remains 90 K. We performed a similar experiment on much smoother YBCO films, specifically grown under slightly off-optimum conditions (substrate temperature: 780 C, energy density: 1.4 J/cm<sup>2</sup>). Although the T<sub>c</sub> in this case was only 7 K lower, the surface morphology was considerably flatter, as shown in the insets of Fig. 1b and 1c.
The resistance of the LSMO/YBCO junction is composed of three parts: the Au/LSMO interface, the LSMO electrode and the LSMO/YBCO interface. According to the literature , even an ex situ grown Au/LSMO interface has a surface resistivity of about 1$`\times `$10<sup>-6</sup> $`\mathrm{\Omega }`$ cm<sup>2</sup> at room temperature, hence the resistance coming from the Au/LSMO interface in our case is smaller than 3.3 m$`\mathrm{\Omega }`$. Taking a very conservative estimate for the resistivity of LSMO of 10 m$`\mathrm{\Omega }`$ cm, we conclude that the LSMO electrode has a resistance of about 2.3 m$`\mathrm{\Omega }`$. Given that the total resistance in our structure is about 300 m$`\mathrm{\Omega }`$, these parasitic resistances can be neglected and the G-V curves below T<sub>c</sub> can be regarded as genuinely representing the properties of the LSMO/YBCO interface.
The conductance data, shown for the “faceted” (Fig. 1b) and smooth (Fig. 1c) interface, display a general “V”-shaped background, similar to the tunneling data in metallic oxide systems . The G-V curves were seen to converge at higher bias in all the data up to 85 K, indicating that there is little contribution from the YBCO channel over this temperature range. Similarly, we can argue that the thermal effects are minimized. The change of curvature in Fig. 1b at higher voltages around 60 K can be explained by the presence of vortices when the channel current approaches the critical current. The most significant difference between the two sets of data is the structure near the zero bias in Fig. 1b. At higher temperatures it evolves into ZBCP, but is completely absent for the smoother surface in Fig. 1c. We attribute the possibility to observe ZBCP in our $`c`$-axis oriented film to the known interfacial roughness of optimally grown YBCO which facilitates $`ab`$ plane transport. Consequently, such surface morphology, due to the sign change of $`d`$-wave order parameter would introduce non-vanishing weight of the $`ab`$ plane Andreev bound states . To investigate this point we perform control measurements with non-magnetic systems in the geometry shown in Fig. 1a where LSMO is replaced by LNO (Fig. 2a) and by Ag (Fig. 2b). In spite of the significantly different electronic properties of LNO and Ag, conductance data in both cases display similar behavior: a ZBCP, present already at lowest examined temperature diminishes at higher temperature, consistent with the effects of thermal smearing. These measurements (in particular the one on LNO which, as a metallic oxide, is similar to LSMO) allow us to investigate the distinguishing conductance features arising from the spin-polarized transport.
In order to reveal the properties of interface transport from Fig. 1b near zero bias with greater clarity, we have to distinguish between the contributions of the two electrodes of the junction, to the G-V curves. Since we are mainly interested in studying the effect of spin polarized tunneling into the HTSC, we would like to remove the effect of the complicated density of states (DOS) of the CMR electrode from the raw conductance data. The DOS of CMR have been studied by tunneling spectroscopy leading to the corresponding conductance contribution at low bias given by $`G(V)=G_0(1+(|eV|/\mathrm{\Delta })^n),`$ where $`G_0`$ is the zero bias value, $`\mathrm{\Delta }`$ is the correlation gap and $`0.5<n<1`$. To determine background conductance (assuming that it is predominantly due to the CMR DOS), which should be removed from the measured G-V curves, we first fit G(V) to the data at lowest T and next, for each higher T, apply thermal smearing to the fitted curve. The resulting curves, after removing background conductance at three different temperatures, are shown in Fig. 3. Each curve is normalized with respect to its value at $`\mathrm{\Delta }_0`$, corresponding to the maximum gap . We compare these results with the theoretical analysis for transport across CMR/HTSC junction by adopting the notation and methods from Ref. , generalized to finite temperature. The strength of interfacial scattering is modeled by parameter Z<sub>0</sub> and the spin polarization is represented by X, the ratio of the exchange and Fermi energies for CMR. The limit X=0 depicts the unpolarized case, while X=1 corresponds to the complete polarization of a half-metallic ferromagnet. To include the effects of different electronic densities in the two materials, we use L<sub>0</sub>, the ratio of Fermi wavevectors in HTSC and CMR. Additionally, to capture the main aspects of surface roughness, we average results over different interface orientations. This is in contrast to the usually studied extreme cases of (110) and (100) planes, corresponding to the maximum and minimum spectral weight of the Andreev bound states, respectively. The calculated conductance is normalized with respect to its value at the maximum gap.
For each G-V curve from Fig. 3a, we plot in Fig. 3b two curves for corresponding temperature with slightly different values of $`X`$, to illustrate the sensitivity of results to the spin polarization of a ferromagnet. We show that the the essential features from Fig. 3a are well reproduced. The overall amplitude is expected to be much smaller from measured values since in the analysis we do not include $`c`$-axis contribution which could be modeled by a parallel conductance channel. Findings from Fig. 3b (which we have also verified in a wide parameter range) show that at a fixed temperature the increase in the exchange energy reduces the amplitude of ZBCP. This can partly explain smaller magnitude of ZBCP observed in LSMO/YBCO as compared to the non-magnetic junctions. Another influence, contributing to this difference in magnitude, could arise from the previously discussed DOS effect of the CMR electrode. Using the theoretical framework from Ref. it is possible to understand various effects on conductance data from the temperature dependent exchange interaction responsible for the spin subband splitting and the spin polarization in the ferromagnetic region. In the limit of low T and almost complete spin polarization it is predicted that there would be a strong conductance suppression at low bias voltage where the transport properties are governed by Andreev reflection. For spin-polarized carriers only a fraction of the incident electrons from a majority spin subband will have partner from a minority spin subband in order to be Andreev reflected, leading to a reduced charge transport across the interface . If the spin-polarization decreases with temperature there would be a smaller difference between population in the two spin subbands resulting in an enhanced Andreev reflection and more pronounced ZBCP. This temperature dependence is qualitatively different from the effects of thermal smearing which reduce and broaden ZBCP. Our raw data, from the blow up in Fig. 1b, provides therefore a strong support for the decreasing spin polarization with temperature which can even be sufficiently fast to off-set the opposite effects of thermal smearing. Observing ZBCP at temperatures close to T<sub>c</sub> further suggests that fabricating our junctions have provided high interfacial quality without degrading the superconducting properties. From the calculated conductance we could infer high but not complete, spin polarization at lowest measured temperature, which is consistent with several other findings using different techniques. A detailed study on the LSMO/YBCO interface by x-ray magnetic circular dichroism has concluded that the polarization of LSMO is further suppressed, compared to the free surface, by the presence of a capping YBCO layer. In the case of LSMO, this should serve as a caution for attempts to get the precise quantitative agreement in the degree of spin polarization obtained by various measurement techniques which often also involve different material fabrication. The temperature dependence of the spin polarization from our data is consistent with the results of spin-resolved photoemission spectroscopy data, which show that the spin polarization within 50 Å of the free surface of LSMO drops faster than the bulk with increasing temperature.
In conclusion, by controlling the surface morphology in a $`c`$-axis grown YBCO we have demonstrated qualitatively different temperature evolution of the measured conductance features in the ferromagnetic and non-magnetic superconducting heterostructures. These differences are attributed to the suppression of the spectral weight of the Andreev bound states by the spin-polarized transport across the interface.
This work is supported by the US ONR Grants $`\mathrm{\#}`$N000140010028, N000149810218, the NSF MRSEC Grant $`\mathrm{\#}`$DMR-96-32521 and by Darpa. We thank I. Takeuchi, P. Fournier, I. I. Mazin and B. Nadgorny for valuable discussions.
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# Parton-Based Gribov-Regge Theory
## Chapter 1 Introduction
The purpose of this paper is to provide the theoretical framework to treat hadron-hadron scattering and the initial stage of nucleus-nucleus collisions at ultra-relativistic energies, in particular with view to RHIC and LHC. The knowledge of these initial collisions is crucial for any theoretical treatment of a possible parton-hadron phase transition, the detection of which being the ultimate aim of all the efforts of colliding heavy ions at very high energies.
It is quite clear that coherence is crucial for the very early stage of nuclear collisions, so a real quantum mechanical treatment is necessary and any attempt to use a transport theoretical parton approach with incoherent quasi-classical partons should not be considered at this point. Also semi-classical hadronic cascades cannot be stretched to account for the very first interactions, even when this is considered to amount to a string excitation, since it is well known that such longitudinal excitation is simply kinematically impossible (see appendix A.2). There is also the very unpleasant feature of having to treat formation times of leading and non-leading particles very different. Otherwise, due to a large gamma factor, it would be impossible for a leading particle to undergo multiple collisions.
So what are the currently used fully quantum mechanical approaches? There are presently considerable efforts to describe nuclear collisions via solving classical Yang-Mills equations, which allows to calculate inclusive parton distributions . This approach is to some extent orthogonal to ours: here, screening is due to perturbative processes, whereas we claim to have good reasons to consider soft processes to be at the origin of screening corrections.
Provided factorization works for nuclear collisions, on may employ the parton model, which allows to calculate inclusive cross sections as a convolution of an elementary cross section with parton distribution functions, with these distribution function taken from deep inelastic scattering. In order to get exclusive parton level cross sections, some additional assumptions are needed, as discussed later.
Another approach is the so-called Gribov-Regge theory . This is an effective field theory, which allows multiple interactions to happen “in parallel”, with the phenomenological object called “Pomeron” representing an elementary interaction. Using the general rules of field theory, on may express cross sections in terms of a couple of parameters characterizing the Pomeron. Interference terms are crucial, they assure the unitarity of the theory. A disadvantage is the fact that cross sections and particle production are not calculated consistently: the fact that energy needs to be shared between many Pomerons in case of multiple scattering is well taken into account when calculating particle production (in particular in Monte Carlo applications), but energy conservation is not taken care of for cross section calculations. This is a serious problem and makes the whole approach inconsistent. Also a problem is the question of how to include in a consistent way hard processes, which are usually treated in the parton model approach. Another unpleasant feature is the fact that different elementary interactions in case of multiple scattering are usually not treated equally, so the first interaction is usually considered to be quite different compared to the subsequent ones.
Here, we present a new approach which we call “Parton-based Gribov-Regge Theory”, where we solve some of the above-mentioned problems: we have a consistent treatment for calculating cross sections and particle production considering energy conservation in both cases; we introduce hard processes in a natural way, and, compared to the parton model, we can deal with exclusive cross sections without arbitrary assumptions. A single set of parameters is sufficient to fit many basic spectra in proton-proton and lepton-nucleon scattering, as well as for electron-positron annihilation (with the exception of one parameter which needs to be changed in order to optimize electron-positron transverse momentum spectra).
The basic guideline of our approach is theoretical consistency. We cannot derive everything from first principles, but we use rigorously the language of field theory to make sure not to violate basic laws of physics, which is easily done in more phenomenological treatments (see discussion above).
There are still problems and open questions: there is clearly a problem with unitarity at very high energies, which should be cured by considering screening corrections due to so-called triple-Pomeron interactions, which we do not treat rigorously at present but which is our next project.
####
### 1.1 Present Status
Before presenting new theoretical ideas, we want to discuss a little bit more in detail the present status and, in particular, the open problems in the parton model approach and in Gribov-Regge theory.
#### Gribov-Regge Theory
Gribov-Regge theory is by construction a multiple scattering theory. The elementary interactions are realized by complex objects called “Pomerons”, who’s precise nature is not known, and which are therefore simply parameterized: the elastic amplitude $`T`$ corresponding to a single Pomeron exchange is given as
$$T(s,t)is^{\alpha _0+\alpha ^{}t}$$
(1.1)
with a couple of parameters to be determined by experiment. Even in hadron-hadron scattering, several of these Pomerons are exchanged in parallel, see fig. 1.1.
Using general rules of field theory (cutting rules), one obtains an expression for the inelastic cross section,
$$\sigma _{\mathrm{inel}}^{h_1h_2}=d^2b\left\{1\mathrm{exp}\left(G(s,b)\right)\right\},$$
(1.2)
where the so-called eikonal $`G(s,b)`$ (proportional to the Fourier transform of $`T(s,t)`$) represents one elementary interaction (a thick line in fig. 1.1). One can generalize to nucleus-nucleus collisions, where corresponding formulas for cross sections may be derived.
In order to calculate exclusive particle production, one needs to know how to share the energy between the individual elementary interactions in case of multiple scattering. We do not want to discuss the different recipes used to do the energy sharing (in particular in Monte Carlo applications). The point is, whatever procedure is used, this is not taken into account in calculation of cross sections discussed above. So, actually, one is using two different models for cross section calculations and for treating particle production. Taking energy conservation into account in exactly the same way will modify the cross section results considerably, as we are going to demonstrate later.
This problem has first been discussed in , . The authors claim that following from the non-planar structure of the corresponding diagrams, conserving energy and momentum in a consistent way is crucial, and therefore the incident energy has to be shared between the different elementary interactions, both real and virtual ones.
Another very unpleasant and unsatisfactory feature of most “recipes” for particle production is the fact, that the second Pomeron and the subsequent ones are treated differently than the first one, although in the above-mentioned formula for the cross section all Pomerons are considered to be identical.
#### The Parton Model
The standard parton model approach to hadron-hadron or also nucleus-nucleus scattering amounts to presenting the partons of projectile and target by momentum distribution functions, $`f_{h_1}`$ and $`f_{h_2}`$, and calculating inclusive cross sections for the production of parton jets with the squared transverse momentum $`p_{}^2`$ larger than some cutoff $`Q_0^2`$ as
$$\sigma _{\mathrm{incl}}^{h_1h_2}=\underset{ij}{}𝑑p_{}^2𝑑x^+𝑑x^{}f_{h_1}^i(x^+,p_{}^2)f_{h_2}^j(x^{},p_{}^2)\frac{d\widehat{\sigma }_{ij}}{dp_{}^2}(x^+x^{}s)\theta \left(p_{}^2Q_0^2\right),$$
(1.3)
where $`d\widehat{\sigma }_{ij}/dp_{}^2`$ is the elementary parton-parton cross section and $`i,j`$ represent parton flavors.
This simple factorization formula is the result of cancelations of complicated diagrams (AGK cancelations) and hides therefore the complicated multiple scattering structure of the reaction. The most obvious manifestation of such a structure is the fact that at high energies ($`\sqrt{s}10`$ GeV) the inclusive cross section in proton-(anti)proton scattering exceeds the total one, so the average number $`\overline{N}_{\mathrm{int}}^{pp}`$ of elementary interactions must be bigger than one:
$$\overline{N}_{\mathrm{int}}^{h_1h_2}=\sigma _{\mathrm{incl}}^{h_1h_2}/\sigma _{\mathrm{tot}}^{h_1h_2}>1$$
(1.4)
The usual solution is the so-called eikonalization, which amounts to re-introducing multiple scattering, based on the above formula for the inclusive cross section:
$$\sigma _{\mathrm{inel}}^{h_1h_2}(s)=d^2b\left\{1\mathrm{exp}\left(A(b)\sigma _{\mathrm{incl}}^{h_1h_2}(s)\right)\right\}=\sigma _m^{h_1h_2}(s),$$
(1.5)
with
$$\sigma _m^{h_1h_2}(s)=d^2b\frac{\left(A(b)\sigma _{\mathrm{incl}}^{h_1h_2}(s)\right)^m}{m!}\mathrm{exp}\left(A(b)\sigma _{\mathrm{incl}}^{h_1h_2}(s)\right)$$
(1.6)
representing the cross section for $`m`$ scatterings; $`A(b)`$ being the proton-proton overlap function (the convolution of the two proton profiles). In this way the multiple scattering is “recovered”. The disadvantage is that this method does not provide any clue how to proceed for nucleus-nucleus ($`AB`$) collisions. One usually assumes the proton-proton cross section for each individual nucleon-nucleon pair of a $`AB`$ system. We are going to demonstrate that this assumption is incorrect.
Another problem, in fact the same one as discussed earlier for the GRT, arises in case of exclusive calculations (event generation), since the above formulas do not provide any information on how to share the energy between the many elementary interactions. The Pythia-method amounts to generating the first elementary interaction according to the inclusive differential cross section, then taking the remaining energy for the second one and so on. In this way, the event generation will reproduce the theoretical inclusive spectrum for hadron-hadron interaction (by construction), as it should be. The method is, however, very arbitrary, and is certainly not a convincing procedure for the multiple scattering aspects of the collisions.
### 1.2 Parton-based Gribov-Regge Theory
In this paper, we present a new approach for hadronic interactions and for the initial stage of nuclear collisions, which is able to solve several of the above-mentioned problems. We provide a rigorous treatment of the multiple scattering aspect, such that questions as energy conservation are clearly determined by the rules of the field theory, both for cross section and particle production calculations. In both (!) cases, energy is properly shared between the different interactions happening in parallel, see fig. 1.2.
for proton-proton and fig. 1.3 for proton-nucleus collisions (generalization to nucleus-nucleus is obvious).
This is the most important and new aspect of our approach, which we consider to be a first necessary step to construct a consistent model for high energy nuclear scattering.
The elementary interactions, shown as the thick lines in the above figures, are in fact a sum of a soft, a hard, and a semi-hard contribution, providing a consistent treatment of soft and hard scattering. To some extend, our approach provides a link between the Gribov-Regge approach and the parton model, we call it “Parton-based Gribov-Regge Theory”.
There are still many problems to be solved: as we are going to show later, a rigorous treatment of energy conservation will lead to unitarity problems (increasingly severe with increasing energy), which is nothing but a manifestation of the fact that screening corrections will be increasingly important. In this paper we will employ a “unitarization procedure” to solve this problem, but this is certainly not the final answer. The next step should be a consistent approach, taking into account both energy conservation and screening corrections due to multi-Pomeron interactions.
Our approach is realized as the Monte Carlo code neXus, which is nothing but the direct implementation of the formalism described in this paper. All our numerical results are calculated with neXus version 2.00.
## Chapter 2 The Formalism
We want to calculate cross sections and particle production in a consistent way, in both cases based on the same formalism, with energy conservation being ensured. The formalism operates with Feynman diagrams of the QCD-inspired effective field theory, such that calculations follow the general rules of the field theory. A graphical representation of a contribution to the elastic amplitude of nucleus-nucleus scattering (related to particle production via the optical theorem) is shown in fig. 2.1:
here the nucleons are split into several “constituents”, each one carrying a fraction of the incident momentum, such that the sum of the momentum fractions is one (momentum conservation). Per nucleon there are one or several “participants” and exactly one “spectator” or “remnant”, where the former ones interact with constituents from the other side via some “elementary interaction” (vertical lines in the figure 2.1). The remnant is just all the rest, i.e. the nucleon minus the participants.
After a technical remark concerning profile functions, we are going to discuss parton-parton scattering, before we develop the multiple scattering theory for hadron-hadron and nucleus-nucleus scattering.
### 2.1 Profile Functions
Profile functions play a fundamental role in our formalism, so we briefly sketch their definition and physical meaning.
Let $`T`$ be the elastic scattering amplitude $`T`$ for the two-body scattering depicted in fig.2.2.
The 4-momenta $`p`$ and $`p^{}`$ are the ones for the incoming particles , $`\stackrel{~}{p}=p+q`$ and $`\stackrel{~}{p}^{}=p^{}q`$ the ones for the outgoing particles, and $`q`$ the 4-momentum transfer in the process. We define as usual the Mandelstam variables $`s`$ and $`t`$ (see appendix A.1). Using the optical theorem, we may write the total cross section as
$$\sigma _{\mathrm{tot}}(s)=\frac{1}{2s}2\mathrm{I}\mathrm{m}T(s,t=0).$$
(2.1)
We define the Fourier transform $`\stackrel{~}{T}`$ of $`T`$ as
$$\stackrel{~}{T}(s,b)=\frac{1}{4\pi ^2}d^2q_{}e^{i\stackrel{}{q}_{}\stackrel{}{b}}T(s,t),$$
(2.2)
using $`t=q_{}^2`$ (see appendix A.2), and a so-called “profile function” $`G(s,b)`$ as
$$G(s,b)=\frac{1}{2s}2\mathrm{I}\mathrm{m}\stackrel{~}{T}(s,b).$$
(2.3)
One can easily verify that
$$\sigma _{\mathrm{tot}}(s)=d^2bG(s,b),$$
(2.4)
which allows an interpretation of $`G(s,b)`$ to be the probability of an interaction at impact parameter $`b`$.
In the following, we are working with partonic, hadronic, and even nuclear profile functions. The central result to be derived in the following sections is the fact that hadronic and also nuclear profile functions can be expressed in terms of partonic ones, allowing a clean formulation of a multiple scattering theory.
### 2.2 Parton-Parton Scattering
We distinguish three types of elementary parton-parton scatterings, referred to as “soft”, “hard” and “semi-hard”, which we are going to discuss briefly in the following. The detailed derivations can be found in appendix B.1.
#### The Soft Contribution
Let us first consider the pure non-perturbative contribution to the process of the fig. 2.2, where all virtual partons appearing in the internal structure of the diagram have restricted virtualities $`Q^2<Q_0^2`$, where $`Q_0^21`$ GeV<sup>2</sup> is a reasonable cutoff for perturbative QCD being applicable. Such soft non-perturbative dynamics is known to dominate hadron-hadron interactions at not too high energies. Lacking methods to calculate this contribution from first principles, it is simply parameterized and graphically represented as a ‘blob’, see fig. 2.3. It is traditionally assumed to correspond to multi-peripheral production of partons (and final hadrons) and is described by the phenomenological soft Pomeron exchange contribution :
$$T_{\mathrm{soft}}(\widehat{s},t)=8\pi s_0\eta (t)\gamma _{\mathrm{part}}^2\left(\frac{\widehat{s}}{s_0}\right)^{\alpha _{\mathrm{soft}}(0)}\mathrm{exp}\left(\lambda _{\mathrm{soft}}^{(2)}\left(\widehat{s}/s_0\right)t\right),$$
(2.5)
with
$$\lambda _{\mathrm{soft}}^{(n)}(z)=nR_{\mathrm{part}}^2+\alpha _{\mathrm{soft}}^{}\mathrm{ln}z,$$
(2.6)
where $`\widehat{s}=(p+p^{})^2`$. The parameters $`\alpha _{\mathrm{soft}}(0)`$, $`\alpha _{\mathrm{soft}}^{}`$ are the intercept and the slope of the Pomeron trajectory, $`\gamma _{\mathrm{part}}`$ and $`R_{\mathrm{part}}^2`$ are the vertex value and the slope for the Pomeron-parton coupling, and $`s_01`$ GeV<sup>2</sup> is the characteristic hadronic mass scale. The so-called signature factor $`\eta `$ is given as
$$\eta (t)=i\mathrm{cot}\frac{\pi \alpha __\mathrm{P}(t)}{2}i.$$
(2.7)
Cutting the diagram of the fig. 2.3 corresponds to the summation over multi-peripheral intermediate hadronic states, connected via unitarity to the imaginary part of the amplitude (2.5),
$`{\displaystyle \frac{1}{i}}\mathrm{disc}_{\widehat{s}}T_{\mathrm{soft}}(\widehat{s},t)`$ $`=`$ $`{\displaystyle \frac{1}{i}}\left[T_{\mathrm{soft}}(\widehat{s}+i0,t)T_{\mathrm{soft}}(\widehat{s}i0,t)\right]`$ (2.8)
$`=`$ $`2\mathrm{I}\mathrm{m}T_{\mathrm{soft}}(\widehat{s},t)`$ (2.9)
$`=`$ $`{\displaystyle \underset{n,\mathrm{spins},\mathrm{}}{}}{\displaystyle 𝑑\tau _nT_{p,p^{}X_n}T_{\stackrel{~}{p},\stackrel{~}{p}^{}X_n}^{}},`$ (2.10)
where $`T_{p,p^{}X_n}`$ is the amplitude for the transition of the initial partons $`p,p^{}`$ into the $`n`$-particle state $`X_n`$, $`d\tau _n`$ is the invariant phase space volume for the $`n`$-particle state $`X_n`$ and the summation is done over the number of particles $`n`$ and over their spins and species, the averaging over initial parton colors and spins is assumed; $`\mathrm{disc}_{\widehat{s}}T_{\mathrm{soft}}(\widehat{s},t)`$ denotes the discontinuity of the amplitude $`T_{\mathrm{soft}}(\widehat{s},t)`$ on the right-hand cut in the variable $`\widehat{s}`$.
For $`t=0`$ one obtains via the optical theorem the contribution $`\sigma _{\mathrm{soft}}`$ of the soft Pomeron exchange to the total parton interaction cross section,
$$\sigma _{\mathrm{soft}}\left(\widehat{s}\right)=\frac{1}{2\widehat{s}}2\mathrm{I}\mathrm{m}T_{\mathrm{soft}}(\widehat{s},0)=8\pi \gamma _{\mathrm{part}}^2\left(\frac{\widehat{s}}{s_0}\right)^{\alpha _{\mathrm{soft}}(0)1},$$
(2.11)
where $`2\widehat{s}`$ defines the initial parton flux.
The corresponding profile function for parton-parton interaction is expressed via the Fourier transform $`\stackrel{~}{T}`$ of $`T`$ divided by the flux $`2s`$,
$$D_{\mathrm{soft}}(\widehat{s},b)=\frac{1}{2\widehat{s}}2\mathrm{I}\mathrm{m}\stackrel{~}{T}_{\mathrm{soft}}(\widehat{s},b),$$
(2.12)
which gives
$`D_{\mathrm{soft}}(\widehat{s},b)`$ $`=`$ $`{\displaystyle \frac{1}{8\pi ^2\widehat{s}}}{\displaystyle d^2q_{}\mathrm{exp}\left(i\stackrel{}{q}_{}\stackrel{}{b}\right)\mathrm{\hspace{0.17em}2}\mathrm{Im}T_{\mathrm{soft}}(\widehat{s},q_{}^2)}`$ (2.13)
$`=`$ $`{\displaystyle \frac{2\gamma _{\mathrm{part}}^2}{\lambda _{\mathrm{soft}}^{(2)}(\widehat{s}/s_0)}}\left({\displaystyle \frac{\widehat{s}}{s_0}}\right)^{\alpha _{\mathrm{soft}}(0)1}\mathrm{exp}\left({\displaystyle \frac{b^2}{4\lambda _{\mathrm{soft}}^{(2)}(\widehat{s}/s_0)}}\right).`$ (2.14)
The external legs of the diagram of fig. 2.3 are “partonic constituents”, which are assumed to be quark-antiquark pairs.
#### The Hard Contribution
Let us now consider the other extreme, when all the processes in the ‘box’ of the fig. 2.2 are perturbative, i.e. all internal intermediate partons are characterized by large virtualities $`Q^2>Q_0^2`$. In that case, the corresponding hard parton-parton scattering amplitude $`T_{\mathrm{hard}}^{jk}(\widehat{s},t)`$ ($`j,k`$ denote the types (flavors) of the initial partons) can be calculated using the perturbative QCD techniques , and the intermediate states contributing to the absorptive part of the amplitude of the fig. 2.2 can be defined in the parton basis. In the leading logarithmic approximation of QCD, summing up terms where each (small) running QCD coupling constant $`\alpha _s(Q^2)`$ appears together with a large logarithm $`\mathrm{ln}(Q^2/\lambda _{\mathrm{QCD}}^2)`$ (with $`\lambda _{QCD}`$ being the infrared QCD scale), and making use of the factorization hypothesis, one obtains the contribution of the corresponding cut diagram for $`t=q^2=0`$ as the cut parton ladder cross section $`\sigma _{\mathrm{hard}}^{jk}(\widehat{s},Q_0^2)`$ <sup>1</sup><sup>1</sup>1 Strictly speaking, one obtains the ladder representation for the process only using axial gauge. , which will correspond to the cut diagram of fig. 2.4,
where all horizontal rungs are the final (on-shell) partons and the virtualities of the virtual $`t`$-channel partons increase from the ends of the ladder towards the largest momentum transfer parton-parton process (indicated symbolically by the ‘blob’ in the middle of the ladder):
$`\sigma _{\mathrm{hard}}^{jk}(\widehat{s},Q_0^2)`$ $`=`$ $`{\displaystyle \frac{1}{2\widehat{s}}}2\mathrm{I}\mathrm{m}T_{\mathrm{hard}}^{jk}\left(\widehat{s},t=0\right)`$
$`=`$ $`K{\displaystyle \underset{ml}{}}{\displaystyle 𝑑x_B^+𝑑x_B^{}𝑑p_{}^2\frac{d\sigma _{\mathrm{Born}}^{ml}}{dp_{}^2}(x_B^+x_B^{}\widehat{s},p_{}^2)}`$
$`\times `$ $`E_{\mathrm{QCD}}^{jm}(x_B^+,Q_0^2,M_F^2)E_{\mathrm{QCD}}^{kl}(x_B^{},Q_0^2,M_F^2)\theta \left(M_F^2Q_0^2\right),`$
Here $`d\sigma _{\mathrm{Born}}^{ml}/dp_{}^2`$ is the differential $`22`$ parton scattering cross section, $`p_{}^2`$ is the parton transverse momentum in the hard process, $`m,l`$ and $`x_B^\pm `$ are correspondingly the types and the shares of the light cone momenta of the partons participating in the hard process, and $`M_F^2`$ is the factorization scale for the process (we use $`M_F^2=p_{}^2/4`$). The ‘evolution function’ $`E_{\mathrm{QCD}}^{jm}(Q_0^2,M_F^2,z)`$ represents the evolution of a parton cascade from the scale $`Q_0^2`$ to $`M_F^2`$, i.e. it gives the number density of partons of type $`m`$ with the momentum share $`z`$ at the virtuality scale $`M_F^2`$, resulted from the evolution of the initial parton $`j`$, taken at the virtuality scale $`Q_0^2`$. The evolution function satisfies the usual DGLAP equation with the initial condition $`E_{\mathrm{QCD}}^{jm}(Q_0^2,Q_0^2,z)=\delta _m^j\delta (1z)`$, as discussed in detail in appendix B.3. The factor $`K1.5`$ takes effectively into account higher order QCD corrections.
In the following we shall need to know the contribution of the uncut parton ladder $`T_{\mathrm{hard}}^{jk}(\widehat{s},t)`$ with some momentum transfer $`q`$ along the ladder (with $`t=q^2`$). The behavior of the corresponding amplitudes was studied in in the leading logarithmic($`1/x`$ ) approximation of QCD. The precise form of the corresponding amplitude is not important for our application; we just use some of the results of , namely that one can neglect the real part of this amplitude and that it is nearly independent on $`t`$, i.e. that the slope of the hard interaction $`R_{\mathrm{hard}}^2`$ is negligible small, i.e. compared to the soft Pomeron slope one has $`R_{\mathrm{hard}}^20`$. So we parameterize $`T_{\mathrm{hard}}^{jk}(\widehat{s},t)`$ in the region of small $`t`$ as
$$T_{\mathrm{hard}}^{jk}(\widehat{s},t)=i\widehat{s}\sigma _{\mathrm{hard}}^{jk}(\widehat{s},Q_0^2)\mathrm{exp}\left(R_{\mathrm{hard}}^2t\right)$$
(2.16)
The corresponding profile function is obtained by calculating the Fourier transform $`\stackrel{~}{T}_{\mathrm{hard}}`$ of $`T_{\mathrm{hard}}`$ and dividing by the initial parton flux $`2\widehat{s}`$,
$$D_{\mathrm{hard}}^{jk}(\widehat{s},b)=\frac{1}{2\widehat{s}}2\mathrm{I}\mathrm{m}\stackrel{~}{T}_{\mathrm{hard}}^{jk}(\widehat{s},b),$$
(2.17)
which gives
$`D_{\mathrm{hard}}^{jk}(\widehat{s},b)={\displaystyle \frac{1}{8\pi ^2\widehat{s}}}{\displaystyle d^2q_{}\mathrm{exp}\left(i\stackrel{}{q}_{}\stackrel{}{b}\right)\mathrm{\hspace{0.17em}2}\mathrm{Im}T_{\mathrm{hard}}^{jk}(\widehat{s},q_{}^2)}`$
$`=\sigma _{\mathrm{hard}}^{jk}(\widehat{s},Q_0^2){\displaystyle \frac{1}{4\pi R_{\mathrm{hard}}^2}}\mathrm{exp}\left({\displaystyle \frac{b^2}{4R_{\mathrm{hard}}^2}}\right),`$ (2.18)
In fact, the above considerations are only correct for valence quarks, as discussed in detail in the next section. Therefore, we also talk about “valence-valence” contribution and we use $`D_{\mathrm{val}\mathrm{val}}`$ instead of $`D_{\mathrm{hard}}`$:
$$D_{\mathrm{val}\mathrm{val}}^{jk}(\widehat{s},b)D_{\mathrm{hard}}^{jk}(\widehat{s},b),$$
(2.19)
so these are two names for one and the same object.
#### The Semi-hard Contribution
The discussion of the preceding section is not valid in case of sea quarks and gluons, since here the momentum share $`x_1`$ of the “first” parton is typically very small, leading to an object with a large mass of the order $`Q_0^2/x_1`$ between the parton and the proton . Microscopically, such ’slow’ partons with $`x_11`$ appear as a result of a long non-perturbative parton cascade, where each individual parton branching is characterized by a small momentum transfer squared $`Q^2<Q_0^2`$ and nearly equal partition of the parent parton light cone momentum . When calculating proton structure functions or high-$`p_t`$ jet production cross sections that non-perturbative contribution is usually included into parameterized initial parton momentum distributions at $`Q^2=Q_0^2`$. However, the description of inelastic hadronic interactions requires to treat it explicitly in order to account for secondary particles produced during such non-perturbative parton pre-evolution, and to describe correctly energy-momentum sharing between multiple elementary scatterings. As the underlying dynamics appears to be identical to the one of soft parton-parton scattering considered above, we treat this soft pre-evolution as the usual soft Pomeron emission, as discussed in detail in appendix B.1.
So for sea quarks and gluons, we consider so-called semi-hard interactions between parton constituents of initial hadrons, represented by a parton ladder with “soft ends”, see fig. 2.5.
As in case of soft scattering, the external legs are quark-antiquark pairs, connected to soft Pomerons. The outer partons of the ladder are on both sides sea quarks or gluons (therefore the index “sea-sea”). The central part is exactly the hard scattering considered in the preceding section. As discussed in length in the appendix B.1, the mathematical expression for the corresponding amplitude is given as
$`iT_{\mathrm{sea}\mathrm{sea}}(\widehat{s},t)`$ $`=`$ $`{\displaystyle \underset{jk}{}}{\displaystyle _0^1}{\displaystyle \frac{dz^+}{z^+}}{\displaystyle \frac{dz^{}}{z^{}}}\mathrm{Im}T_{\mathrm{soft}}^j({\displaystyle \frac{s_0}{z^+}},t)\mathrm{Im}T_{\mathrm{soft}}^k({\displaystyle \frac{s_0}{z^{}}},t)iT_{\mathrm{hard}}^{jk}(z^+z^{}\widehat{s},t),`$ (2.20)
with $`z^\pm `$ being the momentum fraction of the external leg-partons of the parton ladder relative to the momenta of the initial (constituent) partons. The indices $`j`$ and $`k`$ refer to the flavor of these external ladder partons. The amplitudes $`T_{\mathrm{soft}}^j`$ are the soft Pomeron amplitudes discussed earlier, but with modified couplings, since the Pomerons are now connected to the ladder on one side. The arguments $`s_0/z^\pm `$ are the squared masses of the two soft Pomerons, $`z^+z^{}\widehat{s}`$ is the squared mass of the hard piece.
Performing as usual the Fourier transform to the impact parameter representation and dividing by $`2\widehat{s}`$, we obtain the profile function
$$D_{\mathrm{sea}\mathrm{sea}}(\widehat{s},b)=\frac{1}{2\widehat{s}}\mathrm{\hspace{0.17em}2}\mathrm{Im}\stackrel{~}{T}_{\mathrm{sea}\mathrm{sea}}(\widehat{s},b),$$
(2.21)
which may be written as
$`D_{\mathrm{sea}\mathrm{sea}}(\widehat{s},b)`$ $`=`$ $`{\displaystyle \underset{jk}{}}{\displaystyle _0^1}𝑑z^+𝑑z^{}E_{\mathrm{soft}}^j\left(z^+\right)E_{\mathrm{soft}}^k\left(z^{}\right)\sigma _{\mathrm{hard}}^{jk}(z^+z^{}\widehat{s},Q_0^2)`$ (2.22)
$`\times {\displaystyle \frac{1}{4\pi \lambda _{\mathrm{soft}}^{(2)}(1/(z^+z^{}))}}\mathrm{exp}\left({\displaystyle \frac{b^2}{4\lambda _{\mathrm{soft}}^{(2)}\left(1/(z^+z^{})\right)}}\right)`$
with the soft Pomeron slope $`\lambda _{\mathrm{soft}}^{(2)}`$ and the cross section $`\sigma _{\mathrm{hard}}^{jk}`$ being defined earlier. The functions $`E_{\mathrm{soft}}^j\left(z^\pm \right)`$ representing the “soft ends” are defined as
$$\mathrm{E}_{\mathrm{soft}}^j(z^\pm )=\mathrm{Im}T_{\mathrm{soft}}^j\left(\frac{s_0}{z^+},t=0\right),$$
(2.23)
or explicitly
$`E_{\mathrm{soft}}^g\left(z\right)`$ $`=`$ $`8\pi s_0\gamma _{\mathrm{part}}\gamma _gz^{\alpha _{\mathrm{soft}}(0)}(1z)^{\beta _g},`$ (2.24)
$`E_{\mathrm{soft}}^q\left(z\right)`$ $`=`$ $`\gamma _{qg}{\displaystyle _z^1}𝑑\xi P_g^q(\xi )E_{\mathrm{soft}}^g\left({\displaystyle \frac{z}{\xi }}\right),`$ (2.25)
with
$$\gamma _{qg}\gamma _g=w_{\mathrm{split}}\stackrel{~}{\gamma }_g,\gamma _g=\left(1w_{\mathrm{split}}\right)\stackrel{~}{\gamma }_g,$$
(2.26)
and
$$\stackrel{~}{\gamma }_g=\frac{1}{8\pi s_0\gamma _{\mathrm{part}}}\frac{\mathrm{\Gamma }\left(3\alpha _{\mathrm{soft}}(0)+\beta _g\right)}{\mathrm{\Gamma }\left(2\alpha _{\mathrm{soft}}(0)\right)\mathrm{\Gamma }\left(1+\beta _g\right)}$$
(2.27)
(see appendix B.1). We neglected the small hard scattering slope $`R_{\mathrm{hard}}^2`$ compared to the Pomeron slope $`\lambda _{\mathrm{soft}}^{(2)}`$. We call $`E_{\mathrm{soft}}`$ also the “ soft evolution”, to indicate that we consider this as simply a continuation of the QCD evolution, however, in a region where perturbative techniques do not apply any more. As discussed in the appendix B.1, $`E_{\mathrm{soft}}^j\left(z\right)`$ has the meaning of the momentum distribution of parton $`j`$ in the soft Pomeron.
Consistency requires to also consider the mixed semi-hard contributions with a valence quark on one side and a non-valence participant (quark-antiquark pair) on the other one, see fig. 2.6.
We have
$$iT_{\mathrm{val}\mathrm{sea}}^j(\widehat{s},t)=_0^1\frac{dz^{}}{z^{}}\underset{k}{}\mathrm{Im}T_{\mathrm{soft}}^k(\frac{s_0}{z^{}},t)iT_{\mathrm{hard}}^{jk}(z^{}\widehat{s},t)$$
(2.28)
and
$`D_{\mathrm{val}\mathrm{sea}}^j(\widehat{s},b)`$ $`=`$ $`{\displaystyle \underset{k}{}}{\displaystyle _0^1}𝑑z^{}E_{\mathrm{soft}}^k\left(z^{}\right)\sigma _{\mathrm{hard}}^{jk}(z^{}\widehat{s},Q_0^2)`$
$`\times {\displaystyle \frac{1}{4\pi \lambda _{\mathrm{soft}}^{(1)}(1/z^{})}}\mathrm{exp}\left({\displaystyle \frac{b^2}{4\lambda _{\mathrm{soft}}^{(1)}\left(1/z^{}\right)}}\right)`$
where $`j`$ is the flavor of the valence quark at the upper end of the ladder and $`k`$ is the type of the parton on the lower ladder end. Again, we neglected the hard scattering slope $`R_{\mathrm{hard}}^2`$ compared to the soft Pomeron slope. A contribution $`D_{\mathrm{sea}\mathrm{val}}^j(\widehat{s},b)`$, corresponding to a valence quark participant from the target hadron, is given by the same expression,
$$D_{\mathrm{sea}\mathrm{val}}^j(\widehat{s},b)=D_{\mathrm{val}\mathrm{sea}}^j(\widehat{s},b),$$
(2.30)
since eq. (2.2) stays unchanged under replacement $`z^{}z^+`$ and only depends on the total c.m. energy squared $`\widehat{s}`$ for the parton-parton system.
### 2.3 Hadron-Hadron Scattering
Let us now consider hadron-hadron interactions (a more detailed treatment can be found in appendix C). We ignore first contributions involving valence quark scatterings. In the general case, the expression for the hadron-hadron scattering amplitude includes contributions from multiple scattering between different parton constituents of the initial hadrons, as shown in fig. 2.7,
and can be written according to the standard rules as
$`iT_{h_1h_2}(s,t)={\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n!}}{\displaystyle \underset{l=1}{\overset{n}{}}\left[\frac{d^4k_l}{(2\pi )^4}\frac{d^4k_l^{}}{(2\pi )^4}\frac{d^4q_l}{(2\pi )^4}\right]N_{h_1}^{(n)}(p,k_1,\mathrm{},k_n,q_1,\mathrm{},q_n)}`$ (2.31)
$`\times {\displaystyle \underset{l=1}{\overset{n}{}}}\left[iT_{1\mathrm{I}\mathrm{P}}(\widehat{s}_l,q_l^2)\right]N_{h_2}^{(n)}(p^{},k_1^{},\mathrm{},k_n^{},q_1,\mathrm{},q_n)(2\pi )^4\delta ^{(4)}({\displaystyle \underset{k=1}{\overset{n}{}}}q_iq),`$
with $`t=q^2`$, $`s=(p+p^{})^2p^+p^{}`$, with $`p,p^{}`$ being the 4-momenta of the initial hadrons, and with $`\widehat{s}_l=(k_l+k_l^{})^2k_l^+k_l^{}`$. $`T_{1\mathrm{I}\mathrm{P}}`$ is the sum of partonic one-Pomeron-exchange scattering amplitudes, discussed in the preceding section, $`T_{1\mathrm{I}\mathrm{P}}=T_{\mathrm{soft}}+T_{\mathrm{sea}\mathrm{sea}}`$. The momenta $`k_l,k_l^{}`$ and $`q_l`$ denote correspondingly the 4-momenta of the initial partonic constituents (quark-antiquark pairs) for the $`l`$-th scattering and the 4-momentum transfer in that partial process. The factor $`1/n!`$ takes into account the identical nature of the $`n`$ re-scattering contributions. $`N_h^{(n)}(p,k_1,\mathrm{},k_n,q_1,\mathrm{},q_n)`$ denotes the contribution of the vertex for $`n`$-parton coupling to the hadron $`h`$.
As discussed in appendix C, the hadron-hadron amplitude (2.31) may be written as
$`iT_{h_1h_2}(s,t)=8\pi ^2s{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n!}}{\displaystyle _0^1}{\displaystyle \underset{l=1}{\overset{n}{}}}dx_l^+dx_l^{}{\displaystyle \underset{l=1}{\overset{n}{}}}\left[{\displaystyle \frac{1}{8\pi ^2\widehat{s}_l}}{\displaystyle d^2q_l_{}iT_{1\mathrm{I}\mathrm{P}}^{h_1h_2}(x_l^+,x_l^{},s,q_l_{}^2)}\right]`$
$`\times F_{\mathrm{remn}}^{h_1}\left(1{\displaystyle \underset{j=1}{\overset{n}{}}}x_j^+\right)F_{\mathrm{remn}}^{h_2}\left(1{\displaystyle \underset{j=1}{\overset{n}{}}}x_j^{}\right)\delta ^{(2)}\left({\displaystyle \underset{k=1}{\overset{n}{}}}\stackrel{}{q}_k_{}\stackrel{}{q}_{}\right).`$ (2.32)
(see eq. (C.22)), where the partonic amplitudes are defined as $`T_{1\mathrm{I}\mathrm{P}}^{h_1h_2}=T_{\mathrm{soft}}^{h_1h_2}+T_{\mathrm{sea}\mathrm{sea}}^{h_1h_2}`$, with
$`T_{\mathrm{soft}/\mathrm{sea}\mathrm{sea}}^{h_1h_2}(x^+,x^{},s,q_{}^2)=T_{\mathrm{soft}/\mathrm{sea}\mathrm{sea}}(x^+x^{}s,q_{}^2)F_{\mathrm{part}}^{h_1}(x^+)F_{\mathrm{part}}^{h_2}(x^{})`$
$`\times \mathrm{exp}\left(\left[R_{h_1}^2+R_{h_2}^2\right]q_{}^2\right)`$ (2.33)
representing the contributions of “elementary interactions plus external legs”; the functions $`F_{\mathrm{remn}}^h,F_{\mathrm{part}}^h`$ are defined in (C.18-C.21) as
$$F_{\mathrm{part}}^h(x)=\gamma _hx^{\alpha _{\mathrm{part}}},$$
(2.34)
$$F_{\mathrm{remn}}^h(x)=x^{\alpha _{\mathrm{remn}}^h}.$$
(2.35)
Formula (2.32) is also correct if one includes valence quarks (see appendix C) , if one defines
$$T_{1\mathrm{I}\mathrm{P}}^{h_1h_2}=T_{\mathrm{soft}}^{h_1h_2}+T_{\mathrm{sea}\mathrm{sea}}^{h_1h_2}+T_{\mathrm{val}\mathrm{val}}^{h_1h_2}+T_{\mathrm{val}\mathrm{sea}}^{h_1h_2}+T_{\mathrm{sea}\mathrm{val}}^{h_1h_2},$$
(2.36)
with the hard contribution
$`T_{\mathrm{val}\mathrm{val}}^{h_1h_2}(x^+,x^{},s,q_{}^2)`$ $`=`$ $`{\displaystyle _0^{x^+}}𝑑x_v^+{\displaystyle \frac{x^+}{x_v^+}}{\displaystyle _0^x^{}}𝑑x_v^{}{\displaystyle \frac{x^{}}{x_v^{}}}{\displaystyle \underset{j,k}{}}T_{\mathrm{hard}}^{jk}(x_v^+x_v^{}s,q_{}^2)`$
$`\times \overline{F}_{\mathrm{part}}^{h_1,j}(x_v^+,x^+x_v^+)\overline{F}_{\mathrm{part}}^{h_2,k}(x_v^{},x^{}x_v^{})\mathrm{exp}\left(\left[R_{h_1}^2+R_{h_2}^2\right]q_{}^2\right),`$
with the mixed semi-hard “val-sea” contribution
$`T_{\mathrm{val}\mathrm{sea}}^{h_1h_2}(x^+,x^{},s,q_{}^2)`$ $`=`$ $`{\displaystyle _0^{x^+}}𝑑x_v^+{\displaystyle \frac{x^+}{x_v^+}}{\displaystyle \underset{j}{}}T_{\mathrm{val}\mathrm{sea}}^j(x_v^+x^{}s,q_{}^2)`$
$`\times \overline{F}_{\mathrm{part}}^{h_1,j}(x_v^+,x^+x_v^+)F_{\mathrm{part}}^{h_2}(x^{})\mathrm{exp}\left(\left[R_{h_1}^2+R_{h_2}^2\right]q_{}^2\right),`$
and with the contribution “sea-val” obtained from “val-sea” by exchanging variables,
$$T_{\mathrm{sea}\mathrm{val}}^{h_1h_2}(x^+,x^{},s,q_{}^2)=T_{\mathrm{val}\mathrm{sea}}^{h_2h_1}(x^{},x^+,s,q_{}^2).$$
Here, we allow formally any number of valence type interactions (based on the fact that multiple valence type processes give negligible contribution). In the valence contributions, we have convolutions of hard parton scattering amplitudes $`T_{\mathrm{hard}}^{jk}`$ and valence quark distributions $`\overline{F}_{\mathrm{part}}^j`$ over the valence quark momentum share $`x_v^\pm `$ rather than a simple product, since only the valence quarks are involved in the interactions, with the anti-quarks staying idle (the external legs carrying momenta $`x^+`$ and $`x^{}`$ are always quark-antiquark pairs). The functions $`\overline{F}`$ are given as
$$\overline{F}_{\mathrm{part}}^{h,i}(x_v,x_{\overline{q}})=N^1q_{\mathrm{val}}^i(x_v,Q_0^2)(1x_v)^{\alpha _{\mathrm{I}\mathrm{R}}1\alpha _{\mathrm{remn}}}(x_{\overline{q}})^{\alpha _{\mathrm{I}\mathrm{R}}},$$
(2.39)
with the normalization factor
$$N=\frac{\mathrm{\Gamma }\left(1+\alpha _{\mathrm{remn}}\right)\mathrm{\Gamma }\left(1\alpha _{\mathrm{I}\mathrm{R}}\right)}{\mathrm{\Gamma }\left(2+\alpha _{\mathrm{remn}}\alpha _{\mathrm{I}\mathrm{R}}\right)},$$
(2.40)
where $`q_{\mathrm{val}}^i`$ is a a usual valence quark distribution function.
The Fourier transform $`\stackrel{~}{T}_{h_1h_2}`$ of the amplitude (2.32) is given as
$`{\displaystyle \frac{i}{2s}}\stackrel{~}{T}_{h_1h_2}(s,b)`$ $`=`$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n!}}{\displaystyle _0^1}{\displaystyle \underset{l=1}{\overset{n}{}}}dx_l^+dx_l^{}{\displaystyle \underset{l=1}{\overset{n}{}}}{\displaystyle \frac{i}{2\widehat{s}_l}}\stackrel{~}{T}_{1\mathrm{I}\mathrm{P}}^{h_1h_2}(x_l^+,x_l^{},s,b)`$ (2.41)
$`\times `$ $`F_{\mathrm{remn}}^{h_1}\left(1{\displaystyle \underset{j=1}{\overset{n}{}}}x_j^+\right)F_{\mathrm{remn}}^{h_2}\left(1{\displaystyle \underset{j=1}{\overset{n}{}}}x_j^{}\right),`$
with$`\stackrel{~}{T}_{1\mathrm{I}\mathrm{P}}^{h_1h_2}`$ being the Fourier transform of $`T_{1\mathrm{I}\mathrm{P}}^{h_1h_2}`$. The profile function $`\gamma `$ is as usual defined as
$$\gamma _{h_1h_2}(s,b)=\frac{1}{2s}2\mathrm{I}\mathrm{m}\stackrel{~}{\mathrm{T}}_{h_1h_2}(s,b),$$
(2.42)
which may be evaluated using the AGK cutting rules ,
$`\gamma _{h_1h_2}(s,b)`$ $`=`$ $`{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{m!}}{\displaystyle _0^1}{\displaystyle \underset{\mu =1}{\overset{m}{}}}dx_\mu ^+dx_\mu ^{}{\displaystyle \underset{\mu =1}{\overset{m}{}}}{\displaystyle \frac{1}{2\widehat{s}_\mu }}2\mathrm{I}\mathrm{m}\stackrel{~}{T}_{1\mathrm{I}\mathrm{P}}^{h_1h_2}(x_\mu ^+,x_\mu ^{},s,b)`$ (2.43)
$`\times `$ $`{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{l!}}{\displaystyle _0^1}{\displaystyle \underset{\lambda =1}{\overset{l}{}}}d\stackrel{~}{x}_\lambda ^+d\stackrel{~}{x}_\lambda ^{}{\displaystyle \underset{\lambda =1}{\overset{l}{}}}{\displaystyle \frac{1}{2\widehat{s}_\lambda }}(2)\mathrm{Im}\stackrel{~}{T}_{1\mathrm{I}\mathrm{P}}^{h_1h_2}(\stackrel{~}{x}_\lambda ^+,\stackrel{~}{x}_\lambda ^{},s,b)`$
$`\times `$ $`F_{\mathrm{remn}}^{h_1}\left(1{\displaystyle \underset{j=1}{\overset{m}{}}}x_j^+{\displaystyle \underset{k=1}{\overset{l}{}}}\stackrel{~}{x}_k^+\right)F_{\mathrm{remn}}^{h_2}\left(1{\displaystyle \underset{j=1}{\overset{m}{}}}x_j^{}{\displaystyle \underset{k=1}{\overset{l}{}}}\stackrel{~}{x}_k^{}\right),`$
where $`2\mathrm{I}\mathrm{m}\stackrel{~}{\mathrm{T}}^{h_1h_2}`$ represents a cut elementary diagram and $`2\mathrm{I}\mathrm{m}\stackrel{~}{\mathrm{T}}^{h_1h_2}`$ an uncut one (taking into account that the uncut contribution may appear on either side from the cut plane). It is therefore useful to define a partonic profile function $`G`$ via
$`G_{1\mathrm{I}\mathrm{P}}^{h_1h_2}(x_\lambda ^+,x_\lambda ^{},s,b)`$ $`=`$ $`{\displaystyle \frac{1}{2x_\lambda ^+x_\lambda ^{}s}}2\mathrm{I}\mathrm{m}\stackrel{~}{T}_{1\mathrm{I}\mathrm{P}}^{h_1h_2}(x_\lambda ^+,x_\lambda ^{},s,b),`$ (2.44)
which allows to write the integrand of the right-hand-side of eq. (2.43) as a product of $`G`$ and $`(G)`$ terms:
$`\gamma _{h_1h_2}(s,b)`$ $`=`$ $`{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{m!}}{\displaystyle _0^1}{\displaystyle \underset{\mu =1}{\overset{m}{}}}dx_\mu ^+dx_\mu ^{}{\displaystyle \underset{\mu =1}{\overset{m}{}}}G_{1\mathrm{I}\mathrm{P}}^{h_1h_2}(x_\mu ^+,x_\mu ^{},s,b)`$ (2.45)
$`\times `$ $`{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{l!}}{\displaystyle _0^1}{\displaystyle \underset{\lambda =1}{\overset{l}{}}}d\stackrel{~}{x}_\lambda ^+d\stackrel{~}{x}_\lambda ^{}{\displaystyle \underset{\lambda =1}{\overset{l}{}}}G_{1\mathrm{I}\mathrm{P}}^{h_1h_2}(\stackrel{~}{x}_\lambda ^+,\stackrel{~}{x}_\lambda ^{},s,b)`$
$`\times `$ $`F_{\mathrm{remn}}\left(x^{\mathrm{proj}}{\displaystyle \underset{\lambda }{}}\stackrel{~}{x}_\lambda ^+\right)F_{\mathrm{remn}}\left(x^{\mathrm{targ}}{\displaystyle \underset{\lambda }{}}\stackrel{~}{x}_\lambda ^{}\right),`$
see fig. 2.8, with
$$x^{\mathrm{proj}/\mathrm{targ}}=1x_\mu ^\pm $$
being the momentum fraction of the projectile/target remnant.
This is a very important result, allowing to express the total profile function $`\gamma _{h_1h_2}`$ via the elementary profile functions $`G_{1\mathrm{I}\mathrm{P}}^{h_1h_2}`$.
Based on the above definitions, we may write the profile function $`G_{1\mathrm{I}\mathrm{P}}^{h_1h_2}`$ as
$$G_{1\mathrm{I}\mathrm{P}}^{h_1h_2}=G_{\mathrm{soft}}^{h_1h_2}+G_{\mathrm{sea}\mathrm{sea}}^{h_1h_2}+G_{\mathrm{val}\mathrm{val}}^{h_1h_2}+G_{\mathrm{val}\mathrm{sea}}^{h_1h_2}+G_{\mathrm{sea}\mathrm{val}}^{h_1h_2},$$
(2.46)
with
$$G_{\mathrm{soft}/\mathrm{sea}\mathrm{sea}}^{h_1h_2}(x^+,x^{},s,b)=D_{\mathrm{soft}/\mathrm{sea}\mathrm{sea}}^{h_1h_2}(x^+x^{}s,b)F_{\mathrm{part}}^{h_1}(x^+)F_{\mathrm{part}}^{h_2}(x^{})$$
(2.47)
for the soft and semi-hard “sea-sea” contribution, with
$`G_{\mathrm{val}\mathrm{val}}^{h_1h_2}(x^+,x^{},s,b)`$ $`=`$ $`{\displaystyle 𝑑x_{q_v}^+𝑑x_{q_v}^{}\underset{i,j}{}D_{\mathrm{val}\mathrm{val}}^{h_1h_2,ij}(x_{q_v}^+x_{q_v}^{}s,b)}`$ (2.48)
$`\times \overline{F}_{\mathrm{part}}^{h_1,i}(x_{q_v}^+,x^+x_{q_v}^+)\overline{F}_{\mathrm{part}}^{h_2,j}(x_{q_v}^{},x^{}x_{q_v}^{})`$
for the hard “val-val” contribution, and with
$$G_{\mathrm{val}\mathrm{sea}}^{h_1h_2}(x^+,x^{},s,b)=𝑑x_{q_v}^+\underset{i}{}D_{\mathrm{val}\mathrm{sea}}^{h_1h_2,i}(x_{q_v}^+x^{}s,b)\overline{F}_{\mathrm{part}}^{h_1,i}(x_{q_v}^+,x^+x_{q_v}^+)F_{\mathrm{part}}^{h_2}(x^{})$$
(2.49)
and
$$G_{\mathrm{sea}\mathrm{val}}^{h_1h_2}(x^+,x^{},s,b)=𝑑x_{q_v}^{}\underset{i}{}D_{\mathrm{val}\mathrm{sea}}^{h_2h_1,i}(x^+x_{q_v}^{}s,b)\overline{F}_{\mathrm{part}}^{h_2,i}(x_{q_v}^{},x^{}x_{q_v}^{})F_{\mathrm{part}}^{h_1}(x^+),$$
(2.50)
for the mixed semi-hard “val-sea” and “sea-val” contributions. For the soft and “sea-sea” contributions, the $`D`$-functions are given as
$`D_{\mathrm{soft}/\mathrm{sea}\mathrm{sea}}^{h_1h_2}(\widehat{s},b)`$ $`=`$ $`{\displaystyle d^2b^{}D_{\mathrm{soft}/\mathrm{sea}\mathrm{sea}}(\widehat{s},|\stackrel{}{b}\stackrel{}{b}^{}|)\frac{1}{4\pi \left(R_{h_1}^2+R_{h_2}^2\right)}\mathrm{exp}\left[\frac{b^2}{4\left(R_{h_1}^2+R_{h_2}^2\right)}\right]},`$
which means that $`D_{\mathrm{soft}/\mathrm{sea}\mathrm{sea}}^{h_1h_2}`$ has the same functional form as $`D_{\mathrm{soft}/\mathrm{sea}\mathrm{sea}}`$, with $`\lambda _{\mathrm{soft}}^{(2)}(\xi )`$ being replaced by
$`\lambda _{\mathrm{soft}}^{h_1h_2}(\xi )`$ $`=`$ $`\lambda _{\mathrm{soft}}^{(2)}(\xi )+R_{h_1}^2+R_{h_2}^2`$ (2.52)
$``$ $`R_{h_1}^2+R_{h_2}^2+\alpha _{_{\mathrm{soft}}}^{}\mathrm{ln}\xi ,`$ (2.53)
where we neglected the parton slope $`R_{\mathrm{part}}^2`$ compared to the hadron slope $`R_h^2`$.
For the hard contribution, we have correspondingly $`D_{\mathrm{val}\mathrm{val}}^{h_1h_2,ij}(\widehat{s},b)`$ being given in eq. (2.18) with $`R_{\mathrm{hard}}^2`$ being replaced by $`R_{h_1}^2+R_{h_2}^2`$ (neglecting the hard scattering slope $`R_{\mathrm{hard}}^2`$ as compared to the hadron Regge slopes $`R_{h_1}^2,R_{h_2}^2`$). In case of mixed Pomerons, we have $`D_{\mathrm{val}\mathrm{sea}}^{h_1h_2,i}(\widehat{s},b)`$ given in eq. (2.2) with $`R_{\mathrm{part}}^2`$ being replaced by $`R_{h_1}^2+R_{h_2}^2`$.
### 2.4 Nucleus-Nucleus Scattering
We generalize the discussion of the last section in order to treat nucleus-nucleus scattering. In the Glauber-Gribov approach , the nucleus-nucleus scattering amplitude is defined by the sum of contributions of diagrams, shown at fig.2.1, corresponding to multiple scattering processes between parton constituents of projectile and target nucleons. Nuclear form factors are supposed to be defined by the nuclear ground state wave functions. Assuming the nucleons to be uncorrelated, one can make the Fourier transform to obtain the amplitude in the impact parameter representation. Then, for given impact parameter $`\stackrel{}{b}_0`$ between the nuclei, the only formal differences from the hadron-hadron case will be the possibility for a given nucleon to interact with a number of nucleons from the partner nucleus as well as the averaging over nuclear ground states, which amounts to an integration over transverse nucleon coordinates $`\stackrel{}{b}_i^A`$ and $`\stackrel{}{b}_j^B`$ in the projectile and in the target correspondingly. We write this integration symbolically as
$$𝑑T_{AB}:=\underset{i=1}{\overset{A}{}}d^2b_i^AT_A(b_i^A)\underset{j=1}{\overset{B}{}}d^2b_j^BT_B(b_j^B),$$
(2.54)
with $`A,B`$ being the nuclear mass numbers and with the so-called nuclear thickness function $`T_A(b)`$ being defined as the integral over the nuclear density $`\rho _{A(B)}`$,
$$T_A(b):=𝑑z\rho _A(\sqrt{b^2+z^2}).$$
(2.55)
For the nuclear densities, we use a parameterization of experimental data of ,
$$\rho _A(r)\{\begin{array}{ccc}\mathrm{exp}\left\{r^2/1.72^2\right\}\hfill & \mathrm{if}& A=2,\hfill \\ \mathrm{exp}\left\{r^2/(0.9A^{1/3})^2\right\}\hfill & \mathrm{if}& 2<A<10,\hfill \\ 1/\left\{1+\mathrm{exp}(\frac{r0.7A^{0.446}}{0.545})\right\}\hfill & \mathrm{if}& A10.\hfill \end{array}$$
(2.56)
It is convenient to use the transverse distance $`b_k`$ between the two nucleons from the $`k`$-th nucleon-nucleon pair, i.e.
$$b_k=\left|\stackrel{}{b}_0+\stackrel{}{b}_{\pi (k)}^A\stackrel{}{b}_{\tau (k)}^B\right|,$$
(2.57)
where the functions $`\pi (k)`$ and $`\tau (k)`$ refer to the projectile and the target nucleons participating in the $`k^{\mathrm{th}}`$ interaction (pair $`k`$). In order to simplify the notation, we define a vector $`b`$ whose components are the overall impact parameter $`b_0`$ as well as the transverse distances $`b_1,\mathrm{},b_{AB}`$ of the nucleon pairs,
$$b=\{b_0,b_1,\mathrm{},b_{AB}\}.$$
(2.58)
Then the nucleus-nucleus interaction cross section can be obtained applying the cutting procedure to elastic scattering diagrams of fig.2.1 and written in the form
$$\sigma _{\mathrm{inel}}^{AB}(s)=d^2b_0𝑑T_{AB}\gamma _{AB}(s,b),$$
(2.59)
where the so-called nuclear profile function $`\gamma _{AB}`$ represents an interaction for given transverse coordinates of the nucleons.
The calculation of the profile function $`\gamma _{AB}`$ as the sum over all cut diagrams of the type shown in fig.2.9 does not differ from the hadron-hadron case and follows the rules formulated
in the preceding section:
* For a remnant carrying the light cone momentum fraction $`x`$ ($`x^+`$ in case of projectile, or $`x^{}`$in case of target), one has a factor $`F_{\mathrm{remn}}(x)`$, defined in eq. (2.35).
* For each cut elementary diagram (real elementary interaction = dashed vertical line) attached to two participants with light cone momentum fractions $`x^+`$ and $`x^{}`$, one has a factor $`G(x^+,x^{},s,b),`$ given by eqs. (2.46-2.50). Apart of $`x^+`$ and $`x^{}`$, $`G`$ is also a function of the total squared energy $`s`$ and of the relative transverse distance $`b`$ between the two corresponding nucleons (we use $`G`$ as an abbreviation for $`G_{1\mathrm{IP}}^{NN}`$ for nucleon-nucleon scattering).
* For each uncut elementary diagram (virtual emissions = full vertical line) attached to two participants with light cone momentum fractions $`x^+`$ and $`x^{}`$, one has a factor $`G(x^+,x^{},s,b),`$ with the same $`G`$ as used for the cut diagrams.
* Finally one sums over all possible numbers of cut and uncut Pomerons and integrates over the light cone momentum fractions.
So we find
$`\gamma _{AB}(s,b)`$ $`=`$ $`{\displaystyle \underset{m_1l_1}{}}\mathrm{}{\displaystyle \underset{m_{AB}l_{AB}}{}}(1\delta _{0\mathrm{\Sigma }m_k}){\displaystyle \underset{k=1}{\overset{AB}{}}\left\{\underset{\mu =1}{\overset{m_k}{}}dx_{k,\mu }^+dx_{k,\mu }^{}\underset{\lambda =1}{\overset{l_k}{}}d\stackrel{~}{x}_{k,\lambda }^+d\stackrel{~}{x}_{k,\lambda }^{}\right\}}`$ (2.60)
$`\times `$ $`{\displaystyle \underset{k=1}{\overset{AB}{}}}\left\{{\displaystyle \frac{1}{m_k!}}{\displaystyle \frac{1}{l_k!}}{\displaystyle \underset{\mu =1}{\overset{m_k}{}}}G(x_{k,\mu }^+,x_{k,\mu }^{},s,b_k){\displaystyle \underset{\lambda =1}{\overset{l_k}{}}}G(\stackrel{~}{x}_{k,\lambda }^+,\stackrel{~}{x}_{k,\lambda }^{},s,b_k)\right\}`$
$`\times `$ $`{\displaystyle \underset{i=1}{\overset{A}{}}}F_{\mathrm{remn}}\left(x_i^+{\displaystyle \underset{\pi (k)=i}{}}\stackrel{~}{x}_{k,\lambda }^+\right){\displaystyle \underset{j=1}{\overset{B}{}}}F_{\mathrm{remn}}\left(x_j^{}{\displaystyle \underset{\tau (k)=j}{}}\stackrel{~}{x}_{k,\lambda }^{}\right)`$
with
$`x_i^{\mathrm{proj}}`$ $`=`$ $`1{\displaystyle \underset{\pi (k)=i}{}}x_{k,\mu ,}^+`$ (2.61)
$`x_j^{\mathrm{targ}}`$ $`=`$ $`1{\displaystyle \underset{\tau (k)=j}{}}x_{k,\mu }^{}.`$ (2.62)
The summation indices $`m_k`$ refer to the number of cut elementary diagrams and $`l_k`$ to the number of uncut elementary diagrams, related to nucleon pair $`k`$. For each possible pair $`k`$ (we have altogether $`AB`$ pairs), we allow for any number of cut and uncut diagrams. The integration variables $`x_{k,\mu }^\pm `$ refer to the $`\mu ^{\mathrm{th}}`$ elementary interaction of the $`k^{\mathrm{th}}`$ pair for the cut elementary diagrams, the variables $`\stackrel{~}{x}_{k,\lambda }^\pm `$ refer to the corresponding uncut elementary diagrams. The arguments of the remnant functions $`F_{\mathrm{remn}}`$ are the remnant light cone momentum fractions, i.e. unity minus the the momentum fractions of all the corresponding elementary contributions (cut and uncut ones). We also introduce the variables $`x_i^+`$and $`x_j^{}`$, defined as unity minus the momentum fractions of all the corresponding cut contributions, in order to integrate over the uncut ones (see below).
The expression for $`\gamma _{AB}(\mathrm{})`$ sums up all possible numbers of cut Pomerons $`m_k`$ with one exception due to the factor $`(1\delta _{0\mathrm{\Sigma }m_k})`$: one does not consider the case of all $`m_k`$’s being zero, which corresponds to “no interaction” and therefore does not contribute to the inelastic cross section. We may therefore define a quantity $`\overline{\gamma }_{AB}(\mathrm{})`$, representing “no interaction”, by taking the expression for $`\gamma _{AB}(\mathrm{})`$ with $`(1\delta _{0\mathrm{\Sigma }m_k})`$ replaced by $`(\delta _{0\mathrm{\Sigma }m_k})`$:
$`\overline{\gamma }_{AB}(s,b)`$ $`=`$ $`{\displaystyle \underset{l_1}{}}\mathrm{}{\displaystyle \underset{l_{AB}}{}}{\displaystyle \underset{k=1}{\overset{AB}{}}\left\{\underset{\lambda =1}{\overset{l_k}{}}d\stackrel{~}{x}_{k,\lambda }^+d\stackrel{~}{x}_{k,\lambda }^{}\right\}\underset{k=1}{\overset{AB}{}}\left\{\frac{1}{l_k!}\underset{\lambda =1}{\overset{l_k}{}}G(\stackrel{~}{x}_{k,\lambda }^+,\stackrel{~}{x}_{k,\lambda }^{},s,b_k)\right\}}`$ (2.63)
$`\times `$ $`{\displaystyle \underset{i=1}{\overset{A}{}}}F^+\left(1{\displaystyle \underset{\pi (k)=i}{}}\stackrel{~}{x}_{k,\lambda }^+\right){\displaystyle \underset{j=1}{\overset{B}{}}}F^{}\left(1{\displaystyle \underset{\tau (k)=j}{}}\stackrel{~}{x}_{k,\lambda }^{}\right).`$
One now may consider the sum of “interaction” and “no interaction”, and one obtains easily
$$\gamma _{AB}(s,b)+\overline{\gamma }_{AB}(s,b)=1.$$
(2.64)
Based on this important result, we consider $`\gamma _{AB}`$ to be the probability to have an interaction and correspondingly $`\overline{\gamma }_{AB}`$ to be the probability of no interaction, for fixed energy, impact parameter and nuclear configuration, specified by the transverse distances $`b_k`$ between nucleons, and we refer to eq. (2.64) as “unitarity relation”. But we want to go even further and use an expansion of $`\gamma _{AB}`$ in order to obtain probability distributions for individual processes, which then serves as a basis for the calculations of exclusive quantities.
The expansion of $`\gamma _{AB}`$ in terms of cut and uncut Pomerons as given above represents a sum of a large number of positive and negative terms, including all kinds of interferences, which excludes any probabilistic interpretation. We have therefore to perform summations of interference contributions - sum over any number of virtual elementary scatterings (uncut Pomerons) - for given non-interfering classes of diagrams - with given numbers of real scatterings (cut Pomerons) . Let us write the formulas explicitly. We have
$`\gamma _{AB}(s,b)`$ $`=`$ $`{\displaystyle \underset{m_1}{}}\mathrm{}{\displaystyle \underset{m_{AB}}{}}(1\delta _{0{\scriptscriptstyle m_k}}){\displaystyle \underset{k=1}{\overset{AB}{}}\left\{\underset{\mu =1}{\overset{m_k}{}}dx_{k,\mu }^+dx_{k,\mu }^{}\right\}}`$ (2.65)
$`\times `$ $`{\displaystyle \underset{k=1}{\overset{AB}{}}}\left\{{\displaystyle \frac{1}{m_k!}}{\displaystyle \underset{\mu =1}{\overset{m_k}{}}}G(x_{k,\mu }^+,x_{k,\mu }^{},s,b_k)\right\}\mathrm{\Phi }_{AB}(x^{\mathrm{proj}},x^{\mathrm{targ}},s,b),`$
where the function $`\mathrm{\Phi }`$ representing the sum over virtual emissions (uncut Pomerons) is given by the following expression
$`\mathrm{\Phi }_{AB}(x^{\mathrm{proj}},x^{\mathrm{targ}},s,b)`$ $`=`$ $`{\displaystyle \underset{l_1}{}}\mathrm{}{\displaystyle \underset{l_{AB}}{}}{\displaystyle \underset{k=1}{\overset{AB}{}}\left\{\underset{\lambda =1}{\overset{l_k}{}}d\stackrel{~}{x}_{k,\lambda }^+d\stackrel{~}{x}_{k,\lambda }^{}\right\}\underset{k=1}{\overset{AB}{}}\left\{\frac{1}{l_k!}\underset{\lambda =1}{\overset{l_k}{}}G(\stackrel{~}{x}_{k,\lambda }^+,\stackrel{~}{x}_{k,\lambda }^{},s,b_k)\right\}}`$ (2.66)
$`\times `$ $`{\displaystyle \underset{i=1}{\overset{A}{}}}F_{\mathrm{remn}}\left(x_i^{\mathrm{proj}}{\displaystyle \underset{\pi (k)=i}{}}\stackrel{~}{x}_{k,\lambda }^+\right){\displaystyle \underset{j=1}{\overset{B}{}}}F_{\mathrm{remn}}\left(x_j^{\mathrm{targ}}{\displaystyle \underset{\tau (k)=j}{}}\stackrel{~}{x}_{k,\lambda }^{}\right).`$
This summation has to be carried out, before we may use the expansion of $`\gamma _{AB}`$ to obtain probability distributions. This is far from trivial, as we are going to discuss in the next section, but let us assume for the moment that it can be done. To make the notation more compact, we define matrices $`X^+`$ and $`X^{}`$, as well as a vector $`m`$, via
$`X^+`$ $`=`$ $`\left\{x_{k,\mu }^+\right\},`$ (2.67)
$`X^{}`$ $`=`$ $`\left\{x_{k,\mu }^{}\right\},`$ (2.68)
$`m`$ $`=`$ $`\{m_k\},`$ (2.69)
which leads to
$`\gamma _{AB}(s,b)`$ $`=`$ $`{\displaystyle \underset{m}{}}(1\delta _{0m}){\displaystyle 𝑑X^+𝑑X^{}\mathrm{\Omega }_{AB}^{(s,b)}(m,X^+,X^{})},`$ (2.70)
$`\overline{\gamma }_{AB}(s,b)`$ $`=`$ $`\mathrm{\Omega }_{AB}^{(s,b)}(0,0,0),`$ (2.71)
with
$$\mathrm{\Omega }_{AB}^{(s,b)}(m,X^+,X^{})=\underset{k=1}{\overset{AB}{}}\left\{\frac{1}{m_k!}\underset{\mu =1}{\overset{m_k}{}}G(x_{k,\mu }^+,x_{k,\mu }^{},s,b_k)\right\}\mathrm{\Phi }_{AB}(x^{\mathrm{proj}},x^{\mathrm{targ}},s,b).$$
(2.72)
This allows to rewrite the unitarity relation eq. (2.64) in the following form,
$$\underset{m}{}𝑑X^+𝑑X^{}\mathrm{\Omega }_{AB}^{(s,b)}(m,X^+,X^{})=1.$$
(2.73)
This equation is of fundamental importance, because it allows us to interpret $`\mathrm{\Omega }^{(s,b)}(m,X^+,X^{})`$ as probability density of having an interaction configuration characterized by $`m`$, with the light cone momentum fractions of the Pomerons being given by $`X^+`$ and $`X^{}`$.
### 2.5 Diffractive Scattering
We do not have a consistent treatment of diffractive scattering at the moment, this is left to a future project in connection with a complete treatment of enhanced diagrams. For the moment, we introduce diffraction “by hand”: in case of no interaction in $`pp`$ or $`pA`$ scattering, we consider the projectile to be diffractively excited with probability
$$w_{\mathrm{diff}}\frac{\left(1\sqrt{\mathrm{\Phi }(1,x^{\mathrm{targ}},s,b)}\right)^2}{\mathrm{\Phi }(1,x^{\mathrm{targ}},s,b)},$$
(2.74)
with a fit parameter $`w_{\mathrm{diff}}`$. Nucleus-nucleus scattering is here (but only here!) considered as composed of $`pA`$ or $`Ap`$ collisions.
### 2.6 AGK Cancelations in Hadron-Hadron Scattering
As a first application, we are going to prove that AGK cancelations apply perfectly in our model<sup>2</sup><sup>2</sup>2 We speak here about the contribution of elementary interactions (Pomeron exchanges) to the secondary particle production; the AGK cancellations do not hold for the contribution of remnant states (spectator partons) hadronization . .
As we showed above, the description of high energy hadronic interaction requires to consider explicitly a great number of contributions, corresponding to multiple scattering process, with a number of elementary parton-parton interactions happening in parallel. However, when calculating inclusive spectra of secondary particles, it is enough to consider the simplest hadron-hadron (nucleus-nucleus) scattering diagrams containing a single elementary interaction, as the contributions of multiple scattering diagrams with more than one elementary interaction exactly cancel each other. This so-called AGK-cancelation is a consequence of the general Abramovskii-Gribov-Kancheli cutting rules .
Let us consider the most fundamental inclusive distribution, where all other inclusive spectra may be derived from: the distribution $`dn_{\mathrm{Pom}}^{h_1h_2}/dx^+dx^{}`$, with $`dn_{\mathrm{Pom}}^{h_1h_2}`$ being the number of Pomerons with light cone momentum fractions between $`x^+`$ and $`x^++dx^+`$and between $`x^{}`$ and $`x^{}+dx^{}`$ respectively, at a given value of $`b`$ and $`s`$. If AGK cancelations apply, the result for $`dn_{\mathrm{Pom}}^{h_1h_2}/dx^+dx^{}`$ should coincide with the contribution coming from exactly one elementary interaction (see eq. (2.45)):
$$\frac{dn_{\mathrm{Pom}}^{(1)h_1h_2}}{dx^+dx^{}}(x^+,x^{},s,b)=G_{1\mathrm{I}\mathrm{P}}^{h_1h_2}(x^+,x^{},s,b)F_{\mathrm{remn}}^{h_1}(1x^+)F_{\mathrm{remn}}^{h_2}(1x^{}),$$
(2.75)
and the contributions from multiple scattering should exactly cancel. We have per definition
$`{\displaystyle \frac{dn_{\mathrm{Pom}}^{h_1h_2}}{dx^+dx^{}}}(x^+,x^{},s,b)`$ $`=`$ $`{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{\mu =1}{\overset{m}{}}dx_\mu ^+dx_\mu ^{}\underset{\lambda =m+1}{\overset{m+l}{}}dx_\lambda ^+dx_\lambda ^{}}`$
$`\times `$ $`{\displaystyle \frac{1}{m!}}{\displaystyle \frac{1}{l!}}{\displaystyle \underset{\mu =1}{\overset{m}{}}}G_{1\mathrm{I}\mathrm{P}}^{h_1h_2}(x_\mu ^+,x_\mu ^{},s,b){\displaystyle \underset{\lambda =m+1}{\overset{m+l}{}}}G_{1\mathrm{I}\mathrm{P}}^{h_1h_2}(x_\lambda ^+,x_\lambda ^{},s,b)`$
$`\times `$ $`F_{\mathrm{remn}}^{h_1}(1{\displaystyle \underset{\nu =1}{\overset{m+l}{}}}x_\nu ^+)F_{\mathrm{remn}}^{h_2}(1{\displaystyle \underset{\nu =1}{\overset{m+l}{}}}x_\nu ^{})`$
$`\times `$ $`{\displaystyle \underset{\mu ^{}=1}{\overset{m}{}}}\delta (x^+x_\mu ^{}^+)\delta (x^{}x_\mu ^{}^{}).`$
Due to the symmetry of the integrand in the r.h.s. of eq. (2.6) in the variables $`x_1^\pm ,\mathrm{},x_m^\pm `$, the sum of delta functions produces a factor $`mG(x^+,x^{},s,b)`$, and removes one $`dx_\mu ^+dx_\mu ^{}`$ integration. Using
$$m\underset{m=1}{\overset{\mathrm{}}{}}\mathrm{}\frac{1}{m!}\underset{\mu =1}{\overset{m1}{}}\mathrm{}=\underset{m^{}=0}{\overset{\mathrm{}}{}}\mathrm{}\frac{1}{m^{}!}\underset{\mu =1}{\overset{m^{}}{}}\mathrm{},$$
(2.77)
with $`m^{}=m1,`$ we get
$`{\displaystyle \frac{dn_{\mathrm{Pom}}^{h_1h_2}}{dx^+dx^{}}}(x^+,x^{},s,b)`$ $`=`$ $`G_{1\mathrm{I}\mathrm{P}}^{h_1h_2}(x^+,x^{},s,b){\displaystyle \frac{1}{n!}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{\nu =1}{\overset{n}{}}dx_\nu ^+dx_\nu ^{}}`$ (2.78)
$`\times `$ $`\left\{{\displaystyle \underset{m=0}{\overset{n}{}}}\left(\begin{array}{c}n\\ m\end{array}\right){\displaystyle \underset{\mu =1}{\overset{m}{}}}G_{1\mathrm{I}\mathrm{P}}^{h_1h_2}(x_\mu ^+,x_\mu ^{},s,b){\displaystyle \underset{\lambda =m+1}{\overset{n}{}}}G_{1\mathrm{I}\mathrm{P}}^{h_1h_2}(x_\lambda ^+,x_\lambda ^{},s,b)\right\}`$
$`\times `$ $`F_{\mathrm{remn}}^{h_1}(1x^+{\displaystyle \underset{\nu =1}{\overset{n}{}}}x_\nu ^+)F_{\mathrm{remn}}^{h_2}(1x^{}{\displaystyle \underset{\nu =1}{\overset{n}{}}}x_\nu ^{}).`$
The term in curly brackets $`\left\{\mathrm{}\right\}`$ is 1 for $`n=0`$ and zero otherwise, so we get the important final result
$$\frac{dn_{\mathrm{Pom}}^{h_1h_2}}{dx^+dx^{}}(x^+,x^{},s,b)=\frac{dn_{\mathrm{Pom}}^{(1)h_1h_2}}{dx^+dx^{}}(x^+,x^{},s,b),$$
(2.82)
which corresponds to one single elementary interaction; the multiple scattering aspects completely disappeared, so AGK cancelations indeed apply in our approach. AGK cancelations are closely related to the factorization formula for jet production cross section, since as a consequence of eq. (2.82), we may obtain the inclusive jet cross section in a factorized form as
$$\sigma _{\mathrm{jet}}^{h_1h_2}=\underset{j,k}{}𝑑p_{}^2𝑑z^+𝑑z^{}f_j^{h_1}(z^+,M_F^2)f_k^{h_2}(z^{},M_F^2)\frac{d\sigma _{\mathrm{Born}}^{jk}}{dp_{}^2}(z^+z^{}s,p_{}^2),$$
(2.83)
with $`f_j^{h_1}`$ and $`f_k^{h_2}`$ representing the parton distributions of the two hadrons.
### 2.7 AGK Cancelations in Nucleus-Nucleus Scattering
We have shown in the previous section that for hadron-hadron scattering AGK cancelations apply, which means that inclusive spectra coincide with the contributions coming from exactly one elementary interaction. For multiple Pomeron exchanges we have a complete destructive interference, they do not contribute at all. Here, we are going to show that AGK cancelations also apply for nucleus-nucleus scattering, which means that the inclusive cross section for $`A+B`$ scattering is $`AB`$ times the corresponding inclusive cross section for proton-proton interaction.
The inclusive cross section for forming a Pomeron with light cone momentum fractions $`x^+`$ and $`x^{}`$ in nucleus-nucleus scattering is given as
$`{\displaystyle \frac{d\sigma _{\mathrm{Pom}}^{AB}}{dx^+dx^{}}}(x^+,x^{},s)`$ $`=`$ $`{\displaystyle d^2b_0𝑑T_{AB}}`$ (2.84)
$`\times `$ $`{\displaystyle \underset{m_1l_1}{}}\mathrm{}{\displaystyle \underset{m_{AB}l_{AB}}{}}(1\delta _{0{\scriptscriptstyle m_k}}){\displaystyle \underset{k=1}{\overset{AB}{}}\left\{\underset{\mu =1}{\overset{m_k+l_k}{}}dx_{k,\mu }^+dx_{k,\mu }^{}\right\}}`$
$`\times `$ $`{\displaystyle \underset{k=1}{\overset{AB}{}}}\left\{{\displaystyle \frac{1}{m_k!l_k!}}{\displaystyle \underset{\mu =1}{\overset{m_k}{}}}G(x_{k,\mu }^+,x_{k,\mu }^{},s,b_k){\displaystyle \underset{\lambda =m_k+1}{\overset{m_k+l_k}{}}}G(x_{k,\lambda }^+,x_{k,\lambda }^{},s,b_k)\right\}`$
$`\times `$ $`{\displaystyle \underset{i=1}{\overset{A}{}}}F_{\mathrm{remn}}\left(1{\displaystyle \underset{\pi (k)=i}{}}x_{k,\lambda }^+\right){\displaystyle \underset{j=1}{\overset{B}{}}}F_{\mathrm{remn}}\left(1{\displaystyle \underset{\tau (k)=j}{}}x_{k,\lambda }^{}\right)`$
$`\times `$ $`{\displaystyle \underset{k^{}=1}{\overset{AB}{}}}{\displaystyle \underset{\mu ^{}=1}{\overset{m_k}{}}}\delta (x^+x_{k^{}\mu ^{}}^+)\delta (x^{}x_{k^{}\mu ^{}}^{}).`$
The factor $`(1\delta _{0{\scriptscriptstyle m_k}})`$ makes sure that at least one of the indices $`m_k`$ is bigger than zero. Integrating over the variables appearing in the delta functions, we obtain a factor $`_k^{}G(x^+,x^{},s,b_k)m_k^{}`$ which may be written in front of $`_{m_1l_1}\mathrm{}`$ . In the following expression one may rename the integration variables such that the variables $`x_{k^{}m_k^{}}^+`$ and $`x_{k^{}m_k^{}}^{}`$ disappear. This means for the arguments of the functions $`F_{\mathrm{remn}}`$ that for $`i=\pi (k^{})`$ and $`j=\tau (k^{})`$ one replaces $`1`$ by $`1x^+`$ and $`1x^{}`$ respectively. Then one uses the factor $`m_k^{}`$ mentioned above to replace $`m_k^{}!`$ by $`(m_k^{}1)!`$. One finally renames $`(m_k^{}1)`$ by $`m_k^{}`$, as a consequence of which one may drop the factor $`(1\delta _{0{\scriptscriptstyle m_k}})`$. This is crucial, since now we have factors of the form
$$\underset{m_k=0}{\overset{\mathrm{}}{}}\underset{l_k=0}{\overset{\mathrm{}}{}}\mathrm{}\frac{1}{m_k!l_k!}\underset{\mu =1}{\overset{m_k}{}}G(x_{k,\mu }^+,x_{k,\mu }^{},s,b)\underset{\lambda =m_k+1}{\overset{m_k+l_k}{}}G(x_{k,\lambda }^+,x_{k,\lambda }^{},s,b)\mathrm{}.$$
(2.85)
In this sum only the term for $`m_k=0`$ and $`l_k=0`$ is different from zero, namely $`1`$, and so we get
$$\frac{d\sigma _{\mathrm{Pom}}^{AB}}{dx^+dx^{}}(x^+,x^{},s,b)=d^2b_0𝑑T_{AB}\underset{k^{}}{}G(x^+,x^{},s,b_k^{})F_{\mathrm{remn}}(1x^+)F_{\mathrm{remn}}(1x^{}).$$
(2.86)
Using the definition of $`dT_{AB}`$, writing $`b_k^{}`$ explicitly as $`|\stackrel{}{b}_0+\stackrel{}{b}_{\pi (k)}^A\stackrel{}{b}_{\tau (k)}^B|`$, we obtain
$`{\displaystyle \frac{d\sigma _{\mathrm{Pom}}^{AB}}{dx^+dx^{}}}(x^+,x^{},s)`$ $`=`$ $`AB{\displaystyle d^2b_0d^2b_+T(b_+)d^2b_{}T(b_{})G(x^+,x^{},s,|\stackrel{}{b}_0+\stackrel{}{b}_+\stackrel{}{b}_{}|)}`$ (2.87)
$`\times F_{\mathrm{remn}}(1x^+)F_{\mathrm{remn}}(1x^{}).`$
Changing the order of the integrations, we obtain finally
$$\frac{d\sigma _{\mathrm{Pom}}^{AB}}{dx^+dx^{}}(x^+,x^{},s)=AB\frac{d\sigma _{\mathrm{Pom}}^{pp}}{dx^+dx^{}}(x^+,x^{},s)$$
(2.88)
with
$$\frac{d\sigma _{\mathrm{Pom}}^{pp}}{dx^+dx^{}}(x^+,x^{},s)=d^2bG(x^+,x^{},s,b)F_{\mathrm{remn}}(1x^+)F_{\mathrm{remn}}(1x^{}).$$
(2.89)
Since any other inclusive cross section $`d\sigma _{\mathrm{incl}}/dq`$ may be obtained from the inclusive Pomeron distribution via convolution, we obtain the very general result
$$\frac{d\sigma _{\mathrm{incl}}^{AB}}{dq}(q,s,b)=AB\frac{d\sigma _{\mathrm{incl}}^{pp}}{dq}(q,s,b),$$
(2.90)
so nucleus-nucleus inclusive cross sections are just $`AB`$ times the proton-proton ones. So, indeed, AGK cancelations apply perfectly in our approach.
### 2.8 Outlook
What did we achieve so far? We have a well defined model, introduced by using the language of field theory (Feynman diagrams). We were able to prove some elementary properties (AGK cancelations in case of proton-proton and nucleus-nucleus scattering). To proceed further, we have to solve (at least) two fundamental problems:
* the sum over virtual emissions has to be performed,
* tools have to be developed to deal with the multidimensional probability distribution
$`\mathrm{\Omega }_{AB}^{(s,b)}(m,X^+,X^{})`$,
both being very difficult tasks.
Calculating the sum over virtual emissions ($`\mathrm{\Phi }_{AB}`$) is not only technically difficult, there are also conceptual problems. By studying the properties of $`\mathrm{\Phi }_{AB}`$, we find that at very high energies the theory is no longer unitary without taking into account additional screening corrections. In this sense, we consider our work as a first step to construct a consistent model for high energy nuclear scattering, but there is still work to be done.
Concerning the multidimensional probability distribution $`\mathrm{\Omega }_{AB}^{(s,b)}(m,X^+,X^{})`$, we are going to develop methods well known in statistical physics (Markov chain techniques), which we also are going to discuss in detail later. So finally, we are able to calculate the probability distribution $`\mathrm{\Omega }_{AB}^{(s,b)}(m,X^+,X^{})`$, and are able to generate (in a Monte Carlo fashion) configurations $`(m,X^+,X^{})`$ according to this probability distribution.
The two above mentioned problems will be discussed in detail in the following chapters.
## Chapter 3 Virtual Emissions
In order to proceed, we need to calculate the sum over virtual emissions, represented by the function $`\mathrm{\Phi }_{AB}`$. Understanding the behavior of $`\mathrm{\Phi }_{AB}`$ is crucial, since this function is related to $`\overline{\gamma }_{AB}`$ and plays therefore a crucial role in connection with unitarity, the unitarity equation being given as $`\gamma _{AB}+\overline{\gamma }_{AB}=1`$. By studying the properties of $`\mathrm{\Phi }_{AB}`$, we find inconsistencies in the limit of high energies, in the sense that the individual terms appearing in the unitarity equation are not necessarily positive. Attempting to understand this unphysical behavior, we find that any model where AGK cancelations apply (so most of the models used presently) has to run asymptotically into this problem. So eventually one needs to construct models, where AGK cancelations are violated, which is going to be expected when contributions of Pomeron-Pomeron interactions are taken into consideration.
As a first phenomenological solution of the unitarity problem, we are going to “unitarize” the “bare theory” introduced in the preceding chapter “by hand”, such that the theory is changed as little as possible, but the asymptotic problems disappear. The next step should of course be a consistent treatment including contributions of enhanced Pomeron diagrams.
In the following, we are going to present the calculation of $`\mathrm{\Phi }_{AB}`$, we discuss the unitarity problems and the phenomenological solution, as well as properties of the “unitarized theory”.
### 3.1 Parameterizing the Elementary Interaction
The basis for all the calculations which follow is the function $`G_{1\mathrm{I}\mathrm{P}}^{NN}`$, which is the profile function representing a single elementary nucleon-nucleon ($`NN`$) interaction. For simplicity, we write simply $`GG_{1\mathrm{I}\mathrm{P}}^{NN}`$. This function $`G`$ is a sum of several terms, representing soft, semi-hard, valence, and screening contributions. In case of soft and semi-hard, one has $`G=(x^+x^{})^{\alpha _{\mathrm{part}}}D`$, where $`D`$ represents the Pomeron exchange and the factor in front of $`D`$ the “external legs”, the nucleon participants. For the other contributions the functional dependence on $`x^+,x^{}`$ is somewhat more complicated, but nevertheless it is convenient to define a function
$$D(x^+,x^{},s,b)=\frac{G(x^+,x^{},s,b)}{(x^+x^{})^{\alpha _{\mathrm{part}}}}.$$
(3.1)
We obtain $`G`$ and therefore $`D`$ as the result of a quite involved numerical calculation, which means that these functions are given in a discretized fashion. Since this is not very convenient and since the dependence of $`x^+`$ , $`x^{}`$, and $`b`$ are quite simple, we are going to parameterize our numerical results and use this analytical expression for further calculations. This makes the following discussions much easier and more transparent.
We first consider the case of zero impact parameter ($`b=0`$). In fig. 3.1, we plot the function $`D`$ together with the individual contributions as functions of $`x=x^+x^{}`$, for $`b=0`$ and for different values of $`s`$.
To fit the function $`D`$, we make the following ansatz,
$$D(x^+,x^{},s,b=0)=\underset{i=1}{\overset{N}{}}\underset{D_i}{\underset{}{\left\{\alpha _{D_i}+\alpha _{D_i}^{}(x^+x^{})^{\beta _{D_i}^{}}\right\}\left(x^+x^{}s\right)^{\beta _{D_i}}}},$$
(3.2)
where the parameters may depend on $`s`$, and the parameters marked with a star are non-zero for a given $`i`$ only if the corresponding $`\alpha _{D_i}`$ is zero. This parameterization works very well, as shown in fig. 3.2, where we compare the original $`D`$ function with the fit according to eq. (3.2).
Let us now consider the $`b`$-dependence for fixed $`x^+`$ and $`x^{}`$. Since we observe almost a Gaussian shape with a weak (logarithmic) dependence of the width on $`x=x^+x^{}`$, we could make the following ansatz:
$`D(x^+,x^{},s,b)={\displaystyle \underset{i=1}{\overset{N}{}}}\left\{\alpha _{D_i}+\alpha _{D_i}^{}(x^+x^{})^{\beta _{D_i}^{}}\right\}(x^+x^{}s)^{\beta _{D_i}}\mathrm{exp}\left({\displaystyle \frac{b^2}{\delta _{D_i}+ϵ_{D_i}\mathrm{ln}(x^+x^{}s)}}\right).`$ (3.3)
However, we can still simplify the parameterization. We have
$`\mathrm{exp}\left({\displaystyle \frac{b^2}{\delta +ϵ\mathrm{ln}x}}\right)`$ $``$ $`\mathrm{exp}\left({\displaystyle \frac{b^2}{\delta }}\left(1{\displaystyle \frac{ϵ}{\delta }}\mathrm{ln}(x^+x^{}s)\right)\right)`$ (3.4)
$`=`$ $`\mathrm{exp}\left({\displaystyle \frac{b^2}{\delta }}\right)(x^+x^{}s)^{\gamma b^2}`$ (3.5)
So we make finally the ansatz
$$D(x^+,x^{},s,b)=\underset{i=1}{\overset{N}{}}\left\{\alpha _{D_i}+\alpha _{D_i}^{}(x^+x^{})^{\beta _{D_i}^{}}\right\}\left(x^+x^{}s\right)^{\beta _{D_i}+\gamma _{D_i}b^2}e^{\frac{b^2}{\delta _{D_i}}},$$
(3.6)
which provides a very good analytical representation of the numerically obtained function $`D`$, as shown in fig. 3.3.
### 3.2 Calculating $`\mathrm{\Phi }`$ for proton-proton collisions
We first consider proton-proton collisions. To be more precise, we are going to derive an expression for $`\mathrm{\Phi }_{pp}`$ which can be evaluated easily numerically and which will serve as the basis to investigate the properties of $`\mathrm{\Phi }_{pp}`$. We have
$`\mathrm{\Phi }_{pp}(x^+,x^{},s,b)`$ $`=`$ $`{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}{\displaystyle 𝑑x_1^+𝑑x_1^{}\mathrm{}𝑑x_l^+𝑑x_l^{}\left\{\frac{1}{l!}\underset{\lambda =1}{\overset{l}{}}G(x_\lambda ^+,x_\lambda ^{},s,b)\right\}}`$ (3.7)
$`\times `$ $`F_{\mathrm{remn}}\left(x^+{\displaystyle x_\lambda ^+}\right)F_{\mathrm{remn}}\left(x^{}{\displaystyle x_\lambda ^{}}\right).`$
with
$`F_{\mathrm{remn}}(x)`$ $`=`$ $`x^{\alpha _{\mathrm{remn}}}\mathrm{\Theta }(x)\mathrm{\Theta }(1x),`$ (3.8)
where $`\mathrm{\Theta }(x)`$ is the Heavyside function, and
$`G(x_\lambda ^+,x_\lambda ^{},s,b)`$ $`=`$ $`(x_\lambda ^+x_\lambda ^{})^{\alpha _{\mathrm{part}}}D(x_\lambda ^+,x_\lambda ^{},s,b).`$ (3.9)
Using eq. (3.6), we have
$$G(x_\lambda ^+,x_\lambda ^{},s,b)=\underset{i=1}{\overset{N}{}}\underset{G_{i,\lambda }}{\underset{}{\alpha _i(x_\lambda ^+x_\lambda ^{})^{\beta _i}}}$$
(3.10)
with
$`\alpha _i`$ $`=`$ $`\left(\alpha _{D_i}+\alpha _{D_i}^{}\right)s^{\left(\beta _{D_i}+\gamma _{D_i}b^2\right)}e^{\frac{b^2}{\delta _{D_i}}},`$ (3.11)
$`\beta _i`$ $`=`$ $`\beta _{D_i}+\beta _{D_i}^{}+\gamma _{D_i}b^2\alpha _{\mathrm{part}},`$ (3.12)
with $`\alpha _{D_i}^{}0`$ and $`\beta _{D_i}^{}0`$ only if $`\alpha _{D_i}=0`$. Using eq. (3.10), we obtain from eq. (3.7) the following expression
$`\mathrm{\Phi }_{pp}(x^+,x^{}s,b)`$ $`=`$ $`{\displaystyle \underset{r_1=0}{\overset{\mathrm{}}{}}}\mathrm{}{\displaystyle \underset{r_N=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{r_1!}}\mathrm{}{\displaystyle \frac{1}{r_N!}}{\displaystyle \underset{\lambda =1}{\overset{r_1+\mathrm{}+r_N}{}}dx_\lambda ^+dx_\lambda ^{}}`$ (3.13)
$`\times `$ $`{\displaystyle \underset{\rho _1=1}{\overset{r_1}{}}}G_{1,\rho _1}\mathrm{}{\displaystyle \underset{\rho _N=r_1+\mathrm{}+r_{N1}+1}{\overset{r_1+\mathrm{}+r_N}{}}}G_{N,\rho _N}`$
$`\times `$ $`F_{\mathrm{remn}}(x^+{\displaystyle \underset{\lambda }{}}x_\lambda ^+)F_{\mathrm{remn}}(x^{}{\displaystyle \underset{\lambda }{}}x_\lambda ^{}).`$
Using the fact that the functions $`G_{i,\lambda }`$ are separable,
$$G_{i,\lambda }=\alpha _i(x_\lambda ^+)^{\beta _i}(x_\lambda ^{})^{\beta _i},$$
(3.14)
one finds finally (see appendix D.1)
$`\mathrm{\Phi }_{pp}(x^+,x^{},s,b)`$ $`=`$ $`x^{\alpha _{\mathrm{remn}}}{\displaystyle \underset{r_1=0}{\overset{\mathrm{}}{}}}\mathrm{}{\displaystyle \underset{r_N=0}{\overset{\mathrm{}}{}}}\left\{{\displaystyle \frac{\mathrm{\Gamma }(1+\alpha _{\mathrm{remn}})}{\mathrm{\Gamma }(1+\alpha _{\mathrm{remn}}+r_1\stackrel{~}{\beta }_1+\mathrm{}+r_N\stackrel{~}{\beta }_N)}}\right\}^2`$ (3.15)
$`\times `$ $`{\displaystyle \frac{(\alpha _1x^{\stackrel{~}{\beta }_1}\mathrm{\Gamma }^2(\stackrel{~}{\beta }_1))^r}{r_1!}}\mathrm{}{\displaystyle \frac{(\alpha _Nx^{\stackrel{~}{\beta }_N}\mathrm{\Gamma }^2(\stackrel{~}{\beta }_N))^t}{r_N!}}`$
with $`x=x^+x^{}`$and $`\stackrel{~}{\beta }_i=\beta _i+1`$. Since the sums converge very fast, this expression can be easily evaluated numerically.
### 3.3 Unitarity Problems
In this section, we are going to present numerical results for $`\mathrm{\Phi }_{pp}`$, based on equation (3.15). We will observe an unphysical behavior in certain regions of phase space, which amounts to a violation of unitarity. Trying to understand its physical origin, we find that AGK cancelations, which apply in our model, automatically lead to unitarity violations. This means on the other hand that a fully consistent approach requires explicit violation of AGK cancelations, which should occur in case of considering contributions of enhanced Pomeron diagrams.
In which way is $`\mathrm{\Phi }_{pp}`$ related to unitarity? We have shown in the preceding chapter that the inelastic non-diffractive cross section $`\sigma _{\mathrm{inel}}(s)`$ may be written as
$$\sigma _{\mathrm{inel}}(s)=d^2b\gamma _{pp}(s,b),$$
(3.16)
with the profile function $`\gamma _{pp}(s,b)`$ representing all diagrams with at least one cut Pomeron. We defined as well the corresponding quantity $`\overline{\gamma }_{pp}(s,b)`$ representing all diagrams with zero cut Pomerons. We demonstrated that the sum of these two quantities is one (2.64),
$$\gamma _{pp}(s,b)+\overline{\gamma }_{pp}(s,b)=1,$$
(3.17)
which represents a unitarity relation. The function $`\mathrm{\Phi }_{pp}`$ enters finally, since we have the relation
$$\overline{\gamma }_{pp}(s,b)=\mathrm{\Phi }_{pp}(1,1,s,b).$$
(3.18)
Based on these formulas, we interpret $`\overline{\gamma }_{pp}(s,b)=\mathrm{\Phi }_{pp}(1,1,s,b)`$ as the probability of having no interaction, whereas $`\gamma _{pp}(s,b)=1\mathrm{\Phi }_{pp}(1,1,s,b)`$ represents the probability to have an interaction, at given impact parameter and energy. Such an interpretation of course only makes sense as long as any of the $`\gamma `$’s is positive, otherwise unitarity is said to be violated, even if the equation (3.17) still holds.
In fig. 3.4, we plot $`\mathrm{\Phi }_{pp}`$ as a function of $`x=x^+x^{}`$ for $`\sqrt{s}=200`$ GeV for two different values of $`b`$. The curve for $`b=1.5`$ fm (solid curve) is close to one with a minimum of about 0.8 at $`x=1`$. The $`x`$-dependence for $`b=0`$ fm (dashed curve) is much more dramatic: the curve deviates from 1 already at relatively small values of $`x`$ and drops finally to negative values at $`x=1.`$
The values for $`x=1`$ are of particular interest, since $`1\mathrm{\Phi }_{pp}(1,1,s,b)=\gamma _{pp}(s,b)`$ represents the profile function in the sense that the integration over $`b`$ provides the inelastic non-diffractive cross section. Therefore, in fig. 3.5, we plot the $`b`$-dependence of $`1\mathrm{\Phi }_{pp}(1,1,s,b)`$, which increases beyond 1 for small values of $`b`$, since $`\mathrm{\Phi }_{pp}`$ is negative in this region, as discussed above for the case of $`b=0`$ fm.
On the other hand, an upper limit of 1 is really necessary in order to assure unitarity. So the fact that $`\mathrm{\Phi }_{pp}`$ grows to values bigger than one is a manifestation of unitarity violation.
In the following, we try to understand the physical reason for this unitarity problem. We are going to show, that it is intimately related to the fact that in our approach AGK cancelations are fulfilled, as shown earlier, which means that any approach where AGK cancelations apply will have exactly the same problem.
We are going to demonstrate in the following that AGK cancelations imply automatically unitarity violation. The average light cone momentum taken by a Pomeron may be calculated from the Pomeron inclusive spectrum $`dn_{\mathrm{Pom}}/dx^+dx^{}`$ as
$$<x^+>=𝑑x^+x^+\left\{\frac{1}{\sigma _{\mathrm{inel}}(s)}d^2b𝑑x^{}\frac{dn_{\mathrm{Pom}}}{dx^+dx^{}}(x^+,x^{},s,b)\right\}.$$
(3.19)
If AGK cancelations apply, we have
$$dn_{\mathrm{Pom}}/dx^+dx^{}=dn_{\mathrm{Pom}}^{(1)}/dx^+dx^{}=G(x^+,x^{},s,b)F_{\mathrm{remn}}(x^+)F_{\mathrm{remn}}(x^{}),$$
(3.20)
and therefore
$$\frac{dn_{\mathrm{Pom}}}{dx^+dx^{}}(x^+,x^{},s,b)=\mu (s)(x^+x^{}s)^{\mathrm{\Delta }(s)}e^{\frac{b^2}{\lambda (s)}}f(x^+)f(x^{}),$$
(3.21)
where $`\mathrm{\Delta }(s)`$ is bigger than zero and increases with energy and $`\mu (s)`$ and $`\lambda (s)`$ depend weakly (logarithmically) on $`s`$, whereas $`f`$ is an energy independent function. We obtain
$$<x^+>=\frac{\gamma (s)s^{\mathrm{\Delta }(s)}}{\sigma _{\mathrm{inel}}(s)},$$
(3.22)
where $`\gamma (s)`$ depends only logarithmically on $`s`$. Since $`<x^+>`$ must be smaller or equal to one, we find
$$\sigma _{\mathrm{inel}}(s)\gamma (s)s^{\mathrm{\Delta }(s)},$$
(3.23)
which violates the Froissard bound and therefore unitarity. This problem is related to the old problem of unitarity violation in case of single Pomeron exchange. The solution appeared to be the observation that one needs to consider multiple scattering such that virtual multiple emissions provide sufficient screening to avoid the unreasonably fast increase of the cross section. If AGK cancelations apply (as in our model), the problem comes back by considering inclusive spectra, since these are determined by single scattering.
Thus, we have shown that *a consistent application of the eikonal Pomeron scheme both to interaction cross sections and to particle production calculations unavoidably leads to the violation of the unitarity*. This problem is not observed in many models currently used, since there simply no consistent treatment is provided, and the problem is therefore hidden. The solution of the unitarity problem requires to employ the full Pomeron scheme, which includes also so-called enhanced Pomeron diagrams, to be discussed later. The simplest diagram of that kind - so-called Y-diagram, for example, contributes a negative factor to all inclusive particle distributions in the particle rapidity region $`y_0<y<Y`$, where $`Y`$ is the total rapidity range for the interaction and $`y_0`$ corresponds to the rapidity position of the Pomeron self-interaction vertex. Thus, one speaks about breaking of the AGK-cancelations in the sense that one gets corrections to all inclusive quantities calculated from just one Pomeron exchange graph . In particular, presenting the inclusive Pomeron distribution $`d\sigma _{\mathrm{Pom}}(x^+,x^{})/dx^+dx^{}`$ by the formula (2.75) implies that the functions $`f(x^\pm )`$ acquire a dependence on the energy of the interaction $`s`$. It is this dependence which is expected to slow down the energy increase of the Pomeron number and thus to cure the unitarity problem.
So we think it is mandatory to proceed in the following way: first one needs to provide a consistent treatment of cross sections and particle production, which will certainly lead to unitarity problems, and second one has to refine the theory to solve the unitarity problem in a consistent way, via screening corrections. The first part of this program is provided in this paper, the second one will be treated in some approximate fashion later, but a rigorous, self-consistent treatment of this second has still to be done.
### 3.4 A Phenomenological Solution: Unitarization of $`\mathrm{\Phi }`$
As we have seen in the preceding sections, unitarity violation manifests itself by the fact that the virtual emission function $`\mathrm{\Phi }_{pp}`$ appears to be negative at high energies and small impact parameter for large values of $`x^+`$ and $`x^{}`$, particularly for $`x^+=x^{}=1`$. What is the mathematical origin of these negative values? In eq. (3.15), the sums over $`r_i`$ contains terms of the form $`(\mathrm{})^{r_i}/r_i!`$ and an additional factor of the form
$$\left\{\frac{\mathrm{\Gamma }(1+\alpha _{\mathrm{remn}})}{\mathrm{\Gamma }(1+\alpha _{\mathrm{remn}}+r_1\stackrel{~}{\beta }_1+\mathrm{}+r_N\stackrel{~}{\beta }_N)}\right\}.$$
(3.24)
It is this factor which causes the problem, as it strongly suppresses contributions of terms with large $`r_i`$, which are important when the interaction energy $`s`$ increases . Physically it is connected to the reduced phase space in case of too many virtual Pomerons emitted. By dropping this factor one would obtain a simple exponential function which is definitely positive.
Our strategy is to modify the scheme such that $`\mathrm{\Phi }_{pp}`$ stays essentially unchanged for values of $`s`$, $`b`$, $`x^+`$, and $`x^{}`$, where $`\mathrm{\Phi }_{pp}`$ is positive and that $`\mathrm{\Phi }_{pp}`$ is “corrected” to assure positive values in regions where it is negative. We call this procedure “unitarization”, which should not be considered as an approximation, since one is really changing the physical content of the theory. This is certainly only a phenomenological solution of the problem, the correct procedure should amount to taking into account the mentioned screening corrections due to enhanced Pomeron diagrams, which should provide a “natural unitarization”. Nevertheless we consider our approach as a necessary first step towards a consistent formulation of multiple scattering theory in nuclear (including hadron-hadron) collisions at very high energies.
Let us explain our “unitarization” in the following. We define
$$g(z)=\frac{\mathrm{\Gamma }(1+\alpha _{\mathrm{remn}})}{\mathrm{\Gamma }(1+\alpha _{\mathrm{remn}}+z)}$$
(3.25)
such that
$$g(r_1\stackrel{~}{\beta }_1+\mathrm{}+r_N\stackrel{~}{\beta }_N)$$
(3.26)
is the factor causing unitarity problems. This expression should be of the form
$$(\mathrm{})^{r_1}\mathrm{}(\mathrm{})^{r_N},$$
which would make $`\mathrm{\Phi }_{AB}`$ a well behaved exponential function. In order to achieve this, the function $`g`$ should be an exponential. So we replace $`g(z)`$ by $`g_\mathrm{e}(z),`$ where the latter function is defined as
$$g_\mathrm{e}(z)=e^{ϵ_\mathrm{e}z},$$
(3.27)
where the parameter $`ϵ_\mathrm{e}`$ should be chosen such that $`g(z)`$ is well approximated for values of $`z`$ between (say) $`0`$ and $`0.5`$ (see fig. 3.6). The index “e” refers to “exponentiation”.
So, we replace the factor $`g(r_1\stackrel{~}{\beta }_1+\mathrm{}+r_N\stackrel{~}{\beta }_N)`$ by
$`g_\mathrm{e}(r_1\stackrel{~}{\beta }_1+\mathrm{}+r_N\stackrel{~}{\beta }_N)`$ $`=`$ $`e^{ϵ_\mathrm{e}(r_1\stackrel{~}{\beta }_1+\mathrm{}+r_N\stackrel{~}{\beta }_N)}`$ (3.28)
$`=`$ $`\left(e^{ϵ_\mathrm{e}\stackrel{~}{\beta }_1}\right)^{r_1}\mathrm{}\left(e^{ϵ_\mathrm{e}\stackrel{~}{\beta }_N}\right)^{r_N},`$ (3.29)
and obtain correspondingly instead of $`\mathrm{\Phi }_{pp}`$
$`\mathrm{\Phi }_{\mathrm{e}_{pp}}(x^+,x^{},s,b)`$ $`=`$ $`x^{\alpha _{\mathrm{remn}}}{\displaystyle \underset{r_1=0}{\overset{\mathrm{}}{}}}\mathrm{}{\displaystyle \underset{r_N=0}{\overset{\mathrm{}}{}}}\left(e^{2ϵ_\mathrm{e}\stackrel{~}{\beta }_1}\right)^{r_1}\mathrm{}\left(e^{2ϵ_\mathrm{e}\stackrel{~}{\beta }_N}\right)^{r_N}`$ (3.30)
$`\times `$ $`{\displaystyle \frac{(\alpha _1x^{\stackrel{~}{\beta }_1}\mathrm{\Gamma }^2(\stackrel{~}{\beta }_1))^{r_1}}{r_1!}}\mathrm{}{\displaystyle \frac{(\alpha _Nx^{\stackrel{~}{\beta }_N}\mathrm{\Gamma }^2(\stackrel{~}{\beta }_N))^{t_N}}{r_N!}}.`$
Now the sums can be performed and we get
$`\mathrm{\Phi }_{\mathrm{e}_{pp}}(x^+,x^{},s,b)`$ $`=`$ $`x^{\alpha _{\mathrm{remn}}}{\displaystyle \underset{i=1}{\overset{N}{}}}\mathrm{exp}\left\{\alpha _ix^{\stackrel{~}{\beta }_i}\mathrm{\Gamma }^2(\stackrel{~}{\beta }_i)e^{2ϵ_\mathrm{e}\stackrel{~}{\beta }_i}\right\}`$ (3.31)
which may be written as
$$\mathrm{\Phi }_{\mathrm{e}_{pp}}(x^+,x^{},s,b)=(x^+x^{})^{\alpha _{\mathrm{remn}}}\mathrm{exp}\left\{\stackrel{~}{G}(x^+x^{},s,b)\right\}$$
(3.32)
with
$$\stackrel{~}{G}(x,s,b)=\underset{i=1}{\overset{N}{}}\stackrel{~}{\alpha }_ix^{\stackrel{~}{\beta }_i}$$
(3.33)
with
$`\stackrel{~}{\alpha }_i`$ $`=`$ $`\alpha _i\mathrm{\Gamma }^2(\stackrel{~}{\beta }_i)e^{2ϵ_\mathrm{e}\stackrel{~}{\beta }_i}`$ (3.34)
$`\stackrel{~}{\beta }_i`$ $`=`$ $`\beta _i+1,`$ (3.35)
where $`\alpha _i`$ and $`\beta _i`$ are given as
$`\alpha _i`$ $`=`$ $`\left(\alpha _{D_i}+\alpha _{D_i}^{}\right)s^{\left(\beta _{D_i}+\gamma _{D_i}b^2\right)}e^{\frac{b^2}{\delta _{D_i}}},`$ (3.36)
$`\beta _i`$ $`=`$ $`\beta _{D_i}+\beta _{D_i}^{}+\gamma _{D_i}b^2\alpha _{\mathrm{part}},`$ (3.37)
with $`\alpha _{D_i}^{}0`$ and $`\beta _{D_i}^{}0`$ only if $`\alpha _{D_i}0`$.
We are not yet done. We modified $`\mathrm{\Phi }_{pp}`$ such that the new function $`\mathrm{\Phi }_{\mathrm{e}_{pp}}`$ is surely positive. But what happened to our unitarity equation? If we replace $`\mathrm{\Phi }_{pp}`$ by $`\mathrm{\Phi }_{\mathrm{e}_{pp}}`$, we obtain
$$\overline{\gamma }_{pp}+\gamma _{pp}=\underset{m=0}{\overset{\mathrm{}}{}}𝑑X^+𝑑X^{}\mathrm{\Omega }_{\mathrm{e}_{pp}}^{(s,b)}(m,X^+,X^{}),$$
(3.38)
with
$$\mathrm{\Omega }_{\mathrm{e}_{pp}}^{(s,b)}(m,X^+,X^{})=\left\{\frac{1}{m!}\underset{\mu =1}{\overset{m}{}}G(s,x_\mu ^+,x_\mu ^{},b)\right\}\mathrm{\Phi }_{\mathrm{e}_{pp}}(x^+,x^{},s,b),$$
(3.39)
where $`x^+`$ and $`x^{}`$ refer to the remnant light cone momenta. Since $`\mathrm{\Phi }_{\mathrm{e}_{pp}}`$ is always bigger than $`\mathrm{\Phi }_{pp}`$ for small values of $`b`$, the sum $`\gamma _{pp}+\overline{\gamma }_{pp}`$ is bigger than one, so the unitarity equation does not hold any more. This is quite natural, since we modified the virtual emissions without caring about the real ones. In order to account for this, we define
$$Z(s,b)=\underset{m=0}{\overset{\mathrm{}}{}}𝑑X^+𝑑X^{}\left\{\frac{1}{m!}\underset{\mu =1}{\overset{m}{}}G(s,x_\mu ^+,x_\mu ^{},b)\right\}\mathrm{\Phi }_{\mathrm{e}_{pp}}(x^+,x^{},s,b),$$
(3.40)
which is equal to one in the exact case, but which is different from one if we use $`\mathrm{\Phi }_{\mathrm{e}_{pp}}`$ instead of $`\mathrm{\Phi }_{pp}`$. In order to recover the unitarity equation, we have to “renormalize” $`\mathrm{\Phi }_{\mathrm{e}_{pp}}`$, and we define therefore the “unitarized” virtual emission function $`\mathrm{\Phi }_{\mathrm{u}_{pp}}`$ via
$$\mathrm{\Phi }_{\mathrm{u}_{pp}}(x^+,x^{},s,b)=\frac{\mathrm{\Phi }_{\mathrm{e}_{pp}}(x^+,x^{},s,b)}{Z(s,b)}.$$
(3.41)
Now, the unitarity equation holds,
$$\overline{\gamma }_{pp}+\gamma _{pp}=\underset{m=0}{\overset{\mathrm{}}{}}𝑑X^+𝑑X^{}\mathrm{\Omega }_{\mathrm{u}_{pp}}^{(s,b)}(m,X^+,X^{})=1,$$
(3.42)
with
$$\mathrm{\Omega }_{\mathrm{u}_{pp}}^{(s,b)}(m,X^+,X^{})=\left\{\frac{1}{m!}\underset{\mu =1}{\overset{m}{}}G(s,x_\mu ^+,x_\mu ^{},b)\right\}\mathrm{\Phi }_{\mathrm{u}_{pp}}(x^+,x^{},s,b)$$
(3.43)
being strictly positive, which allows finally the probability interpretation.
### 3.5 Properties of the Unitarized Theory
We are now going to investigate the consequences of our unitarization, in other words, how the results are affected by this modification. In fig. 3.7 we compare the exact and the exponentiated version of the virtual emission function ($`\mathrm{\Phi }_{pp}`$ and $`\mathrm{\Phi }_{\mathrm{e}_{pp}}`$) for a large value of the impact parameter ($`b=1.5`$ fm). The exponentiated result (dashed) is somewhat below the exact one (solid curve), but the difference is quite small.
The situation is somewhat different in case of zero impact parameter $`b`$. For small values of $`x`$ the two curves coincide more or less, however for $`x=1`$ the exponentiated result (dashed) is well above the exact one (solid curve). In particular, and this is most important, the dashed curve rests positive and in this sense corrects for the unphysical behavior (negative values) for the exact curve.
The behavior for $`x=1`$ for different values of $`b`$ is summarized in fig. 3.9, where we plot $`1\mathrm{\Phi }_{pp}(1,1,s,b)`$ as a function of $`b`$. We clearly observe that for large $`b`$ exact (solid) and exponentiated (dashed curve) result agree approximately, whereas for small values of $`b`$ they differ substantially, with the exponentiated version always staying below 1, as it should be. So the effect of our exponentiation is essentially to push the function below 1.
Next we calculate explicitly the normalization function $`Z(s,b)`$. We have
$`Z(s,b)`$ $`=`$ $`{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{\mu =1}{\overset{m}{}}dx_\mu ^+dx_\mu ^{}\frac{1}{m!}\underset{\mu =1}{\overset{m}{}}G(x_\mu ^+,x_\mu ^{},s,b)}`$ (3.44)
$`\times \mathrm{\Phi }_{\mathrm{e}_{pp}}(1{\displaystyle \underset{\nu =1}{\overset{m}{}}}x_\nu ^+,1{\displaystyle \underset{\nu =1}{\overset{m}{}}}x_\nu ^{}),`$
which may be written as
$$Z(s,b)=\mathrm{\Phi }_{\mathrm{e}_{pp}}(1,1,s,b)+𝑑z^+𝑑z^{}H(z^+,z^{})\mathrm{\Phi }_{\mathrm{e}_{pp}}(z^+,z^{},s,b),$$
(3.45)
with
$`H(z^+,z^{})`$ $`=`$ $`{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{\mu =1}{\overset{m}{}}dx_\mu ^+dx_\mu ^{}\frac{1}{m!}\underset{\mu =1}{\overset{m}{}}G(x_\mu ^+,x_\mu ^{},s,b)}`$ (3.46)
$`\times `$ $`\delta \left(1z^+{\displaystyle \underset{\mu =1}{\overset{m}{}}}x_\mu ^+\right)\delta \left(1z^{}{\displaystyle \underset{\mu =1}{\overset{m}{}}}x_\mu ^{}\right).`$
Using the analytical form of $`G`$, we obtain
$`H(z^+,z^{})`$ $`=\underset{r_1+\mathrm{}+r_K0}{\underset{}{{\displaystyle \underset{r_1=0}{\overset{\mathrm{}}{}}}\mathrm{}{\displaystyle \underset{r_N=0}{\overset{\mathrm{}}{}}}}}`$ $`{\displaystyle \frac{\left[(1z^+)(1z^{})\right]^{r_1\stackrel{~}{\beta }_1+\mathrm{}+r_N\stackrel{~}{\beta }_N1}}{\mathrm{\Gamma }(r_1\stackrel{~}{\beta }_1+\mathrm{}+r_N\stackrel{~}{\beta }_N)^2}}`$ (3.47)
$`\times `$ $`{\displaystyle \frac{(\alpha _1\mathrm{\Gamma }(\stackrel{~}{\beta }_1)^2)^{r_1}}{r_1!}}\mathrm{}{\displaystyle \frac{(\alpha _N\mathrm{\Gamma }(\stackrel{~}{\beta }_N)^2)^{r_N}}{r_N!}}`$
(see appendix D.2). This can be calculated, and after numerically doing the integration over $`z^+,z^{}`$, we obtain the normalization function $`Z(s,b)`$, as shown in fig. 3.10.
We observe, as expected, a value close to unity at large values of $`b`$, whereas for small impact parameter $`Z(s,b)`$ is bigger than one, since only at small values of $`b`$ the virtual emission function $`\mathrm{\Phi }_{pp}`$ has been changed substantially towards bigger values.
Knowing $`\mathrm{\Phi }_{\mathrm{e}_{pp}}`$ and $`Z`$, we are ready to calculate the unitarized emission functions $`\mathrm{\Phi }_{\mathrm{u}_{pp}}`$ which finally replaces $`\mathrm{\Phi }_{pp}`$ in all formulas for cross section calculations. The results are shown in fig. 3.11,
where we plot 1-$`\mathrm{\Phi }_{\mathrm{u}_{pp}}`$ together with $`1\mathrm{\Phi }_{\mathrm{e}_{pp}}`$ and $`1\mathrm{\Phi }_{pp}`$ for both $`x^+`$ and $`x^{}`$ being one. We observe that compared to $`1\mathrm{\Phi }_{\mathrm{e}_{pp}}`$ the function $`1\mathrm{\Phi }_{\mathrm{u}_{pp}}`$ is somewhat increased at small values of $`b`$ due to the fact that here $`Z(s,b)`$ is bigger than one, whereas for large impact parameters there is no difference.
Since the unitarity equation holds, we may integrate $`1\mathrm{\Phi }_{\mathrm{u}_{pp}}(1,1,s,b)`$ over impact parameter, to obtain the inelastic non-diffractive cross section,
$$\sigma _{\mathrm{inel}}(s)=d^2b\left\{1\mathrm{\Phi }_{\mathrm{u}_{pp}}(1,1,s,b)\right\},$$
(3.48)
the result being shown in fig. 3.12. Here the exact and the unitarized result (using $`\mathrm{\Phi }_{pp}`$ and $`\mathrm{\Phi }_{\mathrm{u}_{pp}}`$ respectively) are quite close due to the fact that one has a two-dimensional $`b`$-integration, and therefore the small values of $`b`$, where we observe the largest differences, do not contribute much to the integral.
We now turn to inclusive spectra. We consider the inclusive $`x`$-spectrum of Pomerons, $`dn_{\mathrm{Pom}}/dx`$, where $`x=x^+x^{}`$ is the squared mass of the Pomeron divided by $`s`$. In the exact theory, we may take advantage of the AGK cancelations, and obtain
$$\frac{dn_{\mathrm{Pom}}}{dx}(x,s,b)=_{+\mathrm{ln}\sqrt{x}}^{\mathrm{ln}\sqrt{x}}𝑑y\frac{dn_{\mathrm{Pom}}^{(1)}}{dx^+dx^{}}(x^+,x^{},s,b)|_{x^+=\sqrt{x}e^y,x^{}=\sqrt{x}e^y},$$
(3.49)
where $`dn_{\mathrm{Pom}}^{(1)}/dx^+dx^{}`$ is the corresponding inclusive distribution for one single elementary interaction, which is given in eq. (2.75). The $`y`$-integration can be easily performed numerically, and we obtain the results shown in fig. 3.13 as solid curves, the upper one for $`b=0`$ fm and the lower one for $`b=1.5`$ fm.
The calculation of the unitarized result is more involved, since now we cannot use the AGK cancelations any more. We have
$`{\displaystyle \frac{dn_{\mathrm{Pom}}}{dx^+dx^{}}}(x^+,x^{},s,b)`$ $`=`$ $`{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{\mu =1}{\overset{m}{}}dx_\mu ^+dx_\mu ^{}}`$ (3.50)
$`\times `$ $`{\displaystyle \frac{1}{m!}}{\displaystyle \underset{\mu =1}{\overset{m}{}}}\left\{G(x_\mu ^+,x_\mu ^{},s,b)\right\}\mathrm{\Phi }_{\mathrm{u}_{pp}}(1{\displaystyle \underset{\mu =1}{\overset{m}{}}}x_\mu ^+,1{\displaystyle \underset{\mu =1}{\overset{m}{}}}x_\mu ^{},s,b)`$
$`\times `$ $`{\displaystyle \underset{\mu ^{}=1}{\overset{m}{}}}\delta (x^+x_\mu ^{}^+)\delta (x^{}x_\mu ^{}^{}),`$
where we used the unitarized version of $`\mathrm{\Phi }_{pp}`$. We find
$`{\displaystyle \frac{dn_{\mathrm{Pom}}}{dx^+dx^{}}}(x^+,x^{},s,b)`$ $`=`$ $`G(x^+,x^{},s,b)`$
$`\times `$ $`[\mathrm{\Phi }_{\mathrm{u}_{pp}}(1x^+,1x^{},s,b)`$
$`+{\displaystyle }dz^+dz^{}H(z^++x^+,z^{}+x^{})\mathrm{\Phi }_{\mathrm{u}_{pp}}(z^+,z^{},s,b)],`$
where $`H`$ is defined in eq. (3.46), with the final result given in eq. (3.47). The integration over $`z^+,z^{}`$ can now be done numerically. Expressing $`x^+`$ and $`x^{}`$ via $`x`$ and $`y`$ and integrating over $`y`$, we finally obtain $`dn_{\mathrm{Pom}}/dx`$, as shown in fig. 3.13 (dashed curves). In fig. 3.14
we show the b-averaged inclusive spectra, which are given as
$$\frac{1}{\sigma _{\mathrm{inel}}}d^2b\frac{dn_{\mathrm{Pom}}}{dx}(x,s,b)$$
(3.52)
for both, the exact and the unitarized version.
### 3.6 Comparison with the Conventional Approach
At this point it is noteworthy to compare our approach with the conventional one . There one neglects the energy conservation effects in the cross section calculation and sums up virtual Pomeron emissions, each one taken with the initial energy of the interaction $`s`$. We can recover the conventional approach by simply considering independent (planar) emission of all the Pomerons, neglecting energy-momentum sharing between them. In the cross section formulas (2.65-3.58) this amounts to perform formally the convolutions of the Pomeron eikonals $`G(x^+,x^{},s,b)`$ with the remnant functions $`F_{\mathrm{remn}}(x^\pm )`$ for all the Pomerons. In case of proton-proton scattering, we then get
$$\mathrm{\Phi }_{\mathrm{conv}_{pp}}(x^+,x^{},s,b)=e^{\chi (s,b)},$$
(3.53)
with
$$\chi (s,b)=𝑑x^+𝑑x^{}G(x^+,x^{},s,b)F_{\mathrm{remn}}(x^+)F_{\mathrm{remn}}(x^{}),$$
(3.54)
where $`\mathrm{\Phi }`$ does not depend on $`x^+`$ and $`x^{}`$ anymore. We obtain a unitarity relation of the form
$$\underset{m=0}{\overset{\mathrm{}}{}}\mathrm{\Omega }_m=1,$$
(3.55)
with
$$\mathrm{\Omega }_m=\frac{\left(\chi (s,b)\right)^m}{m!}e^{\chi (s,b)}$$
(3.56)
representing the probability of having $`m`$ cut Pomerons (Pomeron multiplicity distribution). So in the traditional case, the Pomeron multiplicity distribution is a Poissonian with the mean value given by $`\chi (s,b)`$. As already mentioned above, that approach is not self-consistent as the AGK rules are assumed to hold when calculating interaction cross sections but are violated at the particle production generation. This inconsistency was already mentioned in , where the necessity to develop the correct, Feynman diagram-based scheme, was first argued.
The exact procedure is based on the summation over virtual emissions with the energy-momentum conservation taken into account. This results in the formula (3.7) or, using our parametrization, (3.15) for $`\mathrm{\Phi }_{pp}(x^+,x^{},s,b)`$, explicitly dependent on the momentum, left after cut Pomerons emission, and in the formula
$$\sigma _{\mathrm{inel}}(s)=d^2b\left\{1\mathrm{\Phi }_{pp}(1,1,s,b)\right\}$$
(3.57)
for the inelastic cross section; AGK rules are exactly fulfilled both for the cross sections and for the particle production. But with the interaction energy increasing the approach starts to violate the unitarity and is no longer self-consistent.
The ”unitarized” procedure, which amounts to replacing $`\mathrm{\Phi }_{pp}`$ by $`\mathrm{\Phi }_{\mathrm{u}_{pp}}`$, allows to avoid the unitarity problems. The expressions for cross sections and for inclusive spectra are consistent with each other and with the particle generation procedure. The latter one assures the AGK cancelations validity in the region, where unitarity problems do not appear yet (not too high energies or large impact parameters).
In order to see the effect of energy conservation we calculate $`\chi `$ as given in eq. (3.54) with the same parameters as we use in our approach for different values of $`b`$, and we show the corresponding Pomeron multiplicity distribution in fig. 3.15 as dashed lines. We compare this traditional approach with our full simulation, where energy conservation is treated properly (solid lines in the figures). One observes a huge difference between the two approaches. So energy conservation makes the Pomeron multiplicity distributions much narrower, in other words, the mean number of Pomerons is substantially reduced. The reason is that due to energy conservation the phase space of light cone momenta of the Pomeron ends is considerably reduced.
Of course, in the traditional approach one chooses different parameters in order to obtain reasonable values for the Pomeron numbers in order to reproduce the experimental cross sections. But this only “simulates” in some sense the phase space reduction due to energy conservation in an uncontrolled way.
We conclude that considering energy conservation properly in the cross section formulas has an enormous effect and cannot be neglected.
### 3.7 Unitarization for Nucleus-Nucleus Scattering
In this section, we discuss the unitarization scheme for nucleus-nucleus scattering. The sum over virtual emissions is defined as
$`\mathrm{\Phi }_{AB}(X^+,X^{},s,b)`$ $`=`$ $`{\displaystyle \underset{l_1}{}}\mathrm{}{\displaystyle \underset{l_{AB}}{}}{\displaystyle \underset{k=1}{\overset{AB}{}}\left\{\underset{\lambda =1}{\overset{l_k}{}}d\stackrel{~}{x}_{k,\lambda }^+d\stackrel{~}{x}_{k,\lambda }^{}\right\}\underset{k=1}{\overset{AB}{}}\left\{\frac{1}{l_k!}\underset{\lambda =1}{\overset{l_k}{}}G(\stackrel{~}{x}_{k,\lambda }^+,\stackrel{~}{x}_{k,\lambda }^{},s,b)\right\}}`$ (3.58)
$`\times `$ $`{\displaystyle \underset{i=1}{\overset{A}{}}}F_{\mathrm{remn}}\left(x_i^+{\displaystyle \underset{\pi (k)=i}{}}\stackrel{~}{x}_{k,\lambda }^+\right){\displaystyle \underset{j=1}{\overset{B}{}}}F_{\mathrm{remn}}\left(x_j^{}{\displaystyle \underset{\tau (k)=j}{}}\stackrel{~}{x}_{k,\lambda }^{}\right).`$
where $`X^+=\left\{x_1^+\mathrm{}x_A^+\right\}`$, $`X^{}=\left\{x_1^{}\mathrm{}x_B^{}\right\}`$ and $`\pi (k)`$ and $`\tau (k)`$ represent the projectile or target nucleon linked to pair $`k`$. This calculation is very close to the calculation for proton-proton scattering. Using the expression eq. (3.10) of $`G(\stackrel{~}{x}_{k,\lambda }^+,\stackrel{~}{x}_{k,\lambda }^{},s,b)`$, the definition eq. (3.8) of $`F_{\mathrm{remn}}(x)`$, one finally finds
$`\mathrm{\Phi }_{AB}(X^+,X^{},s,b)=`$
$`{\displaystyle \underset{r_{1,1}\mathrm{}r_{N,1}}{}}\mathrm{}{\displaystyle \underset{r_{1,AB}\mathrm{}r_{N,AB}}{}}{\displaystyle \underset{k=1}{\overset{AB}{}}}{\displaystyle \frac{\left(\alpha _1\right)^{r_{1,k}}}{r_{1,k}!}}\mathrm{}{\displaystyle \frac{\left(\alpha _N\right)^{r_{N,k}}}{r_{N,k}!}}`$ (3.59)
$`{\displaystyle \underset{i=1}{\overset{A}{}}}(x_i^+)^{\alpha _{\mathrm{remn}}}{\displaystyle \underset{\pi (k)=i}{}}\left(\mathrm{\Gamma }(\stackrel{~}{\beta }_1)(x_i^+)^{\stackrel{~}{\beta }_1}\right)^{r_{1,k}}\mathrm{}\left(\mathrm{\Gamma }(\stackrel{~}{\beta }_N)(x_i^+)^{\stackrel{~}{\beta }_N}\right)^{r_{N,k}}g\left({\displaystyle \underset{\pi (k)=i}{}}r_{1,k}\stackrel{~}{\beta }_1+\mathrm{}+r_{N,k}\stackrel{~}{\beta }_N\right)`$
$`{\displaystyle \underset{j=1}{\overset{B}{}}}(x_j^{})^{\alpha _{\mathrm{remn}}}{\displaystyle \underset{\tau (k)=j}{}}\left(\mathrm{\Gamma }(\stackrel{~}{\beta }_1)(x_j^{})^{\stackrel{~}{\beta }_1}\right)^{r_{1,k}}\mathrm{}\left(\mathrm{\Gamma }(\stackrel{~}{\beta }_N)(x_j^{})^{\stackrel{~}{\beta }_N}\right)^{r_{N,k}}g\left({\displaystyle \underset{\tau (k)=j}{}}r_{1,k}\stackrel{~}{\beta }_1+\mathrm{}+r_{N,k}\stackrel{~}{\beta }_N\right)`$
(see appendix D.3), where the function $`g(z)`$ is defined in eq. (3.25), and the parameters $`\alpha _i`$ and $`\stackrel{~}{\beta }_i`$ are the same ones as for proton-proton scattering.
In case of nucleus-nucleus scattering, we use the same unitarization prescription as already applied to proton-proton scattering. The first step amounts to replace the function $`g(z)`$, which appears in the final expression of $`\mathrm{\Phi }_{AB}`$, by the exponential form $`g_\mathrm{e}(z)`$. This allows to perform the sums in eq. (D.70), and we obtain
$`\mathrm{\Phi }_{\mathrm{e}_{AB}}(X^+,X^{},s,b)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{A}{}}}(x_i^+)^{\alpha _{\mathrm{remn}}}{\displaystyle \underset{j=1}{\overset{B}{}}}(x_j^{})^{\alpha _{\mathrm{remn}}}{\displaystyle \underset{k=1}{\overset{AB}{}}}e^{2\stackrel{~}{G}\left(x_{\pi (k)}^+.x_{\tau (k)}^{}\right)}`$ (3.60)
(see appendix D.3), where $`\stackrel{~}{G}(x)`$ is defined in eq. (3.33). Having modified $`\mathrm{\Phi }`$, the unitarity equation
$$\underset{m}{}𝑑X^+𝑑X^{}\mathrm{\Omega }_{\mathrm{e}_{AB}}^{(s,b)}(m,X^+,X^{})=1$$
(3.61)
does not hold any more, since $`\mathrm{\Omega }_\mathrm{e}`$ depends on $`\mathrm{\Phi }_\mathrm{e}`$, and only the exact $`\mathrm{\Phi }`$ assures a correct unitarity relation. So as in proton-proton scattering, we need a second step, which amounts to renormalizing $`\mathrm{\Phi }_\mathrm{e}`$. So we introduce a normalization factor
$$Z_{AB}(s,b)=\underset{m}{}𝑑X^+𝑑X^{}\mathrm{\Omega }_{\mathrm{e}_{AB}}^{(s,b)}(m,X^+,X^{}),$$
(3.62)
with $`\mathrm{\Omega }_{\mathrm{e}_{AB}}`$ defined in the same way as $`\mathrm{\Omega }`$ but with $`\mathrm{\Phi }_{AB}`$ replaced by $`\mathrm{\Phi }_{\mathrm{e}_{AB}}`$, which allows to define the unitarized $`\mathrm{\Phi }_{\mathrm{u}_{AB}}`$ function as
$$\mathrm{\Phi }_{\mathrm{u}_{AB}}(x^+,x^{},s,b)=\frac{\mathrm{\Phi }_{\mathrm{e}_{AB}}(x^+,x^{},s,b)}{Z_{AB}(s,b)}.$$
(3.63)
In this way we recover the unitarity relation,
$$\underset{m}{}𝑑X^+𝑑X^{}\mathrm{\Omega }_{\mathrm{u}_{AB}}^{(s,b)}(m,X^+,X^{})=1,$$
(3.64)
with
$$\mathrm{\Omega }_{\mathrm{u}_{AB}}^{(s,b)}(m,X^+,X^{})=\underset{k=1}{\overset{AB}{}}\left\{\frac{1}{m_k!}\underset{\mu =1}{\overset{m_k}{}}G(s,x_{k,\mu }^+,x_{k,\mu }^{},b)\right\}\mathrm{\Phi }_{\mathrm{u}_{AB}}(x^+,x^{},s,b),$$
(3.65)
and $`\mathrm{\Omega }_\mathrm{u}^{(s,b)}(m,X^+,X^+)`$ may be interpreted as probability distribution for configurations $`(m,X^+,X^+)`$.
### 3.8 Profile Functions in Nucleus-Nucleus Scattering
In case of nucleus-nucleus scattering, the conventional approach represents a “Glauber-type model”, where nucleus-nucleus scattering may be considered as a sequence of nucleon-nucleon scatterings with constant cross sections; the nucleons move through the other nucleus along straight line trajectories. In order to test this picture, we consider all pairs of nucleons, which due to their distributions inside the nuclei provide a more or less flat $`b`$-distribution.
We then simply count, for a given $`b`$-bin, the number of interacting pairs and then divide by the number of pairs in the corresponding bin. The resulting distribution, which we call nucleon-nucleon profile function for nucleus-nucleus scattering, represents the probability density of an interaction of a pair of nucleons at given impact parameter. This may be compared with the proton-proton profile function $`1\mathrm{\Phi }_{\mathrm{u}_{pp}}(1,1,s,b)`$. In the Glauber model, these two distributions coincide. As demonstrated in fig.3.16 for S+S scattering this is absolutely not the case. The profile function in case of S+S scattering is considerably reduced as compared to the proton-proton one. Since integrating the proton-proton profile function represents the inelastic cross section, one may also define the corresponding integral in nucleus-nucleus scattering as “individual nucleon-nucleon cross section”. So we conclude that this cross section is smaller than the proton-proton cross section. This is due to the energy conservation, which reduces the number of Pomerons connected to any nucleon from the projectile and the target and finally affects also the “individual cross section”.
### 3.9 Inclusive Cross Sections in Nucleus-Nucleus Scattering
We have shown in the preceding chapter that in the “bare” theory AGK cancelations apply perfectly, which means that nucleus-nucleus inclusive cross sections are just $`AB`$ times the proton-proton ones,
$$\frac{d\sigma _{\mathrm{incl}}^{AB}}{dq}(q,s,b)=AB\frac{d\sigma _{\mathrm{incl}}^{pp}}{dq}(q,s,b).$$
(3.66)
In the unitarized theory, the results is somewhat different. Unfortunately, we cannot calculate cross sections analytically any more, so we perform a numerical calculation using the Markov chain techniques explained later. In order to investigate the deviation from exact AGK cancelations, we calculate
the inclusive nucleus-nucleus cross section for Pomeron production (being the basic inclusive cross section), divided by $`AB`$,
$$\frac{1}{AB}\frac{d\sigma _{\mathrm{Pom}}^{AB}}{dx}(x,s,b),$$
(3.67)
and compare the result with the corresponding proton-proton cross section, see figs. 3.17 and 3.18. For large and for small values of $`x`$, we still observe AGK cancelations (the two curves agree), but for intermediate values of $`x`$, the AGK cancelations are violated, the nucleus-nucleus cross section is smaller than $`AB`$ times the nucleon-nucleon one. The effect is, however, relatively moderate. If one writes the proton-nucleus cross section as
$$\frac{d\sigma _{\mathrm{Pom}}^{pA}}{dx}(x,s,b)=A^{\alpha (x)}\frac{d\sigma _{\mathrm{Pom}}^{pp}}{dx},$$
(3.68)
we obtain for $`\alpha (x)`$ values between $`0.85`$ and $`1`$. So one may summarize that “AGK cancelations are violated, but not too strongly”.
## Chapter 4 Markov Chain Techniques
In this chapter we discuss how to deal with the multidimensional probability distribution $`\mathrm{\Omega }_{AB}(K)`$ with $`K=\{m,X^+,X^{}\}`$, where the vector $`m`$ characterizes the type of interaction of each pair of nucleons (the number of elementary interactions per pair), and the matrices $`X^+`$, $`X^{}`$ contain the light cone momenta of all Pomerons (energy sharing between the Pomerons).
### 4.1 Probability Distributions for Configurations
In this section we essentially repeat the basic formulas of the preceding chapters which allowed us to derive probability distributions for interaction configurations in a consistent way within an effective theory based on Feynman diagrams.
Our basic formula for the inelastic cross section for a nucleus-nucleus collision (which includes also as a particular case proton-proton scattering) could be written in the following form
$$\sigma _{\mathrm{inel}}(s)=d^2b_0𝑑T_{AB}\gamma _{AB}(s,b_0,b_1\mathrm{}b_{AB}),$$
(4.1)
where $`dT_{AB}`$ represents the integration over the transverse coordinates $`b_i^A`$ and $`b_j^B`$ of projectile and target nucleons, $`b_0`$ is the impact parameter between the two nuclei, and $`b_k=|\stackrel{}{b}+\stackrel{}{b}_{\pi (k)}^A\stackrel{}{b}_{\tau (k)}^B|`$ is the transverse distance between the nucleons of $`k^{th}`$ pair. Using the compact notation
$$b=\{b_k\},m=\{m_k\},X^+=\left\{x_{k,\mu }^+\right\},X^{}=\left\{x_{k,\mu }^{}\right\},$$
(4.2)
the function $`\gamma _{AB}`$ is given as
$$\gamma _{AB}(s,b)=\underset{m}{}(1\delta _{0m})𝑑X^+𝑑X^{}\mathrm{\Omega }_{\mathrm{u}_{AB}}^{(s,b)}(m,X^+,X^{}),$$
(4.3)
which represents all diagrams with at least one cut Pomeron. One may define a corresponding quantity $`\overline{\gamma }_{AB}`$, which represents the configuration with exactly zero cut Pomerons. The latter one can be obtained from (4.3) by exchanging $`1\delta _{0m}`$ by $`\delta _{0m}`$, which leads to
$$\overline{\gamma }_{AB}(s,b)=\mathrm{\Omega }_{\mathrm{u}_{AB}}^{(s,b)}(0,0,0).$$
(4.4)
The expression for $`\mathrm{\Omega }`$ is given as
$$\mathrm{\Omega }_{\mathrm{u}_{AB}}^{(s,b)}(m,X^+,X^{})=\underset{k=1}{\overset{AB}{}}\left\{\frac{1}{m_k!}\underset{\mu =1}{\overset{m_k}{}}G(x_{k,\mu }^+,x_{k,\mu }^{},s,b_k)\right\}\mathrm{\Phi }_{\mathrm{u}_{\mathrm{AB}}}(x^+,x^{},s,b),$$
(4.5)
with
$`\mathrm{\Phi }_{\mathrm{u}_{AB}}(x^+,x^{},s,b)`$ $`=`$ $`{\displaystyle \frac{1}{Z_{AB}}}{\displaystyle \underset{k=1}{\overset{AB}{}}}\mathrm{exp}\left(\stackrel{~}{G}(x_{\pi (k)}^+x_{\tau (k)}^{},s,b_k)\right)`$
$`\times `$ $`{\displaystyle \underset{i=1}{\overset{A}{}}}(x_i^+)^{\alpha _{\mathrm{remn}}}\mathrm{\Theta }(x_i^+)\mathrm{\Theta }(1x_i^+){\displaystyle \underset{j=1}{\overset{B}{}}}(x_j^{})^{\alpha _{\mathrm{remn}}}\mathrm{\Theta }(x_j^{})\mathrm{\Theta }(1x_j^{}).`$
The arguments of $`\mathrm{\Phi }_{\mathrm{u}_{AB}}`$ are the momentum fractions of projectile and target remnants,
$$x_i^+=1\underset{\pi (k)=i}{}x_{k,\mu ,}^+x_j^{}=1\underset{\tau (k)=j}{}x_{k,\mu }^{},$$
(4.7)
where $`\pi (k)`$ and $`\tau (k)`$ point to the remnants linked to the $`k^{\mathrm{th}}`$ interaction.
In the following, we perform the analysis for given $`s`$ and $`b=(b_0,b_1,\mathrm{},b_{AB})`$, so we do not write these variables explicitly. In addition, we always refer to the unitarized functions, so we will also suppress the subscript “u”. Furthermore, we suppress the index $`AB`$.
Crucial for our applications is the probability conservation constraint
$$\gamma +\overline{\gamma }=1,$$
(4.8)
which may be written more explicitly as
$$\underset{m}{}𝑑X^+𝑑X^{}\mathrm{\Omega }(m,X^+,X^{})=1.$$
(4.9)
This allows us to interpret $`\mathrm{\Omega }(m,X^+,X^{})`$ as the probability distribution for a configuration $`(m,X^+,X^{})`$. For any given configuration the function $`\mathrm{\Omega }`$ can be easily calculated using the techniques developed in the chapter 2. The difficulty with the Monte Carlo generation of interaction configurations arises from the fact that the configuration space is huge and rather nontrivial in the sense that it cannot be written as a product of single Pomeron contributions. We are going to explain in the next sections, how we deal with this problem.
### 4.2 The Interaction Matrix
Since $`\mathrm{\Omega }(m,X^+,X^{})`$ is a high-dimensional and nontrivial probability distribution, the only way to proceed amounts to employing dynamical Monte Carlo methods, well known in statistical and solid state physics.
We first need to choose the appropriate framework for our analysis. So we translate our problem into the language of spin systems : we number all nucleon pairs as 1, 2, …, $`AB`$ and for each nucleon pair $`k`$ the possible elementary interactions as 1,2, …, $`m_k`$ Let $`m_{\mathrm{max}}`$ be the maximum number of elementary interactions per nucleon pair one may imagine. We now consider a two dimensional lattice with $`AB`$ lines and $`m_{\mathrm{max}}`$ columns, see fig. 4.1. Lattice sites are occupied $`\left(=1\right)`$ or empty $`\left(=0\right)`$, representing an elementary interaction (1) or the case of no interaction (0), for the $`k^{th}`$ pair.
In order to represent $`m_k`$ elementary interactions for the pair $`k`$, we need $`m_k`$ occupied cells (1’s) in the $`k^{th}`$ line. A line containing only empty cells (0’s) represents a pair without interaction. Any possible interaction may be represented by this “interaction matrix” $`M`$ with elements
$$m_{k\mu }\{0,1\}.$$
(4.10)
Such an “interaction configuration” is exactly equivalent to a spin configuration of the Ising model. Unfortunately the situation is somewhat more complicated in case of nuclear collisions: we need to consider the energy available for each elementary interaction, represented via the momentum fractions $`x_{k\mu }^+`$ and $`x_{k\mu }^{}`$. So we have a “generalized” matrix $`K`$,
$$K=(M,X^+,X^{}),$$
(4.11)
representing an interaction configuration, with elements
$$K_{k\mu }=(m_{k\mu ,}x_{k\mu }^+,x_{k\mu }^{}).$$
(4.12)
It is important to note that a number of matrices $`M`$ represents one and the same vector $`m`$. In fact, $`m`$ is represented by all the matrices $`M`$ with
$$\underset{\mu =1}{\overset{m_{\mathrm{max}}}{}}m_{k\mu }=m_k,$$
(4.13)
for each $`k`$. Since all the corresponding configurations $`(M,X^+,X^{})`$ should have the same weight, and since there are
$$c=\underset{k=1}{\overset{AB}{}}\frac{m_{\mathrm{max}}!}{m_k!(m_{\mathrm{max}}m_k)!}$$
(4.14)
configurations $`(M,X^+,X^{})`$ representing the same configuration $`(m,X^+,X^{})`$, the weight for the former is $`c^1`$ times the weight for the latter, so we obtain the following probability distribution for $`K=(M,X^+,X^{})`$:
$$\mathrm{\Omega }(K)=\underset{k=1}{\overset{AB}{}}\left\{\frac{(m_{\mathrm{max}}m_k)!}{m_{\mathrm{max}}!}\underset{\mu =1}{\overset{m_k}{}}G(s,x_{k,\mu }^+,x_{k,\mu }^{},b)\right\}\mathrm{\Phi }_{\mathrm{u}_{AB}}(x^+,x^{},s,b),$$
(4.15)
or, using the expression for $`\mathrm{\Phi }_{\mathrm{u}_{AB}}`$,
$`\mathrm{\Omega }(K)`$ $`=`$ $`{\displaystyle \frac{1}{Z_{AB}}}{\displaystyle \underset{k=1}{\overset{AB}{}}}\left\{{\displaystyle \frac{(m_{\mathrm{max}}m_k)!}{m_{\mathrm{max}}!}}{\displaystyle \underset{\mu =1}{\overset{m_k}{}}}\left\{G(s,x_{k,\mu }^+,x_{k,\mu }^{},b_k)\right\}\mathrm{exp}\left(\stackrel{~}{G}(x_{\pi (k)}^+x_{\tau (k)}^{},s,b_k)\right)\right\}`$ (4.16)
$`\times {\displaystyle \underset{i=1}{\overset{A}{}}}(x_i^+)^{\alpha _{\mathrm{remn}}}\mathrm{\Theta }(x_i^+)\mathrm{\Theta }(1x_i^+){\displaystyle \underset{j=1}{\overset{B}{}}}(x_j^{})^{\alpha _{\mathrm{remn}}}\mathrm{\Theta }(x_j^{})\mathrm{\Theta }(1x_j^{}).`$
The probability conservation now reads
$$\underset{K}{}\mathrm{\Omega }(K)=\underset{M}{}𝑑X^+𝑑X^{}\mathrm{\Omega }(M,X^+,X^{})=1.$$
(4.17)
In the following, we shall deal with the “interaction matrix” $`K`$, and the probability distribution $`\mathrm{\Omega }(K)`$.
### 4.3 The Markov Chain Method
In order to generate $`K`$ according to the given distribution $`\mathrm{\Omega }\left(K\right)`$, defined earlier, we construct a Markov chain
$$K^{\left(0\right)},K^{\left(1\right)},K^{\left(2\right)},\mathrm{}K^{\left(t_{\mathrm{max}}\right)}$$
(4.18)
such that the final configurations $`K^{\left(t_{\mathrm{max}}\right)}`$ are distributed according to the probability distribution $`\mathrm{\Omega }\left(K\right)`$, if possible for a $`t_{\mathrm{max}}`$ not too large!
Let us discuss how to obtain a new configuration $`K^{(t+1)}=L`$ from a given configuration $`K^{(t)}=K`$. We use Metropolis’ Ansatz for the transition probability
$$p(K,L)=prob\left(K^{\left(t+1\right)}=L|K^{\left(t\right)}=K\right)$$
(4.19)
as a product of a proposition matrix $`w(K,L)`$ and an acceptance matrix $`u(K,L)`$:
$$p(K,L)=\{\begin{array}{ccc}w(K,L)u(K,L)\hfill & \mathrm{if}\hfill & LK\hfill \\ w(K,K)+_{LK}w(K,L)\{1u(K,L)\}\hfill & \mathrm{if}\hfill & L=K\hfill \end{array},$$
(4.20)
where we use
$$u(K,L)=\mathrm{min}(\frac{\mathrm{\Omega }(L)}{\mathrm{\Omega }(K)}\frac{w(L,K)}{w(K,L)},1),$$
(4.21)
in order to assure detailed balance. We are free to choose $`w(K,L)`$, but of course, for practical reasons, we want to minimize the autocorrelation time, which requires a careful definition of $`w`$. An efficient procedure requires $`u(K,L)`$ to be not too small (to avoid too many rejections), so an ideal choice would be $`w(K,L)=\mathrm{\Omega }\left(L\right)`$. This is of course not possible, but we choose $`w(K,L)`$ to be a “reasonable” approximation to $`\mathrm{\Omega }(L)`$ if $`K`$ and $`L`$ are reasonably close, otherwise $`w`$ should be zero. So we define
$$w(K,L)=\{\begin{array}{ccc}\mathrm{\Omega }_0(L)& \mathrm{if}& d(K,L)1\\ 0& \mathrm{otherwise}& \end{array},$$
(4.22)
where $`d(K,L)`$ is the number of lattice sites being different in $`L`$ compared to $`K,`$ and where $`\mathrm{\Omega }^{(0)}`$ is defined by the same formulas as $`\mathrm{\Omega }`$ with one exception : $`\mathrm{\Phi }_{\mathrm{u}_{AB}}`$ is replaced by $`1`$. So we get
$$\mathrm{\Omega }_0(L)\underset{k=1}{\overset{AB}{}}\left\{(m_{\mathrm{max}}m_k)!\underset{\mu =1}{\overset{m_k}{}}G(x_{k,\mu }^+,x_{k,\mu }^{},s,b_k)\right\}.$$
(4.23)
The above definition of $`w(K,L)`$ may be realized by the following algorithm:
* choose randomly a lattice site $`(k,\mu )`$,
* propose a new matrix element $`(m_{k\mu },x_{k\mu }^+,x_{k,\mu }^{})`$ according to the probability distribution $`\rho (m_{k\mu },x_{k\mu }^+,x_{k,\mu }^{})`$,
where we are going to derive the form of $`\rho `$ in the following. From eq. (4.23), we know that $`\rho `$ should be of the form
$$\rho (m,x^+,x^{})m_0!\{\begin{array}{ccc}G(x^+,x^{},s,b)& \mathrm{if}& m=1\\ 1& \mathrm{if}& m=0\end{array},$$
(4.24)
where $`m_0=m_{\mathrm{max}}m`$ is the number of zeros in the row $`k`$. Let us define $`\overline{m}_0`$ as the number of zeros (empty cells) in the row $`k`$ not counting the current site $`(k,\mu )`$. Then the factor $`m_0!`$ is given as $`\overline{m}_0!`$ in case of $`m0`$ and as $`\overline{m}_0!(\overline{m}+1)`$ in case of $`m=0`$, and we obtain
$$\rho (m,x^+,x^{})(\overline{m}_0+1)\delta _{m0}+G(x^+,x^{},s,b)\delta _{m1}.$$
(4.25)
Properly normalized, we obtain
$$\rho (m,x^+,x^{})=p_0\delta _{m0}+(1p_0)\frac{G(x^+,x^{},s,b)}{\chi }\delta _{m1},$$
(4.26)
where the probability $`p_0`$ of proposing no interaction is given as
$$p_0=\frac{\overline{m}_0+1}{\overline{m}_0+1+\chi (s,b)},$$
(4.27)
with $`\chi `$ being obtained by integrating $`G`$ over $`x^+`$ and $`x^{}`$,
$$\chi (s,b)=_0^1𝑑x^+𝑑x^{}G(x^+,x^{},s,b).$$
(4.28)
Having proposed a new configuration $`L`$, which amounts to generating the values $`m_{k\mu },x_{k\mu }^+,x_{k\mu }^{}`$ for a randomly chosen lattice site as described above, we accept this proposal with the probability
$$u(K,L)=\mathrm{min}(z_1z_2,1),$$
(4.29)
with
$$z_1=\frac{\mathrm{\Omega }(L)}{\mathrm{\Omega }(K)},z_2=\frac{w(L,K)}{w(K,L)}.$$
(4.30)
Since $`K`$ and $`L`$ differ in at most one lattice site, say $`(k,\mu )`$, we do not need to evaluate the full formula for the distribution $`\mathrm{\Omega }`$ to calculate $`z_1`$, we rather calculate
$$z_1=\frac{\mathrm{\Omega }^{k\mu }(L)}{\mathrm{\Omega }^{k\mu }(K)},$$
(4.31)
with
$`\mathrm{\Omega }^{k\mu }(K)`$ $`=`$ $`\rho (m_{k\mu },x_{k\mu }^+,x_{k\mu }^{})\mathrm{exp}\left({\displaystyle \underset{l\mathrm{linked}\mathrm{to}k}{}}\stackrel{~}{G}(x_{\pi (l)}^+x_{\tau (l)}^{},s,b_l)\right)`$ (4.32)
$`\times `$ $`(x_{\pi (k)}^+)^{\alpha _{\mathrm{remn}}}\mathrm{\Theta }(x_{\pi (k)}^+)\mathrm{\Theta }(1x_{\pi (k)}^+)(x_{\tau (k)}^{})^{\alpha _{\mathrm{remn}}}\mathrm{\Theta }(x_{\tau (k)}^{})\mathrm{\Theta }(1x_{\tau (k)}^{}),`$
which is technically quite easy. Our final task is the calculation of the asymmetry $`z_2`$. In many applications of the Markov chain method one uses symmetric proposal matrices, in which case this factor is simply one. This is not the case here. We have
$$z_2=\frac{\mathrm{\Omega }_0(K)}{\mathrm{\Omega }_0(L)}=\frac{\mathrm{\Omega }_0^{k\mu }(K)}{\mathrm{\Omega }_0^{k\mu }(L)},$$
(4.33)
with
$$\mathrm{\Omega }_0^{k\mu }(K)=\rho (m_{k\mu },x_{k\mu }^+,x_{k\mu }^{}),$$
(4.34)
which is also easily calculated. So we accept the proposal $`L`$ with the probability $`\mathrm{min}(z_1z_2,1)`$, in which case we have $`K^{(t+1)}=L`$, otherwise we keep the old configuration $`K`$, which means $`K^{(t+1)}=K`$.
### 4.4 Convergence
A crucial item is the question of how to determine the number of iterations, which are sufficient to reach the stationary region. In principle one could calculate the autocorrelation time, or better one could estimate it based on an actual iteration. One could then multiply it with some “reasonable number”, between 10 and 20, in order to obtain the number of iterations. Since this “reasonable number” is not known anyway, we proceed differently. We consider a number of quantities like the number of binary interactions, the number of Pomerons, and other observables, and we monitor their values during the iterations. Simply by inspecting the results for many events, one can quite easily convince oneself if the numbers of iterations are sufficiently large. As a final check one makes sure that the distributions of some relevant observables do not change by doubling the number of iterations. In fig. 4.2,
we show the number of collisions (left) and the number of Pomerons (right) as a function of the iteration step $`t`$ for a S+S collision, where the number of iterations $`t_{\mathrm{max}}`$ has been determined according to some empirical procedure described below. We observe that these two quantities approach very quickly the stationary region. In order to determine the number $`t_{\mathrm{max}}`$ of iterations for a given reaction $`A+B`$, we first calculate the upper limit for the number of possibly interacting nucleon pairs as the number of pairs $`k_{\mathrm{max}}`$ with a transverse distance smaller than some value $`b_{\mathrm{max}}`$ being defined as
$$1e^{\chi (s,b_{\mathrm{max}})}=0.001,$$
(4.35)
and we then define
$$t_{\mathrm{max}}=100\frac{2}{3}k_{\mathrm{max}.}$$
(4.36)
Actually, in the real calculations, we never consider sums of nucleon pairs from 1 to $`AB`$, but only from 1 to $`k_{\mathrm{max}}`$, because for the other ones the chance to be involved in an interaction is so small that one can safely ignore it.
### 4.5 Some Tests for Proton-Proton Scattering
As a first test, we check whether the Monte Carlo procedure reproduces the theoretical profile function $`\gamma _{\mathrm{inel}}`$. So we make a large number of simulations of proton-proton collisions at a given energy $`\sqrt{s}`$, where the impact parameters are chosen randomly between zero and the earlier defined maximum impact parameter $`b_{\mathrm{max}}`$. We then count simply the number $`\mathrm{\Delta }n_{\mathrm{inel}}(b)`$ of inelastic interactions in a given impact parameter bin $`[b\mathrm{\Delta }b/2,b+\mathrm{\Delta }b/2]`$, and divide this by the number $`\mathrm{\Delta }n_{\mathrm{tot}}(b)`$ of simulations in this impact parameter interval. Since the total number $`\mathrm{\Delta }n_{\mathrm{tot}}(b)`$ of simulated configurations for the given $`b`$bin splits into the number $`\mathrm{\Delta }n_{\mathrm{inel}}(b)`$ of “interactions” and the number $`\mathrm{\Delta }n_{\mathrm{nonint}}(b)`$ of “non-interactions”, with $`\mathrm{\Delta }n_{\mathrm{tot}}(b)=`$$`\mathrm{\Delta }n_{\mathrm{inel}}(b)+\mathrm{\Delta }n_{\mathrm{nonint}}(b)`$, the result
$$P(b)=\frac{\mathrm{\Delta }n_{\mathrm{inel}}(b)}{\mathrm{\Delta }n_{\mathrm{tot}}(b)}$$
(4.37)
represents the probability to have an interaction at a given impact parameter $`b`$, which should coincide with the profile function
$$\gamma _{AB}(s,b)=1\mathrm{\Phi }_{\mathrm{u}_{AB}}(1,1,s,b)$$
(4.38)
for the corresponding energy. In fig. 4.3,
we compare the two quantities for a proton-proton collision at $`\sqrt{s}=200`$ GeV and we find an excellent agreement, as it should be.
Another elementary quantity is the inclusive momentum spectrum of Pomerons. Pomerons, representing elementary interactions, are characterized by their light cone momentum fractions $`x^+`$ and $`x^{}`$, so one might study two dimensional distributions, or for example a distribution in $`x=x^+x^{}`$, where the second variable $`y=0.5\mathrm{log}(x^+/x^{})`$ is integrated over. So again we simulate many proton-proton events at a given energy $`\sqrt{s}`$ and we count the number of Pomerons $`\mathrm{\Delta }N_{\mathrm{Pom}}`$ within a certain interval $`[x\mathrm{\Delta }x/2,x+\mathrm{\Delta }x/2]`$ and we calculate
$$\frac{dn_{\mathrm{Pom}}^{\mathrm{MC}}}{dx}=\frac{1}{N_{\mathrm{events}}}\frac{\mathrm{\Delta }N_{\mathrm{Pom}}}{\mathrm{\Delta }x},$$
(4.39)
representing the Monte Carlo Pomeron $`x`$-distribution, which may be compared with the analytical result calculated earlier, as shown in fig. 4.4.
The analytical results of course refer to the unitarized theory. Again we find perfect agreement between Monte Carlo simulations and analytical curves, as it should be. In fig. 4.5,
we compare inclusive Pomeron cross sections (integrated over impact parameter). Here, the impact parameters are generated randomly between 1 and some $`b_{\mathrm{max}}`$, one counts the number of Pomerons $`\mathrm{\Delta }N_{\mathrm{Pom}}`$ within a certain interval of size $`\mathrm{\Delta }x`$, and one calculates
$$\frac{d\sigma _{\mathrm{Pom}}^{\mathrm{MC}}}{dx}=\frac{\pi r_{\mathrm{max}}^2}{N_{\mathrm{events}}}\frac{\mathrm{\Delta }N_{\mathrm{Pom}}}{\mathrm{\Delta }x},$$
(4.40)
which is compared with the analytical result
$$\frac{d\sigma _{\mathrm{Pom}}}{\mathrm{dx}}(x,s)=d^2b\frac{dn_{\mathrm{Pom}}}{dx}(x,s,b),$$
(4.41)
which again show an excellent agreement.
These Pomeron distributions are of particular interest, because they are elementary distributions based on which other inclusive spectra like a transverse momentum distribution of pions may be obtained via convolution.
The two examples of this section provide on one hand a check that the numerical procedures work properly, on the other hand they demonstrate nicely that our Monte Carlo procedure is a very well defined numerical method to solve a particular mathematical problem. In simple cases where analytical results exist, they may be compared with the Monte Carlo results, and they must absolutely agree.
## Chapter 5 Enhanced Pomeron diagrams
The eikonal type diagrams shown at fig. 5.1, considered in the previous chapters, correspond to pair-like scatterings between hadron constituents and form the basis for the description of hadronic interactions at not too high energies.
However, when the interaction energy increases, the contribution of so-called enhanced Pomeron diagrams as, for example, the diagrams shown in fig. 5.2, become more and more important.
The latter ones take into account interactions of Pomerons with each other. The corresponding amplitudes increase asymptotically much faster than the usual eikonal type contributions considered so far. In this paper, we restrict ourselves to the lowest order diagrams ($`Y`$-diagrams and inverted $`Y`$-diagrams).
In the following sections, we discuss the amplitudes corresponding to the lowest order enhanced diagrams and the modification of the hadronic profile function in the presence of these diagrams, before we discuss their most important features.
### 5.1 Calculating lowest order enhanced diagrams
To introduce enhanced type diagrams let us come back to the process of double soft Pomeron exchange, which is a particular case of the diagram of fig. 2.7. The corresponding contribution to the elastic scattering amplitude is given in eqs. (2.31), (C.11) with $`n=2`$ and with $`T_{1\mathrm{I}\mathrm{P}}`$ being replaced by $`T_{\mathrm{soft}}`$:
$`iT_{h_1h_2}^{(2)}(s,t)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{d^4k_1}{(2\pi )^4}\frac{d^4k_1^{}}{(2\pi )^4}\frac{d^4k_2}{(2\pi )^4}\frac{d^4k_2^{}}{(2\pi )^4}\frac{d^4q_1}{(2\pi )^4}\mathrm{\Theta }\left(s_1^+\right)\mathrm{\Theta }\left(s_1^{}\right)\mathrm{\Theta }\left(s_2^+\right)\mathrm{\Theta }\left(s_2^{}\right)\mathrm{\Theta }\left(s_{q_1}^+\right)\mathrm{\Theta }\left(s_{q_1}^{}\right)}`$ (5.1)
$`\times \mathrm{disc}_{s_1^+,s_2^+,s_{q_1}^+}N_{h_1}^{(2)}(p,k_1,k_2,q_1,qq_1)`$
$`\times {\displaystyle \underset{l=1}{\overset{2}{}}}\left[iT_{\mathrm{soft}}(\widehat{s}_l,q_l^2)\right]\mathrm{disc}_{s_1^{},s_2^{},s_{q_1}^{}}N_{h_2}^{(2)}(p^{},k_1^{},k_2^{},q_1,q+q_1)`$
see fig. 5.3.
We are now interested in the contribution with some of the invariants
$`s_1^+`$ $`=`$ $`(pk_1)^2p^+k_1^{},`$ (5.2)
$`s_2^+`$ $`=`$ $`(pk_1k_2)^2p^+(k_1^{}+k_2^{}),`$ (5.3)
$`s_{q_1}^+`$ $`=`$ $`(p+q_1)^2p^+q_1^{},`$ (5.4)
being large, implying $`k_i^{},q_1^{}`$ to be not too small. As shown in appendix C.3, in that case the above amplitude may be written as
$$iT_{h_1h_2}^{3\mathrm{I}\mathrm{P}}(s,t)=_0^1\frac{dx^+}{x^+}\frac{dx^{}}{x^{}}F_{\mathrm{remn}}^{h_1}\left(1x^+\right)F_{\mathrm{remn}}^{h_2}\left(1x^{}\right)iT_{3\mathrm{I}\mathrm{P}}^{h_1h_2}(x^+,x^{},s,t)$$
(5.5)
with
$`iT_{3\mathrm{I}\mathrm{P}}^{h_1h_2}(x^+,x^{},s,t)=\mathrm{\hspace{0.25em}8}\pi ^2x^+x^{}s{\displaystyle \frac{r_{3\mathrm{I}\mathrm{P}}}{2}}{\displaystyle _{s_0/x^{}}^{x^+}}{\displaystyle \frac{dx_{12}^+}{x_{12}^+}}\left[{\displaystyle \frac{1}{2s^+}}\mathrm{Im}T^{h_1}(x^+,s^+,q_{}^2)\right]`$
$`\times {\displaystyle }dz^+{\displaystyle }d^2q_1_{}d^2q_2_{}{\displaystyle _0^x^{}}dx_1^{}dx_2^{}{\displaystyle \underset{l=1}{\overset{2}{}}}\left[{\displaystyle \frac{1}{8\pi ^2\widehat{s}_l}}iT^{h_2}(x_l^{},\widehat{s}_l,q_l_{}^2)\right]`$
$`\times \delta (x^{}x_1^{}x_2^{})\delta ^{(2)}\left(\stackrel{}{q}_{}\stackrel{}{q}_1_{}\stackrel{}{q}_2_{}\right),`$ (5.6)
with
$$T^h(x,s,q_{}^2)=T_{\mathrm{soft}}^h(x,s,q_{}^2)=T_{\mathrm{soft}}(s,q_{}^2)F_{\mathrm{part}}^h(x)\mathrm{exp}\left(R_h^2q_{}^2\right)$$
(5.7)
and
$$\widehat{s}_1=x_{12}^+z^+x_1^{}s,\widehat{s}_2=x_{12}^+(1z^+)x_2^{}s,$$
(5.8)
where the following definitions have been used
$`x^+=k^+/p^+,`$ (5.9)
$`x_{12}^+=k_{12}^+/p^+,`$ (5.10)
$`x_1^+=k_1^+/p^+,`$ (5.11)
$`x_{12}^+x_1^+=(k_{12}^+k_1^+)/p^+=k_2^+/p^+,`$ (5.12)
$`z^+=k_1^+/k_{12}^+=x_1^+/x_{12}^+,`$ (5.13)
$`s^+=(kk_{12})^2k^+k_{12}^{}s_0k^+/k_{12}^+=s_0x^+/x_{12}^+,`$ (5.14)
see fig. 5.4.
The sign “$``$” in “$`3\mathrm{I}\mathrm{P}`$” refers to the Pomeron “splitting” towards the target hadron (reversed $`Y`$-diagram); the lower limit for the integral $`dx_{12}^+`$ is due to $`x_{12}^{}s_0/x_{12}^+<x^{}`$.
The triple-Pomeron contribution (5.6) is by construction expressed via amplitudes $`T_{\mathrm{soft}}^h`$ for parton-parton scattering due to soft Pomeron exchange, each one corresponding to non-perturbative parton dynamics, characterized by restricted parton virtualities $`Q^2<Q_0^2`$. We can also take into account contributions to the triple-Pomeron diagram from semi-hard processes, when some part of the parton cascade mediating the scattering between partons of momenta $`k_l`$ and $`k_l^{}`$ at fig. C.3 ($`k`$ and $`k_{12}`$) develops in the perturbative region $`Q^2>Q_0^2`$. Then, according to the general discussion of chapter 2, the amplitudes $`T^h`$ obtain also contributions from semi-hard sea-type parton-parton scattering $`T_{\mathrm{sea}\mathrm{sea}}^h`$ and from valence quark scattering $`T_{\mathrm{val}\mathrm{sea}}^h`$:
$$T^h=T_{\mathrm{soft}}^h+T_{\mathrm{sea}\mathrm{sea}}^h+T_{\mathrm{val}\mathrm{sea}}^h,$$
(5.15)
with
$$T_{\mathrm{soft}/\mathrm{sea}\mathrm{sea}}^h(x,\widehat{s},q_{}^2)=T_{\mathrm{soft}/\mathrm{sea}\mathrm{sea}}(\widehat{s},q_{}^2)F_{\mathrm{part}}^h(x)\mathrm{exp}\left(R_h^2q_{}^2\right)$$
(5.16)
and
$`T_{\mathrm{val}\mathrm{sea}}^h(x,\widehat{s},q_{}^2)={\displaystyle _0^x}𝑑x_v{\displaystyle \frac{x}{x_v}}{\displaystyle \underset{k}{}}T_{\mathrm{val}\mathrm{sea}}^k({\displaystyle \frac{x_v}{x}}\widehat{s},q_{}^2)\mathrm{exp}\left(R_h^2q_{}^2\right)`$
$`\times \overline{F}_{\mathrm{part}}^{h,k}(x_v,xx_v).`$ (5.17)
We have to stress again that we do not consider the possibility of Pomeron-Pomeron coupling in the perturbative region $`Q^2>Q_0^2`$. Therefore in our scheme hard parton processes can only contribute into internal structure of elementary parton-parton scattering amplitudes but do not influence the triple-Pomeron coupling.
A similar contribution $`T_{3\mathrm{I}\mathrm{P}+}^{h_1h_2}`$ of the $`Y`$-diagram can be obtained via interchanging $`x^+x^{}`$ and $`h_1h_2`$:
$$T_{3\mathrm{I}\mathrm{P}+}^{h_1h_2}(x^+,x^{},s,t)=T_{3\mathrm{I}\mathrm{P}}^{h_2h_1}(x^{},x^+,s,t).$$
(5.18)
One can repeat the above derivation for the case of a general soft multiple scattering process, see eq. (2.31), where some of the energy invariants $`s_{q_l}^+`$, $`s_{q_l}^{}`$ are large . One then obtains finally the general multiple scattering expression (C.1), with the contribution of corresponding pairs of Pomerons being replaced by expressions
$$\frac{1}{8\pi ^2x^+x^{}s}d^2q_{}iT_{3\mathrm{I}\mathrm{P}\pm }^{h_1h_2}(x^+,x^{},s,t),$$
(5.19)
where $`x^\pm `$ refers to the summary light cone momentum share of the constituent partons participating in the triple-Pomeron process. So we get
$`iT_{h_1h_2}(s,t)=8\pi ^2s{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n!}}{\displaystyle _0^1}{\displaystyle \underset{l=1}{\overset{n}{}}}dx_l^+dx_l^{}{\displaystyle \underset{l=1}{\overset{n}{}}}\left[{\displaystyle \frac{1}{8\pi ^2\widehat{s}_l}}{\displaystyle d^2q_l_{}iT_{\mathrm{I}\mathrm{P}}^{h_1h_2}(\widehat{s}_l,q_l_{}^2)}\right]`$
$`\times F_{\mathrm{remn}}^{h_1}\left(1{\displaystyle \underset{j=1}{\overset{n}{}}}x_j^+\right)F_{\mathrm{remn}}^{h_2}\left(1{\displaystyle \underset{j=1}{\overset{n}{}}}x_j^{}\right)\delta ^{(2)}\left({\displaystyle \underset{k=1}{\overset{n}{}}}\stackrel{}{q}_k_{}\stackrel{}{q}_{}\right).`$ (5.20)
with
$$T_{\mathrm{I}\mathrm{P}}^{h_1h_2}=T_{1\mathrm{I}\mathrm{P}}^{h_1h_2}+T_{3\mathrm{I}\mathrm{P}}^{h_1h_2}+T_{3\mathrm{I}\mathrm{P}+}^{h_1h_2}$$
(5.21)
Here, we allow any number of simple triple-Pomeron diagrams; thus we restrict ourselves to the contributions of double Pomeron iteration in the $`t`$-channel rather than to the first order in $`r_{3\mathrm{I}\mathrm{P}}`$.
The Fourier transform $`\stackrel{~}{T}`$ of the amplitude (5.20) is given as
$`{\displaystyle \frac{i}{2s}}\stackrel{~}{T}_{h_1h_2}(s,b)`$ $`=`$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n!}}{\displaystyle _0^1}{\displaystyle \underset{l=1}{\overset{n}{}}}dx_l^+dx_l^{}{\displaystyle \underset{l=1}{\overset{n}{}}}{\displaystyle \frac{i}{2s}}\stackrel{~}{T}_{\mathrm{I}\mathrm{P}}^{h_1h_2}(x_l^+,x_l^{},s,b)`$ (5.22)
$`\times `$ $`F_{\mathrm{remn}}^{h_1}\left(1{\displaystyle \underset{j=1}{\overset{n}{}}}x_j^+\right)F_{\mathrm{remn}}^{h_2}\left(1{\displaystyle \underset{j=1}{\overset{n}{}}}x_j^{}\right),`$
with $`\stackrel{~}{T}_{\mathrm{I}\mathrm{P}}^{h_1h_2}`$ being the Fourier transform of $`T_{\mathrm{I}\mathrm{P}}^{h_1h_2}`$,
$$\stackrel{~}{T}_{\mathrm{I}\mathrm{P}}^{h_1h_2}=\stackrel{~}{T}_{1\mathrm{I}\mathrm{P}}^{h_1h_2}+\stackrel{~}{T}_{3\mathrm{I}\mathrm{P}}^{h_1h_2}+\stackrel{~}{T}_{3\mathrm{I}\mathrm{P}+}^{h_1h_2},$$
(5.23)
where the Fourier transform $`\stackrel{~}{T}_{3\mathrm{I}\mathrm{P}}^{h_1h_2}`$ of the triple Pomeron amplitude $`T_{3\mathrm{I}\mathrm{P}}^{h_1h_2}`$ is given as
$`{\displaystyle \frac{i}{2\widehat{s}}}\stackrel{~}{T}_{3\mathrm{I}\mathrm{P}}^{h_1h_2}(x^+,x^{},s,b)={\displaystyle \frac{r_{3\mathrm{I}\mathrm{P}}}{2}}{\displaystyle d^2b_1_{s_0/x^{}}^{x^+}\frac{dx_{12}^+}{x_{12}^+}\left[\frac{1}{2s^+}\mathrm{Im}\stackrel{~}{T}^{h_1}(x^+,s^+,|\stackrel{}{b}\stackrel{}{b}_1|)\right]}`$
$`\times {\displaystyle }dz^+{\displaystyle _0^x^{}}dx_1^{}dx_2^{}{\displaystyle \underset{l=1}{\overset{2}{}}}\left[{\displaystyle \frac{1}{2\widehat{s}_l}}i\stackrel{~}{T}^{h_2}(x_l^{},\widehat{s}_l,b_1)\right]\delta (x^{}x_1^{}x_2^{}),`$ (5.24)
with $`\widehat{s}=x^+x^{}s`$ (and similarly for $`\stackrel{~}{T}_{3\mathrm{I}\mathrm{P}+}^{h_1h_2}`$).
Here we used
$$\delta ^{(2)}\left(\stackrel{}{q}_{}\stackrel{}{q}_1_{}\stackrel{}{q}_2_{}\right)=\frac{1}{4\pi ^2}d^2b_1\mathrm{exp}\left(i\left(\stackrel{}{q}_{}\stackrel{}{q}_1_{}\stackrel{}{q}_2_{}\right)\stackrel{}{b}_1\right).$$
(5.25)
The profile function $`\gamma `$ for hadron-hadron interaction is as usual defined as
$$\gamma _{h_1h_2}(s,b)=\frac{1}{2s}2\mathrm{I}\mathrm{m}\stackrel{~}{\mathrm{T}}_{h_1h_2}(s,b),$$
(5.26)
which may be evaluated using the AGK cutting rules,
$`\gamma _{h_1h_2}(s,b)`$ $`=`$ $`{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{m!}}{\displaystyle \underset{\mu =1}{\overset{m}{}}dx_\mu ^+dx_\mu ^{}\underset{\mu =1}{\overset{m}{}}G_{\mathrm{I}\mathrm{P}}^{h_1h_2}(x_\mu ^+,x_\mu ^{},s,b)}`$ (5.27)
$`\times `$ $`{\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{l!}}{\displaystyle \underset{\lambda =1}{\overset{l}{}}d\stackrel{~}{x}_\lambda ^+d\stackrel{~}{x}_\lambda ^{}\underset{\lambda =1}{\overset{l}{}}}G_{\mathrm{I}\mathrm{P}}^{h_1h_2}(\stackrel{~}{x}_\lambda ^+,\stackrel{~}{x}_\lambda ^{},s,b)`$
$`\times `$ $`F_{\mathrm{remn}}\left(x^{\mathrm{proj}}{\displaystyle \underset{\lambda }{}}\stackrel{~}{x}_\lambda ^+\right)F_{\mathrm{remn}}\left(x^{\mathrm{targ}}{\displaystyle \underset{\lambda }{}}\stackrel{~}{x}_\lambda ^{}\right),`$
with $`x^{\mathrm{proj}/\mathrm{targ}}=1x_\mu ^\pm `$ being the momentum fraction of the projectile/target remnant, and with
$$G_{\mathrm{I}\mathrm{P}}^{h_1h_{2=}}G_{1\mathrm{I}\mathrm{P}}^{h_1h_2}+G_{3\mathrm{I}\mathrm{P}}^{h_1h_2}+G_{3\mathrm{I}\mathrm{P}+}^{h_1h_2},$$
(5.28)
where $`G_{3\mathrm{I}\mathrm{P}\pm }^{h_1h_2}`$ is twice the imaginary part of the Fourier transformed triple-Pomeron amplitude $`\stackrel{~}{T}_{3\mathrm{I}\mathrm{P}\pm }^{h_1h_2}`$ divided by $`2\widehat{s}`$,
$$G_{3\mathrm{I}\mathrm{P}\pm }^{h_1h_2}(x^+,x^{},s,b)=\frac{1}{2x^+x^{}s}2\mathrm{I}\mathrm{m}\stackrel{~}{T}_{3\mathrm{I}\mathrm{P}\pm }^{h_1h_2}(x^+,x^{},s,b),$$
(5.29)
which gives, assuming imaginary amplitudes,
$`G_{3\mathrm{I}\mathrm{P}}^{h_1h_2}(x^+,x^{},s,b)`$ $`=`$ $`{\displaystyle \frac{r_{3\mathrm{I}\mathrm{P}}}{8}}{\displaystyle d^2b_1_{s_0/x^{}}^{x^+}\frac{dx_{12}^+}{x_{12}^+}G^{h_1}(x^+,s^+,\left|\stackrel{}{b}\stackrel{}{b}_1\right|)}`$ (5.30)
$`\times {\displaystyle _0^1}dz^+{\displaystyle _0^x^{}}dx_1^{}G^{h_2}(x_1^{},\widehat{s}_1,b_1)G^{h_2}(x^{}x_1^{},\widehat{s}_2,b_1),`$
see fig. 5.5,
with
$$s^+=s_0\frac{x^+}{x_{12}^+},\widehat{s}_1=x_{12}^+z^+x_1^{}s,\widehat{s}_2=x_{12}^+(1z^+)(x^{}x_1^{})s.$$
(5.31)
The functions $`G^h`$ are defined as
$$G^h(x,\widehat{s},b)=\frac{1}{2\widehat{s}}\mathrm{\hspace{0.17em}2}\mathrm{Im}\stackrel{~}{T}^h(x,\widehat{s},b),$$
(5.32)
with $`\stackrel{~}{T}^h`$ being the Fourier transform of $`T^h`$, which gives
$$G^h=G_{\mathrm{soft}}^h+G_{\mathrm{sea}\mathrm{sea}}^h+G_{\mathrm{val}\mathrm{sea}}^h,$$
(5.33)
with
$`G_{\mathrm{soft}}^h(x,\widehat{s},b)={\displaystyle \frac{2\gamma _{part}}{\lambda _{\mathrm{soft}}^h(\frac{\widehat{s}}{s_0})}}\left({\displaystyle \frac{\widehat{s}}{s_0}}\right)^{\alpha _{_{\mathrm{soft}}}(0)1}\mathrm{exp}\left({\displaystyle \frac{b^2}{4\lambda _{\mathrm{soft}}^h(\frac{\widehat{s}}{s_0})}}\right)F_{\mathrm{part}}(x)`$ (5.34)
$`G_{\mathrm{sea}\mathrm{sea}}^h(x,\widehat{s},b)={\displaystyle \frac{1}{4\pi }}{\displaystyle \underset{jk}{}}{\displaystyle _0^1}dz^+dz^{}E_{\mathrm{soft}}^j\left(z^+\right)E_{\mathrm{soft}}^k\left(z^{}\right)\sigma _{\mathrm{ladder}}^{jk}(z^+z^{}\widehat{s},Q_0^2)\times `$ (5.35)
$`\times {\displaystyle \frac{1}{\lambda _{\mathrm{soft}}^h\left(1/z^+z^{}\right)}}\mathrm{exp}\left({\displaystyle \frac{b^2}{4\lambda _{\mathrm{soft}}^h\left(1/z^+z^{}\right)}}\right)F_{\mathrm{part}}(x)`$
$`G_{\mathrm{val}\mathrm{sea}}^h(x,\widehat{s},b)={\displaystyle \frac{1}{4\pi }}{\displaystyle \underset{jk}{}}{\displaystyle _0^x}𝑑x_v{\displaystyle _0^1}𝑑z^+E_{\mathrm{soft}}^j\left(z^+\right)\sigma _{\mathrm{ladder}}^{jk}({\displaystyle \frac{x_v}{x}}z^+\widehat{s},Q_0^2)`$ (5.36)
$`\times {\displaystyle \frac{1}{\lambda _{\mathrm{soft}}^h\left(1/z^+\right)}}\mathrm{exp}\left({\displaystyle \frac{b^2}{4\lambda _{\mathrm{soft}}^h\left(1/z^+\right)}}\right)\overline{F}_{\mathrm{part}}^k(x_v,xx_v),`$
with
$$\lambda _{\mathrm{soft}}^h\left(z\right)=R_h^2+\alpha _{_{\mathrm{soft}}}^{}\mathrm{ln}z.$$
(5.37)
### 5.2 Cutting Enhanced Diagrams
To treat particle production, we have to investigate the different cuts of an enhanced diagram. We consider the inverted $`Y`$-diagram here, the same arguments apply to the $`Y`$-diagram. We employ the cutting rules to eq. (5.24),
$$G_{3\mathrm{I}\mathrm{P}}^{h_1h_2}=\frac{1}{2\widehat{s}}2\mathrm{I}\mathrm{m}\stackrel{~}{T}_{3\mathrm{I}\mathrm{P}}^{h_1h_2}=\mathrm{sum}\mathrm{over}\mathrm{cut}\mathrm{diagrams}=\underset{i}{}G_{3\mathrm{I}\mathrm{P}(i)}^{h_1h_2},$$
(5.38)
where the index $`i`$ counts the different cuts. We take into account that the cutting procedure only influences the two Pomerons exchanged “in parallel” in the triple-Pomeron graph (the lower Pomerons) with the third Pomeron already being cut , so that we have three contributions: none of the lower Pomerons cut ($`i=0`$, diffraction), one of these Pomerons cut ($`i=1`$, screening), and both Pomerons cut ($`i=2`$, Pomeron-Pomeron fusion), see fig. 5.6.
So we have
$$G_{3\mathrm{I}\mathrm{P}}^{h_1h_2}=G_{3\mathrm{I}\mathrm{P}(0)}^{h_1h_2}+G_{3\mathrm{I}\mathrm{P}(1)}^{h_1h_2}+G_{3\mathrm{I}\mathrm{P}(2)}^{h_1h_2},$$
(5.39)
with
$`G_{3\mathrm{I}\mathrm{P}(0)}^{h_1h_2}`$ $`=`$ $`\left\{{\displaystyle \frac{i}{2\widehat{s}}}\stackrel{~}{T}_{3\mathrm{I}\mathrm{P}}^{h_1h_2}\right\}\times 2,`$
$`G_{3\mathrm{I}\mathrm{P}(1)}^{h_1h_2}`$ $`=`$ $`\left\{{\displaystyle \frac{i}{2\widehat{s}}}\stackrel{~}{T}_{3\mathrm{I}\mathrm{P}}^{h_1h_2}\right\}\times (2)\times 2\times 2,`$ (5.40)
$`G_{3\mathrm{I}\mathrm{P}(2)}^{h_1h_2}`$ $`=`$ $`\left\{{\displaystyle \frac{i}{2\widehat{s}}}\stackrel{~}{T}_{3\mathrm{I}\mathrm{P}}^{h_1h_2}\right\}\times 2\times 2.`$
Here, we assumed imaginary amplitudes, and we replace as usual a factor $`i\stackrel{~}{T}^h`$ in eq. (5.24) by $`2\mathrm{I}\mathrm{m}\stackrel{~}{T}^h=2i\stackrel{~}{T}^h`$ for a cut Pomeron, and by $`(i\stackrel{~}{T}^h)^{}=i\stackrel{~}{T}^h`$ for an uncut Pomeron being to the left of the cut plane, and by $`i\stackrel{~}{T}^h`$ for an uncut Pomeron being to the right of the cut plane. Using (see eq. (5.29))
$$\frac{i}{2\widehat{s}}\stackrel{~}{T}_{3\mathrm{I}\mathrm{P}}^{h_1h_2}=\frac{1}{2}G_{3\mathrm{I}\mathrm{P}}^{h_1h_2},$$
(5.41)
we get
$`G_{3\mathrm{I}\mathrm{P}(0)}^{h_1h_2}`$ $`=`$ $`1\times G_{3\mathrm{I}\mathrm{P}}^{h_1h_2},`$
$`G_{3\mathrm{I}\mathrm{P}(1)}^{h_1h_2}`$ $`=`$ $`+4\times G_{3\mathrm{I}\mathrm{P}}^{h_1h_2},`$ (5.42)
$`G_{3\mathrm{I}\mathrm{P}(2)}^{h_1h_2}`$ $`=`$ $`2\times G_{3\mathrm{I}\mathrm{P}}^{h_1h_2},`$
which means that each cut contribution is equal to the profile function, up to a pre-factor, see fig. 5.6. The sum of the three contributions is $`G_{3\mathrm{I}\mathrm{P}}^{h_1h_2}`$, as it should be.
There is a substantial difference between the different cuts of triple-Pomeron contributions: in the case of both lower Pomerons being cut ($`i=2`$) all the momentum of the constituent partons, participating in the process, is transferred to secondary hadrons produced, whereas for the cut between these Pomerons ($`i=0`$) only the light cone momentum fractions of the cut Pomeron ($`x^+`$, $`x_{12}^{}=s_0/x_{12}^+`$ in eq. (5.30)) are available for hadron production, the momentum share $`x^{}x_{12}^{}`$ of the partons, connected to the uncut (virtual) Pomerons is given back to the remnant state. Correspondingly, the contribution with one of the two lower Pomerons being cut ($`i=1`$) defines the screening correction to the elementary rescattering with the momentum fractions $`x^+`$, $`x_1^{}`$ (considering the first of the two Pomerons being cut). It is therefore useful to rewrite the expression for the profile function as
$`\gamma _{h_1h_2}(s,b)`$ $`=`$ $`{\displaystyle \underset{m}{}}{\displaystyle \frac{1}{m!}}{\displaystyle \underset{\mu =1}{\overset{m}{}}dx_\mu ^+dx_\mu ^{}d\widehat{x}_\mu ^+d\widehat{x}_\mu ^{}\underset{\mu =1}{\overset{m}{}}\widehat{G}_{\mathrm{I}\mathrm{P}}^{h_1h_2}(x_\mu ^+,x_\mu ^{},\widehat{x}_\mu ^+,\widehat{x}_\mu ^{},s,b)}`$ (5.43)
$`\times `$ $`{\displaystyle \underset{l}{}}{\displaystyle \frac{1}{l!}}{\displaystyle \underset{\lambda =1}{\overset{l}{}}d\stackrel{~}{x}_\lambda ^+d\stackrel{~}{x}_\lambda ^{}\underset{\lambda =1}{\overset{l}{}}}G_{\mathrm{I}\mathrm{P}}^{h_1h_2}(\stackrel{~}{x}_\lambda ^+,\stackrel{~}{x}_\lambda ^{},s,b)`$
$`\times `$ $`F_{\mathrm{remn}}\left(x^{\mathrm{proj}}{\displaystyle \underset{\mu }{}}\widehat{x}_\mu ^+{\displaystyle \underset{\lambda }{}}\stackrel{~}{x}_\lambda ^+\right)F_{\mathrm{remn}}\left(x^{\mathrm{targ}}{\displaystyle \underset{\mu }{}}\widehat{x}_\mu ^{}{\displaystyle \underset{\lambda }{}}\stackrel{~}{x}_\lambda ^{}\right),`$
with $`G_{\mathrm{I}\mathrm{P}}^{h_1h_2}`$ being defined earlier, and with
$`\widehat{G}_{\mathrm{I}\mathrm{P}}^{h_1h_2}(x_\mu ^+,x_\mu ^{},\widehat{x}_\mu ^+,\widehat{x}_\mu ^{},s,b)=G_{1\mathrm{I}\mathrm{P}}(x_\mu ^+,x_\mu ^{},s,b)\delta (\widehat{x}_\mu ^+)\delta (\widehat{x}_\mu ^{})+{\displaystyle \underset{\sigma =\pm }{}}{\displaystyle \underset{i=0}{\overset{2}{}}}\widehat{G}_{3\mathrm{I}\mathrm{P}\sigma (i)}^{h_1h_2}(x_\mu ^+,x_\mu ^{},\widehat{x}_\mu ^+,\widehat{x}_\mu ^{},s,b),`$
(5.44)
where we used
$$x^{\mathrm{proj}/\mathrm{targ}}=1x_\mu ^\pm .$$
(5.45)
The variables $`x_\mu ^\pm `$ are now the momentum fractions for the individual cut contributions, which define the energy for the production of secondary hadrons resulting from a given elementary interaction. The expressions for the functions $`\widehat{G}_{3\mathrm{I}\mathrm{P}\sigma (i)}^{h_1h_2}`$ are obtained from eqs. (5.27, 5.30, 5.55, 5.56), by changing the variables properly. Simplest is the case of all Pomerons being cut, no changing of variables is necessary, we have simply
$`\widehat{G}_{3\mathrm{I}\mathrm{P}(2)}^{h_1h_2}(x^+,x^{},\widehat{x}^+,\widehat{x}^{},s,b)`$ $`=`$ $`2G_{3\mathrm{I}\mathrm{P}}^{h_1h_2}(x^+,x^{},s,b)\delta (\widehat{x}^+)\delta (\widehat{x}^{})`$
$`=`$ $`{\displaystyle \frac{r_{3\mathrm{I}\mathrm{P}}}{4}}{\displaystyle d^2b_1_{s_0/x^{}}^{x^+}\frac{dx_{12}^+}{x_{12}^+}G^{h_1}(x^+,s^+,\left|\stackrel{}{b}\stackrel{}{b}_1\right|)}`$
$`\times `$ $`{\displaystyle _0^1}𝑑z^+{\displaystyle _0^x^{}}𝑑x_1^{}G^{h_2}(x_1^{},\widehat{s}_1,b_1)G^{h_2}(x^{}x_1^{},\widehat{s}_2,b_1)`$
$`\times `$ $`\delta (\widehat{x}^+)\delta (\widehat{x}^{}),`$ (5.47)
For none of the two lower Pomerons being cut, we rename $`x_{12}^{}`$ into $`x^{}`$and $`x^{}x_{12}^{}`$ into $`\widehat{x}^{}`$, so we get
$`\widehat{G}_{3\mathrm{I}\mathrm{P}(0)}^{h_1h_2}(x^+,x^{},\widehat{x}^+,\widehat{x}^{},s,b)`$ $`=`$ $`{\displaystyle \frac{r_{3\mathrm{I}\mathrm{P}}}{8}}{\displaystyle d^2b_1\frac{1}{x^{}}G^{h_1}(x^+,x^+x^{}s,\left|\stackrel{}{b}\stackrel{}{b}_1\right|)}`$
$`\times `$ $`{\displaystyle _0^1}𝑑z^+{\displaystyle _0^{\widehat{x}^{}+x^{}}}𝑑x_1^{}G^{h_2}(x_1^{},x_1^{}{\displaystyle \frac{s_0}{x^{}}}z^+s,b_1)`$
$`\times `$ $`G^{h_2}(\widehat{x}^{}+x^{}x_1^{},(\widehat{x}^{}+x^{}x_1^{}){\displaystyle \frac{s_0}{x^{}}}(1z^+)s,b_1)\delta (\widehat{x}^+).`$
For one of the two lower Pomerons being cut, we rename $`x_1^{}`$ into $`x^{}`$and $`x^{}x_1^{}`$ into $`\widehat{x}^{}`$, and we get
$`\widehat{G}_{3\mathrm{I}\mathrm{P}(1)}^{h_1h_2}(x^+,x^{},\widehat{x}^+,\widehat{x}^{},s,b)`$ $`=`$ $`{\displaystyle \frac{r_{3\mathrm{I}\mathrm{P}}}{2}}{\displaystyle d^2b_1_{s_0/x^{}}^{x^+}\frac{dx_{12}^+}{x_{12}^+}G^{h_1}(x^+,s^+,\left|\stackrel{}{b}\stackrel{}{b}_1\right|)}`$
$`\times `$ $`{\displaystyle _0^1}𝑑z^+G^{h_2}(x^{},x_{12}^+z^+x^{}s,b_1)G^{h_2}(\widehat{x}^{},x_{12}^+(1z^+)\widehat{x}^{}s,b_1)\delta (\widehat{x}^+).`$
The contributions $`G_{3\mathrm{I}\mathrm{P}+(i)}^{h_1h_2}`$ can be obtained via interchanging $`x^+x^{},\widehat{x}^+\widehat{x}^{}`$ and $`h_1h_2`$ in the above formulas:
$$G_{3\mathrm{I}\mathrm{P}+(i)}^{h_1h_2}(x^+,x^{},\widehat{x}^+,\widehat{x}^{},s,b)=G_{3\mathrm{I}\mathrm{P}(i)}^{h_2h_1}(x^{},x^+,\widehat{x}^{},\widehat{x}^+,s,b).$$
(5.50)
We define as in the usual eikonal case the virtual emission function
$`\mathrm{\Phi }_{h_1h_2}(x^{\mathrm{proj}},x^{\mathrm{targ}},s,b)`$ $`=`$ $`{\displaystyle \underset{l}{}}{\displaystyle \underset{\lambda =1}{\overset{l}{}}d\stackrel{~}{x}_\lambda ^+d\stackrel{~}{x}_\lambda ^{}\frac{1}{l!}\underset{\lambda =1}{\overset{l}{}}}G_{\mathrm{I}\mathrm{P}}^{h_1h_2}(\stackrel{~}{x}_\lambda ^+,\stackrel{~}{x}_\lambda ^{},s,b)`$ (5.51)
$`\times `$ $`F_{\mathrm{remn}}\left(x^{\mathrm{proj}}{\displaystyle \underset{\lambda }{}}\stackrel{~}{x}_\lambda ^+\right)F_{\mathrm{remn}}\left(x^{\mathrm{targ}}{\displaystyle \underset{\lambda }{}}\stackrel{~}{x}_\lambda ^{}\right),`$
which allows to write the profile functions eq. (5.43) as
$`\gamma _{h_1h_2}(s,b)`$ $`=`$ $`{\displaystyle \underset{m}{}}{\displaystyle \frac{1}{m!}}{\displaystyle \underset{\mu =1}{\overset{m}{}}dx_\mu ^+dx_\mu ^{}d\widehat{x}_\mu ^+d\widehat{x}_\mu ^{}\underset{\mu =1}{\overset{m}{}}\widehat{G}_{\mathrm{I}\mathrm{P}}^{h_1h_2}(x_\mu ^+,x_\mu ^{},\widehat{x}_\mu ^+,\widehat{x}_\mu ^{},s,b)}`$ (5.52)
$`\times `$ $`\mathrm{\Phi }_{h_1h_2}(x^{\mathrm{proj}}{\displaystyle \underset{\mu }{}}\widehat{x}_\mu ^+,x^{\mathrm{targ}}{\displaystyle \underset{\mu }{}}\widehat{x}_\mu ^{},s,b),`$
which may be approximated as
$`\gamma _{h_1h_2}(s,b)`$ $`=`$ $`{\displaystyle \underset{m}{}}{\displaystyle \frac{1}{m!}}{\displaystyle \underset{\mu =1}{\overset{m}{}}dx_\mu ^+dx_\mu ^{}\underset{\mu =1}{\overset{m}{}}\widehat{\widehat{G}}_{\mathrm{I}\mathrm{P}}^{h_1h_2}(x_\mu ^+,x_\mu ^{},x^{\mathrm{proj}},x^{\mathrm{targ}},s,b)}`$ (5.53)
$`\times \mathrm{\Phi }_{h_1h_2}(x^{\mathrm{proj}},x^{\mathrm{targ}},s,b),`$
with
$`\widehat{\widehat{G}}_{\mathrm{I}\mathrm{P}}^{h_1h_2}(x^+,x^{},x^{\mathrm{proj}},x^{\mathrm{targ}},s,b)={\displaystyle \frac{1}{F_{\mathrm{remn}}\left(x^{\mathrm{proj}}\right)F_{\mathrm{remn}}\left(x^{\mathrm{targ}}\right)}}`$ (5.54)
$`\times {\displaystyle }d\widehat{x}^+d\widehat{x}^{}\widehat{G}_{\mathrm{I}\mathrm{P}}^{h_1h_2}(x^+,x^{},\widehat{x}^+,\widehat{x}^{},s,b)F_{\mathrm{remn}}(x^{\mathrm{proj}}\widehat{x}^+)F_{\mathrm{remn}}(x^{\mathrm{targ}}\widehat{x}^{}).`$
Based on eq. (5.53), we proceed as in the eikonal case. We unitarize the theory by replacing $`\mathrm{\Phi }_{h_1h_2}`$ by $`\mathrm{\Phi }_{\mathrm{u}_{h_1h_2}}`$ in a complete analogy to the eikonal model. The numerical Markov chain procedures have to be modified slightly due to the fact that the $`\widehat{\widehat{G}}`$ functions contain $`x^{\mathrm{proj}}`$ and $`x^{\mathrm{targ}}`$ as arguments. One can no longer restrict oneself to considering one single site of the interaction matrix, since changing $`x^+`$ and $`x^{}`$modifies as well $`x^{\mathrm{proj}}`$ and $`x^{\mathrm{targ}}`$ and affects therefore the other sites as well (in case of nucleus-nucleus all the sites related to the same projectile and target nucleon). But this does not pose major problems.
### 5.3 Important Features of Enhanced Diagrams
The amplitude corresponding to a single Pomeron exchange as shown at fig. 5.7(a) behaves as a function of energy approximately as $`s^\mathrm{\Delta }`$, where $`s`$ is the c.m. energy squared for hadron-hadron interaction and $`\mathrm{\Delta }`$ is some effective exponent. At the same moment, the amplitude of the so-called $`Y`$-diagram shown in fig. 5.7(b) increases asymptotically as $`s^{2\mathrm{\Delta }}`$ , as can be seen from eq. (5.30). The amplitude corresponding to the diagram in fig. 5.7(c) behaves as $`s^{3\mathrm{\Delta }}`$.
This indicates very important property of enhanced diagrams, namely that they increase with energy much faster than the usual eikonal ones.
In the following, we discuss exclusively enhanced $`Y`$-diagrams for given hadron types $`h_1`$ and $`h_2`$, and to simplify the notation, we use simply $`G`$ to refer to the corresponding profile function, omitting all indices referring to the initial hadron types and the Pomeron type.
In order to calculate contributions of enhanced graphs to the total interaction cross section, one has to consider different cuts of elastic scattering diagrams. We have
$$G\frac{1}{2s}2\mathrm{I}\mathrm{m}\stackrel{~}{T}(s,t=0)=G_{(0)}+G_{(1)}+G_{(2)},$$
(5.55)
where $`G_{(i)}`$ refers to the different cut diagrams, as shown in fig. 5.8.
We use our convention employed already earlier to plot cut Pomerons as dashed and uncut ones as full vertical lines. The diagram in fig. 5.8(a) gives rise to the process of high mass target diffraction ($`G_{(0)}`$), the diagram in fig. 5.8(b) represents the screening correction to the one cut Pomeron process ($`G_{(1)}`$), and the diagram in fig 5.8(c) - the cut Pomeron fusion process ($`G_{(2)}`$). As discussed in the preceding section, we have
$`G_{(0)}`$ $`=`$ $`1G,`$
$`G_{(1)}`$ $`=`$ $`+4G,`$ (5.56)
$`G_{(2)}`$ $`=`$ $`2G,`$
the sum being therefore equal to $`G`$, as it should be. Since $`G`$ is negative, the first order contribution to the inelastic cross section is negative, so another very important property of enhanced diagrams is the suppression of the increase of the inelastic cross section with energy.
A remarkable feature of enhanced diagrams is connected to their effect on the inclusive particle spectra. In particular, if one assumes (for a qualitative discussion) that each cut Pomeron gives rise to a flat rapidity distribution of secondary hadrons, $`dn/dy=\rho _0`$, the sum of all three contributions gives the screening correction to the inclusive particle spectrum as
$$\frac{d\mathrm{\Delta }n}{dy}1G\times \mathrm{\hspace{0.17em}0}\rho _0+4G\times \mathrm{\hspace{0.17em}1}\rho _02G\times \mathrm{\hspace{0.17em}2}\rho _0=0,y>y_0$$
(5.57)
and
$$\frac{d\mathrm{\Delta }n}{dy}1G\times \rho _0+4G\times \rho _02G\times \rho _0=G\rho _0<0,y<y_0,$$
(5.58)
where $`y_0`$ is the rapidity position of the triple Pomeron vertex, see fig. 5.9. The contribution is negative, because $`G`$ is negative.
Thus enhanced diagrams give rise to screening corrections to secondary hadron spectra only in restricted regions of the kinematical phase space; contributions of different cuts exactly cancel each other in the region of rapidity space where two or more Pomerons are exchanged in parallel (in case of the diagram of fig. 5.9 for $`y>y_0`$) . Therefore another important effect of Pomeron-Pomeron interactions is the modification of secondary hadron spectra, being mainly suppressed in the fragmentation regions of the interaction: close to $`y=0`$ from the $`Y`$-diagrams and close to $`y_{\mathrm{max}}`$ for the inverted $`Y`$-diagrams. That explains the great importance of enhanced diagrams for the solution of the unitarity problems inherent for the pure eikonal scheme and for the construction of a consistent unitary approach to hadronic interactions at very high energies (see the discussion in chapter 3).
Another important effect is the considerable increase of fluctuations of hadronic interactions. Let us for a moment consider the indices indicating Pomeron types: $`G_{3\mathrm{I}\mathrm{P}+}`$ for the $`Y`$-diagram as discussed above, $`G_{3\mathrm{I}\mathrm{P}}`$ for the corresponding inverted diagram, and $`G_{1\mathrm{I}\mathrm{P}}(s,x^+,x^{})`$ for the normal Pomeron. The full contribution (so far) is
$$G_{\mathrm{I}\mathrm{P}}=G_{1\mathrm{I}\mathrm{P}}+G_{3\mathrm{I}\mathrm{P}}+G_{3\mathrm{I}\mathrm{P}+}.$$
(5.59)
Using the fact that $`G_{3\mathrm{I}\mathrm{P}\pm }`$ can be written as the sum of the different cut contributions $`G_{3\mathrm{I}\mathrm{P}\pm (i)}`$, we get
$`G_{\mathrm{I}\mathrm{P}}`$ $`=`$ $`G_{1\mathrm{I}\mathrm{P}}+{\displaystyle \underset{i}{}}G_{3\mathrm{I}\mathrm{P}(i)}+{\displaystyle \underset{i}{}}G_{3\mathrm{I}\mathrm{P}+(i)},`$
which may be written as
$`G_{\mathrm{I}\mathrm{P}}`$ $`=`$ $`\left\{G_{1\mathrm{I}\mathrm{P}}+G_{3\mathrm{I}\mathrm{P}(1)}+G_{3\mathrm{I}\mathrm{P}+(1)}\right\}+\left\{G_{3\mathrm{I}\mathrm{P}(0)}+G_{3\mathrm{I}\mathrm{P}+(0)}\right\}+\left\{G_{3\mathrm{I}\mathrm{P}(2)}+G_{3\mathrm{I}\mathrm{P}+(2)}\right\}.`$
This means that we have three contributions: a modified one cut Pomeron exchange, with probability $`w_{\mathrm{one}}`$, the high mass target (see fig. 5.8(a)) and projectile diffraction, with probability $`w_{\mathrm{diff}}`$, and the process of Pomeron fusion of fig. 5.8(c), with probability $`w_{\mathrm{fusion}}`$, where the probabilities are given as
$`w_{\mathrm{one}}={\displaystyle \frac{G_{1\mathrm{I}\mathrm{P}}+G_{3\mathrm{I}\mathrm{P}(1)}+G_{3\mathrm{I}\mathrm{P}+(1)}}{G_{\mathrm{I}\mathrm{P}}}},w_{\mathrm{diff}}={\displaystyle \frac{G_{3\mathrm{I}\mathrm{P}(0)}+G_{3\mathrm{I}\mathrm{P}+(0)}}{G_{\mathrm{I}\mathrm{P}}}},w_{\mathrm{fusion}}={\displaystyle \frac{G_{3\mathrm{I}\mathrm{P}(2)}+G_{3\mathrm{I}\mathrm{P}+(2)}}{G_{\mathrm{I}\mathrm{P}}}}.`$
The two latter processes result in correspondingly much smaller and much larger values of the hadron multiplicity than for the one Pomeron process.
The problem of consistent treatment of Pomeron-Pomeron interactions was addressed already in . The number of diagrams which contribute essentially to the interaction characteristics increases fast with the energy. Therefore, one has to develop a suitable method to take into account the necessary contributions to the forward scattering amplitude, the latter one being related via the optical theorem to the total cross section of the reaction and to the weights for particular configurations of the interaction (via the “cutting” procedure). Such a scheme is still under development and our current goal was the proper treatment of some lowest order enhanced diagrams. Thus we proposed a minimal modification of the standard eikonal scheme, which allowed us to obtain a consistent description of hadronic interactions in the range of c.m. energies from some ten GeV till few thousand TeV. Already this minimal scheme allows to achieve partly the goals mentioned above: the slowing down of the energy increase of the interaction cross section and the non-AGK-type modification of particle spectra, as well as the improvement of the description of the multiplicity and inelasticity fluctuations in hadron-hadron interaction.
An important question exists concerning the nature of the triple-Pomeron coupling. As discussed above, we used the perturbative treatment for the part of a parton cascade developing in the region of parton virtualities bigger than some cutoff $`Q_0^2`$, whereas the region of smaller virtualities is treated phenomenologically, based on the soft Pomeron. There was an argumentation in that the triple Pomeron coupling is perturbative and therefore can be described on the basis of the QCD techniques. At the same time, it was shown in that such perturbative coupling, corresponding to non-small parton virtualities, would result in negligible contribution to the basic interaction characteristics, in particular, to the proton structure function $`F_2`$. The latter result was confirmed experimentally by HERA measurements, where no real shoulder in the behavior of $`F_2(x,Q^2)`$ (predicted in ) was found in the limit $`x0`$. Another argument in favor of the smallness of the perturbative Pomeron-Pomeron coupling comes from the HERA diffractive data, where the proportion of diffractive type events (with a large rapidity gap in secondary hadron spectra) appeared to be nearly independent on the virtuality of the virtual photon probe. This implies that that the Pomeron self-interaction is rather inherent to the non-perturbative initial condition for the QCD evolution than to the dynamical evolution itself. Therefore we assumed that the Pomerons interact with each other in the non-perturbative region of parton virtualities $`Q^2<Q_0^2`$ and considered it as the interaction between soft Pomerons. In our scheme the relatively big value of the soft triple-Pomeron coupling results in the screening corrections which finally prevent the large increase of parton densities in the small $`x`$ limit and restore the unitarity, thus leaving a little room for higher twist effects in the perturbative part of the interaction.
## Chapter 6 Parton Configurations
In this section, we consider the generation of parton configurations in nucleus-nucleus (including proton-proton) scattering for a given interaction configuration, which has already been determined, as discussed above. So, the numbers $`m_k`$ of elementary interactions per nucleon-nucleon pair $`k`$ are known, as well as the light cone momentum fractions $`x_{k\mu }^+`$ and $`x_{k\mu }^{}`$ of each elementary interaction of the pair $`k`$. A parton configuration is specified by the number of partons, their types and momenta.
### 6.1 General Procedure of Parton Generation
We showed earlier that the inelastic cross section may be written as
$$\sigma _{\mathrm{inel}}=d^2b\underset{m}{}𝑑X^+𝑑X^{}\mathrm{\Omega }(m,X^+,X^{}),$$
(6.1)
where $`\{m,X^+,X^{}\}`$ represents an interaction configuration. The function $`\mathrm{\Omega }`$ is known (see eq. (2.72)) and is interpreted as the probability distribution for interaction configurations. For each individual elementary interaction a term $`G_{1\mathrm{I}\mathrm{P}}^{h_1h_2}`$ appears in the formula for $`\mathrm{\Omega }`$, where the function $`G_{1\mathrm{I}\mathrm{P}}^{h_1h_2}`$ itself can be expressed in terms of contributions of different parton configurations. Namely the $`QCD`$ evolution function $`E_{\mathrm{QCD}}^{ij}`$, which enters into the formula for the elementary interaction contribution $`G_{1\mathrm{I}\mathrm{P}}^{h_1h_2}`$, is the solution of a ladder equation, where adding a ladder rung corresponds to an integration over the momenta of the corresponding resolvable parton emitted. The complete evolution function is therefore a sum over $`n`$-rung ladder contributions, where the latter one can be written as an integration over $`n`$ parton momenta. So we have
$$\mathrm{\Omega }(m,X^+,X^{})=\underset{k=1}{\overset{AB}{}}\underset{\mu =1}{\overset{m_k}{}}\underset{\tau =1}{\overset{t}{}}\underset{\nu =1}{\overset{n_{k\mu \tau }}{}}d^3p_{k\mu \tau \nu }\mathrm{\Psi }(\{p_{k\mu \tau \nu }\}),$$
(6.2)
where $`t`$ is the number of Pomeron types (soft, sea-sea, …), and $`n_{k\mu \tau }`$ the number of partons for the $`\mu ^{\mathrm{th}}`$ interaction of the pair $`k`$ in case of Pomeron type $`\tau `$. We interpret
$$\frac{\mathrm{\Psi }(\{p_{k\mu \tau \nu }\})}{\mathrm{\Omega }(m,X^+,X^{})}$$
(6.3)
as the probability distribution for parton configurations for a given interaction configuration $`\{m,X^+,X^{}\}`$. The Monte Carlo method provides a convenient tool for treating such multidimensional distributions: with $`\mathrm{\Theta }`$ being known (see chapter 2 and the discussion below), one generates parton configurations according to this distribution.
We want to stress that the parton generation is also based on the master formula eq. (2.72), no new elements enter.
In the following, we want to sketch the generation of parton configurations, technical details are provided in the next section. Let us consider a particular elementary interaction with given light cone momentum fractions $`x^+`$ and $`x^{}`$ and given impact parameter difference $`b`$ between the corresponding pair of interacting nucleons, for a fixed primary energy squared $`s`$.
For the sake of simplicity, we discuss here the procedure without the triple-Pomeron contribution, with the corresponding generalization being done in appendix C.4. We have to start with specifying the type of elementary interaction (soft, semi-hard, or valence type). The corresponding probabilities are
$`G_{\mathrm{soft}}^{h_1h_2}(x^+,x^{},s,b)/G_{1\mathrm{I}\mathrm{P}}^{h_1h_2}(x^+,x^{},s,b),`$
$`G_{\mathrm{sea}\mathrm{sea}}^{h_1h_2}(x^+,x^{},s,b)/G_{1\mathrm{I}\mathrm{P}}^{h_1h_2}(x^+,x^{},s,b),`$
$`G_{\mathrm{val}\mathrm{val}}^{h_1h_2}(x^+,x^{},s,b)/G_{1\mathrm{I}\mathrm{P}}^{h_1h_2}(x^+,x^{},s,b),`$ (6.4)
$`G_{\mathrm{val}\mathrm{sea}}^{h_1h_2}(x^+,x^{},s,b)/G_{1\mathrm{I}\mathrm{P}}^{h_1h_2}(x^+,x^{},s,b),`$
$`G_{\mathrm{sea}v\mathrm{al}}^{h_1h_2}(x^+,x^{},s,b)/G_{1\mathrm{I}\mathrm{P}}^{h_1h_2}(x^+,x^{},s,b).`$
In the case of a soft elementary interaction, no perturbative parton emission takes place. Therefore we are left with the trivial parton configuration, consisting from the initial active partons – hadron constituents – to which the Pomeron is attached.
Let us now consider a semi-hard contribution. We obtain the desired probability distributions from the explicit expressions for $`G_{\mathrm{sea}\mathrm{sea}}^{h_1h_2}`$. For given $`x^+`$, $`x^{}`$, and $`b`$ we have
$`G_{\mathrm{sea}\mathrm{sea}}^{h_1h_2}(x^+,x^{},s,b)`$ $``$ $`{\displaystyle _0^1}𝑑z^+{\displaystyle _0^1}𝑑z^{}{\displaystyle \underset{ij}{}}E_{\mathrm{soft}}^i\left(z^+\right)E_{\mathrm{soft}}^j\left(z^{}\right)\sigma _{\mathrm{hard}}^{ij}(z^+z^{}\widehat{s},Q_0^2)`$
$`\times `$ $`{\displaystyle \frac{1}{4\pi \lambda _{\mathrm{soft}}^{h_1h_2}(1/(z^+z^{}))}}\mathrm{exp}\left({\displaystyle \frac{b^2}{4\lambda _{\mathrm{soft}}^{h_1h_2}\left(1/(z^+z^{})\right)}}\right)F_{\mathrm{part}}^{h_1}(x^+)F_{\mathrm{part}}^{h_2}(x^{})`$
with $`\widehat{s}=x^+x^{}s`$, $`\lambda _{\mathrm{soft}}^{h_1h_2}(\xi )=R_{h_1}^2+R_{h_2}^2+\alpha _{_{\mathrm{soft}}}^{}\mathrm{ln}\xi ,`$ and
$`\sigma _{\mathrm{hard}}^{ij}(\widehat{s},Q_0^2)`$ $`=`$ $`K{\displaystyle \underset{kl}{}}{\displaystyle 𝑑x_B^+𝑑x_B^{}𝑑p_{}^2\frac{d\sigma _{\mathrm{Born}}^{kl}}{dp_{}^2}(x_B^+x_B^{}\widehat{s},p_{}^2)}`$ (6.6)
$`\times `$ $`E_{\mathrm{QCD}}^{ik}(x_B^+,Q_0^2,M_F^2)E_{\mathrm{QCD}}^{jl}(x_B^{},Q_0^2,M_F^2)\theta \left(M_F^2Q_0^2\right),`$
representing the perturbative parton-parton cross section, where both initial partons are taken at the virtuality scale $`Q_0^2`$; we choose the factorization scale as $`M_\mathrm{F}^2=p_{}^2/4`$. The integrand of eq. (6.1) serves as the probability distribution to generate the momentum fractions $`x_1^\pm =x^\pm z^\pm `$ and the flavors $`i`$ and $`j`$ of the initial partons for the parton ladder.
The knowledge of the initial conditions – the momentum fractions $`x_1^\pm `$ and the starting virtuality $`Q_0^2`$ for the “first partons” as well as the flavors $`i`$ and $`j`$ – allows us to reconstruct the complete ladder, based on the eq. (6.6) and on the evolution equations (B.28-B.29) for $`E_{\mathrm{QCD}}^{ij}`$. To do so, we generalize the definition of the parton-parton cross section $`\sigma _{\mathrm{hard}}^{ij}`$ to arbitrary virtualities of the initial partons, defining
$`\sigma _{\mathrm{hard}}^{ij}(\widehat{s},Q_1^2,Q_2^2)`$ $`=`$ $`K{\displaystyle \underset{kl}{}}{\displaystyle 𝑑x_B^+𝑑x_B^{}𝑑p_{}^2\frac{d\sigma _{\mathrm{Born}}^{kl}}{dp_{}^2}(x_B^+x_B^{}\widehat{s},p_{}^2)}`$
$`\times `$ $`E_{\mathrm{QCD}}^{ik}(x_B^+,Q_1^2,M_\mathrm{F}^2)E_{\mathrm{QCD}}^{jl}(x_B^{},Q_2^2,M_\mathrm{F}^2)\mathrm{\Theta }\left(M_\mathrm{F}^2\mathrm{max}[Q_1^2,Q_2^2]\right)`$
and
$`\sigma _{\mathrm{ord}}^{ij}(\widehat{s},Q_1^2,Q_2^2)`$ $`=`$ $`K{\displaystyle \underset{k}{}}{\displaystyle 𝑑x_B^+𝑑x_B^{}𝑑p_{}^2\frac{d\sigma _{\mathrm{Born}}^{kj}}{dp_{}^2}(x_B^+x_B^{}\widehat{s},p_{}^2)}`$
$`\times `$ $`E_{\mathrm{QCD}}^{ik}(x_B^+,Q_1^2,M_\mathrm{F}^2,w^+)\mathrm{\Delta }^j(Q_2^2,M_\mathrm{F}^2)\mathrm{\Theta }\left(M_\mathrm{F}^2\mathrm{max}[Q_1^2,Q_2^2]\right)`$
representing the full ladder contribution ($`\sigma _{\mathrm{hard}}`$) and the contribution, corresponding to the ordering of parton virtualities towards the end of the ladder, i.e. to the case of parton $`j`$, involved into the highest virtuality Born process ($`\sigma _{\mathrm{ord}}`$). We calculate and tabulate $`\sigma _{\mathrm{hard}}`$ and $`\sigma _{\mathrm{ord}}`$ initially, so that we can use them via interpolation to generate partons. The generation of partons is done in an iterative fashion based on the following equations:
$`\sigma _{\mathrm{hard}}^{ij}(\widehat{s},Q_1^2,Q_2^2)`$ $`=`$ $`{\displaystyle \underset{k}{}}{\displaystyle \frac{dQ^2}{Q^2}𝑑\xi \mathrm{\Delta }^i(Q_1^2,Q^2)\frac{\alpha _s}{2\pi }P_i^k(\xi )\sigma _{\mathrm{hard}}^{kj}(\xi \widehat{s},Q^2,Q_2^2)}`$
$`+`$ $`\sigma _{\mathrm{ord}}^{ji}(\widehat{s},Q_2^2,Q_1^2)`$
and
$`\sigma _{\mathrm{ord}}^{ij}(\widehat{s},Q_1^2,Q_2^2)`$ $`=`$ $`{\displaystyle \underset{k}{}}{\displaystyle \frac{dQ^2}{Q^2}𝑑\xi \mathrm{\Delta }^i(Q_1^2,Q^2)\frac{\alpha _s}{2\pi }P_i^k(\xi )\sigma _{\mathrm{ord}}^{kj}(\xi \widehat{s},Q^2,Q_2^2)}`$
$`+`$ $`\sigma _{\mathrm{Born}}^{ij}(\widehat{s},Q_1^2,Q_2^2)`$
Here, $`\sigma _{\mathrm{Born}}^{ij}`$ gives the contribution of the configuration without any resolvable emission before the highest virtuality Born process:
$`\sigma _{\mathrm{Born}}^{ij}(\widehat{s},Q_1^2,Q_2^2)`$ $`=`$ $`K{\displaystyle 𝑑p_{}^2\frac{d\sigma _{\mathrm{Born}}^{ij}}{dp_{}^2}(\widehat{s},p_{}^2)}`$
$`\times `$ $`\mathrm{\Delta }^i(Q_1^2,M_\mathrm{F}^2)\mathrm{\Delta }^j(Q_2^2,M_\mathrm{F}^2)\mathrm{\Theta }\left(M_\mathrm{F}^2\mathrm{max}[Q_1^2,Q_2^2]\right)`$
The procedure is described in detail in the next section.
In the case of elementary interactions involving valence quarks, the method is almost identical. In that case, the corresponding momentum fractions $`x_1^\pm `$ are the ones of valence quarks, $`x_1^\pm =x_q^\pm `$, to be determined according to the corresponding integrands in the expressions for $`G_{\mathrm{val}}^{h_1h_2}`$, $`G_{\mathrm{val}\mathrm{sea}}^{h_1h_2}`$, $`G_{\mathrm{sea}\mathrm{val}}^{h_1h_2}`$, see eqs. (2.48-2.50). For example, in the case of both hadron constituents being valence quarks, one generates momentum fractions $`x_1^\pm `$ and valence quark flavors $`i`$, $`j`$ with the distribution (up to a normalization constant)
$$D_{\mathrm{val}\mathrm{val}}^{h_1h_2,ij}(x_{q_v}^+x_{q_v}^{}s,b)\overline{F}_{\mathrm{part}}^{h_1,i}(x_{q_v}^+,x^+x_{q_v}^+)\overline{F}_{\mathrm{part}}^{h_2,j}(x_{q_v}^{},x^{}x_{q_v}^{}),$$
(6.12)
see eqs. (2.48), (2.18). One then proceeds to generate parton emissions as discussed above.
### 6.2 Generating the Parton Ladder
We now discuss in detail the generation of the partons in a ladder, starting from the initial partons (“leg partons”) with flavors $`i`$ and $`j`$ and light cone momentum fractions $`x_1^+`$ and $`x_1^{}`$. To simplify the discussion, we will neglect the effects of finite virtualities and transverse momenta of initial partons in the kinematical formulas so that the 4-momenta $`k_1`$ and $`k_1^{}`$ of the two leg partons are purely longitudinal. In the hadron-hadron (nucleus-nucleus) center of mass frame we have:
$$\begin{array}{cc}k_1^+=x_1^+\sqrt{s}/2,& k_1^{}=0,k_1_{}=0,\\ k_1^{}=x_1^{}\sqrt{s}/2,& k_1^+=0,k_1_{}^{}=0.\end{array}$$
(6.13)
The invariant mass squared of the ladder is $`\widehat{s}=(k_1+k_1^{})^2`$.
One first generates all resolvable partons emitted at one side of the ladder before the hardest process (for the definiteness we start with the leg parton $`i`$ moving in the forward direction). At each step one decides whether there is any resolvable emission at the forward end of the ladder before the hardest process. An emission is done with the probability
$$\mathrm{prob}(\mathrm{forward}\mathrm{emission})=\left(\sigma _{\mathrm{hard}}^{ij}(\widehat{s},Q_1^2,Q_2^2)\sigma _{\mathrm{ord}}^{ji}(\widehat{s},Q_2^2,Q_1^2)\right)/\sigma _{\mathrm{hard}}^{ij}(\widehat{s},Q_1^2,Q_2^2).$$
(6.14)
In case of an emission, the generation of light cone momentum fraction $`\xi `$ and momentum transfer squared $`Q^2`$ for the current parton branching is done – up to a normalization constant – according to the integrand of $`(\sigma _{\mathrm{hard}}^{ij}\sigma _{\mathrm{ord}}^{ji})`$,
$$\mathrm{prob}(\xi ,Q^2)\frac{1}{Q^2}\mathrm{\Delta }^i(Q_1^2,Q^2)\frac{\alpha _s}{2\pi }\underset{i^{}}{}P_i^i^{}(\xi )\sigma _{\mathrm{hard}}^{i^{}j}(\xi \widehat{s},Q^2,Q_2^2)$$
(6.15)
see eq. (6.1). Here the emitted $`s`$-channel parton gets the share $`1\xi `$ of the parent (leg) parton light cone momentum $`k^+`$ and the transverse momentum squared $`p_{}^2(1\xi )Q^2`$ . To leading logarithmic accuracy, the initial parton virtuality is neglected in the branching probability eq. (6.15), because of $`Q_1^2Q^2`$. Generating randomly the polar angle $`\phi `$ for the emission, one reconstructs the 4-vector $`p`$ of the final $`s`$-channel parton as
$$p^+=(1\xi )k_1^+,p^{}=p_{}^2/((1\xi )k_1^+),\stackrel{}{p}_{}=\left(\begin{array}{c}p_{}\mathrm{cos}\phi \\ p_{}\mathrm{sin}\phi \end{array}\right).$$
(6.16)
The remaining ladder after the parton emission is now characterized by the mass squared $`\widehat{s}^{}=(k_1+k_1^{}p)^2\xi \widehat{s}`$ and the initial virtualities $`Q_{}^{}{}_{1}{}^{2}=Q^2`$ and $`Q_2^2`$. The flavor $`i^{}`$ of the new leg parton is generated according to the corresponding weights in eq. (6.15), properly normalized given as
$$\mathrm{prob}(i^{})=\frac{P_i^i^{}(\xi )\sigma _{\mathrm{hard}}^{i^{}j}(\widehat{s}^{},Q^2,Q_2^2)}{_lP_i^l(\xi )\sigma _{\mathrm{hard}}^{lj}(\widehat{s}^{},Q^2,Q_2^2)},$$
(6.17)
where $`\sigma _{\mathrm{hard}}^{i^{}j}(\xi \widehat{s},Q^2,Q_2^2)`$ is the parton cross section (6.1) for the new ladder. One then renames $`\widehat{s}^{}`$, $`i^{}`$, and $`Q_1^2`$ into $`\widehat{s}`$, $`i`$, and $`Q_1^2`$ and repeats the above procedure.
In case of no forward emission, the generation of all resolvable parton emissions at the forward side of the ladder has been completed.
One then proceeds to generate all resolvable parton emissions for the backward side of the ladder, starting from the original leg parton $`j`$ of virtuality $`Q_2^2=Q_0^2`$, by using a corresponding recursive algorithm, now based on eq. (6.1). On the other end of the ladder, we have (after renaming) parton $`i`$ with the virtuality $`Q_1^2`$. One decides whether there is any resolvable emission before the hardest process, where the probability of an emission is given as
$$\mathrm{prob}(\mathrm{backward}\mathrm{emission})=\left(\sigma _{\mathrm{ord}}^{ji}(\widehat{s},Q_2^2,Q_1^2)\sigma _{\mathrm{Born}}^{ij}(\widehat{s},Q_1^2,Q_2^2)\right)/\sigma _{\mathrm{ord}}^{ji}(\widehat{s},Q_2^2,Q_1^2).$$
(6.18)
In case of an emission, the generation of the fraction $`\xi `$ of the light cone momentum $`k_1^{}`$, and of the momentum transfer squared $`Q^2`$ is done – up to a normalization constant – according to the integrand of $`(\sigma _{\mathrm{ord}}^{ji}`$ \- $`\sigma _{\mathrm{Born}}^{ij})`$,
$$\mathrm{prob}(\xi ,Q^2)\frac{1}{Q^2}\mathrm{\Delta }^j(Q_2^2,Q^2)\frac{\alpha _s}{2\pi }\underset{j^{}}{}P_j^j^{}(\xi )\sigma _{\mathrm{ord}}^{j^{}i}(\xi \widehat{s},Q^2,Q_1^2),$$
(6.19)
see eq. (6.1). The flavor $`j^{}`$ of the new leg parton is defined according to the partial contributions in (6.19). The generation of resolvable parton emissions is completed when the iterative procedure stops, with the probability
$$1\mathrm{prob}(\mathrm{backward}\mathrm{emission}).$$
Note, that all parton emissions are simulated in the original Lorentz frame, where the original leg partons (the initial partons for the perturbative evolution) are moving along the $`z`$-axis.
The final step is the generation of the hardest $`22`$ parton scattering process. In the center of mass system of two partons $`i`$ and $`j`$ with center-of-mass energy squared $`\widehat{s}`$, we simulate the transverse momentum $`p_{}^2`$ for the scattering within the limits (given by the condition $`M_\mathrm{F}^2=p_{}^2/4>\mathrm{max}[Q_1^2,Q_2^2]`$)
$$4\mathrm{max}[Q_1^2,Q_2^2]<p_{}^2<\widehat{s}/4$$
(6.20)
according to
$$\mathrm{prob}(p_{}^2)\frac{d\sigma _{\mathrm{Born}}^{ij}}{dp_{}^2}(\widehat{s},p_{}^2)\mathrm{\Delta }^i(Q_1^2,p_{}^2/4)\mathrm{\Delta }^j(Q_2^2,p_{}^2),$$
(6.21)
where the differential parton-parton cross section is
$$\frac{d\sigma _{\mathrm{Born}}^{ij}}{dp_{}^2}(\widehat{s},p_{}^2)=\frac{1}{16\pi \widehat{s}^2\sqrt{14p_{}^2/\widehat{s}}}\underset{k,l}{}\left|M^{ijkl}(\widehat{s},p_{}^2)\right|^2$$
(6.22)
with $`|M^{ijk,l}(\widehat{s},p_{}^2)|^2`$ being the squared matrix elements of the parton subprocesses .
Then we choose a particular subprocess $`ijkl`$ according to its contribution to the differential cross section (6.22), and reconstruct the 4-momenta $`p_1`$ and $`p_2`$ of the final partons in their center of mass system as
$$\begin{array}{c}p_1^+=z\sqrt{\widehat{s}},p_1^{}=p_{}^2/(z\sqrt{\widehat{s}}),\stackrel{}{p}_1=\left(\begin{array}{c}p_{}\mathrm{cos}\phi \\ p_{}\mathrm{sin}\phi \end{array}\right),\\ p_2^+=(1z)\sqrt{\widehat{s}},p_2^{}=p_{}^2/((1z)\sqrt{\widehat{s}}),\stackrel{}{p}_2=\left(\begin{array}{c}p_{}\mathrm{cos}\phi \\ p_{}\mathrm{sin}\phi \end{array}\right),\end{array}$$
(6.23)
with
$$z=\frac{1}{2}\left(1+\sqrt{14p_{}^2/\widehat{s}}\right)$$
(6.24)
and a random polar angle $`\phi `$. We finally boost the momenta to the original Lorentz frame.
### 6.3 The Time-like Parton Cascade
The above discussion of how to generate parton configurations is not yet complete: the emitted partons are in general off–shell and can therefore radiate further partons. This so-called time-like radiation is taken into account using standard techniques , to be discussed in the following.
The parton emission from an off-shell parton is done using the so-called DGLAP evolution equations, which describes the process with the leading logarithmic accuracy. The splitting probability for the initial parton of type $`j`$ is then given as
$$\frac{d𝒫}{𝒫}=\frac{dQ^2}{Q^2}𝑑z\underset{k}{}\frac{\alpha _s(p_{}^2)}{2\pi }P_j^k(z),$$
(6.25)
with the usual Altarelli-Parisi splitting functions $`P_j^k(z)`$. Here $`Q^2=Q_j^2`$ is the virtuality of the parent parton $`j`$
and $`z`$ is interpreted as the energy fraction carried away by the daughter parton $`k`$. The maximum possible virtuality $`q_{j\mathrm{max}}^2`$ of the parton $`j`$ is given by the virtuality of the parent of $`j`$. One imagines now to decrease the virtuality of $`j`$, starting from the maximum value, such that the DGLAP evolution equations give then the probability $`dP`$ that during a change $`dQ^2`$ of the virtuality, a parton splits into two daughter partons $`k`$ and $`l`$ – see fig. 9.1. For the energies of the daughter partons one has
$$E_k=zE_j,E_l=(1z)E_j.$$
(6.26)
Choosing a frame where $`p_j_{}=0,`$ we have
$$p_k_{}=p_l_{}=p_{},$$
(6.27)
with
$$p_{}^2=\frac{E_j^2\left(z(1z)Q_j^2zQ_l^2(1z)Q_k^2\right)\frac{1}{4}\left(Q_l^2Q_j^2Q_k^2\right)^2+Q_j^2Q_k^2}{E_j^2Q_j^2}.$$
(6.28)
As usual, a cutoff parameter terminates the cascade of parton emissions. We introduce a parameter $`p_{\mathrm{fin}}^2`$, which represents the lower limit for $`p_{}^2`$ during the evolution, from which we obtain the lower limit for the virtualities as $`Q_{\mathrm{min}}^2=4p_{\mathrm{fin}}^2`$.
Let us provide some technical details on the splitting procedure. Using eq. (6.25) and applying the rejection method proposed in we determine the variables $`Q_j^2`$ and $`z`$ as well as the flavors $`k,l`$ of the daughter partons – see appendix B.3. To determine the 4-momentum $`p_j`$ of the parent parton we have to distinguish between three different modes of branching (see figure 6.2):
1. If $`j`$ is the initial parton for the time-like cascade, the energy $`E_j`$ is given, and with the obtained value $`Q_j^2`$, we can calculate $`|\stackrel{}{p}_j|=\sqrt{E_j^2Q_j^2}`$. The direction of $`\stackrel{}{p}_j`$ is obtained from the momentum conservation constraint for the summary momentum of all partons produced in the current time-like cascade.
2. If $`j`$ is the initial parton for the time-like cascade resulted from the Born process, together with a partner parton $`j^{}`$, then we know the total energy $`E_j+E_j^{}`$ for the two partons. After obtaining the virtualities $`Q_j^2`$ and $`Q_j^{}^2`$ from the secondary splittings $`jk,l`$ and $`j^{}k^{},l^{}`$ we can use the momentum conservation constraint in the parton-parton center of mass frame, which gives $`\stackrel{}{p}_j=\stackrel{}{p}_j^{}=\stackrel{}{p}`$ and allows to determine $`|\stackrel{}{p}|`$. The direction of $`\stackrel{}{p}`$ can be chosen randomly.
3. If parton $`j`$ is produced in a secondary time-like branching, its energy $`E_j=zE_{\mathrm{parent}}`$ as well as the energy of the second “daughter” $`E_j^{}=(1z)E_{\mathrm{parent}}`$ are known from the previous branching, as well as $`Q_j^2`$ and $`Q_j^{}^2`$ – from the splittings $`jk,l`$ and $`j^{}k^{},l^{}`$, which allows to determine $`\stackrel{}{p}_j_{}=\stackrel{}{p}_j_{}^{}=\stackrel{}{p}_{}`$ using eq. (6.28) in the frame where the parent parton moves along the $`z`$ -axis. In some cases one gets $`p_{_{}}^2<0`$; then the current splitting of the parton $`j`$ is rejected and its evolution continues.
The described leading order algorithm is known to be not accurate enough for secondary hadron production, in particular it gives too high multiplicities of secondaries in $`e^+e^{}`$-annihilation. The method can be corrected if one takes into account the phenomenon of color coherence. The latter one appears if one considers some higher order corrections to the simplest leading logarithmic contributions, the latter ones being the basis for the usual Altarelli-Parisi evolution equations. In the corresponding treatment – so-called modified leading logarithmic approach – one essentially recovers the original scheme for the time-like parton cascading supplemented by the additional condition, the strict ordering of the emission angles in successive parton branchings. The appearance of the angular ordering can be explained in a qualitative way: (see ): if a transverse wavelength of an emitted gluon ($`f`$ or $`e`$ on figure 6.3) is larger than the separation between the two, this gluon cannot see the color charge of the parent ($`l`$) but only the total (much smaller) charge of the two partons ($`k+l=j`$) and the radiation is suppressed. In a Monte-Carlo model this can be easily realized by imposing the angular ordering condition via rejection . Thus, for each branching we check whether the angular ordering $`\theta _2<\theta _1`$ is valid, where $`\theta _1`$ is the angle of the previous branching, and reject the current splitting otherwise.
## Chapter 7 Hadronization
Till now, our discussion concerned exclusively partons, whereas “in the real world” one observes finally hadrons. It is the purpose of this chapter, to provide the link, i.e. to discuss how to calculate hadron production starting from partonic configurations, discussed in the previous chapters.
Hadron production is related to the structure of cut Pomerons. A cut Pomeron is in principle a sum over squared amplitudes of the type $`a+b\mathrm{hadrons},`$ integrated over phase space, with $`a`$ and $`b`$ being the nucleon constituents involved in the interaction. So far there was no need to talk about the details of the hadron production, which could be considered to be “integrated out”. In case of soft Pomerons, we used a parameterization of the whole object, based on general asymptotic considerations, which means that all the hadron production is hidden in the few parameters characterizing the soft Pomeron. In case of hard Pomerons, we discussed the explicit partonic structure of the corresponding diagram without talking about hadrons. This is justified based on the assumption that summing over hadronic final states is identical to summing over partonic final states, both representing complete sets of states. But although our ignorance of hadronic states so far was well justified, we finally have to be specific about the hadronic structure of the cut Pomerons, because these are hadronic spectra which are measured experimentally, and not parton configurations.
Lacking a rigorous theoretical treatment, we are going to use the same strategy as we used already for treating the soft Pomeron: we are going to present a “parameterization” of the hadronic structure of the cut Pomerons, as simple as possible with no unnecessary details, in agreement with basic laws of physics and basic experimental observations. We do not claim at all to understand the microscopic mechanism, so our parameterization, called “string model”, should not be considered as a microscopic hadronization model.
### 7.1 Hadronic Structure of Cut Pomerons
In order to develop our multiple scattering theory, we used a simple graphical representation of a cut Pomeron, namely a thick vertical line connecting the external legs representing nucleon components,
as shown in fig. 7.1. This simple diagram hides somewhat the fact that there is a complicated structure hidden in this Pomeron, and the purpose of this section is to discuss in particular the hadronic content of the Pomeron.
Let us start our discussion with the soft Pomeron. Based on Veneziano’s topological expansion one may consider a soft Pomeron as a “cylinder” i.e. the sum of all possible QCD diagrams having a cylindrical topology, see fig. 7.2.
As discussed in detail in chapter 2.2, the “nucleon components” mentioned earlier, representing the external legs of the diagram, are always quark-anti-quark pairs, indicated by a dashed line (anti-quark) and a full line (quark) in fig. 7.2. Important for the discussion of particle production are of course cut diagrams, therefore we show in fig. 7.2 a cut cylinder representing a cut Pomeron: the cut plane is shown as two vertical dotted lines. Let us consider the half-cylinder, for example, the one to the right of the cut, representing an inelastic amplitude.
We can unfold this object in order to have a planar representation, as shown in fig.7.3. Here, the dotted vertical lines indicate the cuts of the previous figure, and it is here where the hadronic final state hadrons appear. Lacking a theoretical understanding of this hadronic structure, we simply apply a phenomenological procedure, essentially a parameterization. We require the method to be as simple as possible, with a minimum of necessary parameters. A solution coming close to these demands is the so-called string model: each cut line is identified with a classical relativistic string, a Lorentz invariant string breaking procedure provides the transformation into a hadronic final state, see fig. 7.4.
The phenomenological microscopic picture which stays behind this procedure was discussed in a number of reviews : the string end-point partons resulted from the interaction appear to be connected by a color field. With the partons flying apart, this color field is stretched into a tube, which finally breaks up giving rise to the production of hadrons and to the neutralization of the color field.
We now consider a semi-hard Pomeron of the “sea-sea” type, where we have a hard pQCD process in the middle and a soft evolution at the end, see fig. 7.5.
We generalize the picture introduced above for the soft Pomeron. Again, we assume a cylindrical structure. For the example of fig. 7.5, we have the picture shown in fig. 7.6:
the shaded areas on the cylinder ends represent the soft Pomerons, whereas in the middle part we draw explicitly the gluon lines on the cylinder surface. We apply the same procedure as for the soft Pomeron: we cut the diagram and present a half-cylinder in a planar fashion, see fig. 7.6. We observe one difference compared to the soft case: there are three partons (dots) on each cut line: apart from the quark and the anti-quark at the end, we have a gluon in the middle. We again apply the string picture, but here we identify a cut line with a so-called kinky string, where the internal gluons correspond to internal kinks. The underlying microscopic picture will be presented by three color-connected partons - the gluon connected by the color field to the quark and to the anti-quark. The string model provides then a “parameterization” of hadron production, see fig. 7.7.
The procedure described above can be easily generalized to the case of complicated parton ladders involving many gluons and quark-anti-quark pairs. One should note that the treatment of semi-hard Pomerons is just a straightforward generalization of the string model for soft Pomerons, or one might see it the other way round: the soft string model is a natural limiting case of the kinky string procedure for semi-hard Pomerons.
We now need to discuss Pomerons of valence type. In case of “valence-valence” the first partons of the parton ladder are valence quarks,
there is no soft Pomeron between the parton ladder and the nucleon. The nucleon components representing the external legs are, as usual, quark-anti-quark pairs, but the anti-quark plays in fact just the role of a spectator. The simplest possible interaction is the exchange of two gluons, as shown in fig.7.8. We follow the scheme used for soft Pomerons and “sea-sea” type semi-hard Pomerons: we draw the diagram on a cylinder, see fig. 7.9. There is no soft region, the gluons couple directly to the external partons. We cut the cylinder, one gluon being to the right and one gluon to the left of the cut, and then we consider the corresponding half-cylinder presented in a planar fashion, see fig. 7.9 (right). Here, we have only internal gluons, on the cut line we observe just the external partons, the corresponding string is therefore just an ordinary quark-anti-quark string without internal kinks, as in the case of the soft Pomerons. We apply the usual string breaking procedure to obtain hadrons, see fig. 7.10.
Let us consider a more complicated valence-type diagram, as shown in fig. 7.11.
It is again a contribution to the Pomeron of the “valence-valence” type: the external partons of the parton ladders are the valence quarks of the nucleons. In contrast to the previous example, we have here an emission of s-channel gluons, traversing the cut. As usual, we present the diagram on the cylinder, as shown in fig. 7.12, where we also show the corresponding planar half-cylinder. In addition to internal gluons, we now observe also external ones, presented as dots on the cut line. As usual, we identify the cut line with a relativistic kinky string, where each external (s-channel) gluon represents a kink. We then employ the usual string procedure to produce hadrons, as sketched in fig. 7.13.
The general procedure should be clear from the above examples: in any case, no matter what type of Pomeron, we have the following procedure:
1. drawing of a cylinder diagram;
2. cutting the cylinder;
3. planar presentation of the half-cylinder;
4. identification of a cut line with a kinky string;
5. kinky string hadronization
The last point, the string hadronization procedure, will be discussed in detail in the following. This work is basically inspired by and further developed in . The main differences to are that hadrons are directly obtained from strings instead from low mass clusters, and an intrinsic transverse momentum is added to a string break-up. The main difference to the Lund-model is to use the area-law instead of a fragmentation function.
### 7.2 Lagrange Formalism for Strings
A string can be considered as a point particle with one additional space-like dimension. The trajectory in Minkowski space depends on two parameters:
$$x^\mu =x^\mu (\sigma ,\tau ),\sigma =0\mathrm{}\pi ,$$
(7.1)
with $`\sigma `$ being a space-like and $`\tau `$ a time-like parameter. In order to obtain the equation of motion, we need a Lagrangian. It is obtained by demanding the invariance of the trajectory with respect to gauge transformations of the parameters $`\sigma `$ and $`\tau `$. This way we find the Lagrangian of Nambu-Goto:
$$=\kappa \sqrt{(x^{}\dot{x})^2x^2\dot{x}^2},$$
(7.2)
with $`\dot{x}^\mu =dx^\mu /d\tau `$, $`x^\mu {}_{}{}^{}=dx^\mu /d\sigma `$ and $`\kappa `$ being the energy density or string tension. With this Lagrangian we write down the action
$$S=_0^\pi 𝑑\sigma _{\tau _0}^{\tau _1}𝑑\tau ,$$
(7.3)
which leads to the Euler-Lagrange equation:
$$\frac{}{d\tau }\frac{}{\dot{x}_\mu }+\frac{}{\sigma }\frac{}{x_\mu ^{}}=0,$$
(7.4)
with the initial conditions
$$\frac{}{x_\mu ^{}}=0,\sigma =0,\pi ,$$
(7.5)
since we have $`\delta x=0`$ for $`\tau =\tau _0`$ and $`\tau =\tau _1`$. This equation can be solved most easily by a partial gauge fixing. We have this freedom, since the result is independent on the choice of the parameters. This is done indirectly by imposing the following conditions:
$$\dot{x}^2+x^2=0,\dot{x}x^{}=0.$$
(7.6)
The Euler-Lagrange equation gives us a simple solution, the wave equation:
$$\frac{^2x_\mu }{\tau ^2}\frac{^2x_\mu }{\sigma ^2}=0,$$
(7.7)
with the following boundary conditions:
$$\frac{x_\mu }{\sigma }=0,\sigma =0,\pi .$$
(7.8)
The total momentum of a string is given by
$$p_{\mathrm{string}}^\mu =_C\frac{L}{\dot{x}_\mu }𝑑\sigma +\frac{L}{x_\mu ^{}}d\tau $$
(7.9)
with $`C`$ being a curve between the two ends of the string ($`\sigma =0`$ and $`\sigma =\pi `$). This gives for (7.6) and for $`d\tau =0`$
$$p_{\mathrm{string}}^\mu =_0^\pi \kappa \dot{x}^\mu 𝑑\sigma .$$
(7.10)
We still have to fix completely the gauge since it has been fixed partially only. This can be done with the following conditions for the parameter $`\tau `$:
$$n^\mu x_\mu =\lambda \tau $$
(7.11)
Different choices for $`n`$ and $`\lambda `$ are possible, like $`n=(1,1,0,0)`$ which is called the transverse gauge. We will use $`n=(1,0,0,0)`$ which leads to $`\lambda =E/\pi \kappa `$ and another choice $`\pi =E/\kappa `$ will identify $`\tau `$ with the time $`x_0`$, whereas $`E=_0^\pi \kappa \dot{x}_0𝑑\sigma `$ is the total energy of the string. We define “string units” via $`\kappa =1`$; $`\sigma `$ and $`\tau `$ have thereby the dimension of energy and $`\pi =E`$. In “ordinary” units, one has $`\kappa =\stackrel{~}{\kappa }`$ GeV/fm, with $`\stackrel{~}{\kappa }`$ being approximately 1, so a length of 1 GeV corresponds to 1 fm$`/\stackrel{~}{\kappa }`$ $``$1 fm.
The solution of a wave equation is a function which depends on the sum or the difference of the two parameters $`\sigma `$ and $`\tau `$. As the second derivative shows up, we have two degrees of freedom to impose the initial conditions on the space-like extension and the speed of the string at $`\tau =0`$. One can easily verify that the following Ansatz fulfills the wave equation (7.7):
$`x^\mu (\sigma ,\tau )`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[f^\mu \left(\sigma +\tau \right)+f^\mu \left(\sigma \tau \right)+{\displaystyle _{\sigma \tau }^{\sigma +\tau }}g^\mu \left(\xi \right)𝑑\xi \right]`$ (7.12)
$`f(\sigma )`$ $`=`$ $`x^\mu (\sigma ,\tau )|_{\tau =0}`$ (7.13)
$`g(\sigma )`$ $`=`$ $`\dot{x}^\mu (\sigma ,\tau )|_{\tau =0}`$ (7.14)
We identify the function $`f(\sigma )`$ with the initial spatial extension and $`g(\sigma )`$ with the initial speed of the string at the time $`\tau =0`$.
We will consider here a special class of strings, namely those with $`f=0`$ (initially point-like) and with a piecewise constant function $`g`$,
$$g(\sigma )=v_k\mathrm{for}E_{k1}\sigma E_k,\mathrm{\hspace{0.33em}1}kn$$
(7.15)
with some integer $`n`$. The set $`\{E_k\}`$ is a partition of the $`\sigma `$-range $`[0,E]`$,
$$0=E_0<E_1<\mathrm{}<E_{n1}<E_n=E,$$
(7.16)
and $`\{v_k\}`$ represents $`n`$ constant 4-vectors. Such strings are called kinky strings, with $`n`$ being the number of kinks, and the $`n`$ vectors $`v_k`$ being called kink velocities.
In order to use eq. (7.12), we have to extend the function $`g`$ beyond the physical range between $`0`$ and $`\pi `$. This can be done by using the boundary conditions, which gives
$`g(\tau )`$ $`=`$ $`g(\tau ),`$ (7.17)
$`g(\tau +2\pi )`$ $`=`$ $`g(\tau ),`$ (7.18)
So $`g`$ is a symmetric periodic function, with the period $`2\pi `$. This defines $`g`$ everywhere, and the eq. (7.12) is the complete solution of the string equation, expressed in terms of the initial condition $`g`$ ($`f`$ is taken to be zero). In case of kinky strings the latter is expressed in terms of the kink velocities $`\{v_k\}`$ and the energy partition $`\{E_k\}`$.
### 7.3 Identifying Partons and Kinks
We discussed earlier that a cut Pomeron may be identified with two sequences of partons of the type
$$qgg\mathrm{}g\overline{q},$$
(7.19)
representing all the partons on a cut line. We identify such a sequence with a kinky string, by requiring
$$\mathrm{parton}=\mathrm{kink},$$
(7.20)
which means that we identify the partons of the above sequence with the kinks of a kinky string, such that the partition of the energy is given by the parton energies,
$$E_k=\mathrm{energy}\mathrm{of}\mathrm{parton}k$$
(7.21)
and the kink velocities are just the parton velocities,
$$v_k=\frac{\mathrm{momentum}\mathrm{of}\mathrm{parton}k}{E_k}.$$
(7.22)
We consider massless partons, so that the energy is equal to the absolute value of the parton momentum.
Fig. 7.14 shows as an example the evolution of a kinky string representing three partons: a quark, an anti-quark, and a gluon, as a function of the time $`\tau `$. One sees that the partons start to move along their original direction with the speed of light. After some time which corresponds to their energy they take the direction of the gluon. One could say that they lose energy to the string itself. The gluon loses energy in two directions, to the quark and to the anti-quark and therefore in half the time. The ends of the string move continuously with the speed of each of the partons until the whole string is contracted again in one point. The cycle starts over.
Another example is shown on fig. 7.15, where realistic partons coming from a simulation of a $`e^+e^{}`$ annihilation process at 14 GeV c.m.s. energy are considered. We will see later how to generate these partons. We observe 6 partons, 2 quarks and 4 gluons, symbolically displayed in the first sub-figure. As the total energy is 14 GeV the cycle has a periodicity of 28 GeV. But one sees that the perturbative gluons play an important role in the beginning of the movement, and later from 2 GeV on, the longitudinal character dominates. As we will see later, a string breaks typically after $`1\text{ GeV}/\kappa `$ which gives much importance to the perturbative gluons.
### 7.4 Momentum Bands in Parameter Space
As we will see later, it is not necessary for a fragmentation model to know the spatial extension at each instant. Therefore, we concentrate on a description in momentum space which simplifies the model even more. By using formula (7.12) we can express the derivatives of $`x(\sigma ,\tau )`$ in terms of the initial conditions $`g(\sigma )`$ as
$`\dot{x}(\sigma ,\tau )`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[g(\sigma +\tau )+g(\sigma \tau )\right]`$ (7.23)
$`x^{}(\sigma ,\tau )`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[g(\sigma +\tau )g(\sigma \tau )\right].`$ (7.24)
Since the function $`g`$ is stepwise constant, we easily identify regions in the parameter space $`(\sigma ,\tau )`$, where $`g(\sigma +\tau )`$ is constant or where $`g(\sigma \tau )`$ is constant, as shown in fig. 7.16,7.17.
These regions are called momentum bands, more precisely R-bands and L-bands, being of great importance for the string breaking. If we overlay the two figures of 7.16,7.17, we get fig. 7.18,
which allows us to identify regions, where $`g(\sigma +\tau )`$ and $`g(\sigma \tau )`$ are constant at the same time, namely the intersections of R-bands and L-bands. In these areas $`\dot{x}`$ and $`x^{}`$ are constant, given as
$`\dot{x}(\sigma ,\tau )`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[v^{}+v^+\right]`$ (7.25)
$`x^{}(\sigma ,\tau )`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[v^{}v^+\right],`$ (7.26)
with $`v^+`$ and $`v^{}`$ being the velocities of the partons corresponding to the two intersecting bands.
Rather than considering a $`\sigma `$-range between $`\mathrm{}`$ and $`+\mathrm{}`$, one may simply consider the physical range from 0 to $`\pi `$, and construct the bands via reflection. As an example, let us follow the L-band corresponding to the parton $`i`$, starting at $`\tau =0`$. With increasing $`\tau `$ one reaches at some stage the border $`\sigma =0`$. Here, we have an intersection with the R-band, corresponding to the same parton $`i`$, coming from the unphysical region $`\sigma <0`$. We now follow this R-band, which corresponds to a reflection of the above-mentioned L-band, till we hit the border $`\sigma =\pi `$, … .
In the regions where $`g(\sigma +\tau )`$ and $`g(\sigma \tau )`$ have the same value, corresponding to collinear partons or to an overlap of the momentum bands of one and the same parton, one finds $`x^{}=0`$, i.e. there is no spatial extension in the dependence of the parameter $`\sigma `$. Therefore the coordinates $`x^\mu `$ stay unchanged and we recover the speed of the original partons. In particular, this is the case for the whole string at $`\tau =0`$, due to $`f=0`$. With the string evolving in time, more and more bands of non collinear partons overlap, which gives $`x^{}0`$; the string is extending as we have seen in fig. 7.14 until $`\tau =7`$.
### 7.5 Area Law
In order to consider string breaking, we are going to extend the model in a covariant fashion. We use the method proposed by Artru and Menessier , which is based on a simple extension of the decay law of unstable particles, where the probability $`dP`$ to decay within a time interval $`dt`$ is given as
$$dP=\lambda dt,$$
(7.27)
with some decay constant $`\lambda `$. For strings, we use the same formula by replacing the proper time by proper surface in Minkowski space,
$$dP=\lambda dA.$$
(7.28)
By construction, this method is covariant. Since we work in parameter space it is useful to express this dependence as a function of $`\sigma `$ and $`\tau `$,
$$dA=\sqrt{(\dot{x}x^{})^2x^2\dot{x}^2}d\sigma d\tau .$$
(7.29)
By using the expressions for $`\dot{x}`$ and $`x^{}`$ and $`\dot{x}x^{}=0`$ and $`g^2=0`$ , we find
$`dA`$ $`=`$ $`\sqrt{{\displaystyle \frac{1}{4}}\left(2g(\sigma +\tau )g(\sigma \tau )\right){\displaystyle \frac{1}{4}}\left(2g(\sigma +\tau )g(\sigma \tau )\right)}d\sigma d\tau `$ (7.30)
$`=`$ $`\left[{\displaystyle \frac{1}{2}}g(\sigma +\tau )g(\sigma \tau )\right]d\sigma d\tau `$ (7.31)
$`=`$ $`{\displaystyle \frac{1}{2}}\left(1\mathrm{cos}\varphi \right)d\sigma d\tau ,`$ (7.32)
with $`\varphi `$ being the angle between the partons. Consequently, a string cannot break at a point where the momentum bands of the same parton overlap, because in this case the angle $`\varphi `$ is zero, which leads to $`dA=0`$. The maximal contribution is obtained for partons moving in opposite directions.
We still have to define how a string breaks and how the sub-strings evolve. At each instant, one knows exactly the momenta of the string by eq. (7.25) and (7.26). The configuration of $`g(\sigma +\tau )`$ and $`g(\sigma \tau )`$ at the time $`\tau _1`$ of the break point is used as initial condition for the two substrings. The function $`g(\sigma +\tau )`$ is cut into two pieces between $`0`$ and $`\sigma _1`$ and between $`\sigma _1`$ and $`\pi `$. The two resulting functions are continued beyond their physical ranges $`[0,\sigma _1]`$ and $`[\sigma _1,\pi ]`$ by taking them to be symmetric and periodic with periods $`2\sigma _1`$ and $`2(\pi \sigma _1)`$. Fig. 7.19 shows this for a breaking at $`(\sigma _1,\tau _1)`$ and a second break point at $`(\sigma _2,\tau _2)`$.
In principle, the cycle starts over with the two sub-strings breaking each until the resulting pieces are light enough to form hadrons. However, it is easier to look for many break-points at once. If they are space-like separated, they do not interfere with each other. For the coordinates in the parameter space this translates into the condition
$$|\sigma _1\sigma _2|>|\tau _1\tau _2|.$$
(7.33)
### 7.6 Generating Break Points
Having assumed that string breaking occurs according to the area law,
$$dP=\lambda dA,$$
(7.34)
we now need an algorithm to accomplish this in the framework of the Monte-Carlo method. The most simple way is to sub-divide a given surface into sufficiently small pieces and then to decide according to formula (7.34) if there is a break point or not. This is what we refer to as the naive method, which is of course not efficient. We will therefore construct another algorithm (the direct one) which is based on
$$P_0(A)=e^{\lambda A}$$
(7.35)
being the probability of having no break point within the area $`A`$. One generates surfaces $`A_1`$, $`A_2`$, … according to $`P_0(A)`$ as
$$A_i=\mathrm{log}(r_i)/\lambda ,$$
(7.36)
with random numbers $`r_i`$ between $`0`$ and $`1`$. The formula does not say anything about the form of the surfaces $`A_i`$. Actually, several choices are possible as long as they do not violate causality. In the case of a simple string without any gluons, it is easiest to place the surfaces $`A_i`$ from left to right such that the break points $`P_i`$ are the left upper corners of $`A_{i+1}`$, as shown in fig. 7.20: one first takes $`A_1`$, which defines the line $`L_1`$. The first break point $`P_1`$ is generated randomly on this line. The next surface $`A_2`$ has to be placed in a way that does not violate causality. The first break point is therefore used as a constraint for the next one, etc. Finally, if the last surface obtained is too large to be placed on the rest of the string, the procedure is finished. The advantage of this method is that no break-points are rejected because the causality principle is obeyed constantly throughout the whole procedure.
We generalize the method to work for any number of perturbative gluons in the following way. Since the elementary invariant area $`dA`$ is proportional to the scalar product of the momenta of two partons, we can easily calculate the area $`A_{ij}`$ corresponding to a sub-region $`S_{ij}`$ of the $`(\sigma ,\tau )`$-space, representing the intersection of the momentum bands of the partons $`i`$ and $`j`$. We find
$`A_{ij}`$ $`=`$ $`{\displaystyle _{S_{ij}}}\left({\displaystyle \frac{1}{2}}g(\sigma +\tau )g(\sigma \tau )\right)𝑑\sigma 𝑑\tau ={\displaystyle \frac{1}{4}}p_ip_j,`$ (7.37)
with $`p_i`$ and $`p_j`$ being the 4-momenta of the two partons. We now construct the break points in the parameter space rather than in Minkowski space. One first defines the area in the parameter space of allowed breakpoints as
$$S_{\mathrm{break}}=S_{ij,}$$
(7.38)
with the indices running as indicated in fig. 7.21. To obtain a unique way of counting the regions, we mark bands which come from a left-moving band at $`\tau =0`$ with a . We further observe that the outer bands $`1`$ and $`1^{}`$ as well as $`5`$ and $`5^{}`$, which come from the (anti-)quark-momenta, are neighboring. It is therefore useful to redefine them as one band $`1`$ and $`5^{}`$ (with double momentum), see fig. 7.21.
For each of these sub-areas $`S_{ij}`$ the corresponding area in Minkowski space $`A_{ij}`$ is known $`(=p_ip_j/4)`$. One then generates areas $`A_1`$, $`A_2`$, $`A_3`$, …(in Minkowski space) according to eq. (7.36), and places the corresponding areas $`S_1`$, $`S_2`$, $`S_3`$, …(in parameter space) into $`S_{\mathrm{break}}`$ from left to right such that the break points $`P_i`$ are the left upper corners of $`S_{i+1}`$.
Let us consider an example of five partons (1,2,3,4,5), see fig. 7.22. Suppose that we have sampled a surface $`A_1`$. If it is smaller than the first region $`A_{12^{}}`$, we determine $`S_1=S_{12^{}}A_1/A_{12^{}}`$ and we place $`S_1`$ into the left side of $`S_{12^{}}`$ and generate the break point $`P_1`$ randomly on the right upper border of $`S_{12^{}}`$, see fig. 7.22.
If $`A_1`$ is greater than $`A_{12^{}}`$, we subtract $`A_{12^{}}`$ from $`A_1`$:
$$A_1^{}=A_1A_{12^{}}.$$
(7.39)
In the case of the sum of the three areas $`A_s=A_{2^{}3^{}}+A_{13^{}}+A_{23^{}}`$ being greater than $`A_1^{}`$, the first coordinate $`x`$ of the break point $`P_1`$ (see fig. 7.23)
is determined by
$$x=\frac{A_1^{}}{A_s}.$$
(7.40)
Otherwise we continue the procedure correspondingly. The $`y`$ coordinate is determined as
$$y=\{\begin{array}{cccc}r\frac{A_s}{A_{2^{}3^{}}}& \text{if}& 0<r<\frac{A_{2^{}3^{}}}{A_s}& (\text{region }S_{2^{}3^{}})\hfill \\ \left(r\frac{A_{2^{}3^{}}}{A_s}\right)\left(\frac{A_s}{A_{13^{}}}\right)& \text{if}& \frac{A_{2^{}3^{}}}{A_s}<r<\frac{A_{2^{}3^{}}+A_{13^{}}}{A_s}& (\text{region }S_{13^{}})\hfill \\ \left(r\frac{A_{2^{}3^{}}+A_{13^{}}}{A_s}\right)\left(\frac{A_s}{A_{23^{}}}\right)& \text{if}& \frac{A_{2^{}3^{}}+A_{13^{}}}{A_s}<r<1& (\text{region }S_{23^{}})\hfill \end{array},$$
(7.41)
with $`r`$ being a random number between $`0`$ and $`1`$. This means that after having determined in which of the regions we find the break point, it is placed randomly on the world-line which points to the future. After having obtained the break point $`P_1`$ we continue the procedure in the same way by obeying to the principle of causality. The area to sweep over is then limited by the first break point as shown on fig. 7.23.
In fig. 7.24, we apply our hadronization procedure, referred to as *direct method*, as discussed above, to calculate the distributions of $`xy`$, $`x`$, $`\eta =\frac{1}{2}\mathrm{log}\left(\frac{x}{y}\right),`$ and the multiplicity $`n`$ of break points for a quark-anti-quark string of $`E_q=E_{\overline{q}}=8GeV`$. We compare our results with the *the naive method,* where the area of the string is divided into small elements $`\mathrm{\Delta }A=88\text{ }\mathrm{GeV}^2/N^2`$, with $`N`$ sufficiently large to not change the results any more. In each of these elements, a break point is found with the probability $`\lambda \mathrm{\Delta }A`$. The points which are in the future of another one are rejected. The latter method is a literal realization of the area-law. As one can easily see in fig. 7.24, the two methods agree within statistical errors.
### 7.7 From String Fragments to Hadrons
So far, we discussed how to break a string into small pieces, i.e. string fragments with invariant masses between 0 and about two GeV. In order to identify string fragments and hadrons, we first have to define the flavors (= quark content) of the fragments, and then we have to discuss the question of fragment masses.
#### Flavors of String Fragments
A string as a whole has some flavor, carried by the partons at its two extremities. Additional flavor is created (by definition) at each break point in the form of a quark-anti-quark or a diquark-anti-diquark pair of a certain flavor. The corresponding probabilities are free parameters of the model. In case of quark-anti-quark formation, we introduce the parameter $`p_{\mathrm{ud}}`$, which gives the probability of flavor $`u`$ or $`d`$. The probability to get an $`s`$-quark is therefore $`12p_{\mathrm{ud}}`$ which is smaller than $`p_{\mathrm{ud}}`$ because on the larger mass of the $`s`$ quark. For diquark-anti-diquark production, we introduce the corresponding probability $`p_{\mathrm{diquark}}`$.
#### Masses of String Fragments
In the following, we show how to determine the masses of string fragments, characterized by break points in the parameter space.
Fig. 7.25 shows an example of two break points for a string with 3 inner kinks. The momentum bands and the regions of the their overlaps are shown: in case of the inner bands, we have three R-bands (2,3,4) and three L-bands (6,7,8). The bands at the extremities play a special role, since we may have the corresponding R- and L-band as just one band, due to the fact that one of the bands is reflected immediately. So, we consider two “double bands” (1 and 5).
The string momentum is given as
$$p_{\mathrm{string}}=_C\left[\dot{x}d\sigma +x^{}d\tau \right],$$
(7.42)
where $`C`$ is an arbitrary curve from one border ($`\sigma =0`$) to the other ($`\sigma =\pi `$) in the parameter space. This leads to
$$p_{\mathrm{string}}=\frac{1}{2}_C\left(g(\sigma +\tau )+g(\sigma \tau )\right)𝑑\sigma +\frac{1}{2}_C\left(g(\sigma +\tau )g(\sigma \tau )\right)𝑑\tau $$
(7.43)
The momenta of the bands are by definition
$$p_{(i)}=\{\begin{array}{c}\frac{1}{2}_{\mathrm{band}i}g(\sigma \tau )𝑑\sigma \frac{1}{2}_{\mathrm{band}i}g(\sigma \tau )𝑑\tau \mathrm{if}\mathrm{R}\mathrm{band}\hfill \\ \frac{1}{2}_{\mathrm{band}i}g(\sigma +\tau )𝑑\sigma +\frac{1}{2}_{\mathrm{band}i}g(\sigma +\tau )𝑑\tau \mathrm{if}\mathrm{L}\mathrm{band}\hfill \end{array},$$
(7.44)
where one integrates along an arbitrary curve from one border of the band to the other. An important property: an integration path parallel to a band provides zero contribution. One has to pay attention for the bands at the extremities: integrating only along $`\tau =0`$ represents only half the band. We have
$$\underset{i}{}p_{(i)}=p_{\mathrm{string}.}$$
(7.45)
The momenta of the bands are related to the corresponding parton momenta as
$$p_{(i)}=\{\begin{array}{ccc}\frac{1}{2}p_{\mathrm{parton}}\hfill & \mathrm{if}\hfill & \mathrm{inner}\mathrm{band}\hfill \\ p_{\mathrm{parton}}\hfill & \mathrm{if}\hfill & \mathrm{outer}\mathrm{band}\hfill \end{array},$$
(7.46)
which one verifies easily by expressing $`g`$ in terms of the parton momenta. The difference between inner and outer bands is due to the fact that the outer ones (at the extremities) represent in reality two bands. For the example of fig. 7.25, we have
$`p_{(1)}`$ $`=`$ $`p_q`$
$`p_{(2)}=p_{(8)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}p_{g_1}`$
$`p_{(3)}=p_{(7)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}p_{g_2}`$ (7.47)
$`p_{(4)}=p_{(6)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}p_{g_3}`$
$`p_{(5)}`$ $`=`$ $`p_{\overline{q}}`$
Summing over the bands, we get
$`{\displaystyle \underset{i=1}{\overset{8}{}}}p_{(i)}=p_q+p_{g_1}+p_{g_2}+p_{g_3}+p_{\overline{q}}=p_{\mathrm{string}},`$ (7.48)
which is the total momentum of the string.
For a fragment of the string, the momentum is given as
$$p_{\mathrm{fragm}}=\frac{1}{2}_C^{}\left(g(\sigma +\tau )+g(\sigma \tau )\right)𝑑\sigma +\frac{1}{2}_C^{}\left(g(\sigma +\tau )g(\sigma \tau )\right)𝑑\tau ,$$
(7.49)
where the path of the integration $`C^{}`$ is an arbitrary curve between two breakpoints, or between one break point and a boundary. One may write
$$p_{\mathrm{fragm}}=p_\mathrm{R}+p_\mathrm{L},$$
(7.50)
with
$`p_\mathrm{R}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _C^{}}g(\sigma \tau )𝑑\sigma {\displaystyle \frac{1}{2}}{\displaystyle _C^{}}g(\sigma \tau )𝑑\tau ,`$ (7.51)
$`p_\mathrm{L}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _C^{}}g(\sigma +\tau )𝑑\sigma +{\displaystyle \frac{1}{2}}{\displaystyle _C^{}}g(\sigma +\tau )𝑑\tau ,`$ (7.52)
where $`p_\mathrm{R}`$ and $`p_\mathrm{L}`$ represent sums of momenta of R-bands and L-bands. For the example of fig. 7.25, we can choose the path $`(1B),(B2)`$ for the string fragment between the break points 1 and 2. Since the first path is parallel to all L-bands (only R-bands contribute) and the second one is parallel to all R-bands (only L-bands contribute), we find
$$p_{\mathrm{fragm}\mathrm{\hspace{0.17em}1}2}=p_{1B}+p_{B2}$$
(7.53)
with
$`p_{1B}`$ $`=`$ $`x_1p_{(1)}+p_{(2)}+(1x_2)p_{(3)},`$ (7.54)
$`p_{B2}`$ $`=`$ $`(1y_1)p_{(7)}^\mu +p_{(6)}^\mu +y_2p_{(5)}^\mu ,`$ (7.55)
where the factors $`x_1`$, $`(1x_1)`$, $`y_1`$, $`(1y_1)`$ represent the fact that the bands at the extremities are only partially integrated over. The other string fragments are treated correspondingly. For the left string fragment, we may chose the integration path $`(A1)`$, for the right one $`(2C)`$. So we find for the three string fragments (referred to as $`a`$, $`b`$, $`c`$) of fig. 7.25 the momenta
$`p_\mathrm{a}`$ $`=`$ $`y_1p_{(7)}+p_{(8)}+(1x_1)p_{(1)}`$ (7.56)
$`p_\mathrm{b}`$ $`=`$ $`x_1p_{(1)}+p_{(2)}+(1x_2)p_{(3)}+y_2p_{(5)}+p_{(6)}+(1y_1)p_{(7)}`$ (7.57)
$`p_\mathrm{c}`$ $`=`$ $`x_2p_{(3)}+p_{(4)}+(1y_2)p_{(5)}.`$ (7.58)
It is easy to verify that the sum of the three sub-strings gives the total momentum of the string,
$$p_a+p_b+p_c=\underset{i=1}{\overset{8}{}}p_{(i)}=P_{\mathrm{string}}.$$
(7.59)
The mass squared of the string fragments is finally given as
$$m_{\mathrm{fragm}}^2=p_{\mathrm{fragm}}^2,$$
(7.60)
for example
$$m_a^2=2\left(p_{(1)}p_{(8)}x_1p_{(8)}p_{(1)}+y_1\left(p_{(7)}p_{(8)}+p_{(7)}p_{(1)}\right)x_1y_1p_{(7)}p_{(1)}\right),$$
(7.61)
where we took advantage of the light-cone character of the momenta of the bands $`(p_{(i)}^2=0).`$
#### Determination of Hadrons
So far, we have determined the flavor $`f`$ and the mass $`m`$ of each string fragment. In order to identify string fragments with hadrons, we construct a mass table, which defines the hadron type as a function of the mass and the flavor of the fragment. For a given flavor $`f`$ of the fragment, we introduce a sequence $`m_1^f<m_2^f<\mathrm{}`$ of masses, such that in case of a fragment mass being within an interval $`[m_{i1}^f,m_i^f]`$, one assigns a certain hadron $`h_i`$. The masses $`m_i^f`$ are determined by the masses of the neighboring particles. So we decide for a $`u\overline{u}`$ pair to be a pion if its mass is between $`0`$ and $`(140+770)/2=455`$ MeV. So the particle masses give a natural parameterization. This works, however, only up to strange flavor. For charm and bottom flavor we choose with a fraction $`1:3`$ between pseudo-scalar and vector-mesons.
#### Mass Corrections
An unrealistic feature of our approach, so far, is the fact that stable particles are in general off-mass-shell. In order to correct for this, we employ a slight modification of the break point such that the on-shell mass is imposed.
Let us again consider the example of fig. 7.25. For a given mass, the parameters $`x_1`$ and $`y_1`$ describe hyperbolas in the regions of overlapping bands (different ones in different regions). Fig. 7.26 shows
for our example some curves of constant mass for the left sub-string (between the left side and the first break point). In the same way we find hyperbolas of constant mass for the right sub-string (between the break points 1 and 2).
If two neighboring substrings are stable particles, one needs to impose on-shell masses to both fragments, which amounts to find the intersection of the two corresponding hyperbolas.
If one has to modify the break point according to only one mass condition, with the mass of the second sub-string being still large enough not to represent a stable hadron, a possible break point must lie on the corresponding hyperbola. To completely determine the point, we need a second condition. Apart from the squared mass, another Lorentz invariant variable available is the squared proper time of the break point, defined as
$$\mathrm{\Gamma }^2=\left(x(\sigma _{\mathrm{break}},\tau _{\mathrm{break}})\right)^2.$$
(7.62)
So the second condition is the requirement that the proper time of the new break point should coincide with the proper time of the original one. To calculate the proper time, we use eq. (7.12), to obtain
$$\mathrm{\Gamma }^2=\left(\frac{1}{2}_{\sigma \tau }^{\sigma +\tau }g(\xi )d\xi \right)^2.$$
(7.63)
The integration is done in the same way as for the masses, it is a summation of the momenta of the bands or of fractions of them. In the case of our example, we find in the region $`(1,7)`$
$`\mathrm{\Gamma }^2`$ $`=`$ $`{\displaystyle \frac{1}{4}}\left(x_1p_{(1)}+p_{(2)}+y_1p_{(3)}\right)^2`$ (7.64)
$`=`$ $`{\displaystyle \frac{1}{2}}\left(x_1p_{(1)}p_{(2)}+y_1p_{(2)}p_{(3)}+x_1y_1p_{(1)}p_{(3)}\right),`$ (7.65)
which represents again a hyperbola in the parameter space, as shown in fig. 7.27.
So finding a new break point amounts to finding the intersection of the two curves (hyperbolas) representing constant mass and proper time.
### 7.8 Transverse Momentum
Inspired by the uncertainty principle, a transverse momentum is generated at each breaking, which means that 4-vectors $`p_{}`$ and $`p_{}`$ are assigned to the string ends at both sides of the break point. First we choose the absolute value $`k=|\stackrel{}{p}_{}|`$ of the transverse momentum according to the distribution
$$f\left(k\right)e^{\frac{k}{2p_{}^{\mathrm{break}}}},$$
(7.66)
with the parameter $`p_{}^{\mathrm{break}}`$ to be fixed. We require $`p_{}`$ to be orthogonal to the momenta $`p_{(i)}`$ and $`p_{(j)}`$ of the two intersecting bands where the break point is located. So we have
$`p_{}p_{(i)}`$ $`=`$ $`0`$ (7.67)
$`p_{}p_{(j)}`$ $`=`$ $`0`$ (7.68)
$`p_{}^2`$ $`=`$ $`k^2.`$ (7.69)
Technically this is most easily done, if we perform a Lorentz-boost into the center of mass system of the two momenta $`p_{(i)}`$ and $`p_{(j)}`$ followed by a rotation such that $`p_{(i)}`$ is oriented along the z-axis. One defines a vector $`p_{}^{_{^{}}}`$, having the components
$`(p_{}^{_{^{}}})_o`$ $`=`$ $`0`$ (7.70)
$`(p_{}^{_{^{}}})_x`$ $`=`$ $`k\mathrm{cos}\alpha `$ (7.71)
$`(p_{}^{_{^{}}})_y`$ $`=`$ $`k\mathrm{sin}\alpha `$ (7.72)
$`(p_{}^{_{^{}}})_z`$ $`=`$ $`0,`$ (7.73)
$`\alpha `$ being a random angle between $`0`$ and $`2\pi `$. The transformation back to the original system gives the 4-vector $`p_{}`$.
This operation modifies, however, the mass of the string. In order to account for this, we consider the transverse momentum as an additional band of the string. It is treated in the same way as the others with the only exception that we do not look for break points in this region. For our example of fig. 7.25, we obtain for the left string fragment the momentum
$$p_\mathrm{a}=p_{}+y_1p_{(7)}+p_{(8)}+(1x_1)p_{(1)},$$
(7.74)
rather than eq. (7.56). The modification of the coefficients of the corresponding hyperbola for the mass correction procedure is obvious. In the case where we have to pass to another region to find a modified break point for the mass correction, we have to perform a rotation such that the vector $`p_{}`$ is transverse to the two momenta of the new region.
### 7.9 The Fragmentation Algorithm
In the following, we describe the fragmentation algorithm which is used to obtain a complete set of particles from one string.
1. For a given string, we look for break points. Let $`n`$ be the number of break points.
2. For each break point, we generate a flavor and a transverse momentum.
3. We choose one break point by random and calculate the masses of the two neighboring substrings.
4. If there is at least one mass in the region of the resonances, we try to modify the break point as discussed to get exactly this mass. If this is not possible, we reject (delete) this break point and go to step 3.
5. If the mass of a sub-string is bigger than the upper limit in the mass table, we fragment this sub-string (go to step 1).
In this way, we can deal in an elegant manner with the kinematical constraints. Often, break points are rejected when a sampled transverse momentum is too high, which results in a negative mass squared for a final particle. In this case we look for another break point with another transverse momentum until a valid configuration is found.
## Chapter 8 Parameters
We discuss in this chapter the parameters of the model, how they are determined, and also their values. Parameter fixing is done in a systematic way, starting with the hadronization parameters and the ones determining the time-like cascade, before considering parton-parton-scattering and hadron-hadron scattering.
### 8.1 Hadronization
The breaking probability $`p_{\mathrm{break}}`$ is the essential parameter in the hadronization model to determine the multiplicity and the form of the rapidity distribution. For $`\gamma p,pp,pA,AA`$ , we use a a fixed value, fitted to reproduce the pion multiplicity in $`\gamma p`$ scattering. For $`e^+e^{}`$annihilation, this parameter is considered to be $`Q`$ dependent as
$$p_{\mathrm{break}}(Q)=0.14+\frac{7.16\mathrm{GeV}}{Q}\frac{111.1\mathrm{GeV}^2}{Q^2}+\frac{855\mathrm{G}\mathrm{e}\mathrm{V}^3}{Q^3}$$
(8.1)
in the region $`14Q91.2`$ GeV, and to be constant outside this interval.
Figure 8.1 shows the total multiplicity of charged particles in $`e^+e^{}`$ annihilation as a function of energy. The solid line corresponds to the $`Q^2`$ dependent $`p_{\mathrm{break}}`$, the dashed line is for $`p_{\mathrm{break}}=0.21`$ (the value for $`Q=91`$GeV. In fig. 8.2, we show the corresponding rapidity distributions. The effect of making $`p_{\mathrm{break}}`$ Q-dependent shows up more in the shape of the rapidity distribution, not in the total multiplicity. The parameter $`p_{}^{\mathrm{break}}`$, which determines the transverse momenta of the partons at a string break, is determined by investigating transverse momentum spectra of charged particles. A value of 0.50 GeV provides the best fit to data concerning $`\gamma p,pp,pA,AA`$, whereas for $`e^+e^{}`$ a value of 0.35 GeV is more favorable. The parameter $`p_{\mathrm{ud}}`$ affects strongly kaon production, we use $`p_{\mathrm{ud}}=0.44`$ adjusted to the multiplicity of kaons. Baryon production is determined by $`p_{\mathrm{diquark}}`$, we use 0.08 adjusted to proton production in $`e^+e^{}`$-annihilation. In table 8.4, we give the complete list of hadronization parameters, with their default values.
So we use absolutely the same parameters for all the reactions $`\gamma p,pp,pA,AA`$. A perfect fit for $`e^+e^{}`$requires a modification of two parameters, $`p_{\mathrm{break}}`$ and $`p_{}^{\mathrm{break}}`$.
### 8.2 Time-Like Cascade
Let us discuss the parameters which determine the time-like cascade, all fixed via studying $`e^+e^{}`$ annihilation. For the pQCD parameter, we use the usual leading order value $`\mathrm{\Lambda }=0.2\text{ GeV}`$. We have also a technical parameter $`q_{\mathrm{fin}}^2`$, which determines the lower mass limit of partons in the time-like cascade. In figure 8.3, we analyze how certain spectra depend on this parameter. We show rapidity, transverse momentum and multiplicity distributions for partons and for charged particles for an $`e^+e^{}`$ annihilation at 34 GeV. We show results for different values of $`q_{\mathrm{fin}}^2`$, namely $`0.25`$ GeV<sup>2</sup>, $`1.0`$ GeV<sup>2</sup>and $`4.0`$ GeV<sup>2</sup>, which means a lower mass limit of $`1`$ GeV, $`2`$ GeV or $`4`$ GeV, respectively. One sees that only parton distributions are sensitive to the choice of this parameter, whereas the corresponding charged particle spectra exhibit rather week dependence on it. This can be explained from the fact that decreasing $`q_{\mathrm{fin}}^2`$ mainly results in production of additional partons with transverse momenta $`p_{}^2q_{\mathrm{fin}}^2`$ and such soft collinear partons hardly affect the fragmentation procedure, which is in this sense “infrared stable”. The number $`N_f`$ of active flavors is taken to be 5 for $`e^+e^{}`$ annihilation, whereas for $`\gamma p,pp,pA,AA`$ we use for the moment $`N_f=3`$. In table 8.2
we show the cascade parameters and their default values.
### 8.3 Parton-Parton Scattering
There is first of all the parameter $`Q_0`$ which defines the borderline between soft and hard processes, where one has to choose a reasonable value (say between 1 and 2 GeV<sup>2</sup>).
Then we have a couple of parameters characterizing the soft Pomeron: the intercept $`\alpha _{\mathrm{soft}}(0)`$ and the slope $`\alpha _{\mathrm{soft}}^{}`$ of the Pomeron trajectory, the vertex value $`\gamma _{\mathrm{part}}`$ and the slope $`R_{\mathrm{part}}`$ for the Pomeron-parton coupling, and the characteristic hadronic mass scale $`s_0`$. We have two parameters, $`\beta _g`$ and $`w_{\mathrm{split}}`$, characterizing the coupling between the soft Pomeron and the parton ladder. Whereas for $`s_0`$ one just chooses some “reasonable” value and $`R_{\mathrm{part}}`$ is taken to be zero, one fixes the other parameters by trying to get a good fit for the total cross section and the slope parameter for proton-proton scattering as a function of the energy as well as the structure function $`F_2`$ of deep inelastic lepton-proton scattering.
Concerning the hard scattering part, the resolutions scale $`p_{\mathrm{res}}^2`$ and the K-factor $`K`$ are fixed such that the standard parton evolution is reproduced.
Finally, we have the triple Pomeron coupling weight $`r_{3\mathrm{I}\mathrm{P}}`$, which is fixed by as well checking the energy dependence of the proton-proton total cross section.
The values of these parameters are shown in table 8.3.
### 8.4 Hadron-Hadron Scattering
Let us first discuss the parameters related to the partonic wave function of the hadron $`h`$ (for numerical applications we only consider nucleons: $`h=N`$). The transverse momentum distribution is characterized by the hadronic Regge radius squared $`R_N^2`$, the longitudinal momentum distributions are given in terms of two exponents, $`\alpha _{\mathrm{remn}}^N`$ and $`\alpha _{\mathrm{part}}`$. The latter one is taken to be independent of the hadron type as $`2\alpha _{\mathrm{I}\mathrm{R}}(0)1`$, with the usual Reggeon intercept $`\alpha _{\mathrm{I}\mathrm{R}}(0)=1/2.`$ The parameter $`R_N`$ also affects the proton-proton total cross sections, whereas $`\alpha _{\mathrm{remn}}^N`$ can be determined by investigating baryon spectra (but it also influences the total cross section).
There are several more parameters, which have not been mentioned so far: the remnant excitation probability $`p_{\mathrm{remn}.\mathrm{ex}}`$ and the exponent $`\alpha _{\mathrm{remn}.\mathrm{ex}}`$, which gives a remnant mass distribution as
$$\left(M^2\right)^{\alpha _{\mathrm{remn}.\mathrm{ex}}}.$$
(8.2)
The minimum string mass $`m_{\mathrm{string}\mathrm{min}}`$ assures that Pomerons with string masses less than this minimal mass are ignored. The partons defining the string ends are assumed to have transverse momenta according a Gaussian distribution with a mean value $`p_{\mathrm{SE}}`$. Diffractive scattering is assumed to transfer transverse momentum according to a Gaussian distribution, with a mean value $`p_{\mathrm{diff}}`$. The parameter $`m_{\mathrm{string}\mathrm{min}}`$ is taken to be just slightly bigger than two times the pion mass to allow at least string fragmentation into two pions. The other parameters can be fixed by comparing with experimental inclusive spectra.
In table 8.4, we show the numerical values of the parameters.
## Chapter 9 Testing Time-like Cascade and Hadronization: Electron-Positron Annihilation
Electron-positron annihilation is the simplest possible system to test the time-like cascade as well as the model of fragmentation, since the decay of a virtual photon in electron-positron annihilation gives a quark and an anti-quark, both emitting a cascade of time-like partons, which finally hadronize.
Electron-positron annihilation is therefore used to test both the time-like cascade and the hadronization model, and in particular to fix parameters.
The simulation of an electron-positron annihilation event can be divided into three different stages:
1. The annihilation into a virtual photon or a $`Z`$ boson and its subsequent decay into a quark-anti-quark pair (the basic diagram).
2. The evolution of the quark and the anti-quark into on-shell partons by radiation of perturbative partons (time-like cascade).
3. The transition of the partonic system into hadrons via a fragmentation model (hadronization).
These stages are discussed in the following sections.
After having described the three stages of electron-positron annihilation, we will be able to test the model against numerous data available. We will show comparisons with experimental results at low energies at PETRA (DESY), by the TASSO collaboration . The center-of-mass energies are 14, 22 and 34 GeV. Higher energies are reached at LEP, where we compare especially with results for $`91.2`$ GeV, the $`Z^0`$ mass, where a big number of events has been measured.
By comparing with data, we will be able to fix the essential parameters of the hadronization model, namely $`p_{\mathrm{break}}`$, $`p_{\mathrm{ud}}`$, $`p_{\mathrm{diquark}}`$ and $`p_{}^{\mathrm{break}}`$. The free parameters in the parton cascade are the pQCD scaling parameter $`\mathrm{\Lambda }`$ and $`q_{\mathrm{fin}}^2`$, representing the minimum transverse momentum for a branching in the cascade. For the pQCD parameter, we use the usual leading order value $`\mathrm{\Lambda }=0.2\text{ GeV}`$. The influence of the technical parameter $`q_{\mathrm{fin}}^2`$ has been investigated in detail.
### 9.1 The Basic Diagram
The first order differential cross section for the process
$$e^+e^{}\gamma ^{}\mathrm{or}Zq\overline{q}$$
(9.1)
is given as
$`{\displaystyle \frac{d\sigma }{d\mathrm{cos}\theta }}`$ $`=`$ $`{\displaystyle \frac{\pi \alpha ^2}{2s}}[\begin{array}{c}\end{array}(1+\mathrm{cos}^2\theta )\{q_f^22q_fV_eV_f\chi _1(s)+(A_e^2+V_e^2)(A_f^2+V_f^2)\chi _2(s)\}`$ (9.7)
$`+\mathrm{cos}\theta \{4q_fA_eA_f\chi _1(s)+8A_eV_eA_fV_f\chi _2(s)\}\begin{array}{c}\end{array}]`$
with
$`\chi _1(s)`$ $`=`$ $`\kappa {\displaystyle \frac{s(sM_Z^2)}{(sM_Z^2)^2+\mathrm{\Gamma }_Z^2M_Z^2}},`$ (9.8)
$`\chi _2(s)`$ $`=`$ $`\kappa ^2{\displaystyle \frac{s^2}{(sM_Z^2)^2+\mathrm{\Gamma }_Z^2M_Z^2}},`$ (9.9)
where $`\kappa `$ is given as
$$\kappa =\frac{\sqrt{2}G_FM_Z^2}{16\pi \alpha }.$$
Here, $`\alpha `$ is the fine structure constant, $`q_f`$ the quark flavor, $`M_z`$ the mass of the $`Z`$ boson, $`\mathrm{\Gamma }_Z`$ its decay width. The vector and axial coupling factors are
$$V_f=T_f^32q_f\mathrm{sin}^2\theta _W,A_f=T_f^3$$
(9.10)
and
$$T_f^3=\{\begin{array}{ccc}\frac{1}{2}\hfill & \mathrm{if}\hfill & f=\nu ,u,c\hfill \\ \frac{1}{2}\hfill & \mathrm{if}\hfill & f=e,d,s,b\hfill \end{array}.$$
(9.11)
$`\theta _W`$ is the Weinberg mixing angle. The Fermi constant
$$G_F=\frac{\sqrt{2}g_W^2}{8M_W^2}$$
(9.12)
is expressed via the weak coupling $`g_W^2=4\pi \alpha /\mathrm{sin}^2\theta _W`$ and the W-boson mass $`M_w`$. At low energies, $`sM_Z^2`$, we recover the well known formula
$$\sigma =\frac{4\pi \alpha ^2}{3s}q_f^2.$$
(9.13)
The factors $`\chi _1(s)`$ and $`\chi _2(s)`$ correspond to the intermediate $`Z`$-boson state and to the photon-Z-boson interference, respectively. Formula (9.7) can now be used to generate an initial quark-anti-quark pair.
### 9.2 The Time-like Parton Cascade and String Formation
If one considers a quark and an anti-quark coming from the decay of a virtual photon or Z-boson to be on-shell, this amounts to a lowest order treatment. At high energies, a perturbative correction has to be done. Since it is difficult to calculate the higher order Feynman diagrams exactly, one uses the so-called DGLAP evolution equations, which describes the evolution of a parton system with leading logarithmic accuracy. This amounts to successively emitting partons (time-like parton cascade, see fig. 9.2).
This has been discussed in detail in connection with the parton production in proton-proton or nucleus-nucleus collisions. We use exactly the same method here. Even for determining three momenta of the primary quark and anti-quark, we do not need any new input, since this corresponds exactly to the case of the time-like cascade of the two partons involved in the Born scattering in hadronic collisions.
The fact that we use one and the same procedure for the time-like parton cascade in all the different reactions, allows us to test elements of hadronic interactions in a much simpler context of elementary electron-positron interactions.
The final step is the hadronization of the above-mentioned parton configuration. Here we use the string model, as in case of $`pp`$ or $`\gamma p`$ scattering. The string formation in $`e^+e^{}`$ is much simpler than in proton-proton scattering, where the cut Pomeron is represented as a cylinder. The structure of an $`e^+e^{}`$ event is planar in the sense that the whole event can be represented on a plane. So we simply plot the diagram on a plane with only one cut line. In fig. 9.3(a), we present a half-plane (on one side of the cut) for the amplitude shown in fig. 9.2. The dotted line represents the cut. There are a couple of partons crossing the cut, indicated by dots. As in the case of proton-proton or photon-proton scattering, we identify the cut line as a kinky relativistic string, with the partons representing the kinks. So in our example, we have a kinky string with six kinks, two external ones and four internal ones. We then apply the usual hadronization procedure, discussed earlier in detail, in order to calculate hadron production from a fragmenting string, see fig. 9.3(b).
### 9.3 Event Shape Variables
We start our presentation of results by considering the so-called event-shape variables, which describe the form of an $`e^+e^{}`$ event in general. For example, one is interested in knowing whether the particle momenta are essentially aligned along a certain axis, distributed isotropically over the phase space, or lying more or less in a plane. In the figs. 9.4 to 9.8, we are going to compare our calculated distributions of several event-variables with data. Let us first discuss the different event-shape variables, one after the other.
#### Sphericity
The sphericity is defined by the eigenvalues of the sphericity tensor,
$$S^{\alpha \beta }=\underset{i}{}(p_i)^\alpha (p_i)^\beta ,$$
(9.14)
where $`i`$ sums all particles and $`p^\alpha `$ is the particle four-momentum. One finds three eigenvalues $`\lambda _i`$ with $`\lambda _1<\lambda _2<\lambda _3`$ and $`\lambda _1+\lambda _2+\lambda _3=1`$. The sphericity is then defined as
$$S=\frac{3}{2}(\lambda _1+\lambda _2).$$
(9.15)
For a perfectly isotropic event, one finds $`\lambda _1=\lambda _2=\lambda _3=1/3`$ and therefore $`S=1`$. An event oriented along one axis gives $`S=3/2(0+0)=0`$ .
To test whether an event has planar geometry, one defines the aplanarity
$$A=\frac{3}{2}\lambda _{1.}$$
(9.16)
For events in a plane we will find $`\lambda _1=A=0`$. The maximum of this value is $`A=3/21/3=1/2`$ for an isotropic event, since the eigenvalues are ordered.
The three eigenvectors $`\stackrel{}{v}_{1,2,3}`$ of the matrix $`S^{\alpha \beta }`$ can be used to define a coordinate system.
#### C and D Parameters
The C-parameter is defined by
$$C=3\left(\lambda _1\lambda _2+\lambda _1\lambda _3+\lambda _2\lambda _3\right),$$
(9.17)
with $`\lambda _{1,2,3}`$ being the eigenvalues of the tensor
$$M^{\alpha \beta }=\frac{_i\frac{p_i^\alpha p_i^\beta }{|p_i|}}{_i|p_i|}.$$
(9.18)
The $`D`$ -parameter is
$$D=27\lambda _1\lambda _2\lambda _3.$$
(9.19)
These values measure the multiple jet-structure of events. For small values of C two of the eigenvalues are close to zero, we have a two-jet event. If one of the three eigenvalues is close to zero, the D-parameter is approaching zero as well, we have at least a planar event.
#### Thrust
The thrust of an event is defined as
$$T=\underset{\stackrel{}{n}}{\mathrm{max}}\frac{_j\left|\stackrel{}{n}\stackrel{}{p}_j\right|}{_j\left|\stackrel{}{p}_j\right|}.$$
(9.20)
The vector $`\stackrel{}{n}_{\mathrm{thrust}}`$, which maximizes this expression, defines the thrust axis. A two-jet event will give a thrust value of 1 and a thrust axis along the two jets. An isotropic event gives $`T1/2`$.
One can repeat the same algorithm with the imposed condition $`\stackrel{}{n}\stackrel{}{n}_{\mathrm{thrust}}`$; this gives an expression for the major $`M`$ with the axis $`\stackrel{}{n}_{\mathrm{Major}}`$.
A third variable, the minor $`m`$, is obtained by evaluating the above expression with $`\stackrel{}{n}\stackrel{}{n}_{\mathrm{thrust}}`$ and $`\stackrel{}{n}\stackrel{}{n}_{\mathrm{Major}}`$, the axis being already given. Each of these values describes the extension of the event perpendicular to the thrust axis. Similar values for $`M`$ and $`m`$ describe therefore a cylindrical event. For this, the oblateness is defined as $`O=Mm`$ which, as the aplanarity, describes a cylindrical event for $`O=0`$ an a planar event for higher values.
#### Jet Broadening
In each hemisphere, the sum of the transverse momenta of the particles relative to the thrust axis is divided by the sum of the absolute values of the momenta.
$$B_\pm =\frac{_{\pm \stackrel{}{p}.\stackrel{}{n}_{\mathrm{thrust}}>0}|\stackrel{}{p}_i\times \stackrel{}{n}_{\mathrm{thrust}}|}{2_i|\stackrel{}{p}_i|}$$
(9.21)
The greater $`B_+`$ is, the greater is the mean transverse momentum. One defines in addition the following variables,
$$B_{\mathrm{wide}}=\mathrm{max}(B_+,B_{}),B_{\mathrm{narrow}}=\mathrm{min}(B_+,B_{}),$$
(9.22)
$$B_{\mathrm{total}}=B_++B_{},B_{\mathrm{diff}}=|B_+B_{}|,$$
(9.23)
to compare jet broadening in both hemispheres. For small $`B_{\mathrm{wide}}`$, one finds a longitudinal event, $`B_{\mathrm{diff}}`$ measures the asymmetry between the two hemispheres.
#### Heavy Jet or Hemisphere Mass
The variable $`M_\mathrm{h}^2`$ is defined as
$$M_\mathrm{h}^2=\mathrm{max}(\left(\underset{\stackrel{}{p}_i\dot{}\stackrel{}{n}_{\mathrm{thrust}}>0}{}p_i\right)^2,\left(\underset{\stackrel{}{p}_i\dot{}\stackrel{}{n}_{\mathrm{thrust}}<0}{}p_i\right)^2).$$
(9.24)
and corresponds to the maximal invariant mass squared of the hemispheres. The corresponding formula with “$`\mathrm{min}`$” instead of “$`\mathrm{max}`$” defines the variable $`M_\mathrm{l}`$. One defines as well $`M_{\mathrm{diff}}=M_\mathrm{h}M_\mathrm{l}`$. Usually one analyzes distributions of $`\frac{M_\mathrm{h}^2}{E_{\mathrm{vis}}^2}`$, $`\frac{M_\mathrm{l}^2}{E_{\mathrm{vis}}^2}`$ or $`\frac{M_{\mathrm{diff}}^2}{E_{\mathrm{vis}}^2}`$ , which describe the squared masses normalized to the visible energy.
#### Some Comments
Fig. 9.4 shows the distributions of sphericity and thrust for the lower energies 14, 22 et 34 GeV. Even though one expects an increasing contribution of perturbative gluons, which is confirmed by the inclusive hadron distributions (see next section), the events are more longitudinal at higher energies, corresponding to the values of thrust close to 1 and to the sphericity $`S0`$ . This can be explained by the fact that the leading quarks dominate the event shape. The results for higher energies (figs. 9.5, 9.6, 9.7) confirm the above statements.
In figs. 9.5, 9.6, 9.7, we show also the distributions of other event variables, like heavy jet mass, Major, etc. In general, our model describes quite well all these event shape variables.
### 9.4 The Charged Particle Distributions
We will now consider the distributions of charged particles, which by definition contain all the particles with a decay time smaller than $`10^9`$s, i.e. the spectra contain, for example, products of decay of $`K_{\mathrm{short}}^0`$, while $`K^\pm `$ are considered stable. The decay products of strange baryons are also included in the distributions.
In fig. 9.9, we plot multiplicity distributions of charged particles for three different energies, where one observes an obvious increase of the multiplicity with energy. In fig. 9.10, we show the distributions of the absolute value of the rapidity for the energies 14 GeV, 22 GeV, 34 GeV, 91.2 GeV, 133 GeV and 161 GeV. The rapidity is defined as $`y=0.5\mathrm{ln}((E+p_z)/(Ep_z))`$, where the variable $`p_z`$ may be defined along the thrust axis or along the sphericity axis. For both, multiplicity and rapidity distributions, the theoretical curves agree well with the data.
The multiplicity increases faster than $`<n_{\mathrm{ch}}>=a+b\mathrm{ln}s`$ as a function of $`s`$ , which is due to the fact that the maximal height of the rapidity distribution increases with energy, as seen in fig. 9.10. This comes from radiated gluons, leading to kinky strings, since a flat string without gluons shows an increasing width but a constant height, as one can see in
fig. 9.11. Here, the rapidity distributions of charged particles are plotted for the fragmentation of a flat $`d\overline{d}`$ string for the energies 14, 22 and 34 GeV. One observes that the width of the distributions increases whereas its height does not change, giving rise to a proportionality to $`\mathrm{ln}s`$. Additional hard gluons with non-collinear momenta increase the multiplicity in the mid-rapidity region.
Rather than the rapidity, one may consider the scaled momentum $`x_\mathrm{p}=2|\stackrel{}{p}_i|/E`$ or the scaled energy $`x_\mathrm{E}=2E_i/E`$, as well as the “rapidity-like” variable $`\xi =\mathrm{ln}x_\mathrm{p}`$. Concerning the $`\xi `$ distributions shown in fig. 9.12,
one sees that the value $`\xi _{\mathrm{max}}`$ corresponding to the maximum of the curves increases with energy. The $`x_\mathrm{p}`$ distributions (see fig. 9.13) show the development of a more and more pronounced peak at $`x_\mathrm{p}`$ close to zero, with increasing energy.
Having discussed in detail the variables describing the longitudinal phase space, we now turn to transverse momentum,
which can be defined according the the sphericity axis or to the thrust axis. One writes
$$p_{}^{\mathrm{in}}=\{\begin{array}{ccc}|\stackrel{}{v}_2\stackrel{}{p}|& \text{for sphericity}& \\ |\stackrel{}{n}_{\mathrm{major}}\stackrel{}{p}|& \text{for thrust}& \end{array},$$
(9.25)
and
$$p_{}^{\mathrm{out}}=\{\begin{array}{ccc}|\stackrel{}{v}_3\stackrel{}{p}|& \text{for sphericity}& \\ |\stackrel{}{n}_{\mathrm{minor}}\stackrel{}{p}|& \text{for thrust}& \end{array},$$
(9.26)
There are mainly two “sources” of transverse momentum. The first one is the transverse momentum created at each string break. The second one is the transverse momentum from hard gluon radiation, which can be much larger than the first one. So we find large values of $`p_{}`$ in the event plane, and smaller ones out of the event plane. Here the event plane is essentially defined by the direction of the hardest gluon emitted.
Let us have a look at transverse momenta of charged particles coming from a string decay for different energies (see fig. 9.14). As expected, the curves show the same behavior. In the same figure, we show the results for an $`e^+e^{}`$ annihilation at 14 GeV. Already at 14 GeV the influence of parton radiation is important. Fig. 9.15 shows the results for 14-34 GeV . Our results agree well with data from the TASSO collaboration . One can see how transverse momenta increase with the energy as an indication of more hard gluon radiation: the $`p_{}^{\mathrm{out}}`$ distributions change little, whereas the $`p_{}^{\mathrm{in}}`$ distributions get much harder at high energies as compared to lower ones.
### 9.5 Identified Particles
In this section we consider inclusive spectra of identified hadrons. This provides a crucial test of the fragmentation model and allows to fix the two hadronization parameters $`p_{\mathrm{ud}}`$ et $`p_{\mathrm{diquark}}`$. The first one gives the probability to find a pair $`u\overline{u}`$ or $`d\overline{d}`$, fixing so the strangeness probability to be $`(12p_{\mathrm{ud}})`$. The parameter $`p_{\mathrm{diquark}}`$ determines the multiplicity of baryons.
Let us look at spectra at 29 GeV obtained at SLAC (, ) (figs. 9.17, 9.18) and at 91 GeV at LEP (figs. 9.19, 9.20, 9.21). Since the total multiplicity is dominated by small $`x`$-values, we show this regions separately for some figures. The results are in general quite good, however, $`K^{}`$’s are underestimated.
For charmed particles, there is no production from the string decay due to the large mass of the $`c\overline{c}`$ pair. The corresponding probability $`p_{\mathrm{charm}}`$ is taken to be zero. Charmed quarks come therefore directly from the decay of the virtual photon as well as from perturbative parton cascade. This explains as well the drop of $`D`$ spectrum at small $`x_p`$.
### 9.6 Jet Rates
Jet multiplicities play an important role in $`e^+e^{}`$ physics since their measurements proved the validity of perturbative QCD.
The first 3-jet event was found in 1979. Fig. 9.22 shows this historical event.
There are several methods to determine the number of jets in an event, but they are all based on some distance $`y_{ij}`$ between two particles $`i`$ and $`j`$ in momentum space, something like an invariant mass. For the JADE algorithm one defines
$$y_{ij}=\frac{2E_iE_j(1\mathrm{cos}\theta _{ij})}{E_{vis}^2},$$
(9.27)
$`\theta _{ij}`$ being the angle between the two particles, and for the algorithm DURHAM one has
$$y_{ij}=\frac{2\mathrm{min}(E_i^2,E_j^2)(1\mathrm{cos}\theta _{ij})}{E_{vis}^2}.$$
(9.28)
$`E_{vis}`$ is the total visible energy of all the particles which contribute to the jet finding. The algorithm works as follows: one determines the pair with the lowest distance $`y`$ and replaces it with one pseudo-particle having the sum of the momenta of the two particles $`i`$ and $`j`$ : $`p^\mu =p_i^\mu +p_j^\mu `$. One repeats this until all pairs of pseudo-particles have a distance greater than $`y_{\mathrm{cut}}`$. The number of jets is then the total number of pseudo-particles. Of course, this depends on the choice of $`y_{\mathrm{cut}}`$. Therefore one displays often the jet multiplicity distribution as a function of $`y_{\mathrm{cut}}`$ .
Let us compare the jet rates for different energies. Fig. 9.23 shows the jet rates for 91.2 GeV, fig. 9.24 for 133 and 161 GeV. The greater is $`y_{\mathrm{cut}}`$ the smaller is the number of jets.
## Chapter 10 Testing the Semi-hard Pomeron: Photon-Proton Scattering
It is well known that both photon-proton ($`\gamma ^{}p`$) scattering and hard processes in proton-proton ($`pp`$) collisions can be treated on the basis of perturbative QCD, using the same evolution equations. In both cases, the perturbative partons are finally coupled softly to the proton(s). This provides a very useful consistency check: any model for (semi)hard proton-proton collisions should be applied to photon-proton scattering, where a wealth of data exists, mainly from deep inelastic electron-proton scattering (DIS). In particular the soft coupling to the protons is not calculable from first principles, so photon-proton scattering provides a nice opportunity to test the scheme.
Let us discuss the relation between $`\gamma ^{}p`$ and $`pp`$ scattering.
In figure 10.1, we show the cut diagram (integrated squared amplitude), representing a contribution to photon-proton scattering: a photon couples to a quark of the proton, where this quark represents the last one in a “cascade” of partons emitted from the nucleon. In the leading logarithmic approximation (LLA) the virtualities of the partons are ordered such that the largest one is close to the photon .
Comparing the cut diagram for $`\gamma ^{}p`$ (figure 10.1) with the cut diagram representing a semi-hard elementary proton-proton scattering (figure 10.2), we see immediately that the latter one is essentially made of two $`\gamma ^{}p`$ diagrams, glued together by a Born process. So, understanding $`\gamma ^{}p`$ implies understanding an elementary nucleon-nucleon interaction as well. Actually, probably everybody agrees with this statement, which is nothing but the factorization hypothesis, proved in QCD , and the standard procedure to calculate inclusive cross sections in proton-proton scattering amounts to using input from the DIS structure functions. But one can profit much more from studying $`\gamma ^{}p`$, for example, concerning the production of hadrons. Not being calculable from first principles, the hadronization of parton configurations is a delicate issue in any model for proton-proton (or nucleus-nucleus) scattering. So studying $`\gamma ^{}p`$ provides an excellent possibility to “gauge” the hadronization procedure, such that there is no freedom left on the level of nucleon-nucleon (or nucleus-nucleus) scattering.
The simple picture, depicted at the fig. 10.1, is correct for large virtualities, but it fails when the photon virtuality becomes small. In that case a virtual photon behaves to a large extent as a hadron and is characterized by some parton content instead of interacting with a proton just as a point-like object.
Then the contribution of so-called resolved photon interactions - see fig. 10.3 \- is important and has to be taken into account properly for the description of hadron production in DIS. Only then one may deduce the parton momentum distributions of the proton from the measured virtual photon-proton cross section $`\sigma ^{\gamma ^{}p}`$.
### 10.1 Kinematics
In the following, we consider photon-proton collisions in the context of electron-proton scattering. We first recall the basic kinematic variables, see fig. 10.4. We use standard conventions: $`k`$, $`k^{}`$, and $`p`$ are the four-momenta of incoming and outgoing lepton and the target nucleon, $`q=kk^{}`$ is the four-momentum of the photon.
Then the photon virtuality is $`Q^2=q^2`$, and one defines the Bjorken $`x`$-variable as
$$x_\mathrm{B}=\frac{Q^2}{2(pq)}.$$
(10.1)
The $`y`$-variable is defined as
$$y=\frac{(pq)}{(pk)}=\frac{2(pq)}{s}=\frac{Q^2}{x_\mathrm{B}s},$$
(10.2)
which gives the energy fraction of the photon relative to the incident electron in the proton rest frame. In the above formula we used
$$s=(k+p)^2,$$
(10.3)
being the total center-of-mass squared energy and neglected the proton and electron masses, $`p^20`$, $`k^20`$. For the center-of-mass squared energy of photon-proton interaction, we use
$$\widehat{s}=(p+q)^2,$$
(10.4)
and we finally define the variable
$$\stackrel{~}{s}=\widehat{s}+Q^2=2(pq)=ys,$$
(10.5)
which allows to write
$$x_\mathrm{B}=\frac{Q^2}{\stackrel{~}{s}}.$$
(10.6)
It is often convenient to take $`\stackrel{~}{s}`$ and $`Q^2`$ as the basic kinematical variables instead of $`x_\mathrm{B}`$ and $`Q^2`$.
### 10.2 Cross Sections
The differential cross section for deep inelastic electron-proton scattering in the one-photon approximation can be written as
$$\frac{d\sigma ^{ep}(s,x_B,Q^2)}{dQ^2dx_\mathrm{B}}=\frac{\alpha }{\pi Q^2x_\mathrm{B}}\left[L_T^\gamma (y)\sigma _T^{\gamma ^{}p}(\stackrel{~}{s},Q^2)+L_L^\gamma (y)\sigma _L^{\gamma ^{}p}(\stackrel{~}{s},Q^2)\right],$$
(10.7)
where $`\sigma _{T(L)}^{\gamma ^{}p}(\stackrel{~}{s},Q^2)`$ are the cross sections for interactions of transversely (longitudinally) polarized photons of virtuality $`Q^2`$ with a proton, $`\alpha `$ is the fine structure constant and factors $`L_{T(L)}^\gamma ^{}(y)`$ define the flux of transversely (longitudinally) polarized photons,
$$L_T^\gamma (y)=\frac{1+(1y)^2}{2}L_L^\gamma (y)=(1y).$$
(10.8)
The cross sections $`\sigma _{T(L)}^{\gamma ^{}p}(\stackrel{~}{s},Q^2)`$ are related to the proton structure functions $`F_2`$, $`F_L`$, describing the proton structure as seen by a virtual photon probe, as
$`F_2(x_\mathrm{B},Q^2)`$ $`=`$ $`{\displaystyle \frac{Q^2}{4\pi ^2\alpha }}\left[\sigma _T^{\gamma ^{}p}(\stackrel{~}{s},Q^2)+\sigma _L^{\gamma ^{}p}(\stackrel{~}{s},Q^2)\right]`$ (10.9)
$`F_L(x_\mathrm{B},Q^2)`$ $`=`$ $`{\displaystyle \frac{Q^2}{4\pi ^2\alpha }}\sigma _L^{\gamma ^{}p}(\stackrel{~}{s},Q^2)`$ (10.10)
To the leading logarithmic accuracy one has to take into account a number of processes, contributing to $`\sigma _{T(L)}^{\gamma ^{}p}(\stackrel{~}{s},Q^2)`$. In the case of transverse polarization of the photon, we have three contributions: the direct coupling of the virtual photon $`\gamma ^{}`$ to a light quark from the proton (“light”), the direct coupling to a charm quark (“charm”), and finally we have a “resolved” contribution, i.e.
$$\sigma _T^{\gamma ^{}p}(\stackrel{~}{s},Q^2)=\sigma _{T(\mathrm{light})}^{\gamma ^{}p}(\stackrel{~}{s},Q^2)+\sigma _{T(\mathrm{charm})}^{\gamma ^{}p}(\stackrel{~}{s},Q^2)+\sigma _{T(\mathrm{resolved})}^{\gamma ^{}p}(\stackrel{~}{s},Q^2).$$
(10.11)
The latter one is becoming essential at small $`Q^2`$ and large $`\stackrel{~}{s}`$. For our study, contributions of beauty and top quarks can be neglected. The longitudinal photon cross section receives leading order contributions only from the direct $`\gamma ^{}`$-coupling to a charm quark. So we have
$$\sigma _L^{\gamma ^{}p}(\stackrel{~}{s},Q^2)=\sigma _{L(\mathrm{charm})}^{\gamma ^{}p}(\stackrel{~}{s},Q^2).$$
(10.12)
Again, contributions of beauty and top quarks can be neglected.
Let us list the different contributions. The leading order light quark-$`\gamma ^{}`$ coupling contribution (“light”) can be expressed via the quark momentum distributions in the proton $`f_{q/p}(x_B,Q^2)`$ as
$$\sigma _{T(\mathrm{light})}^{\gamma ^{}p}(\stackrel{~}{s},Q^2)=\frac{4\pi ^2\alpha }{Q^2}\underset{i\{u,d,s,\overline{u},\overline{d},\overline{s}\}}{}e_i^2x_Bf_{i/p}(x_B,Q^2),$$
(10.13)
where $`e_q^2`$ is the quark $`q`$ electric charge squared.
The contributions of heavy quarks can be taken into account via photon-gluon fusion (PGF) process ,
$$\sigma _{T/L(\mathrm{charm})}^{\gamma ^{}p}(\stackrel{~}{s},Q^2)=e_c^2𝑑x^{}𝑑p_{}^2\frac{d\sigma _{T/L}^{\gamma ^{}gc\overline{c}}(x^{}\stackrel{~}{s},Q^2,p_{}^2)}{dp_{}^2}f_{g/p}(x^{},M_F^2),$$
(10.14)
where $`f_{g/p}(x^{},M_F^2)`$ is the gluon momentum distribution in the proton at the factorization scale $`M_F^2`$, and where the photon-gluon cross section in lowest order is given as
$$\frac{d\sigma _{T/L}^{\gamma ^{}gc\overline{c}}(\stackrel{~}{s},Q^2,p_{}^2)}{dp_{}^2}=\frac{1}{16\pi \stackrel{~}{s}\widehat{s}\sqrt{14(p_{}^2+m_c^2)/\widehat{s}}}\left|M_{T/L}^{\gamma ^{}gc\overline{c}}(\stackrel{~}{s},Q^2,p_{}^2)\right|^2$$
(10.15)
with $`\widehat{s}=\stackrel{~}{s}Q^2`$, and with the corresponding matrix elements squared given as
$`\left|M_T^{\gamma ^{}gc\overline{c}}(\stackrel{~}{s},Q^2,p_{}^2)\right|^2`$ $`=`$ $`\pi \alpha \alpha _s[({\displaystyle \frac{t^{}}{u^{}}}+{\displaystyle \frac{u^{}}{t^{}}}){\displaystyle \frac{Q^4+\widehat{s}^2}{\stackrel{~}{s}^2}}`$
$`+`$ $`{\displaystyle \frac{2m_c^2\stackrel{~}{s}^2}{t^2u^2}}(Q^22m_c^2)+{\displaystyle \frac{4m_c^2}{t^{}u^{}}}(\stackrel{~}{s}2Q^2)]`$
$`\left|M_L^{\gamma ^{}gc\overline{c}}(\stackrel{~}{s},Q^2,p_{}^2)\right|^2`$ $`=`$ $`8\pi \alpha \alpha _s\left({\displaystyle \frac{\widehat{s}Q^2}{\stackrel{~}{s}^2}}{\displaystyle \frac{m_c^2Q^2}{t^{}u^{}}}\right)`$ (10.17)
Here, the variables $`t^{},u^{}`$ are expressed via standard Mandelstam variables for parton-parton scattering as $`t^{}=tm_c^2`$, $`u^{}=um_c^2`$. According to , the factorization scale $`M_F^2`$ has to be chosen irrespectively to the photon virtuality $`Q^2`$ to assure the perturbative stability of the result; we use $`M_F^2=p_{}^2+m_c^2`$, which coincides at large $`\widehat{s}`$ with the virtuality (off-shellness) of the intermediate $`t`$-channel $`c`$-quark $`|t^{}|`$, with $`m_c=1.6`$ GeV being the $`c`$-quark mass.
In addition, at small $`Q^2`$ and large $`\stackrel{~}{s}`$, the contribution of resolved photon processes becomes important for the production of parton jets of transverse momenta $`p_{}^2>Q^2`$. Here,
$`\sigma _{T(\mathrm{resolved})}^{\gamma ^{}p}(\stackrel{~}{s},Q^2)`$ $`=`$ $`{\displaystyle 𝑑x^+𝑑x^{}𝑑p_{}^2\underset{i,j}{}\frac{d\sigma _{\mathrm{Born}}^{ij}(x^+x^{}\widehat{s},p_{}^2)}{dp_{}^2}}`$
$`\times `$ $`f_{i/\gamma ^{}}(x^+,M_\gamma ^2,Q^2)f_{j/p}(x^{},M_p^2)\theta \left(p_{}^2Q^2\right)`$
where $`\widehat{s}=\stackrel{~}{s}Q^2`$ is the c.m. energy squared for $`\gamma ^{}`$-proton interaction, $`d\sigma _{\mathrm{Born}}^{ij}/dp_{}^2`$ is the differential partonic cross section, $`f_{i/\gamma ^{}}`$ is the parton momentum distribution in the photon, $`M_p^2`$, $`M_\gamma ^2`$ are the factorization scales for the proton and photon correspondingly. As in hadron-hadron scattering, we use $`M_p^2=p_{}^2/4`$, whereas for the photon we take $`M_\gamma ^2=4p_{}^2`$. This requires some more explanation. To the leading order accuracy, the factorization scales are rather undefined as the difference between the results, obtained for different scale choices, is due to higher order corrections. The scheme would be scale independent only after summing up all order contributions both in the structure functions and in the partonic cross section. High $`p_{}`$ jet production in $`\gamma ^{}`$-proton interaction is known to obtain essential contributions from next to leading order (NLO) direct processes . As it was shown in , the sum of the leading order resolved $`\gamma ^{}`$-proton cross section and the NLO direct one exhibits a remarkable independence on the scale $`M_\gamma ^2`$ for the production of parton jets with $`p_{}>Q`$, where $`Q`$ is the photon virtuality. So our strategy is to choose $`M_\gamma ^2`$ such that it allows to represent effectively the full contribution by the leading order resolved cross section.
### 10.3 Parton Momentum Distributions
The cross sections mentioned in the preceding section are all expressed via the so-called parton distribution functions $`f_{i/a}`$, representing the momentum fraction distribution of parton $`i`$ inside particle $`a`$ (proton or photon). In this section, we are going to discuss these distribution functions.
We are first discussing parton distribution functions of the proton. They are represented by the hadronic part of the photon-proton diagram, i.e. the diagram without the external photon line. As mentioned before, this diagram is also a building block of one of the the elementary diagrams of $`pp`$ scattering, and one can therefore repeat literally the argumentation of the chapter 2. In $`pp`$ scattering, we have (apart of the soft one) four contributions, since on each side the parton ladder couples to the nucleon either via a soft Pomeron or it connects directly to a valence quark. In addition, there is a triple Pomeron diagram. Corresponding we have here three contributions, referred to as “sea”, “triple Pomeron”, and “valence” .
#### The Sea Contribution
The sea contribution contains the perturbative parton cascade, described as a parton ladder with strictly ordered virtualities, and the non-perturbative soft block, dominated by the soft Pomeron asymptotics, see fig. 10.5.
We obtain the momentum distribution $`f_{i/p(\mathrm{sea})}^{1\mathrm{I}\mathrm{P}}(x,M_F^2)`$ of the parton $`i`$ at the virtuality scale $`M_F^2`$ in the proton as the convolution of three distributions (see fig. 10.5): the inclusive parton Fock state distribution in the proton $`\stackrel{~}{F}_p^{(1)}(x_0)`$ (see eq. (C.18)),
$$\stackrel{~}{F}_p^{(1)}(x_0)=F_{\mathrm{remn}}^p(1x_0)F_{\mathrm{part}}^p(x_0),$$
(10.19)
the distribution for the parton momentum share in the soft Pomeron $`E_{\mathrm{soft}}^j(z)`$ (see eqs. (B.21-B.22)), and the QCD evolution function $`E_{\mathrm{QCD}}^{ji}(z,Q_0^2,M_F^2)`$:
$`f_{i/p(\mathrm{sea})}^{1\mathrm{I}\mathrm{P}}(x,M_F^2)`$ $`=`$ $`{\displaystyle \underset{j}{}}{\displaystyle _x^1}{\displaystyle \frac{dx_0}{x_0}}{\displaystyle _x^{x_0}}{\displaystyle \frac{dx_1}{x_1}}`$ (10.20)
$`\times \stackrel{~}{F}_p^{(1)}(x_0)E_{\mathrm{soft}}^j\left({\displaystyle \frac{x_1}{x_0}}\right)E_{\mathrm{QCD}}^{ji}({\displaystyle \frac{x}{x_1}},Q_0^2,M_F^2),`$
see fig. 10.5. This equation may be written as
$$f_{i/p(\mathrm{sea})}^{1\mathrm{I}\mathrm{P}}(x,M_F^2)=\underset{j}{}_x^1\frac{dx_1}{x_1}\phi _{j/p(\mathrm{sea})}^{1\mathrm{I}\mathrm{P}}\left(x_1\right)E_{\mathrm{QCD}}^{ji}(\frac{x}{x_1},Q_0^2,M_F^2),$$
(10.21)
where $`\phi _{j/p(\mathrm{sea})}^{1\mathrm{I}\mathrm{P}}\left(x_1\right)`$ by construction corresponds to the distribution at the initial scale $`Q_0^2`$,
$$\phi _{j/p(\mathrm{sea})}^{1\mathrm{I}\mathrm{P}}\left(x_1\right)=_{x_1}^1\frac{dx_0}{x_0}\stackrel{~}{F}_p^{(1)}(x_0)E_{\mathrm{soft}}^j\left(\frac{x_1}{x_0}\right)$$
(10.22)
Here we use the same expressions for $`F_{\mathrm{remn}}^p(x)`$, $`F_{\mathrm{part}}^p(x)`$, and $`E_{\mathrm{soft}}^j(z)`$ as in the case of proton-proton or nucleus-nucleus scattering, see eqs. (C.20-C.21), (B.21-B.22).
#### The Triple Pomeron Contribution
We have to take also into account triple-Pomeron contributions $`f_{i/p(\mathrm{sea})}^{3\mathrm{I}\mathrm{P}}(x,M_F^2)`$ to gluon and sea quark momentum distributions. The latter ones are defined by the diagrams of fig. 5.6 with the upper cut Pomeron being replaced by the “half” of the sea-sea type semihard Pomeron which consists from a soft Pomeron coupled to the triple-Pomeron vertex and the parton ladder describing the perturbative parton evolution from the initial scale $`Q_0^2`$ to the final scale $`M_F^2`$ , see fig. 10.6.
The lower legs are (un)cut Pomerons.
Let us consider the part $`\mathrm{\Delta }\sigma _{pp(\mathrm{sea}\mathrm{sea})}^{3\mathrm{I}\mathrm{P}}(s)`$ of the contribution of cut triple-Pomeron diagram
$$\mathrm{\Delta }\sigma _{pp}^{3\mathrm{I}\mathrm{P}}(s)=\frac{1}{2s}2\mathrm{I}\mathrm{m}T_{pp}^{3\mathrm{I}\mathrm{P}}(s,t=0),$$
corresponding to the semihard sea-sea type parton-parton scattering in the upper Pomeron, which is according to eqs. (5.5-5.6) given as
$`\mathrm{\Delta }\sigma _{pp(\mathrm{sea}\mathrm{sea})}^{3\mathrm{I}\mathrm{P}}(s)`$ $`=`$ $`{\displaystyle \frac{r_{3\mathrm{I}\mathrm{P}}}{2}}\mathrm{Im}[{\displaystyle _0^1}dx^+{\displaystyle _0^1}dx_1^{}{\displaystyle _0^{1x_1^{}}}dx_2^{}F_{\mathrm{remn}}^p(1x^+)F_{\mathrm{remn}}^p(1x_1^{}x_2^{})`$ (10.23)
$`\times 8\pi ^2{\displaystyle \frac{dx_{12}^{}}{x_{12}^{}}\left[\frac{1}{2s^+}\mathrm{Im}T_{\mathrm{sea}\mathrm{sea}}^p(x^+,s^+,0)\right]}`$
$`\times {\displaystyle }dz^+{\displaystyle }d^2q_1_{}{\displaystyle \underset{l=1}{\overset{2}{}}}\left[{\displaystyle \frac{1}{8\pi ^2\widehat{s}_l}}iT^p(x_l^{},\widehat{s}_l,q_1_{}^2)\right]],`$
where we used $`x_{12}^{}=s_0/(x_{12}^+s)`$, $`s^+=x^+x_{12}^{}s`$, $`\widehat{s}_1=z^+s_0x_1^{}/x_{12}^{}`$, $`\widehat{s}_2=(1z^+)s_0x_2^{}/x_{12}^{}`$, and $`T^p`$, $`T_{\mathrm{sea}\mathrm{sea}}^p`$ are defined in (5.15-5.17). Applying the AGK cutting rules, the contribution (10.23) can be written as a sum of three terms,
$$\mathrm{\Delta }\sigma _{pp(\mathrm{sea}\mathrm{sea})}^{3\mathrm{I}\mathrm{P}}(s)=\mathrm{\Delta }\sigma _{pp(\mathrm{sea}\mathrm{sea})}^{3\mathrm{I}\mathrm{P}(0)}(s)+\mathrm{\Delta }\sigma _{pp(\mathrm{sea}\mathrm{sea})}^{3\mathrm{I}\mathrm{P}(1)}(s)+\mathrm{\Delta }\sigma _{pp(\mathrm{sea}\mathrm{sea})}^{3\mathrm{I}\mathrm{P}(2)}(s),$$
(10.24)
corresponding to the diffractive type semihard interaction, screening correction to the usual semihard interaction, and double cut Pomeron contribution, with the weights
$`\mathrm{\Delta }\sigma _{pp(\mathrm{sea}\mathrm{sea})}^{3\mathrm{I}\mathrm{P}(0)}(s)`$ $`=1`$ $`\times \mathrm{\Delta }\sigma _{pp(\mathrm{sea}\mathrm{sea})}^{3\mathrm{I}\mathrm{P}}(s)`$ (10.25)
$`\mathrm{\Delta }\sigma _{pp(\mathrm{sea}\mathrm{sea})}^{3\mathrm{I}\mathrm{P}(1)}(s)`$ $`=+4`$ $`\times \mathrm{\Delta }\sigma _{pp(\mathrm{sea}\mathrm{sea})}^{3\mathrm{I}\mathrm{P}}(s)`$ (10.26)
$`\mathrm{\Delta }\sigma _{pp(\mathrm{sea}\mathrm{sea})}^{3\mathrm{I}\mathrm{P}(2)}(s)`$ $`=2`$ $`\times \mathrm{\Delta }\sigma _{pp(\mathrm{sea}\mathrm{sea})}^{3\mathrm{I}\mathrm{P}}(s)`$ (10.27)
Making use of (5.16), (B.20), (2.2), we have
$`{\displaystyle \frac{1}{2s^+}}\mathrm{Im}T_{\mathrm{sea}\mathrm{sea}}^p(x^+,s^+,0)`$ $`=`$ $`{\displaystyle \frac{1}{2}}F_{\mathrm{part}}^p(x^+){\displaystyle \underset{jk}{}}{\displaystyle _0^1}{\displaystyle \frac{dz_1^+}{z_1^+}}{\displaystyle \frac{dz_1^{}}{z_1^{}}}E_{\mathrm{soft}}^j\left(z_1^+\right)E_{\mathrm{soft}}^k\left(z_1^{}\right)`$ (10.28)
$`\times `$ $`{\displaystyle \underset{ml}{}}{\displaystyle _{z_1^+}^1}𝑑z_B^+{\displaystyle _{z_1^{}}^1}𝑑z_B^{}{\displaystyle 𝑑p_{}^2E_{\mathrm{QCD}}^{jm}(z_B^+/z_1^+,Q_0^2,M_F^2)E_{\mathrm{QCD}}^{kl}(z_B^{}/z_1^{},Q_0^2,M_F^2)}`$
$`\times `$ $`K{\displaystyle \frac{d\sigma _{\mathrm{Born}}^{ml}}{dp_{}^2}}(z_B^+z_B^{}s^+,p_{}^2)\theta \left(M_F^2Q_0^2\right)`$
Now, using (10.19-10.20), we can rewrite (10.23) as
$`\mathrm{\Delta }\sigma _{pp(\mathrm{sea}\mathrm{sea})}^{3\mathrm{I}\mathrm{P}}(s)`$ $`=`$ $`{\displaystyle \underset{ml}{}}{\displaystyle 𝑑x_B^+𝑑x_B^{}𝑑p_{}^2f_{m/p(\mathrm{sea})}^{1\mathrm{I}\mathrm{P}}(x_B^+,M_F^2)f_{l/p(\mathrm{sea})}^{3\mathrm{I}\mathrm{P}}(x_B^{},M_F^2)}`$ (10.29)
$`\times `$ $`K{\displaystyle \frac{d\sigma _{\mathrm{Born}}^{ml}}{dp_{}^2}}(x_B^+x_B^{}s,p_{}^2)\theta \left(M_F^2Q_0^2\right),`$
where we denoted
$$x_B^+=x^+z_B^+,x_B^{}=z_B^{}\frac{s_0}{x_{12}^+s},$$
and defined
$$f_{i/p(\mathrm{sea})}^{3\mathrm{I}\mathrm{P}}(x,M_F^2)=\underset{j}{}_x^1\frac{dx_1}{x_1}\phi _{j/p(\mathrm{sea})}^{3\mathrm{I}\mathrm{P}}\left(x_1\right)E_{\mathrm{QCD}}^{ji}(\frac{x}{x_1},Q_0^2,M_F^2),$$
(10.30)
with
$`\phi _{j/p(\mathrm{sea})}^{3\mathrm{I}\mathrm{P}}\left(x\right)`$ $`=`$ $`{\displaystyle \frac{r_{3\mathrm{I}\mathrm{P}}}{2}}{\displaystyle _x^1}{\displaystyle \frac{dx_{12}}{x_{12}}}E_{\mathrm{soft}}^j\left({\displaystyle \frac{x}{x_{12}}}\right){\displaystyle 𝑑x_1𝑑x_2F_{\mathrm{remn}}^p\left(1x_1x_2\right)}`$ (10.31)
$`\times `$ $`4\pi ^2{\displaystyle 𝑑zd^2q_{}\mathrm{Im}\left[\underset{l=1}{\overset{2}{}}\left[\frac{1}{8\pi ^2\widehat{s}_l}iT^p(x_l,\widehat{s}_l,q_{}^2)\right]\right]}`$
$`=`$ $`{\displaystyle \frac{r_{3\mathrm{I}\mathrm{P}}}{8}}{\displaystyle _x^1}{\displaystyle \frac{dx_{12}}{x_{12}}}E_{\mathrm{soft}}^j\left({\displaystyle \frac{x}{x_{12}}}\right){\displaystyle 𝑑x_1𝑑x_2F_{\mathrm{remn}}^p\left(1x_1^{}x_2^{}\right)}`$
$`\times `$ $`{\displaystyle d^2b𝑑zG^p(x_1,\widehat{s}_1,b)G^p(x_2,\widehat{s}_2,b)},`$
with $`\widehat{s}_1=zs_0x_1/x_{12}`$, $`\widehat{s}_2=(1z)s_0x_2/x_{12}`$, where $`G^p(x,\widehat{s},b)`$ is given in (5.33-5.36 ).
Now, replacing in (10.29) the interaction with the projectile parton $`m`$ by the interaction with a virtual photon probe of virtuality $`M_F^2`$ or by the interaction with a hypotetical probe which couples directly to a gluon, we immediately see that $`f_{i/p(\mathrm{sea})}^{3\mathrm{I}\mathrm{P}}(x,M_F^2)`$ defines the (negative) contribution of the triple-Pomeron diagram to parton structure functions.
#### The Valence Contribution
The third contribution, referred to as “valence”, amounts to the case where a valence quark is the first parton of the ladder, see fig. 10.7.
Here, the soft pre-evolution, governed by the secondary Reggeon, is typically short and therefore not considered explicitly. One simply uses the parameterized input for valence quark momentum distributions at the initial scale $`Q_0^2`$,
$$\phi _{i/p(\mathrm{val})}(x_1)=q_{\mathrm{val}}^i(x_1,Q_0^2)$$
(10.32)
with the Gluck-Reya-Vogt parameterization for $`q_{\mathrm{val}}^i(x,Q_0^2)`$ .
#### The Complete Proton Distribution Function
The total parton distribution in the proton at the initial scale $`Q_0^2`$ is then defined as
$$\phi _{i/p}(x_1)=\phi _{i/p(\mathrm{sea})}(x_1)+\phi _{i/p(\mathrm{val})}(x_1),$$
(10.33)
with
$$\phi _{i/p(\mathrm{sea})}(x_1)=\phi _{i/p(\mathrm{sea})}^{1\mathrm{I}\mathrm{P}}(x_1)+\phi _{i/p(\mathrm{sea})}^{3\mathrm{I}\mathrm{P}}(x_1)$$
(10.34)
which results for an arbitrary scale $`M_F^2`$ in
$$f_{i/p}(x,M_F^2)=\underset{j}{}\frac{dx_1}{x_1}\phi _{i/p}\left(x_1\right)E_{\mathrm{QCD}}^{ji}(\frac{x}{x_1},Q_0^2,M_F^2).$$
(10.35)
For $`M_F^2=Q_0^2`$, the semi-hard contribution is a function which peaks at very small values of $`x`$ and then decreases monotonically towards zero for $`x=1`$. The valence contribution, on the other hand, has a maximum at large values of $`x`$ and goes towards zero for small values of $`x`$. For moderate values of $`M_F^2`$, the precise form of $`f`$ depends crucially on the exponent for the Pomeron-nucleon coupling $`\alpha _{\mathrm{part}}`$, whereas for large $`M_F^2`$ it is mainly defined by the QCD evolution and depends weekly on the initial conditions at the scale $`Q_0^2`$.
#### The Photon Distribution Functions
To calculate the resolved photon cross section (10.2), one needs also to know parton momentum distributions of a virtual photon $`f_{i/\gamma ^{}}(x,M_\gamma ^2,Q^2)`$. According to , $`f_{i/\gamma ^{}}`$ gets contributions both from vector meson states of the photon and from perturbative point-like photon splitting into a quark-anti-quark pair,
$$f_{i/\gamma }(x,M_\gamma ^2,Q^2)=f_{i/\gamma }^{\mathrm{VDM}}(x,M_\gamma ^2,Q^2)+f_{i/\gamma }^{\mathrm{point}}(x,M_\gamma ^2,Q^2).$$
(10.36)
For the former one, one has
$$f_{i/\gamma }^{\mathrm{VDM}}(x,M_\gamma ^2,Q^2)=\eta (Q^2)\alpha \left[G_i^2f_{i/\pi ^0}(x,M_\gamma ^2)+\frac{1}{2}\delta _i\left(G_u^2G_d^2\right)f_{s/\pi ^0}(x,M_\gamma ^2)\right],$$
(10.37)
where the function $`\eta `$ is given as
$$\eta (Q^2)=\left(1+Q^2/m_\rho ^2\right)^2,$$
(10.38)
with $`m_\rho ^2=0.59`$ GeV<sup>2</sup> and with
$$\begin{array}{cccc}\delta _u=1,\hfill & \delta _d=1,\hfill & \delta _s=0,\hfill & \delta _g=0,\hfill \\ G_u^2=0.836,\hfill & G_d^2=0.250,\hfill & G_s^2=0.543,\hfill & G_g^2=0.543.\hfill \end{array}$$
(10.39)
Pion structure functions $`f_{i/\pi ^0}`$ are defined in the same way as proton ones, namely
$$f_{i/\pi ^0}(x,M_F^2)=\underset{j}{}\frac{dx_1}{x_1}\phi _{i/\pi ^0}\left(x_1\right)E_{\mathrm{QCD}}^{ji}(\frac{x}{x_1},Q_0^2,M_F^2),$$
(10.40)
with
$$\phi _{i/\pi ^0}(x_1)=\phi _{i/\pi (\mathrm{sea})}(x_1)+\phi _{i/\pi ^0(\mathrm{val})}(x_1),$$
(10.41)
Here $`\phi _{i/\pi ^0(\mathrm{val})}`$ is a parameterized initial distribution for the valence component from
$$\phi _{i/\pi ^0(\mathrm{val})}(x_1)=q_{\mathrm{val}/\pi ^0}^i(x_1,Q_0^2),$$
(10.42)
and the sea component $`\phi _{j/\pi (\mathrm{sea})}`$ is given by the formulas (10.33-10.34), (10.22), (10.31) , (10.32) with the subscript $`p`$ being replaced by $`\pi `$ and using the appropriate parameters $`\alpha _{\mathrm{remn}}^\pi `$, $`\gamma _\pi `$ in $`F_{\mathrm{part}}^\pi ,F_{\mathrm{remn}}^\pi `$, but keeping all other parameters, characterizing the Pomeron trajectory, unchanged compared to proton case.
The point-like contribution $`f_{i/\gamma }^{\mathrm{point}}`$ is given as a convolution of the photon splitting into a quark-anti-quark pair (with the Altarelli-Parisi splitting function $`P^{\gamma q\overline{q}}(z)=N_c/2(z^2+(1z)^2)`$, $`N_c=3`$ being the number of colors), followed by the QCD evolution of a (anti-)quark from the initial virtuality $`q^2`$ of the splitting till the scale $`M_\gamma ^2`$
$`f_{i/\gamma }^{\mathrm{point}}(x,M_\gamma ^2,Q^2)`$ $`=`$ $`e_q^2{\displaystyle \frac{dq^2}{q^2}_x^1\frac{dx_\gamma }{x_\gamma }\frac{\alpha }{2\pi }P^{\gamma q\overline{q}}(x_\gamma )}`$ (10.43)
$`\times `$ $`{\displaystyle \underset{j\{u,d,s,\overline{u},\overline{d},\overline{s}\}}{}}E_{\mathrm{QCD}}^{ji}({\displaystyle \frac{x}{x_\gamma }},q^2,M_\gamma ^2)\mathrm{\Theta }\left(q^2\mathrm{max}[Q_0^2,x_\gamma Q^2]\right)`$
with $`x_\gamma `$ being the share of the virtual photon light cone momentum taken by the (anti-)quark and with
$$e_q^2=\frac{1}{3}\left(e_u^2+e_d^2+e_s^2\right)$$
(10.44)
being the average light quark charge squared. Here we have chosen the initial scale for the QCD evolution of a $`t`$-channel (anti-)quark to be equal to its initial virtuality $`q^2=x_\gamma Q^2+p_{}^2/(1x_\gamma )`$, with $`p_{}^2`$ being the transverse momentum squared for the splitting in the $`\gamma ^{}p`$ center of mass system .
### 10.4 The Structure Function $`F_2`$
We have now all elements to calculate the structure function $`F_2`$, based on the formula (10.9), with all leading order contributions to $`\sigma _{T(L)}^{\gamma ^{}p}`$ as given in eqs. (10.11-10.2) and with the parton momentum distributions in the proton and photon as discussed in the preceding section. The results for $`F_2(x,Q^2)`$ are shown in fig. 10.8 together with experimental data from H1 , ZEUS and NMC . The parameters affecting the results for $`F_2`$ are actually the same ones which affect parton-parton and hadron-hadron scattering. So we fix them in order to have an overall good fit for $`F_2`$ and at the same time the energy dependence of the total cross section and of the slope parameter. It is possible to get a reasonable agreement, which is of course not perfect due to the fact that enhanced diagrams are only treated to lowest order. In fig. 10.9, we show separately the direct light quarks contribution (10.13), as well as the ones of charm quarks (10.14) and of resolved photons (10.2), for $`Q^2=`$1.5 GeV<sup>2</sup>, 5 GeV<sup>2</sup>, and 25 GeV<sup>2</sup>. It is easy to see that the resolved photon cross section contributes significantly to $`F_2(x,Q^2)`$ for small photon virtualities $`Q^2`$.
### 10.5 Parton Configurations: Basic Formulas
In order to have a coherent approach, we base our treatment of particle production on exactly the same formulas as derived earlier for the cross sections. To be more precise, we take the formulas for $`d\sigma /dx_BdQ^2`$ as a basis for treating particle production, which means first of all parton production.
The differential cross section for lepton-nucleon scattering is given in eq. (10.7). Using (B.28), (B.30), the $`QCD`$ evolution function $`E_{\mathrm{QCD}}^{ij}`$, which enters into the formulas for the cross sections $`\sigma _{T(L)}^{\gamma ^{}p}(\stackrel{~}{s},Q^2)`$, can be expanded as a sum over $`n`$-rung ladder contributions, where the latter ones can be written as an integration over the momenta $`p_1,p_2,\mathrm{},p_n`$ of $`n`$ resolvable final partons. Introducing a multidimensional variable
$$P=\{p_1,p_2,\mathrm{},p_n\},$$
(10.45)
and considering the symbol $`_P`$ representing $`_n𝑑P_n`$, with $`dP_n`$ being the invariant phase space volume for $`n`$-parton state, we may write
$$\frac{d\sigma _{lp}}{dx_BdQ^2}=\underset{P}{}\sigma (x_B,Q^2,P).$$
(10.46)
After normalization, $`\sigma (x_B,Q^2,P)`$ may be interpreted as the probability distribution for a parton configuration $`P`$ for given values of $`x_B`$ and $`Q^2`$. The Monte Carlo method provides a convenient tool for treating such multidimensional distributions: with $`\sigma (x_B,Q^2,P)`$ being known (see preceding sections), one generates parton configurations $`P`$ according to this distribution. In addition to $`x_B`$, $`Q^2`$, and $`P`$, additional variables occur, specifying a particular contribution to the DIS cross section. One essentially follows the structure of the formula for the cross section. Let us discuss the procedure to generate parton configurations in detail.
We start with some useful definitions. Using the relation (B.28) for the evolution function $`E_{\mathrm{QCD}}^{ij}`$, any parton momentum distribution at the scale $`Q^2`$ can be decomposed into two contributions, corresponding to the case of no resolvable emission in the range of virtualities between $`Q_0^2`$ and $`Q^2`$ and to at least one resolvable emission:
$$f_j(x,Q^2)=f_j(x,Q_0^2)\mathrm{\Delta }^j(Q_0^2,Q^2)+\underset{i}{}_x^{1ϵ}\frac{dz}{z}\overline{E}_{\mathrm{QCD}}^{ij}(z,Q_0^2,Q^2)f_i(\frac{x}{z},Q_0^2)$$
(10.47)
In case of “$`i`$” and “$`j`$” being quarks, we split the sum $`_j\overline{E}_{\mathrm{QCD}}^{ij}`$ into two components,
$$\underset{j}{}\overline{E}_{\mathrm{QCD}}^{ij}=\overline{E}_{\mathrm{NS}}+\overline{E}_\mathrm{S}(i,j=\mathrm{quark}s),$$
(10.48)
with the so-called non-singlet evolution $`\overline{E}_{\mathrm{NS}}`$, where only gluons are emitted as final $`s`$-channel partons, and the singlet evolution $`\overline{E}_\mathrm{S}`$ representing all the other contributions, see fig. 10.10a,c.
The non-singlet evolution (compare with eq. (B.29)) satisfies the evolution equation
$`\overline{E}_{\mathrm{NS}}(x,Q_0^2,Q^2)=`$ (10.49)
$`{\displaystyle _{Q_0^2}^{Q^2}}{\displaystyle \frac{dQ_1^2}{Q_1^2}}\left[{\displaystyle _x^{1ϵ}}{\displaystyle \frac{dz}{z}}{\displaystyle \frac{\alpha _s}{2\pi }}P_q^q(z)\overline{E}_{\mathrm{NS}}({\displaystyle \frac{x}{z}},Q_0^2,Q_1^2)\mathrm{\Delta }^q(Q_1^2,Q^2)+{\displaystyle \frac{\alpha _s}{2\pi }}P_q^q(x)\mathrm{\Delta }^q(Q_0^2,Q^2)\right],`$
and the singlet one
$`\overline{E}_\mathrm{S}(x,Q_0^2,Q^2)={\displaystyle _{Q_0^2}^{Q^2}}{\displaystyle \frac{dQ_1^2}{Q_1^2}}{\displaystyle _x^{1ϵ}}{\displaystyle \frac{dz}{z}}{\displaystyle \frac{\alpha _s}{2\pi }}[P_q^q(z)\overline{E}_\mathrm{S}({\displaystyle \frac{x}{z}},Q_0^2,Q_1^2)`$
$`+\mathrm{\hspace{0.33em}2}n_fP_g^q(z)\overline{E}_{\mathrm{QCD}}^{qg}({\displaystyle \frac{x}{z}},Q_0^2,Q_1^2)]\mathrm{\Delta }^q(Q_1^2,Q^2),`$ (10.50)
with $`n_f=3`$ being the number of active quark flavors.
Now it is convenient to define parton level cross sections, corresponding to different contributions to the deep inelastic scattering process and to different partons, entering the perturbative evolution at the initial scale $`Q_0^2`$, see fig. 10.10. Essentially, we include into the cross sections the perturbative part of parton evolution, whereas the initial conditions, given by parton momentum densities at the initial scale $`Q_0^2`$, are factorized out. The non-singlet and singlet contributions to the direct (light) photon-quark interaction are defined as
$$\sigma _{T(\mathrm{NS})}^{\gamma ^{}q}(\stackrel{~}{s},Q^2,Q_0^2)=\frac{4\pi ^2\alpha }{\stackrel{~}{s}}\overline{E}_{\mathrm{NS}}(\frac{Q^2}{\stackrel{~}{s}},Q_0^2,Q^2)$$
(10.51)
$$\sigma _{T(\mathrm{S})}^{\gamma ^{}q}(\stackrel{~}{s},Q^2,Q_0^2)=\frac{4\pi ^2\alpha }{\stackrel{~}{s}}\overline{E}_\mathrm{S}(\frac{Q^2}{\stackrel{~}{s}},Q_0^2,Q^2)$$
(10.52)
For the direct (light) photon-gluon interaction, we define
$$\sigma _{T(\mathrm{light})}^{\gamma ^{}g}(\stackrel{~}{s},Q^2,Q_0^2)=\frac{4\pi ^2\alpha }{\stackrel{~}{s}}\overline{E}_{\mathrm{QCD}}^{gq}(\frac{Q^2}{\stackrel{~}{s}},Q_0^2,Q^2)$$
(10.53)
The photon-parton charm production cross section is defined as
$$\sigma _{T/L(\mathrm{charm})}^{\gamma ^{}i}(\stackrel{~}{s},Q^2,Q_0^2)=𝑑x^{}𝑑p_{}^2E_{\mathrm{QCD}}^{ig}(x^{},Q_0^2,M_F^2)\frac{d\sigma _{T/L}^{\gamma ^{}gc\overline{c}}(x^{}\stackrel{~}{s},Q^2,p_{}^2)}{dp_{}^2},$$
(10.54)
where the parton ($`i`$) may be a quark or a gluon. Finally, we define the parton-parton cross section for resolved processes similarly to (6.1) as
$`\sigma _{T(\mathrm{resolved})}^{ij}(\widehat{s},Q^2,q^2,Q_0^2)`$ $`=`$ $`{\displaystyle 𝑑x^+𝑑x^{}𝑑p_{}^2\underset{k,l}{}\frac{d\sigma _{\mathrm{Born}}^{kl}(x^+x^{}\widehat{s},p_{}^2)}{dp_{}^2}}`$
$`\times `$ $`E_{\mathrm{QCD}}^{ik}(x^+,q^2,M_\gamma ^2)E_{\mathrm{QCD}}^{jl}(x^{_{}},Q_0^2,M_p^2)\theta \left(p_{}^2Q^2\right).`$
Based on the above partial photon-parton cross sections, we define now the total photon-parton cross sections, summed over the quark flavors of the quark coupling to the photon with the appropriate squared charge ($`e^2`$) factor. We obtain for the photon-gluon cross section
$`\sigma _T^{\gamma ^{}g}(\stackrel{~}{s},Q^2,Q_0^2)`$ $`=`$ $`e_q^2\sigma _{T(\mathrm{light})}^{\gamma ^{}g}(\stackrel{~}{s},Q^2,Q_0^2)+e_c^2\sigma _{T(\mathrm{charm})}^{\gamma ^{}g}(\stackrel{~}{s},Q^2,Q_0^2)`$
$`+`$ $`{\displaystyle \underset{j}{}}{\displaystyle 𝑑x_\gamma f_{j/\gamma }^{\mathrm{VDM}}(x_\gamma ,Q_0^2,Q^2)\sigma _{T(\mathrm{resolved})}^{jg}(x_\gamma (\stackrel{~}{s}Q^2),Q^2,Q_0^2,Q_0^2)}`$
$`+`$ $`e_q^2{\displaystyle \frac{dq^2}{q^2}𝑑x_\gamma \frac{\alpha }{2\pi }P^{\gamma q\overline{q}}(x_\gamma )}`$
$`\times `$ $`{\displaystyle \underset{j\{u,d,s,\overline{u},\overline{d},\overline{s}\}}{}}\sigma _{T(\mathrm{resolved})}^{jg}(x_\gamma \stackrel{~}{s}q^2,Q^2,q^2,Q_0^2)\mathrm{\Theta }\left(q^2\mathrm{max}[Q_0^2,x_\gamma Q^2]\right).`$
The photon-quark cross section for a quark with flavor “$`i`$” is given as
$`\sigma _T^{\gamma ^{}i}(\stackrel{~}{s},Q^2,Q_0^2)`$ $`=`$ $`e_i^2\sigma _{T(\mathrm{NS})}^{\gamma ^{}q}(\stackrel{~}{s},Q^2,Q_0^2)+e_q^2\sigma _{T(\mathrm{S})}^{\gamma ^{}q}(\stackrel{~}{s},Q^2,Q_0^2)`$
$`+`$ $`e_c^2\sigma _{T(\mathrm{charm})}^{\gamma ^{}g}(\stackrel{~}{s},Q^2,Q_0^2)`$
$`+`$ $`{\displaystyle \underset{j}{}}{\displaystyle 𝑑x_\gamma f_{j/\gamma }^{\mathrm{VDM}}(x_\gamma ,Q_0^2,Q^2)\sigma _{T(\mathrm{resolved})}^{ji}(x_\gamma (\stackrel{~}{s}Q^2),Q^2,Q_0^2,Q_0^2)}`$
$`+`$ $`e_q^2{\displaystyle \frac{dq^2}{q^2}𝑑x_\gamma \frac{\alpha }{2\pi }P^{\gamma q\overline{q}}(x_\gamma )}`$
$`\times `$ $`{\displaystyle \underset{j\{u,d,s,\overline{u},\overline{d},\overline{s}\}}{}}\sigma _{T(\mathrm{resolved})}^{ji}(x_\gamma \stackrel{~}{s}q^2,Q^2,q^2,Q_0^2)\mathrm{\Theta }\left(q^2\mathrm{max}[Q_0^2,x_\gamma Q^2]\right)`$
The longitudinal photon-parton cross section for a parton of flavor $`i`$ (quark or gluon) is finally given as
$$\sigma _L^{\gamma ^{}i}(\stackrel{~}{s},Q^2,Q_0^2)=e_c^2\sigma _{L(\mathrm{charm})}^{\gamma ^{}i}(\stackrel{~}{s},Q^2,Q_0^2).$$
(10.58)
Finally, we may express the photon-proton cross sections in terms of the above photon-parton cross sections and the parton distribution functions $`\phi (x)=f(x,Q_0^2)`$ at the scale $`Q_0^2`$. The transverse cross section is given as
$`\sigma _T^{\gamma ^{}p}(\stackrel{~}{s},Q^2)`$ $`=`$ $`{\displaystyle 𝑑x_1\phi _{g/p}\left(x_1\right)\sigma _T^{\gamma ^{}g}(x_1\stackrel{~}{s},Q^2,Q_0^2)}`$
$`+`$ $`{\displaystyle \underset{i\{u,d,s,\overline{u},\overline{d},\overline{s}\}}{}}\left({\displaystyle \frac{4\pi ^2\alpha }{Q^2}}e_i^2x_B\phi _{i/p}(x_B)\mathrm{\Delta }^q(Q_0^2,Q^2)+{\displaystyle 𝑑x_1\phi _{i/p}(x_1)\sigma _T^{\gamma ^{}i}(x_1\stackrel{~}{s},Q^2,Q_0^2)}\right),`$
the longitudinal cross section is given as
$`\sigma _L^{\gamma ^{}p}(\stackrel{~}{s},Q^2)`$ $`=`$ $`{\displaystyle 𝑑x_1\phi _{g/p}\left(x_1\right)\sigma _L^{\gamma ^{}g}(x_1\stackrel{~}{s},Q^2,Q_0^2)}`$
$`+`$ $`{\displaystyle \underset{i\{u,d,s,\overline{u},\overline{d},\overline{s}\}}{}}{\displaystyle 𝑑x_1\phi _{i/p}(x_1)\sigma _L^{\gamma ^{}i}(x_1\stackrel{~}{s},Q^2,Q_0^2)},`$
with the quark momentum distributions being a sum of two terms,
$$\phi _{i/p}(x_1)=\phi _{i/p(\mathrm{sea})}(x_1)+\phi _{i/p(\mathrm{val})}(x_1),$$
(10.61)
see eqs. (10.22-10.33).
The above formulas together with eq. (10.7) serve as the basis to generate all main variables for the description of deep inelastic scattering. After modeling $`Q^2`$ and $`x_B`$ we simulate types (a valence quark, a sea quark, or a gluon) and kinematical characteristics for first partons, entering the perturbative evolution (at the initial scale $`Q_0^2)`$, and then, for given initial conditions, generate corresponding parton configurations, based on the particular structure of perturbative cross sections (10.51-10.5). The detailed description of this procedure is given in the next sections.
The triple Pomeron contributions are included here in the definition of the parton distribution $`\phi _{i/p(\mathrm{sea})}`$. At HERA energies, the triple Pomeron contribution is dominated by the process where the two Pomerons exchanged in parallel are soft ones, and therefore no additional parton production needs to be considered in that case.
### 10.6 Generating Initial Conditions for the Perturbative Evolution
We start with the generation of the kinematical variables $`x_B`$ and $`Q^2`$ according to the differential cross section eq. (10.7) together with the explicit form for the photon-proton cross sections eqs. (10.5-10.5). Then we choose an interaction with the transverse or with the longitudinal polarization component of the photon, with the corresponding weights
$$L_{T/L}^\gamma ^{}(y)\sigma _{T/L}^{\gamma ^{}p}(\stackrel{~}{s},Q^2)/\left(L_T^\gamma ^{}(y)\sigma _T^{\gamma ^{}p}(\stackrel{~}{s},Q^2)+L_L^\gamma ^{}(y)\sigma _L^{\gamma ^{}p}(\stackrel{~}{s},Q^2)\right).$$
(10.62)
After that, we consider virtual photon-proton interaction for given photon virtuality $`Q^2`$ and polarization ($`T`$, $`L`$), and for given c.m. energy squared $`\widehat{s}=\stackrel{~}{s}Q^2`$ for the interaction; we use the photon-proton center of mass system.
Let us first consider the case of transverse photon polarization. We have to choose between “sea” and “valence ” contribution, where the latter one is chosen with the probability
$`\mathrm{prob}(\mathrm{val})={\displaystyle \frac{1}{\sigma _T^{\gamma ^{}p}(\stackrel{~}{s},Q^2)}}\{{\displaystyle \frac{4\pi ^2\alpha }{Q^2}}{\displaystyle \underset{i\{u,d\}}{}}e_i^2x_B\phi _{i/p(\mathrm{val})}(x_B)\mathrm{\Delta }^q(Q_0^2,Q^2)`$ $`+`$
$`+{\displaystyle \underset{i\{u,d\}}{}}{\displaystyle }dx_1\phi _{i/p(\mathrm{val})}(x_1)\sigma _T^{\gamma ^{}i}(x_1\stackrel{~}{s},Q^2,Q_0^2)\}`$ (10.63)
and the former one with the probability $`1\mathrm{prob}(\mathrm{val})`$, which is the sum of the contributions, corresponding to a gluon or a sea quark from the proton being the first ladder parton,
$`\mathrm{prob}(\mathrm{sea})`$ $`=`$ $`{\displaystyle \frac{1}{\sigma _T^{\gamma ^{}p}(\stackrel{~}{s},Q^2)}}\{{\displaystyle \frac{4\pi ^2\alpha }{Q^2}}{\displaystyle \underset{i\{u,d,s,\overline{u},\overline{d},\overline{s}\}}{}}e_i^2x_B\phi _{i/p(\mathrm{sea})}(x_B)\mathrm{\Delta }^q(Q_0^2,Q^2)+`$ (10.64)
$`+{\displaystyle \underset{i\{g,u,d,s,\overline{u},\overline{d},\overline{s}\}}{}}{\displaystyle }dx_1\phi _{i/p(\mathrm{sea})}(x_1)\sigma _T^{\gamma ^{}i}(x_1\stackrel{~}{s},Q^2,Q_0^2)\}.`$
The next step consists of defining the type (flavor) of the first ladder parton and its momentum share $`x_1`$ in the proton. Here one has to distinguish two possible parton configurations for the interaction: parton cascading without any resolvable parton emission in the ladder (corresponding to a valence or a sea quark of the proton scattered back in the Breit frame), with the relative weights given by the first term in the curly brackets in eqs. (10.63, 10.64), i.e.
$$\mathrm{prob}(\mathrm{no}\mathrm{emission})=\frac{\frac{4\pi ^2\alpha }{Q^2}_ie_i^2x_B\phi _{\mathrm{i}/p(\mathrm{val}/\mathrm{sea})}(x_B)\mathrm{\Delta }^q(Q_0^2,Q^2)}{\mathrm{prob}(\mathrm{val}/\mathrm{sea})\sigma _T^{\gamma ^{}p}(\stackrel{~}{s},Q^2)},$$
(10.65)
and the configurations with at least one resolvable parton emission, with the weight $`1\mathrm{prob}(\mathrm{no}\mathrm{emission})`$. In the case of no resolvable emission we have $`x_1=x_B=Q^2/\stackrel{~}{s}`$ and the flavor of the quark is generated according to the weights
$$\{\begin{array}{ccc}e_i^2\phi _{\mathrm{i}/p(\mathrm{val})}(x_B)\hfill & \mathrm{if}\hfill & \mathrm{val}\hfill \\ e_i^2\phi _{i/p(\mathrm{sea})}(x_B)\hfill & \mathrm{if}\hfill & \mathrm{sea}\hfill \end{array}.$$
(10.66)
Then we are left with a trivial parton configuration. For a valence quark contribution it consists of an anti-quark, moving along the original proton direction, and the quark, scattered back.
In the case of at least one resolvable emission, one generates the type (flavor) $`i`$ and the light cone momentum fraction $`x_1`$ of the first parton of the QCD cascade according to the distributions, given by the expressions in the curly brackets in (10.63), (10.64),
$$\mathrm{prob}(i,x_1)=\{\begin{array}{ccc}\delta _q^i\phi _{\mathrm{i}/p(\mathrm{val})}(x_1)\sigma _T^{\gamma ^{}i}(x_1\stackrel{~}{s},Q^2,Q_0^2)\hfill & \mathrm{if}\hfill & \mathrm{val}\hfill \\ \phi _{i/p(\mathrm{sea})}(x_1)\sigma _T^{\gamma ^{}i}(x_1\stackrel{~}{s},Q^2,Q_0^2)\hfill & \mathrm{if}\hfill & \mathrm{sea}\hfill \end{array},$$
(10.67)
where $`\delta _q^i`$ is zero if $`i=g`$ and otherwise one.
Then one chooses between different types of interactions, contributing to the photon-parton cross section, according to their partial weights in eqs. (10.5, 10.5). The direct photon-parton contribution is chosen with the weight
$$\mathrm{prob}(\mathrm{direct})=\{\begin{array}{ccc}\left[e_i^2\sigma _{T(\mathrm{NS})}^{\gamma ^{}q}+e_q^2\sigma _{T(\mathrm{S})}^{\gamma ^{}q}+e_c^2\sigma _{T(\mathrm{charm})}^{\gamma ^{}q}\right]/\sigma _T^{\gamma ^{}i}\hfill & \mathrm{if}\hfill & i=\mathrm{quark}\hfill \\ \left[e_q^2\sigma _{T(\mathrm{light})}^{\gamma ^{}g}+e_c^2\sigma _{T(\mathrm{charm})}^{\gamma ^{}g}\right]/\sigma _T^{\gamma ^{}g}\hfill & \mathrm{if}\hfill & i=\mathrm{gluon}\hfill \end{array}.$$
(10.68)
For the resolved photon contributions, the weight is given by the last two terms in eqs. (10.5, 10.5). The probability for the VDM part is
$$\mathrm{prob}(\mathrm{VDM})=\underset{j}{}𝑑x_\gamma f_{j/\gamma }^{\mathrm{VDM}}(x_\gamma ,Q_0^2,Q^2)\sigma _{T(\mathrm{resolved})}^{ji}(x_\gamma (x_1\stackrel{~}{s}Q^2),Q^2,Q_0^2,Q_0^2)/\sigma _T^{\gamma ^{}i},$$
(10.69)
where $`x_1`$ is the light cone momentum share and $`i`$ is the type (flavor) of the first ladder parton on the proton side (already determined), whereas the probability for the point-like resolved contribution is
$`\mathrm{prob}(\mathrm{point})`$ $`=`$ $`e_q^2{\displaystyle \frac{dq^2}{q^2}𝑑x_\gamma \frac{\alpha }{2\pi }P^{\gamma q\overline{q}}(x_\gamma )}`$
$`\times `$ $`{\displaystyle \underset{j\{u,d,s,\overline{u},\overline{d},\overline{s}\}}{}}\sigma _{T(\mathrm{resolved})}^{ji}(x_\gamma x_1\stackrel{~}{s}q^2,Q^2,q^2,Q_0^2)\mathrm{\Theta }\left(q^2\mathrm{max}[Q_0^2,x_\gamma Q^2]\right)/\sigma _T^{\gamma ^{}i}.`$
In case of a direct light contribution, one has to generate the configuration for a parton ladder, strictly ordered in parton virtualities towards the virtual photon. The method is quite analogous to the one of section 2 in chapter 5 and is described in the next section.
In case of a resolved contribution, we need to define the initial conditions for the other end of the parton ladder, on the photon side, as well as the parton type (a quark of some flavor or a gluon) and the share of the light cone momentum fraction $`x_\gamma `$, taken by the parton from the photon. For the direct resolved contribution, corresponding to the point-like photon splitting into a quark-anti-quark pair, the flavor $`j`$, the share $`x_\gamma `$, and the virtuality $`q^2`$ of the (anti-)quark, being the first ladder parton, are generated according to
$$\mathrm{prob}(j,x_\gamma ,q^2)\frac{1}{q^2}\frac{\alpha }{2\pi }P^{\gamma q\overline{q}}(x_\gamma )\sigma _{T(\mathrm{resolved})}^{ji}(x_\gamma x_1\stackrel{~}{s}q^2,Q^2,q^2,Q_0^2)\mathrm{\Theta }\left(q^2\mathrm{max}[Q_0^2,x_\gamma Q^2]\right)\delta _q^j$$
(10.71)
For the VDM contribution, the first ladder parton is taken at the initial virtuality $`Q_0^2`$ and its type $`j`$ and momentum share $`x_\gamma `$ are chosen according to the distribution
$$\mathrm{prob}(j,x_\gamma )f_{j/\gamma }^{\mathrm{VDM}}(x_\gamma ,Q_0^2,Q^2)\sigma _{T(\mathrm{resolved})}^{ji}(x_\gamma (x_1\stackrel{~}{s}Q^2),Q^2,Q_0^2,Q_0^2),$$
(10.72)
with the VDM parton momentum distributions in the photon defined in (10.37). In both cases for resolved photon interactions, the simulation of parton configurations, corresponding to the ladder of given mass squared $`\widehat{s}^{}`$ ($`\widehat{s}^{}=x_\gamma x_1\stackrel{~}{s}q^2`$ for the direct resolved contribution and $`\widehat{s}^{}=x_\gamma (x_1\stackrel{~}{s}Q^2)`$ for the VDM one), and of given types and virtualities of the leg partons, is done exactly in the same way as for proton-proton (nucleus-nucleus) interactions, as described in the chapter 5. The only difference comes from the presence of two different scales $`M_p^2`$, $`M_\gamma ^2`$ and the cutoff $`p_{}^2>Q^2`$ in the parton-parton cross section $`\sigma _{T(\mathrm{resolved})}^{ij}`$ for resolved DIS processes as given in eq. (10.5), when compared to the cross section in eq. (6.1).
In the case of interaction with the longitudinal photon component, the procedure simplifies considerably, as one only has to consider direct photon-parton interactions via the parton-gluon fusion process. One starts by choosing between the coupling of the parton ladder to a valence quark (“val”) or to a soft Pomeron (“sea”), the weights are
$`\mathrm{prob}(\mathrm{val})`$ $`=`$ $`{\displaystyle \frac{𝑑x_1_i\phi _{i/p(\mathrm{val})}(x_1)\sigma _L^{\gamma ^{}i}(x_1\stackrel{~}{s},Q^2,Q_0^2)}{\sigma _L^{\gamma ^{}p}(\stackrel{~}{s},Q^2)}};`$ (10.73)
$`\mathrm{prob}(\mathrm{sea})`$ $`=`$ $`{\displaystyle \frac{𝑑x_1_i\phi _{i/p(\mathrm{sea})}(x_1)\sigma _L^{\gamma ^{}i}(x_1\stackrel{~}{s},Q^2,Q_0^2)}{\sigma _L^{\gamma ^{}p}(\stackrel{~}{s},Q^2)}}.`$ (10.74)
The weight for the parton type (flavor) $`i`$ and the distribution for the light cone momentum fraction $`x_1`$ of the first parton of the QCD cascade is given by the integrands of (10.73-10.74). The final step amounts to generating the configuration for the parton ladder, strictly ordered in parton virtualities towards the virtual photon, with the largest momentum transfer parton process of parton-gluon fusion type, as discussed in the next section.
### 10.7 Generating the Ladder Partons
In this section we describe the procedure to generate parton configurations, corresponding to direct photon-parton interaction with at least one resolvable emission in the parton cascade. In that case, a parton ladder is strictly ordered in parton virtualities towards the virtual photon and the ladder cross section is given as a sum of the contributions eqs. (10.51-10.54) in the case of transverse photon polarization,
$$\sigma _{T(\mathrm{direct})}^{\gamma ^{}i}(\stackrel{~}{s},Q^2,Q_0^2)=\{\begin{array}{ccc}e_i^2\sigma _{T(\mathrm{NS})}^{\gamma ^{}q}(\stackrel{~}{s},Q^2,Q_0^2)+e_q^2\sigma _{T(\mathrm{S})}^{\gamma ^{}q}(\stackrel{~}{s},Q^2,Q_0^2)\hfill & & \\ +e_c^2\sigma _{T(\mathrm{charm})}^{\gamma ^{}q}(\stackrel{~}{s},Q^2,Q_0^2)\hfill & i=q,\hfill & \\ e_q^2\sigma ^{\gamma ^{}g}(\stackrel{~}{s},Q^2,Q_0^2)+e_c^2\sigma _T^{\gamma ^{}g(c\overline{c})}(\stackrel{~}{s},Q^2,Q_0^2)\hfill & i=g,\hfill & \end{array}$$
(10.75)
or by the cross section eq. (10.54) for the longitudinal one,
$$\sigma _{L(\mathrm{direct})}^{\gamma ^{}i}(\stackrel{~}{s},Q^2,Q_0^2)=e_c^2\sigma _{T(\mathrm{charm})}^{\gamma ^{}i}(\stackrel{~}{s},Q^2,Q_0^2).$$
(10.76)
All the photon-parton cross sections are expressed in terms of the QCD evolution functions $`\overline{E}_{\mathrm{QCD}}`$. Using the explicit representation eqs. (B.30-B.32) for $`\overline{E}_{\mathrm{QCD}}`$, one can rewrite the recursive relations eqs. (B.29), (10.5-10.50) in a form, such that the first (lowest virtuality) emission in the ladder is treated explicitly, multiplied by a weight factor, given by the contribution of the rest of the ladder (the sum of any number of additional resolvable emissions),
$`\overline{E}_{\mathrm{QCD}}^{ij}(x,Q_0^2,Q^2)`$ $`=`$ $`{\displaystyle _{Q_0^2}^{Q^2}}{\displaystyle \frac{dQ_1^2}{Q_1^2}}{\displaystyle \underset{k}{}}{\displaystyle _x^{1ϵ}}{\displaystyle \frac{dz}{z}}{\displaystyle \frac{\alpha _s}{2\pi }}P_i^k(z)\mathrm{\Delta }^i(Q_0^2,Q_1^2)\overline{E}_{\mathrm{QCD}}^{kj}({\displaystyle \frac{x}{z}},Q_1^2,Q^2)`$ (10.77)
$`+`$ $`{\displaystyle _{Q_0^2}^{Q^2}}{\displaystyle \frac{dQ_1^2}{Q_1^2}}\mathrm{\Delta }^i(Q_0^2,Q_1^2)\mathrm{\Delta }^j(Q_1^2,Q^2){\displaystyle \frac{\alpha _s}{2\pi }}P_i^j(x)`$
$`\overline{E}_{NS}(x,Q_0^2,Q^2)`$ $`=`$ $`{\displaystyle _{Q_0^2}^{Q^2}}{\displaystyle \frac{dQ_1^2}{Q_1^2}}{\displaystyle _x^{1ϵ}}{\displaystyle \frac{dz}{z}}{\displaystyle \frac{\alpha _s}{2\pi }}P_q^q(z)\mathrm{\Delta }^q(Q_0^2,Q_1^2)\overline{E}_{NS}({\displaystyle \frac{x}{z}},Q_1^2,Q^2)`$ (10.78)
$`+`$ $`\mathrm{\Delta }^q(Q_0^2,Q^2){\displaystyle _{Q_0^2}^{Q^2}}{\displaystyle \frac{dQ_1^2}{Q_1^2}}{\displaystyle \frac{\alpha _s}{2\pi }}P_q^q(x)`$
$`\overline{E}_S(x,Q_0^2,Q^2)`$ $`=`$ $`{\displaystyle _{Q_0^2}^{Q^2}}{\displaystyle \frac{dQ_1^2}{Q_1^2}}{\displaystyle \underset{k}{}}{\displaystyle _x^{1ϵ}}{\displaystyle \frac{dz}{z}}\mathrm{\Delta }^q(Q_0^2,Q_1^2){\displaystyle \frac{\alpha _s}{2\pi }}[P_q^q(z)\overline{E}_\mathrm{S}({\displaystyle \frac{x}{z}},Q_1^2,Q^2)`$ (10.79)
$`+`$ $`P_q^g(z)\overline{E}_{\mathrm{QCD}}^{gq}({\displaystyle \frac{x}{z}},Q_1^2,Q^2)]`$
With the help of eqs. (10.77-10.79), (10.51-10.54), one can obtain the recursive relations for the cross sections eqs. (10.75, 10.76) for an arbitrary virtuality $`Q_1^2`$ of the initial parton $`i`$,
$`\sigma _{T/L(\mathrm{direct})}^{\gamma ^{}i}(\stackrel{~}{s},Q^2,Q_1^2)`$ $`=`$ $`{\displaystyle \underset{j}{}}{\displaystyle _{Q_1^2}^{Q^2}}{\displaystyle \frac{d\stackrel{~}{Q}^2}{\stackrel{~}{Q}^2}}\mathrm{\Delta }^i(Q_1^2,\stackrel{~}{Q}^2){\displaystyle 𝑑z\frac{\alpha _s}{2\pi }P_i^j(z)\sigma _{T/L(\mathrm{direct})}^{\gamma ^{}j}(z\stackrel{~}{s},Q^2,\stackrel{~}{Q}^2)}+`$ (10.80)
$`+`$ $`\sigma _{T/L(\mathrm{direct}\mathrm{\hspace{0.17em}2}2)}^{\gamma ^{}i}(\stackrel{~}{s},Q^2,Q_1^2),`$
where $`\sigma _{T/L(\mathrm{direct}\mathrm{\hspace{0.17em}2}2)}^{\gamma ^{}i}`$ represents the contribution of parton configurations with only one resolvable parton emission in the ladder, or with the photon-gluon fusion process without any additional resolvable parton emissions,
$`\sigma _{T(\mathrm{direct}\mathrm{\hspace{0.17em}2}2)}^{\gamma ^{}i}(\stackrel{~}{s},Q^2,Q_1^2)`$ $`=`$ $`{\displaystyle \frac{4\pi ^2\alpha e_i^2}{\stackrel{~}{s}}}E_{\mathrm{QCD}}^{(1)qq}({\displaystyle \frac{Q^2}{\stackrel{~}{s}}},Q_1^2,Q^2)\mathrm{\Theta }(Q^2Q_1^2)(i=\mathrm{quark})`$ (10.81)
$`\sigma _{T(\mathrm{direct}\mathrm{\hspace{0.17em}2}2)}^{\gamma ^{}g}(\stackrel{~}{s},Q^2,Q_1^2)`$ $`=`$ $`{\displaystyle \frac{4\pi ^2\alpha e_q^2}{\stackrel{~}{s}}}E_{\mathrm{QCD}}^{(1)gq}({\displaystyle \frac{Q^2}{\stackrel{~}{s}}},Q_1^2,Q^2)\mathrm{\Theta }(Q^2Q_1^2)`$ (10.82)
$`+`$ $`e_c^2{\displaystyle 𝑑p_{}^2\frac{d\sigma _T^{\gamma ^{}gc\overline{c}}(\stackrel{~}{s},Q^2,p_{}^2)}{dp_{}^2}\mathrm{\Delta }^g(Q_1^2,M_F^2)}`$
$`\sigma _{L(\mathrm{direct}\mathrm{\hspace{0.17em}2}2)}^{\gamma ^{}i}(\stackrel{~}{s},Q^2,Q_1^2)`$ $`=`$ $`\delta _i^ge_c^2{\displaystyle 𝑑p_{}^2\frac{d\sigma _L^{\gamma ^{}gc\overline{c}}(\stackrel{~}{s},Q^2,p_{}^2)}{dp_{}^2}\mathrm{\Delta }^g(Q_1^2,M_F^2)}`$ (10.83)
It is noteworthy that to the leading logarithmic accuracy one should not use in $`\sigma _{T(\mathrm{direct}\mathrm{\hspace{0.17em}2}2)}^{\gamma ^{}i}`$ the Born process matrix elements $`d\sigma _T^{\gamma ^{}gq\overline{q}}/dp_{}^2`$ and $`d\sigma _T^{\gamma ^{}qgq}/dp_{}^2`$; the contribution of just one resolvable parton emission is proportional to $`E_{\mathrm{QCD}}^{(1)iq}`$, eq. (B.32), defined by the corresponding Altarelli-Parisi kernel $`P_i^q(z)`$.
The formulas (10.80-10.83) allow us to generate the cascade of partons, corresponding to the direct $`\gamma ^{}`$-quark (gluon) interaction of energy squared $`\widehat{s}=\stackrel{~}{s}Q^2`$ and photon virtuality $`Q^2`$, starting from an initial parton with a flavor $`i`$, taken at a scale $`Q_1^2=Q_0^2`$. We use an iterative procedure, similar to the one of chapter 5. At each step one checks whether there is any additional resolvable parton emission before the last one , with the probability
$$\mathrm{prob}(\mathrm{forward}\mathrm{emission})=\frac{\sigma _{T/L(\mathrm{direct})}^{\gamma ^{}i}(\stackrel{~}{s},Q^2,Q_1^2)\sigma _{T/L(\mathrm{direct}\mathrm{\hspace{0.17em}2}2)}^{\gamma ^{}i(22)}(\stackrel{~}{s},Q^2,Q_1^2)}{\sigma _{T/L(\mathrm{direct})}^{\gamma ^{}i}(\stackrel{~}{s},Q^2,Q_1^2)}$$
(10.84)
In case of an emission, the flavor $`j`$ of the new ladder leg parton, the light cone momentum fraction $`z`$, taken from the parent parton, and the virtuality $`\stackrel{~}{Q}^2`$ are generated according to the integrand of $`\sigma _{T/L(\mathrm{direct})}^{\gamma ^{}i}\sigma _{T/L(\mathrm{direct}\mathrm{\hspace{0.17em}2}2)}^{\gamma ^{}i}`$ in (10.80),
$$\mathrm{prob}(j,z,\stackrel{~}{Q}^2)\frac{1}{\stackrel{~}{Q}^2}\mathrm{\Delta }^i(Q_1^2,\stackrel{~}{Q}^2)\frac{\alpha _s}{2\pi }P_i^j(z)\sigma _{T/L(\mathrm{direct})}^{\gamma ^{}j}(z\stackrel{~}{s},Q^2,\stackrel{~}{Q}^2).$$
(10.85)
The process is repeated for the new ladder of energy squared $`\widehat{s}^{}=z\stackrel{~}{s}Q^2`$, with the initial parton $`i^{}=j`$, and with virtuality $`Q_1^2=\stackrel{~}{Q}^2`$ and so on. At each step one decides about emission or not, which finally terminates the iteration.
Having done the iterative parton emission, we finally generate the last (highest virtuality) resolvable parton emission or the photon-gluon fusion process (if $`i^{}=g`$) in the photon-parton center-of-mass system. The photon-gluon fusion (PGF) process is chosen for $`i^{}=g`$ with the probability
$$\mathrm{prob}(\mathrm{PGF})=\frac{1}{\sigma _{T(\mathrm{direct}\mathrm{\hspace{0.17em}2}2)}^{\gamma ^{}g}(\stackrel{~}{s}^{},Q^2,\stackrel{~}{Q}^2)}e_c^2𝑑p_{}^2\frac{d\sigma _T^{\gamma ^{}gc\overline{c}}(\stackrel{~}{s}^{},Q^2,p_{}^2)}{dp_{}^2}\mathrm{\Delta }^g(\stackrel{~}{Q}^2,M_F^2)$$
(10.86)
for the transverse photon polarization and always for the longitudinal photon polarization. For $`i^{}=q`$, we have $`\mathrm{prob}(\mathrm{PGF})=0`$. In case of PGF, with our choice $`M_F^2=m_c^2+p_{}^2`$, we generate the transverse momentum squared of final charm quarks in the region $`\stackrel{~}{Q}^2m_c^2<p_{}^2<\frac{1}{4}\widehat{s}m_c^2`$ according to
$$\mathrm{prob}(p_{}^2)\frac{d\sigma _T^{\gamma ^{}gc\overline{c}}(\stackrel{~}{s}^{},Q^2,p_{}^2)}{dp_{}^2}\mathrm{\Delta }^g(\stackrel{~}{Q}^2,M_F^2)$$
(10.87)
In case of no PGF, we generate the momentum transfer squared for the last resolvable parton emission in the range $`\stackrel{~}{Q}^2<Q^2<Q^2`$ according to
$$\mathrm{prob}(Q^2)\frac{1}{Q^2}\mathrm{\Delta }^i(\stackrel{~}{Q}^2,Q^2)\mathrm{\Delta }^q(Q^2,Q^2)\frac{\alpha _s}{2\pi }P_i^q\left(z^{}\right),$$
(10.88)
with
$$z^{}=1\left(1\frac{Q^2}{\stackrel{~}{s}}\right)\left(1\frac{Q^2}{\stackrel{~}{s}}\right),$$
(10.89)
and find the parton transverse momentum squared as $`p_{}^2=Q^2(1z^{})`$.
We then reconstruct final parton 4-momenta in their center of mass system with a random polar angle for the transverse momentum vector $`\stackrel{}{p}_{}`$ and boost them to the original Lorentz frame. This completed the description of the algorithm to generate parton configurations, based on exactly the same formulas as for calculatins of $`F_2`$ before.
The above discussion of how to generate parton configurations is not yet complete: the emitted partons are in general off–shell and can therefore radiate further partons. This so called time-like radiation is taken into account using standard techniques , as discussed already in chapter 5.
### 10.8 Hadron Production
For the hadronization, we use exactly the same philosophy and even the same procedure as in case of proton-proton ($`pp`$) scattering. Hadronization is not considered as a dynamical procedure, rather we consider the hadronic states as being integrated out in the considerations of cross section calculations of the preceding sections. Hadronization means simply a phenomenological procedure to explicitly reintroduce these hadronic states. The procedure employed for $`pp`$ scattering and to be used here as well is as follows:
1. drawing a cylinder diagram;
2. cutting the cylinder;
3. planar presentation of half-cylinder;
4. identification of cut line with kinky string;
5. kinky string hadronization (as explained in chapter 6).
We are going to explain the steps (1-4) for a concrete example of a diagram contributing to photon-proton scattering, where the photon interacts directly with a light quark (contribution “light”), and where the first parton of the ladder on the proton side couples to the proton via a soft Pomeron (contribution “sea”), as shown in fig. 10.11.
The external legs on the lower (proton) side are a quark (full line) and an anti-quark (dashed), representing together the proton constituent participating in the interaction. In fig. 10.12(left), we show the result of plotting the diagram on a cylinder. The shaded area on the lower part of the cylinder indicates the soft Pomeron, a complicated non-resolved structure, where we do not specify the microscopic content. The two space-like gluons emerge out of this soft structure. The cut is represented by the two vertical dotted lines on the cylinder. We now consider one of the two half-cylinders, say the left one, and we plot it in a planar fashion, as shown in fig. 10.12(right). We observe one internal gluon, and one external one, appearing on the cut line. We now identify the two cut lines with kinky strings such that a parton on the cut line corresponds to a kink: we have one kinky string with one internal kink (gluon) in addition to the two end kinks, and we have one flat string with just two end kinks, but no internal one. The strings are then hadronized according to the methods explained in chapter 6, see fig. 10.13.
In principle, we have also triple Pomerons contributing to the event topology. However, due to AGK cancellations, such contributions to the inclusive spectra cancel each other in the kinematical region where the two Pomerons are in parallel. Therefore, the average charteristics are correctly described by considering the simple cylinder-type topology corresponding to one Pomeron exchange.
### 10.9 Results
We are now capable to simulate events from deep inelastic scattering. When fixing the parameters, we found that all the ones found in $`e^+e^{}`$ can be kept with the exception of the mean transverse momentum of string breaking $`p_{}^f`$and the break probability $`p_{\mathrm{break}}`$, see the discussion in chapter 8.
We show results from $`ep`$ scattering and compare to the data of the experiments accomplished at HERA. Electrons of an energy 26.7 GeV collide with protons of 820 GeV, which gives a center of mass energy 296 GeV. We made the analysis for the kinematical region $`10^4<x<10^2`$ and $`10<Q^2<100\text{ GeV}^2`$. The distribution of the events calculated using our model is shown on figure 10.14.
We recall the principal variables
$$xx_\mathrm{B}=\frac{Q^2}{2(pq)}=\frac{Q^2}{ys};y=\frac{Q^2}{xs},$$
(10.90)
which gives straight lines for $`y=\text{const}.`$ in fig. 10.14. The sharp borders are due to imposing the experimental cuts $`0.03<y`$, $`E^{}>12\text{ GeV}`$ and $`\theta _e>7.5`$. On figure 10.15, we plot charged particle distributions for different values of $`W=\widehat{s}=2(pq)Q^2`$. The particle spectra look very similar for different $`W`$ values; the distributions decrease rapidly with $`x_F`$. The dependence of the average $`p_{}^2`$on $`x_F`$ shows an overall good agreement with the data from H1 collaboration .
In table 10.1 the bins in $`x`$ and $`Q^2`$ are given for the experimental data points on figures 10.15-10.20 (see ). The bin 0 is the sum of all the others. First, we compare the $`p_{}`$ distribution in the photon-proton center of mass system for the different bins - fig. 10.16. Fig. 10.16(a) shows the comparison of the results for low and high $`x`$ values for the values of $`Q^21020\text{ GeV}^2`$(bins 6 and 3). We find harder spectra for smaller $`x`$, which is the consequence of the larger kinematical space (in $`x`$) for the initial state radiation. Next, on figures 10.16(b,c), the spectra are compared for two different values of $`Q^2`$ and either the energy $`WQ^2/x`$ being fixed (bins 2 and 7, fig. 10.16(b)), or for a given value of $`x`$ (bins 6 and 8, figure 10.16(c)). The spectra in $`p_{}`$ are always harder for larger $`Q^2`$, which is now the consequence of the larger kinematical space in $`p_{}^2`$ for the initial state radiation. Two cuts in pseudo-rapidity $`\eta `$ for bin 3 are considered on figure 10.16(d). We see a harder distribution for $`1.5<\eta <2.5`$. Around mid-rapidity, where $`\eta `$ is maximal, one finds higher transverse momenta as this region is dominated by the contribution of the largest virtuality photon process.
Let us now consider pseudo-rapidity distributions of charged particles. Figures 10.17, 10.18 show the $`\eta `$-distributions for the 9 bins of table 10.1 for different values of $`Q^2`$ and $`x`$. On figure 10.18, a cut for $`p_{}>1\text{ GeV}`$ has been made to extract the contribution of hard processes. The latter one results in approximately 10% of the total hadron multiplicity. At fig. 10.18 we find fewer particles at small $`\eta `$ which is the consequence of smaller kinematical space (in $`x`$) for the initial state radiation and the reduced influence of the largest virtuality photon process.
The transverse momenta for all particles are generally well described by the model - figs. 10.19, 10.20. The fact that we find harder spectra for higher $`Q^2`$ and lower $`x`$ is best seen for smaller values of $`\eta `$ \- see fig. 10.20.
## Chapter 11 Results for Proton-Proton Scattering
In this section we are going to discuss our results for proton-proton interactions in the energy range between roughly 10 and 2000 GeV, which represents the range of validity of our approach. The lower limit is a fundamental limitation due to the fact that our approach requires hadron production to start after the primary interactions are finished, which is no longer fulfilled at low energies. The upper limit is due to the fact that above 2000 GeV higher order screening corrections need to be taken into account.
### 11.1 Energy dependence
We first consider the energy dependence of some elementary quantities in $`pp`$ scattering in the above mentioned energy range. In fig. 11.1, the results for the total cross section are shown.
The cross section is essentially used to fit the soft Pomeron parameters. Also shown in the figure is the energy dependence of the number of soft and semi-hard Pomerons. Over the whole energy range shown in the figure, soft physics is dominating. So for example at RHIC, soft physics dominates by far.
Our results for hadron production are based on the Pomeron parameters, defined from the cross sections fitting, and on the fragmentation procedure, adjusted on the basis of $`e^+e^{}`$data. In fig. 11.2, average multiplicities of different hadron species are given as a function of the energy. In fig. 11.3, we show the energy dependence of the pseudo-rapidity plateau $`dn_C/d\eta (0)`$ and of the mean squared transverse momentum $`<p_t^2>`$.
### 11.2 Charged Particle and Pion Spectra
In fig. 11.4, we present rapidity distributions of pions at 100 GeV, in fig. 11.5 rapidity distributions of pions and charged particles at 200 GeV. The values following the Symbol “I=” represent the integrals, i.e. the average multiplicity. The first number is the simulation, the second number (in brackets) represents data.
In fig. 11.6, we show rapidity distributions for positively and negatively charged particles at 53 GeV (cms), where we adopted also for the simulations the experimental definition of the rapidity by always taking the pion mass. In fig. 11.7, we show pseudo-rapidity distributions of charged particles at 200 and 1800 GeV (cms). In figs. 11.8 and 11.9, we finally show transverse momentum spectra at different energies between 100 GeV (lab) and 1800 GeV (cms).
### 11.3 Proton spectra
In fig. 11.10, we plot longitudinal momentum fraction distributions of protons for different values of $`t`$ at 200 GeV, in figs. 11.11 and 11.12 as well longitudinal momentum fraction distributions at 100-200 GeV, for given values of $`p_t`$ or integrated over $`p_t`$. In fig. 11.13, we show transverse momentum spectra of protons for different values of the longitudinal momentum fraction $`x`$ at 100 and 205 GeV.
### 11.4 Strange particle spectra
In figs. 11.14 and 11.15, we show transverse momentum and rapidity spectra of lambdas (including $`\mathrm{\Sigma }_0`$), anti-lambdas (including $`\overline{\mathrm{\Sigma }}_0`$), and kaons. The numbers represent the integrals, i.e. the average multiplicity. The first number is the simulation, the second number (in brackets) represents data.
## Chapter 12 Results for Collisions Involving Nuclei
It is well known that secondary interactions play an important role in collisions involving nuclei. Nevertheless, in this report, we do no want to consider any rescattering procedure, we just present bare neXus simulations. This seems to us the most honest way to present results.
### 12.1 Proton-Nucleus Scattering
In fig. 12.1, we show rapidity spectra of negatively charged particles for different target nuclei. Missing particles in the backward region are certainly due to rescattering. In the forward region, the model works well, except for $`p+`$Au, which represents the heaviest target, but in addition one has here a centrality trigger, in contrast to the other reactions. Here, we expect some reduction due to nuclear screening effects. The transverse momentum spectra are well reproduced in case of $`p+S`$, as shown in fig. 12.2, whereas for $`p+Au`$ one sees some deviations for small values of $`p_t`$, see fig. 12.3.
Let us turn to proton spectra. In fig. 12.4, we show rapidity spectra of net protons (protons minus anti-protons) for different target nuclei. Since secondary interactions are not considered, we are missing the pronounced peak around rapidity zero (not visible in the figure, since we have chosen the range for the y-axis to be ). Apart of the target fragmentation region, the model works well. The transverse momentum spectra shown in fig. 12.5 refer to the target fragmentation region, and therefore the absolute number is too small, whereas the shape of the spectra is quite good.
Strange particle spectra are shown in figs. 12.6 and 12.7. Whereas the simulations agree with the data for the kaons, the spectra are largely underestimated in case of lambdas, in particular in the backward region, where rescattering plays a dominant role.
### 12.2 Nucleus-Nucleus Scattering
Again, we show results of the bare neXus model, without any secondary interactions. Considering rapidity distributions of negatively charged particles, as shown in fig. 12.8, we observe a strong excess at central rapidities compared to the data. Rescattering will partly cure this, but not completely. For asymmetric systems like for example S+Ag, we observe in addition a missing asymmetry in the shape of the rapidity spectrum, in other words, the simulated spectrum is too symmetric. This is not surprising, since in our approach AGK cancelations apply, which make $`A+B`$ spectra identical to the $`p+p`$ ones, up to a factor. Rescattering will not cure this, since it acts essentially at central rapidities. But we expect another important effect effect due to additional screening effects coming from contributions of ehnanced Pomeron diagrams.
In fig. 12.9, we show transverse momentum or transverse mass spectra of negatively charged particles for different $`A+B`$ collisions at 200 GeV (lab). Several rapidity windows are shown; from top to bottom: 3.15-3.65, 3.65-4.15, 4.15-4.65, 4.65-5.15, 5.15-5.65 in case of Pb+Pb and 0.8-2, 2-3, 3-4, 4-4.4 for the other reactions. The lowest curves are properly normalized, the next ones are multiplied by ten, etc. We again observe an excess at certain rapidities, as already discussed above. In addition, in particular for heavy systems, the slopes are too steep, there is clearly some need of secondary interactions to “heat up” the system. This is consistent with the fact that the multiplicity is too high: collective motion should reduce the multiplicity but instead increase the transverse energy per particle.
In fig. 12.10, we show rapidity distributions of net protons (protons minus anti-protons) for different $`A+B`$ collisions at 200 GeV. For the asymmetric systems one observes clearly the effect of missing target nucleons, which should be cured by rescattering. The simulated results for symmetric systems are close to the data, rescattering does not contribute much here. But as for pion production, we expect also some changes due to screening effects. Transverse momentum spectra, as shown in fig. 12.11, show a similar behavior as for pions but even more pronounced: the theoretical spectra are much too steep, in particular for heavy systems.
In fig. 12.12-12.15, we show rapidity spectra of strange particles. $`\mathrm{K}^{}`$ seem to be correct, whereas $`\mathrm{K}^+`$ are in general somewhat too low compared to the data. Lambdas and anti-lambdas are way too low. For these particles, rescattering has to provide most of the particles which are finally observed.
## Chapter 13 Summary
We presented a new approach for hadronic interactions and for the initial stage of nuclear collisions, which solves several conceptual problems of certain classes of models, which are presently widely used in order to understand experimental data.
The main problem of these models is the fact that energy is not conserved in a consistent fashion: the fact that energy needs to be shared between many elementary interactions in case of multiple scattering is well taken into account when calculating particle production, but energy conservation is not taken care of in cross section calculations. Related to this problem is the fact that different elementary interactions in case of multiple scattering are usually not treated equally, so the first interaction is usually considered to be quite different compared to the subsequent ones.
We provided a rigorous treatment of the multiple scattering aspect, such that questions as energy conservation are clearly determined by the rules of field theory, both for cross section and particle production calculations. In both cases, energy is properly shared between the different interactions happening in parallel. This is the most important and new aspect of our approach, which we consider to be a first necessary step to construct a consistent model for high energy nuclear scattering.
Another important aspect of our approach is the hypothesis that particle production is a universal process for all the elementary interactions, from $`e^+e^{}`$ annihilation to nucleus-nucleus scattering. That is why we also carefully study $`e^+e^{}`$ annihilation and deep inelastic scattering. This allows to control reasonably well for example the hadronization procedure, which is not treatable theoretically from first principles.
This work has been funded in part by the IN2P3/CNRS (PICS 580) and the Russian Foundation of Basic Researches (RFBR-98-02-22024).
## Appendix A Kinematics of Two Body Collisions
### A.1 Conventions
We consider a scattering of a projectile $`P`$ on a target $`T`$ (hadron-hadron or parton-parton). We define the incident 4-momenta to be $`p`$ and $`p^{}`$ and the transferred momentum $`q`$, so that the outgoing momenta are $`\stackrel{~}{p}=p+q`$ and $`\stackrel{~}{p}^{}=p^{}q`$, see fig. A.1.
We define as usual the Mandelstam variables
$$s=(p+p^{})^2,t=(\stackrel{~}{p}p)^2.$$
(A.1)
Usually, we employ light cone momentum variables, connected to the energy and $`z`$-component of the particle momentum, as
$$p^\pm =p_0\pm p_z,$$
(A.2)
and we denote the particle 4-vector as
$$p=(p^+,p^{},\stackrel{}{p}_{}).$$
(A.3)
### A.2 Proof of the Impossibility of Longitudinal Excitations
Here, we present a mathematical proof of the well known fact that longitudinal excitations are impossible at high energies. Consider a collision between two hadrons $`h`$ and $`h^{}`$ which leads to two hadrons $`\stackrel{~}{h}`$ and $`\stackrel{~}{h}^{}`$,
$$h(p)+h^{}(p^{})\stackrel{~}{h}(\stackrel{~}{p})+\stackrel{~}{h}^{}(\stackrel{~}{p}^{})$$
(A.4)
with four-momenta $`p`$, $`p^{}`$, $`\stackrel{~}{p}`$, $`\stackrel{~}{p}^{}`$. As usual, we define $`s=(p+p^{})^2`$ and $`t=(\stackrel{~}{p}p)^2`$. For the following, we consider always the limit $`s\mathrm{}`$ and ignore terms of the order $`p^2/s`$. We expand $`q=\stackrel{~}{p}p=p^{}\stackrel{~}{p}^{}`$ as
$$q=\alpha p+\beta p^{}+q_{},$$
(A.5)
and obtain
$$\alpha =\frac{2qp^{}}{s},\beta =\frac{2qp}{s},$$
(A.6)
where we used
$$p^2=0,p^2=0,pq_{}=p^{}q_{}=0,s=2pp^{}.$$
(A.7)
We get
$$q=\frac{q^2}{s}(pp^{})+q_{},$$
(A.8)
having used
$$q=\stackrel{~}{p}p=p^{}\stackrel{~}{p}^{},p^2=p^2=0,p^{}\stackrel{~}{p}^{}=p\stackrel{~}{p}=q^2/2.$$
(A.9)
This proves
$$q=q_{}$$
(A.10)
for $`s\mathrm{}`$ and limited $`q^2`$. In other words, momentum transfer is purely transversal at high energies.
## Appendix B Partonic Interaction Amplitudes
### B.1 Semi-hard Parton-Parton Scattering
Let us derive the mathematical expression corresponding to the contribution of so-called semi-hard parton-parton scattering, see fig. B.1.
Here $`p,p^{}`$ are the 4-momenta of the constituent partons participating in the process. We denote by $`k,k^{}`$ the 4-momenta of the first partons entering the perturbative evolution, i.e. the initial partons for the perturbative parton cascade (characterized by parton virtualities $`Q^2>Q_0^2`$). Further, we define light cone momentum fractions
$$x^+=\frac{p^+}{p_0^+},x^{}=\frac{p^{}}{p_0^{}},x_1^+=\frac{k^+}{p_0^+},x_1^{}=\frac{k^{}}{p_0^{}},$$
(B.1)
with $`p_0^\pm `$ being the total light-cone momenta for the interaction.
At high energies, the dominant contribution to the process comes from the kinematical region where these partons are slow, i.e. $`x_1^\pm x^\pm `$, so that a relatively small contribution of the perturbative parton cascade (of the ladder part of the diagram of the fig. B.1) is compensated by the large density of such partons, resulted from the soft pre-evolution . Since the initial partons $`k,k^{}`$ are gluons or sea quarks (contrary to valence quarks) we talk about “sea-sea” contribution.
Let us first consider the case where the intermediate partons $`k,k^{}`$ are gluons. Then the amplitude for the diagram of fig. B.1 can be written as
$`iT_{\mathrm{sea}\mathrm{sea}}^{gg}(\widehat{s},t)`$ $`=`$ $`{\displaystyle \frac{d^4k}{(2\pi )^4}\frac{d^4k^{}}{(2\pi )^4}\underset{\lambda \lambda ^{}\gamma \gamma ^{}\delta \delta ^{}\tau \tau ^{}}{}iM_{\mathrm{soft}}^g(p,k,p+q,kq)_{\lambda \gamma }}`$
$`\times `$ $`D_{\lambda \delta }^g\left(k^2\right)D_{\gamma \tau }^g\left((k+q)^2\right)iM_{\mathrm{hard}}^{gg}(k,k^{},k+q,k^{}q,Q_0^2)_{\delta \tau \delta ^{}\tau ^{}}`$
$`\times `$ $`D_{\lambda ^{}\delta ^{}}^g\left(k^2\right)D_{\gamma ^{}\tau ^{}}^g\left((k^{}q)^2\right)iM_{\mathrm{soft}}^g(p^{},k^{},p^{}q,k^{}+q)_{\lambda ^{}\gamma ^{}},`$
where the amplitude $`M_{\mathrm{hard}}^{gg}(k,k^{},k+q,k^{}q,Q_0^2)_{\delta \tau \delta ^{}\tau ^{}}`$ represents the perturbative contribution of the parton ladder with the initial partons of momenta $`k,k^{}`$ and with the momentum transfer along the ladder $`q`$ (hard parton-parton scattering), the amplitude $`M_{\mathrm{soft}}^g(p,k,p+q,kq)_{\lambda \gamma }`$ corresponds to the non-perturbative soft interaction between the constituent parton with the 4-momentum $`p`$ and the gluon with the 4-momentum $`k,`$ and $`D_{\lambda \delta }^g(k)`$ is the non-perturbative gluon propagator, $`D_{\lambda \delta }^g(k)=i\stackrel{~}{D}^g\left(k^2\right)\epsilon _{\lambda \delta }(k)`$ with $`\epsilon _{\lambda \delta }(k)`$ being the usual gluon polarization tensor in the axial gauge; $`\lambda ,\gamma ,\delta ,\mathrm{}`$ denotes symbolically the combination of color and Lorentz indexes for the intermediate gluons.
As by construction partons of large virtualities $`Q^2>Q_0^2`$ can only appear in the parton ladder part of the diagram of the fig. B.1, we assume that the integral $`d^4k`$ converges in the region of restricted virtualities $`k^2s_0`$ with $`s_01`$ GeV<sup>2</sup> being the typical hadronic mass scale, i.e. the combination
$$M_{\mathrm{soft}}^g(p,k,p+q,kq)_{\lambda \gamma }D_{\lambda \delta }^g\left(k^2\right)D_{\gamma \tau }^g\left((k+q)^2\right)$$
(B.3)
drops down fast with increasing $`\left|k^2\right|`$; this implies that the transverse momentum $`k_{}`$ is also restricted to the region $`k_{}s_0`$. Similar arguments apply for $`k^{}`$. Further we make the usual assumption that longitudinal polarizations in the gluon propagators $`D_{\lambda \delta }^g\left(k^2\right)`$ are canceled in the convolution with the soft contribution $`M_{\mathrm{soft}}^g`$ even for finite gluon virtualities $`k^2`$ . Finally we assume that in the considered limit $`x_1^\pm x^\pm `$ the amplitude $`M_{\mathrm{soft}}^g(p,k,p+q,kq)_{\lambda \gamma }`$ is governed by the non-perturbative soft Pomeron exchange between the constituent parton $`p`$ and the gluon $`k`$, which implies in particular that it has the singlet structure in the color and Lorentz indexes:
$$M_{\mathrm{soft}}^g(p,k,p+q,kq)_{\lambda \gamma }\delta _\gamma ^\lambda .$$
(B.4)
Then, for small momentum transfer $`q`$ in the process of fig. B.1 the intermediate gluons of momenta $`k,k^{},k+q,k^{}q`$ can be considered as real (on-shell) ones with respect to the perturbative parton evolution in the ladder, characterized by large momentum transfers $`Q^2>Q_0^2`$. Then we obtain
$`{\displaystyle \frac{1}{K_g^2}}{\displaystyle \underset{\lambda \lambda ^{}\delta \delta ^{}\tau \tau ^{}}{}}M_{\mathrm{hard}}^{gg}(k,k^{},k+q,k^{}q,Q_0^2)_{\delta \tau \delta ^{}\tau ^{}}`$ (B.5)
$`\times \epsilon _{\lambda \delta }(k)\epsilon _{\delta \tau }(k+q)\epsilon _{\lambda ^{}\delta ^{}}(k^{})\epsilon _{\delta ^{}\tau ^{}}(k^{}q)`$ $``$ $`T_{\mathrm{hard}}^{gg}(\widehat{s}_{\mathrm{hard}},q^2,Q_0^2),`$
where the averaging over the spins and the colors of the initial gluons $`k,k^{}`$ is incorporated in the factor $`K_g^2`$, $`T_{\mathrm{hard}}^{gg}(\widehat{s}_{\mathrm{hard}},t,Q_0^2)`$ is defined in (2.16), and $`\widehat{s}_{\mathrm{hard}}=(k+k^{})^2k^+k^{}=x_1^+x_1^{}s`$.
Now, using (B.4-B.5) we can rewrite (B.1) as
$`iT_{\mathrm{sea}\mathrm{sea}}^{gg}(\widehat{s},t)`$ $`=`$ $`{\displaystyle \frac{dk^+dk^{}d^2k_{}}{2(2\pi )^4}\frac{dk^+dk^{}d^2k_{}^{}}{2(2\pi )^4}iT_{\mathrm{hard}}^{gg}(\widehat{s}_{\mathrm{hard}},q^2,Q_0^2)}`$
$`\times `$ $`\left[i{\displaystyle \underset{\lambda }{}}M_{\mathrm{soft}}^g(p,k,p+q,kq)_{\lambda \lambda }\stackrel{~}{D}^g\left(k^2\right)\stackrel{~}{D}^g\left((k+q)^2\right)\right]`$
$`\times `$ $`\left[i{\displaystyle \underset{\lambda ^{}}{}}M_{\mathrm{soft}}^g(p^{},k^{},p^{}q,k^{}+q)_{\lambda ^{}\lambda ^{}}\stackrel{~}{D}^g\left(k^2\right)\stackrel{~}{D}^g\left((k^{}q)^2\right)\right].`$
It is convenient to perform separately the integrations over $`k^{},k^+,k_{},k_{}^{}`$ keeping in mind that the only dependence on those variables in the eq. (B.1) appears in the non-perturbative contributions in the square brackets. Let us consider the first of those contributions, corresponding to the upper soft blob at fig. B.1 together with the intermediate gluon propagators (for the lower blob the derivation is identical). Being described by the soft Pomeron exchange, the amplitude $`M_{\mathrm{soft}}^g(p,k,p+q,kq)_{\lambda \lambda }`$ is an analytical function of the energy invariants $`\widehat{s}_{\mathrm{soft}}=(pk)^2`$ $`p^+k^{}`$ and $`\widehat{u}_{\mathrm{soft}}=(p+k+q)^2`$ $`p^+k^{}`$with the singularities in the complex $`k^{}`$-plane, corresponding to the values $`\widehat{s}_{\mathrm{soft}}=s^{}i0`$ for all real $`s^{}`$ which are greater than some threshold value $`s_{\mathrm{thr}}`$ for the Pomeron asymptotics to be applied, as well as for $`\widehat{u}_{\mathrm{soft}}=u^{}i0`$, where $`u^{}>u_{\mathrm{thr}}`$ with some threshold value $`u_{\mathrm{thr}}`$ . Thus one has the singularities in the variable $`k^{}`$in the upper half of the complex plane at $`k^{}(s^{}i0)/p^+`$ and in the lower half of the complex plane at $`k^{}(u^{}i0)/p^+`$. Then one can use the standard trick to rotate the integration contour $`C`$ in the variable $`k^{}`$such that the new contour $`C^{}`$ encloses the left-hand singularities in $`k^{}`$, corresponding to the right-hand singularities in the variable $`\widehat{s}_{\mathrm{soft}}`$ . Then one ends up with the integral over the discontinuity of the amplitude $`M_{\mathrm{soft}}^g`$ on the left-hand cut in the variable $`k^{}`$, which is up to a minus sign equal to the discontinuity on the right-hand cut in $`\widehat{s}_{\mathrm{soft}}`$:
$`{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑k^{}\left[{\displaystyle \underset{\lambda }{}}M_{\mathrm{soft}}^g(p,k,p+q,kq)_{\lambda \lambda }\stackrel{~}{D}^g\left(k^2\right)\stackrel{~}{D}^g\left((k+q)^2\right)\right]`$
$`={\displaystyle _{\mathrm{}}^{s_{\mathrm{thr}}/k^+}}𝑑k^{}\mathrm{disc}_{\widehat{s}_{\mathrm{soft}}}\left[{\displaystyle \underset{\lambda }{}}M_{\mathrm{soft}}^g(p,k,p+q,kq)_{\lambda \lambda }\stackrel{~}{D}^g\left(k^2\right)\stackrel{~}{D}^g\left((k+q)^2\right)\right]`$ (B.7)
$`={\displaystyle _{\mathrm{}}^{s_{\mathrm{thr}}/k^+}}𝑑k^{}\mathrm{\hspace{0.17em}2}i\mathrm{Im}\left[{\displaystyle \underset{\lambda }{}}M_{\mathrm{soft}}^g(p,k,p+q,kq)_{\lambda \lambda }\stackrel{~}{D}^g\left(k^2\right)\stackrel{~}{D}^g\left((k+q)^2\right)\right].`$ (B.8)
Now, using $`𝑑k^{}=𝑑k^2/k^+`$ and recalling our assumption that the integral over $`k^2`$ gets dominant contribution from the region $`k^2s_0`$, we may write
$`{\displaystyle \frac{dk^{}d^2k_{}}{(2\pi )^4}\mathrm{Im}\left[\underset{\lambda }{}M_{\mathrm{soft}}^g(p,k,p+q,kq)_{\lambda \lambda }\stackrel{~}{D}^g\left(k^2\right)\stackrel{~}{D}^g\left((k+q)^2\right)\right]}`$
$`={\displaystyle \frac{1}{k^+}}\mathrm{Im}T_{\mathrm{soft}}^g(\widehat{s}_{\mathrm{soft}},q^2)`$ (B.9)
with
$$\widehat{s}_{\mathrm{soft}}=s_0\frac{p^+}{k^+}=s_0\frac{x^+}{x_1^+}.$$
(B.10)
The integrations over $`k^2`$ in the vicinity of $`k^2=s_0`$ and over $`k_{}s_0`$ are supposed to just contribute to the redetermination of the Pomeron-gluon coupling of the usual soft Pomeron amplitude, and therefore we parameterize the amplitude $`T_{\mathrm{soft}}^g`$ as (compare with eq. (2.5))
$$T_{\mathrm{soft}}^g(\widehat{s},t)=8\pi s_0\eta (t)\gamma _{\mathrm{part}}\gamma _g\left(\frac{\widehat{s}}{s_0}\right)^{\alpha __\mathrm{P}(0)}\mathrm{exp}\left(\lambda _{\mathrm{soft}}^{(1)}(\widehat{s}/s_0)t\right)\left(1\frac{s_0}{\widehat{s}}\right)^{\beta _g},$$
(B.11)
with
$$\lambda _{\mathrm{soft}}^{(1)}(z)=R_{\mathrm{part}}^2+\alpha _{\mathrm{soft}}^{}\mathrm{ln}z,$$
(B.12)
where we used $`\gamma _g`$ for the Pomeron-gluon coupling and we neglected the radius of the Pomeron-gluon vertex assuming that the coupling is local in the soft Pomeron; the factor $`\left(1s_0/\widehat{s}\right)^{\beta _g}`$ is included to ensure that the Pomeron has sufficiently large mass, which is the necessary condition for applying Regge description for the soft evolution. As we shall see below the parameter $`\beta _g`$ determines the gluon momentum distribution in the Pomeron at $`x_1^\pm x^\pm `$.
Finally, using the above results we obtain
$`iT_{\mathrm{sea}\mathrm{sea}}^{gg}(\widehat{s},t)`$ $`=`$ $`{\displaystyle \frac{dk^+}{k^+}\frac{dk^{}}{k^{}}\mathrm{Im}T_{\mathrm{soft}}^g(\widehat{s}_{\mathrm{soft}},t)\mathrm{Im}T_{\mathrm{soft}}^g(\widehat{s}_{\mathrm{soft}}^{},t)iT_{\mathrm{hard}}^{gg}(\widehat{s}_{\mathrm{hard}},t,Q_0^2)}`$ (B.13)
$`=`$ $`{\displaystyle _0^1}{\displaystyle \frac{dz^+}{z^+}}{\displaystyle \frac{dz^{}}{z^{}}}\mathrm{Im}T_{\mathrm{soft}}^g({\displaystyle \frac{s_0}{z^+}},t)\mathrm{Im}T_{\mathrm{soft}}^g({\displaystyle \frac{s_0}{z^{}}},t)iT_{\mathrm{hard}}^{gg}(z^+z^{}\widehat{s},t,Q_0^2),`$ (B.14)
where the following definitions have been used:
$$z^\pm =\frac{x_1^\pm }{x^\pm },\widehat{s}_{\mathrm{soft}}=s_0\frac{x^+}{x_1^+}=\frac{s_0}{z^+},\widehat{s}_{\mathrm{soft}}^{}=s_0\frac{x^{}}{x_1^{}}=\frac{s_0}{z^{}},\widehat{s}_{\mathrm{hard}}=x_1^+x_1^{}s=z^+z^{}\widehat{s}.$$
(B.15)
In the case of the intermediate parton $`k`$ being a (anti-)quark, we assume that it originates from local gluon splitting in the soft Pomeron. Thus we neglect the slope of the (non-perturbative) gluon-quark vertex. Using the usual Altarelli-Parisi kernel $`P_g^q(z)=\frac{1}{2}(z^2+(1z)^2)`$ for the gluon light cone momentum partition between the quark and the anti-quark, we get the corresponding amplitude $`T_{\mathrm{sea}\mathrm{sea}}^{qg}(\widehat{s},t)`$ as
$$iT_{\mathrm{sea}\mathrm{sea}}^{qg}(\widehat{s},t)=_0^1\frac{dz^+}{z^+}\frac{dz^{}}{z^{}}\mathrm{Im}T_{\mathrm{soft}}^q(\frac{s_0}{z^+},t)\mathrm{Im}T_{\mathrm{soft}}^g(\frac{s_0}{z^{}},t)iT_{\mathrm{hard}}^{qg}(z^+z^{}\widehat{s},t,Q_0^2),$$
(B.16)
where the imaginary part of the amplitude $`\mathrm{Im}T_{\mathrm{soft}}^q`$ for the soft Pomeron exchange between the constituent parton and the quark $`q`$ $``$ $`\{u,d,s,\overline{u},\overline{d},\overline{s}\}`$ is expressed via $`\mathrm{Im}T_{\mathrm{soft}}^g`$ as
$$\mathrm{Im}T_{\mathrm{soft}}^q(\widehat{s}_{\mathrm{soft}},t)=\gamma _{qg}𝑑\xi P_g^q(\xi )\mathrm{Im}T_{\mathrm{soft}}^g(\xi \widehat{s}_{\mathrm{soft}},t),$$
(B.17)
with $`\gamma _{qg}`$ representing the quark-gluon vertex value and $`\xi `$ being the ratio of the quark and the parent gluon light cone momentum, $`\xi =k^+/k_g^+`$; the mass squared of the Pomeron between the initial constituent parton and the gluon is then
$$(pk_g)^2s_0\frac{p^+}{k_g^+}=\xi \widehat{s}_{\mathrm{soft}}.$$
(B.18)
The full amplitude for the semi-hard scattering is the sum of the different terms discussed above, i.e.
$`iT_{\mathrm{sea}\mathrm{sea}}(\widehat{s},t)`$ $`=`$ $`{\displaystyle \underset{jk}{}}iT_{\mathrm{sea}\mathrm{sea}}^{jk}(\widehat{s},t)`$ (B.19)
$`=`$ $`{\displaystyle \underset{jk}{}}{\displaystyle _0^1}{\displaystyle \frac{dz^+}{z^+}}{\displaystyle \frac{dz^{}}{z^{}}}\mathrm{Im}T_{\mathrm{soft}}^j({\displaystyle \frac{s_0}{z^+}},t)\mathrm{Im}T_{\mathrm{soft}}^k({\displaystyle \frac{s_0}{z^{}}},t)iT_{\mathrm{hard}}^{jk}(z^+z^{}\widehat{s},t,Q_0^2),`$
where $`j,k`$ denote the types (flavors) of the initial partons for the perturbative evolution (quarks or gluons).
The discontinuity of the amplitude $`T_{\mathrm{sea}\mathrm{sea}}(\widehat{s},t)`$ on the right-hand cut in the variable $`\widehat{s}`$ defines the contribution of the cut diagram of the fig. B.1. Cutting procedure amounts here to replace the hard parton-parton scattering amplitude $`iT_{\mathrm{hard}}^{jk}(\widehat{s}_{\mathrm{hard}},t,Q_0^2)`$ in (B.19) by $`2\mathrm{I}\mathrm{m}T_{\mathrm{hard}}^{jk}(\widehat{s}_{\mathrm{hard}},t,Q_0^2)`$, whereas the contributions of the soft parton cascades $`\mathrm{Im}T_{\mathrm{soft}}^j`$ stay unchanged as they are already defined by the discontinuities in the corresponding variables $`\widehat{s}_{\mathrm{soft}}`$ and $`\widehat{s}_{\mathrm{soft}}^{}`$. So the cut diagram contribution is just $`2\mathrm{I}\mathrm{m}T_{\mathrm{sea}\mathrm{sea}}(\widehat{s},t)`$. At $`t=0`$ it defines the cross section for the semi-hard parton-parton scattering:
$$\sigma _{\mathrm{sea}\mathrm{sea}}\left(\widehat{s}\right)=\frac{1}{2\widehat{s}}2\mathrm{I}\mathrm{m}T_{\mathrm{sea}\mathrm{sea}}(\widehat{s},0)=\underset{jk}{}_0^1𝑑z^+𝑑z^{}E_{\mathrm{soft}}^j\left(z^+\right)E_{\mathrm{soft}}^k\left(z^{}\right)\sigma _{\mathrm{hard}}^{jk}(z^+z^{}\widehat{s},Q_0^2),$$
(B.20)
where we used the explicit expressions (2.16), (B.11), (B.17) for $`T_{\mathrm{hard}}^{jk}`$, $`T_{\mathrm{soft}}^j`$, and the functions $`E_{\mathrm{soft}}^j`$ are defined as
$`E_{\mathrm{soft}}^g\left(z\right)`$ $`=`$ $`8\pi s_0\gamma _{\mathrm{part}}\gamma _gz^{\alpha _{\mathrm{soft}}(0)}(1z)^{\beta _g}`$ (B.21)
$`E_{\mathrm{soft}}^q\left(z\right)`$ $`=`$ $`\gamma _{qg}{\displaystyle _z^1}𝑑\xi P_g^q(\xi )E_{\mathrm{soft}}^g\left({\displaystyle \frac{z}{\xi }}\right).`$ (B.22)
It is easy to see that $`E_{\mathrm{soft}}^j\left(z\right)`$ has the meaning of the momentum distribution of parton $`j`$ at the scale $`Q_0^2`$ for an elementary interaction, i.e. parton distribution in the soft Pomeron. Introducing the gluon splitting probability $`w_{\mathrm{split}}`$ and the coupling $`\stackrel{~}{\gamma }_g`$ via
$$\gamma _{qg}\gamma _g=w_{\mathrm{split}}\stackrel{~}{\gamma }_g,\gamma _g=\left(1w_{\mathrm{split}}\right)\stackrel{~}{\gamma }_g,$$
(B.23)
the light cone momentum conservation reads
$$1=_0^1𝑑z\underset{j}{}zE_{\mathrm{soft}}^j\left(z\right)=8\pi s_0\gamma _{\mathrm{part}}\stackrel{~}{\gamma }_g_0^1𝑑zz^{1\alpha _{\mathrm{soft}}(0)}(1z)^{\beta _g}$$
(B.24)
and therefore
$$\stackrel{~}{\gamma }_g=\frac{1}{8\pi s_0\gamma _{\mathrm{part}}}\frac{\mathrm{\Gamma }\left(3\alpha _{\mathrm{soft}}(0)+\beta _g\right)}{\mathrm{\Gamma }\left(2\alpha _{\mathrm{soft}}(0)\right)\mathrm{\Gamma }\left(1+\beta _g\right)}$$
(B.25)
### B.2 Parton Evolution
In this appendix, we discuss the properties of the evolution function $`E_{\mathrm{QCD}}`$, describing the perturbative evolution of partons.
The evolution function $`E_{\mathrm{QCD}}^{jm}(z,Q_0^2,Q^2)`$ satisfies the usual DGLAP equation
$$\frac{dE_{\mathrm{QCD}}^{jm}(Q_0^2,Q^2,x)}{d\mathrm{ln}Q^2}=\underset{k}{}_x^1\frac{dz}{z}\frac{\alpha _s}{2\pi }\stackrel{~}{P}_k^m(z)E_{\mathrm{QCD}}^{jk}(\frac{x}{z},Q_0^2,Q^2)$$
(B.26)
with the initial condition $`E_{\mathrm{QCD}}^{jm}(Q_0^2,Q_0^2,z)=\delta _j^m\delta (1z).`$ Here $`\stackrel{~}{P}_k^m(z)`$ are the usual Altarelli-Parisi splitting functions, regularized at $`z1`$ by the contribution of virtual emissions.
One can introduce the concept of “resolvable” parton emission, i.e. an emission of a final ($`s`$-channel) parton with a finite share of the parent parton light cone momentum $`(1z)>ϵ=p_{\text{res }}^{2\text{ }}/Q^2`$ (with finite relative transverse momentum $`p_{}^2=`$ $`Q^2(1z)`$ $`>p_{\mathrm{res}\text{ }}^{2\text{ }}`$) and use the so-called Sudakov form factor, corresponding to the contribution of any number of virtual and unresolvable emissions (i.e. emissions with $`(1z)<ϵ`$) \- see fig. B.2.
$$\mathrm{\Delta }^k(Q_0^2,Q^2)=\mathrm{exp}\left\{_{Q_0^2}^{Q^2}\frac{dq^2}{q^2}_{1ϵ}^1𝑑z\frac{\alpha _s}{2\pi }\stackrel{~}{P}_k^k(z)\right\}$$
(B.27)
This can also be interpreted as the probability of no resolvable emission between $`Q_0^2`$ and $`Q^2`$.
Then $`E_{\mathrm{QCD}}^{jm}`$ can be expressed via $`\overline{E}_{\mathrm{QCD}}^{jm}`$, corresponding to the sum of any number (but at least one) resolvable emissions, allowed by the kinematics:
$$E_{\mathrm{QCD}}^{jm}(z,Q_0^2,Q^2)=\delta _j^m\delta (1z)\mathrm{\Delta }^j(Q_0^2,Q^2)+\overline{E}_{\mathrm{QCD}}^{jm}(z,Q_0^2,Q^2),$$
(B.28)
where $`\overline{E}_{\mathrm{QCD}}^{jm}(z,Q_0^2,Q^2)`$ satisfies the integral equation
$`\overline{E}_{\mathrm{QCD}}^{jm}(x,Q_0^2,Q^2)`$ $`=`$ $`{\displaystyle _{Q_0^2}^{Q^2}}{\displaystyle \frac{dQ_1^2}{Q_1^2}}[{\displaystyle \underset{k}{}}{\displaystyle _x^{1ϵ}}{\displaystyle \frac{dz}{z}}{\displaystyle \frac{\alpha _s}{2\pi }}P_k^m(z)\overline{E}_{\mathrm{QCD}}^{jk}({\displaystyle \frac{x}{z}},Q_0^2,Q_1^2)+`$ (B.29)
$`+`$ $`\mathrm{\Delta }^j(Q_0^2,Q_1^2){\displaystyle \frac{\alpha _s}{2\pi }}P_j^m(x)]\mathrm{\Delta }^m(Q_1^2,Q^2)`$
Here $`P_j^k(z)`$ are the Altarelli-Parisi splitting functions for real emissions, i.e. without $`\delta `$-function and regularization terms at $`z1`$.
Eq. (B.29) can be solved iteratively, expressing $`\overline{E}_{\mathrm{QCD}}^{jm}`$ as the contribution of at most $`n`$ ($`n\mathrm{}`$) resolvable emissions (of an ordered ladder with at most $`n`$ ladder rungs) - see fig. B.3:
$$\overline{E}_{\mathrm{QCD}}^{jm}(Q_0^2,Q^2,x)=\underset{n\mathrm{}}{lim}E_{\mathrm{QCD}}^{(n)jm}(Q_0^2,Q^2,x),$$
(B.30)
with
$`E_{\mathrm{QCD}}^{(n)jm}(x,Q_0^2,Q^2)`$ $`=`$ $`{\displaystyle _{Q_0^2}^{Q^2}}{\displaystyle \frac{dQ_1^2}{Q_1^2}}\left[{\displaystyle \underset{k}{}}{\displaystyle _x^{1ϵ}}{\displaystyle \frac{dz}{z}}{\displaystyle \frac{\alpha _s}{2\pi }}P_k^m(z)E_{\mathrm{QCD}}^{(n1)jk}({\displaystyle \frac{x}{z}},Q_0^2,Q_1^2)\right]`$ (B.31)
$`\times `$ $`\mathrm{\Delta }^m(Q_1^2,Q^2)+E_{\mathrm{QCD}}^{(1)jm}(x,Q_0^2,Q^2)`$
$`E_{\mathrm{QCD}}^{(1)jm}(x,Q_0^2,Q^2)`$ $`=`$ $`{\displaystyle _{Q_0^2}^{Q^2}}{\displaystyle \frac{dQ_1^2}{Q_1^2}}\mathrm{\Delta }^j(Q_0^2,Q_1^2)\mathrm{\Delta }^m(Q_1^2,Q^2){\displaystyle \frac{\alpha _s}{2\pi }}P_j^m(z)`$ (B.32)
So the procedure amounts to only considering resolvable emissions, but to multiply each propagator with $`\mathrm{\Delta }^j`$, which is the reason for the appearance of $`\mathrm{\Delta }^j`$ in eqs. (B.29), (B.31-B.32).
### B.3 Time-Like Parton Splitting
We discuss here the algorithm for Monte Carlo generation of time-like parton emission on the basis of the eq. (6.25).
The standard procedure is to apply the Monte Carlo rejection method . We consider the splitting of a parton $`j`$ with a maximal virtuality $`Q_{j\mathrm{max}}^2`$ given by the preceding process. For the proposal function, we define the limits in $`z`$ for given $`Q_j^2`$ using an approximate formula
$$p_{}^2z(1z)Q_j^2zQ_l^2(1z)Q_k^2$$
(B.33)
instead of the exact one, eq. (6.28), with the lowest possible virtualities for daughter partons $`Q_k^2=Q_l^2=p_{\mathrm{fin}}^2`$, which gives
$$z_{\mathrm{min}/\mathrm{max}}(Q_j^2)=\frac{1}{2}\pm \frac{1}{2}\sqrt{1\frac{4p_{\mathrm{fin}}^2}{Q_j^2}}.$$
(B.34)
Further, we define upper limits,
$`\overline{P}_g^g(z)`$ $`=`$ $`3\left\{{\displaystyle \frac{1}{z}}+{\displaystyle \frac{1}{1z}}\right\},`$
$`\overline{P}_g^q(z)`$ $`=`$ $`{\displaystyle \frac{N_f}{2}},`$ (B.35)
$`\overline{P}_q^g(z)`$ $`=`$ $`{\displaystyle \frac{4}{3}}{\displaystyle \frac{2}{1z}},`$
for the splitting functions,
$`{\displaystyle \frac{1}{2}}P_g^g(z)=3{\displaystyle \frac{(1z(1z))^2}{z(1z)}},`$
$`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i\{u,d,s,\overline{u},\overline{d},\overline{s}\}}{}}P_g^i(z)={\displaystyle \frac{N_f}{2}}(z^2+(1z)^2),`$ (B.36)
$`P_q^g(z)={\displaystyle \frac{4}{3}}{\displaystyle \frac{1+z^2}{1z}},`$
with $`N_f`$ being the number of active quark flavors. Integrating these three functions $`\overline{P}_j^k(z)`$ over $`z`$ from $`z_{\mathrm{min}}=z_{\mathrm{min}}(Q_j^2)`$ to $`z_{\mathrm{max}}=z_{\mathrm{max}}(Q_j^2)`$ as
$$I_j^k(Q_j^2)=_{z_{\mathrm{min}}}^{z_{\mathrm{max}}}𝑑z\overline{P}_j^k(z),$$
(B.37)
one obtains
$`I_g^g(Q_j^2)`$ $`=`$ $`3\left(\mathrm{ln}\left({\displaystyle \frac{z_{\mathrm{max}}}{z_{\mathrm{min}}}}\right)+\mathrm{ln}\left({\displaystyle \frac{1z_{\mathrm{min}}}{1z_{\mathrm{max}}}}\right)\right)`$ (B.38)
$`I_g^q(Q_j^2)`$ $`=`$ $`{\displaystyle \frac{n_f}{2}}\left(z_{\mathrm{max}}z_{\mathrm{min}}\right)`$ (B.39)
$`I_q^q(Q_j^2)`$ $`=`$ $`{\displaystyle \frac{8}{3}}\mathrm{ln}\left({\displaystyle \frac{1z_{\mathrm{min}}}{1z_{\mathrm{max}}}}\right).`$ (B.40)
Defining $`I_j(Q_j^2)`$ as
$$I_j(Q_j^2)=\left\{I_g^g(Q_j^2)+I_g^q(Q_j^2)\right\}\delta _j^g+I_q^q(Q_j^2)\delta _j^q,$$
(B.41)
we propose the value $`Q_j^2`$ based upon the probability distribution
$$f_j(Q_j^2)=\frac{\alpha _{\mathrm{max}}}{2\pi }I_j(Q_j^2)\frac{1}{Q_j^2},$$
(B.42)
with $`\alpha _{\mathrm{max}}=\alpha _s(p_{\mathrm{min}}^2)=\alpha _s(p_{\mathrm{fin}}^2)`$. The flavor $`k`$ of the daughter parton is then chosen according to partial contributions $`I_j^k(Q_j^2)`$ in (B.41), and the value of $`z`$ according to the functions $`\overline{P}_j^k(z)`$.
The proposed values for $`Q_j^2=Q^2`$, $`k`$, and $`z`$ are accepted according to the probability
$$\frac{\alpha _s(p_{}^2)}{\alpha _{\mathrm{max}}}w_j^k,$$
(B.43)
with $`p_{}^2=z(1z)Q_j^2`$ and
$`w_g^g=`$ $`(1z(1z))^2,`$ (B.44)
$`w_g^q=`$ $`z^2+(1z)^2,`$ (B.45)
$`w_q^g=`$ $`(1+z^2)/2.`$ (B.46)
Otherwise, the proposal is rejected and one looks for another splitting with $`Q_{j\mathrm{max}}^2=Q_j^2`$.
## Appendix C Hadron-Hadron Amplitudes
In this appendix, we discuss the hadron-hadron scattering amplitude $`T_{h_1h_2}`$, where $`h_1`$ and $`h_2`$ represent any pair of hadrons.
### C.1 Neglecting Valence Quark Scatterings
We start with the general expression for hadron-hadron scattering amplitude, eq. (2.31),
$`iT_{h_1h_2}(s,t)={\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n!}}{\displaystyle \underset{l=1}{\overset{n}{}}\left[\frac{d^4k_l}{(2\pi )^4}\frac{d^4k_l^{}}{(2\pi )^4}\frac{d^4q_l}{(2\pi )^4}\right]N_{h_1}^{(n)}(p,k_1,\mathrm{},k_n,q_1,\mathrm{},q_n)}`$
$`\times {\displaystyle \underset{l=1}{\overset{n}{}}}\left[iT_{1\mathrm{I}\mathrm{P}}(\widehat{s}_l,q_l^2)\right]N_{h_2}^{(n)}(p^{},k_1^{},\mathrm{},k_n^{},q_1,\mathrm{},q_n)(2\pi )^4\delta ^{(4)}({\displaystyle \underset{k=1}{\overset{n}{}}}q_iq),`$ (C.1)
with $`t=q^2`$, $`s=(p+p^{})^2p^+p^{}`$, and with $`p,p^{}`$ being the 4-momenta of the initial hadrons. We consider for simplicity identical parton constituents (neglecting valence quark scatterings) and take $`T_{1\mathrm{I}\mathrm{P}}`$ to be the sum of the soft Pomeron exchange amplitude (see eq. (2.5)) and the semi-hard sea-sea scattering amplitude (see eq. (LABEL:t-sea-sea)):
$$T_{1\mathrm{I}\mathrm{P}}(\widehat{s}_l,q_l^2)=T_{\mathrm{soft}}(\widehat{s}_l,q_l^2)+T_{\mathrm{sea}\mathrm{sea}}(\widehat{s}_l,q_l^2),$$
(C.2)
with $`\widehat{s}_l=(k_l+k_l^{})^2k_l^+k_l^{}`$. The momenta $`k_l,k_l^{}`$ and $`q_l`$ denote correspondingly the 4-momenta of the initial partons for the $`l`$-th scattering and the 4-momentum transfer in that partial process. The factor $`1/n!`$ takes into account the identical nature of the $`n`$ scattering contributions. $`N_h^{(n)}(p,k_1,\mathrm{},k_n,q_1,\mathrm{},q_n)`$ denotes the contribution of the vertex for $`n`$-parton coupling to the hadron $`h`$.
We assume that the initial partons $`k_l,k_l^{}`$ are characterized by small virtualities $`k_l^2s_0`$, $`k_l^2s_0`$, and therefore by small transverse momenta $`k_l_{}^2<s_0`$, $`k_l_{}^2<s_0`$, so that the general results of the analysis made in are applicable for the hadron-parton vertices $`N_h^{(n)}`$. Using
$$d^4k_l=\frac{1}{2}dk_l^+dk_l^{}d^2k_l_{},d^4q_l=\frac{1}{2}dq_l^+dq_l^{}d^2q_l_{},$$
(C.3)
we can perform the integrations over $`k_l^{},k_l_{},q_l^{}`$ and $`k_l^+,k_l_{}^{},q_l^+`$ separately for the upper and the lower vertexes by making use of
$$q_l^2q_l_{}^2,k_l^{},q_l^{}k_l^{},k_l^+,q_l^+k_l^+,$$
(C.4)
as well as the fact that the integrals $`dk_l^+,dk_l^{}`$ are restricted by the physical region
$$0<k_l^+p^+,\underset{l}{}k_l^+p^+,$$
(C.5)
(similar for $`k_l^{}`$) . We shall consider explicitly the upper vertex as for the lower one the derivation is identical. The integrals over $`k_l^{},q_l^{}`$ are defined by the discontinuities of the analytic amplitude $`N_h^{(n)}`$ with respect to the singularities in the corresponding energy invariants
$`s_1^+=(pk_1)^2p^+k_1^{},`$
$`\mathrm{}`$ (C.6)
$`s_n^+=(pk_1\mathrm{}k_n)^2,`$
and
$`s_{q_1}^+=(p+q_1)^2p^+q_1^{},`$
$`\mathrm{}`$ (C.7)
$`s_{q_{n1}}^+=(p+q_1+\mathrm{}+q_{n1})^2.`$
As the processes corresponding to large values of $`s_l^+,s_{q_l}^+`$ need an explicit treatment (the so-called enhanced diagrams, see chapter 5), we only get contributions from the region of large $`k_l^+p^+`$, so that
$$s_l^+p^+(k_1^{}+\mathrm{}+k_l^{})s_0\left(\frac{p^+}{k_1^+}+\mathrm{}+\frac{p^+}{k_l^+}\right)<M_0^2,$$
(C.8)
where $`M_0^2`$ is some minimal mass for the Pomeron asymptotics to be applied. The similar argument holds for the momenta $`q_l^{}`$, such that
$$s_{q_l}^+p^+(q_1^{}+\mathrm{}+q_l^{})<M_0^2.$$
(C.9)
Using
$$dq_l^{}=\frac{ds_{q_l}^+}{p^+},dk_l^{}=\frac{dk_l^2}{k_l^+},$$
(C.10)
we get
$`{\displaystyle \underset{l=1}{\overset{n}{}}\left[\frac{d^4k_l}{(2\pi )^4}\right]\underset{l=1}{\overset{n1}{}}\left[\frac{dq_l^{}}{2\pi }\right]N_{h_1}^{(n)}(p,k_1,\mathrm{},k_n,q_1,\mathrm{},q_n)}`$
$`={\displaystyle \underset{l=1}{\overset{n}{}}\left[\frac{dk_l^2dk_l^+d^2k_l_{}}{2(2\pi )^4k_l^+}\mathrm{\Theta }\left(s_l^+\right)\right]\underset{l=1}{\overset{n1}{}}\left[\frac{ds_{q_l}^+}{2\pi p^+}\mathrm{\Theta }\left(s_{q_l}^+\right)\right]}`$
$`\times disc_{s_1^+,\mathrm{},s_n^+,s_{q_1}^+,\mathrm{},s_{q_{n1}}^+}N_{h_1}^{(n)}(p,k_1,\mathrm{},k_n,q_1,\mathrm{},q_n)\mathrm{\Theta }\left(1{\displaystyle \underset{j=1}{\overset{n}{}}}x_j^+\right)`$
$`{\displaystyle \frac{1}{\left(p^+\right)^{n1}}}{\displaystyle _0^1}{\displaystyle \underset{l=1}{\overset{n}{}}}{\displaystyle \frac{dx_l^+}{x_l^+}}F_{h_1}^{(n)}(x_1^+,\mathrm{}x_n^+,q_1_{}^2,\mathrm{},q_n_{}^2)\mathrm{\Theta }\left(1{\displaystyle \underset{j=1}{\overset{n}{}}}x_j^+\right).`$ (C.11)
The only difference of the formula eq. (C.11) from the traditional expression for $`F_h^{(n)}`$ in is the fact that we keep explicitly the integrations over the light cone momentum shares of the constituent partons $`x_l^+=k_l^+/p^+`$. Further we assume that the dependences on the light cone momentum fractions $`x_l^+`$ and on the momentum transfers along the Pomerons $`q_l^2q_l_{}^2`$ factorize in $`F_h^{(n)}`$, and we use the Gaussian parameterization for the latter one,
$$F_h^{(n)}(x_1^+,\mathrm{}x_n^+,q_1_{}^2,\mathrm{},q_n_{}^2)=\stackrel{~}{F}_h^{(n)}(x_1^+,\mathrm{}x_n^+)\mathrm{exp}\left(R_h^2\underset{j=1}{\overset{n}{}}q_j_{}^2\right),$$
(C.12)
where the parameter $`R_h^2`$ is known as the hadron Regge slope . Based on the above discussion and a corresponding treatment of the lower part of the diagram, eq. (B.1) can be rewritten as
$`iT_{h_1h_2}(s,t)={\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n!}}{\displaystyle \underset{l=1}{\overset{n1}{}}\left[\frac{d^2q_l_{}}{8\pi ^2s}\right]_0^1\underset{l=1}{\overset{n}{}}\frac{dx_l^+}{x_l^+}\frac{dx_l^{}}{x_l^{}}\stackrel{~}{F}_{h_1}^{(n)}(x_1^+,\mathrm{}x_n^+)\stackrel{~}{F}_{h_2}^{(n)}(x_1^{},\mathrm{}x_n^{})}`$
$`\times {\displaystyle \underset{l=1}{\overset{n}{}}}\left[iT_{1\mathrm{I}\mathrm{P}}(\widehat{s}_l,q_l_{}^2)\mathrm{exp}([R_{h_1}^2+R_{h_2}^2]q_l_{}^2)\right]\mathrm{\Theta }(1{\displaystyle \underset{j=1}{\overset{n}{}}}x_j^+)\mathrm{\Theta }(1{\displaystyle \underset{j=1}{\overset{n}{}}}x_j^{})`$ (C.13)
The formula (C.13) can be also obtained using the parton momentum Fock state expansion of the hadron eigenstate
$$|h=\underset{k=1}{\overset{\mathrm{}}{}}\frac{1}{k!}_0^1\underset{l=1}{\overset{k}{}}dx_lf_k(x_1,\mathrm{}x_k)\delta \left(1\underset{j=1}{\overset{k}{}}x_j\right)a^+(x_1)\mathrm{}a^+(x_k)|0,$$
(C.14)
where $`f_k(x_1,\mathrm{}x_k)`$ is the probability amplitude for the hadron $`h`$ to consist of $`k`$ constituent partons with the light cone momentum fractions $`x_1,\mathrm{},x_k`$ and $`a^+\left(x\right)`$ is the creation operator for a parton with the fraction $`x`$. $`f_k(x_1,\mathrm{}x_k)`$ satisfies the normalization condition
$$\underset{k=1}{\overset{\mathrm{}}{}}\frac{1}{k!}_0^1\underset{l=1}{\overset{k}{}}dx_l\left|f_k(x_1,\mathrm{}x_k)\right|^2\delta \left(1\underset{j=1}{\overset{k}{}}x_j\right)=1$$
(C.15)
Then, for the contribution of $`n`$ pair-like scatterings between the parton constituents of the projectile and target hadrons one obtains eq. (C.13), as shown in , with
$$\frac{1}{n!}\stackrel{~}{F}_h^{(n)}(x_1,\mathrm{}x_n)=\underset{k=n}{\overset{\mathrm{}}{}}\frac{1}{k!}\frac{k!}{n!(kn)!}_0^1\underset{l=n+1}{\overset{k}{}}dx_l\left|f_k(x_1,\mathrm{}x_k)\right|^2\delta \left(1\underset{j=1}{\overset{k}{}}x_j\right)$$
(C.16)
representing the “inclusive” momentum distributions of $`n`$ “participating” parton constituents, involved in the scattering process. From the normalization condition (C.15) follows the momentum conservation constraint
$$_0^1𝑑xx\stackrel{~}{F}_h^{(1)}\left(x\right)=1$$
(C.17)
To get further simplifications, we assume that $`\stackrel{~}{F}_{h_1(h_2)}^{(n)}(x_1,\mathrm{}x_n)`$ can be represented in a factorized form as a product of the contributions $`F_{\mathrm{part}}^h(x_l)`$, depending on the momentum shares $`x_l`$ of the “participating” or “active” parton constituents, and on the function $`F_{\mathrm{remn}}^h\left(1_{j=1}^nx_j\right)`$, representing the contribution of all “spectator” partons, sharing the remaining share $`1_jx_j`$ of the initial light cone momentum (see fig. C.1). So we have
$$\stackrel{~}{F}_h^{(n)}(x_1,\mathrm{}x_n)=\underset{l=1}{\overset{n}{}}F_{\mathrm{part}}^h(x_l)F_{\mathrm{remn}}^h\left(1\underset{j=1}{\overset{n}{}}x_j\right)$$
(C.18)
The participating parton constituents are assumed to be quark-anti-quark pairs (not necessarily of identical flavors), such that the baryon numbers of the projectile and of the target are conserved. So we have $`x=x_q+x_{\overline{q}}`$ with $`x_q`$ and $`x_{\overline{q}}`$ being the light-cone momentum fractions of the quark and the anti-quark. The function $`F_{\mathrm{part}}^h`$ may thus be written as
$$F_{\mathrm{part}}^h(x)=𝑑x_q𝑑x_{\overline{q}}\overline{F}_{\mathrm{part}}^h(x_q,x_{\overline{q}})\delta (xx_qx_{\overline{q}}).$$
(C.19)
In case of soft or semi-hard Pomerons, $`\overline{F}_{\mathrm{part}}^h`$ is taken as a product of two asymptotics $`x_i^{\alpha _q},i=q,\overline{q}`$, so we have
$$F_{\mathrm{part}}^h(x)=\gamma _hx^{\alpha _{\mathrm{part}}},$$
(C.20)
with $`\alpha _{\mathrm{part}}=2\alpha _q1`$. The parameter $`\alpha _q`$ defines the probability to slow down the constituent (“dressed”) (anti-)quark; it is related to the Regge intercept of the $`q\overline{q}`$-trajectory : $`\alpha _q=\alpha _{\mathrm{I}\mathrm{R}}(0)0.5`$. The remnant function $`F_{\mathrm{remn}}^h`$ defines the probability to slow down the initial hadron quark configuration; it is assumed to be of the form
$$F_{\mathrm{remn}}^h(x)=x^{\alpha _{\mathrm{remn}}^h},$$
(C.21)
with an adjustable parameter $`\alpha _{\mathrm{remn}}^h`$. Using (C.18-C.21), the eq. (C.13) can be rewritten as
$`iT_{h_1h_2}(s,t)=8\pi ^2s{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n!}}{\displaystyle _0^1}{\displaystyle \underset{l=1}{\overset{n}{}}}dx_l^+dx_l^{}{\displaystyle \underset{l=1}{\overset{n}{}}}\left[{\displaystyle \frac{1}{8\pi ^2\widehat{s}_l}}{\displaystyle d^2q_l_{}iT_{1\mathrm{I}\mathrm{P}}^{h_1h_2}(x_l^+,x_l^{},s,q_l_{}^2)}\right]`$
$`\times F_{\mathrm{remn}}^{h_1}\left(1{\displaystyle \underset{j=1}{\overset{n}{}}}x_j^+\right)F_{\mathrm{remn}}^{h_2}\left(1{\displaystyle \underset{j=1}{\overset{n}{}}}x_j^{}\right)\delta ^{(2)}\left({\displaystyle \underset{k=1}{\overset{n}{}}}\stackrel{}{q}_k_{}\stackrel{}{q}_{}\right).`$ (C.22)
with
$$T_{1\mathrm{I}\mathrm{P}}^{h_1h_2}(x_l^+,x_l^{},s,q_l_{}^2)=T_{1\mathrm{I}\mathrm{P}}(\widehat{s}_l,q_l_{}^2)F_{\mathrm{part}}^{h_1}(x_l^+)F_{\mathrm{part}}^{h_2}(x_l^{})\mathrm{exp}\left(\left[R_{h_1}^2+R_{h_2}^2\right]q_l_{}^2\right)$$
(C.23)
representing the contributions of “elementary interactions plus external legs”.
### C.2 Including Valence Quark Hard Scatterings
To include valence quark hard scatterings one has to replace the inclusive parton momentum distributions
$$\frac{1}{n!}\stackrel{~}{F}_h^{(n)}(x_1,\mathrm{}x_n)$$
(C.24)
in (C.13) by the momentum distributions
$$\frac{1}{n_v!(nn_v)!}\stackrel{~}{F}_h^{(n,n_v)i_1,\mathrm{}i_{n_v}}(x_{v_1},\mathrm{},x_{v_{n_v}},x_{n_v+1},\mathrm{},x_n),$$
(C.25)
corresponding to the case of $`n_v`$ partons being valence quarks with flavors $`i_1,\mathrm{},i_{n_v}`$ (taken at the virtuality scale $`Q_0^2`$) and other $`nn_v`$ partons being usual non-valence participants (quark-anti-quark pairs). One has as well to take into account different contributions for scatterings between a pair of valence quarks or between a valence quark and a non-valence participant. In particular, for a single hard scattering on a valence quark we have to use
$$\stackrel{~}{F}_h^{(1,1)}\left(x_v\right)=q_{\mathrm{val}}^i(x_\nu ),$$
(C.26)
where $`q_{\mathrm{val}}^i`$ is the momentum distribution of a valence quark of flavor $`i`$ at the scale $`Q_0^2`$. In order to conserve the initial hadron baryon content and to keep the simple factorized structure (C.18), we associate a “quasi-spectator” anti-quark to each valence quark interaction, defining the joint contribution $`\overline{F}_{\mathrm{part}}^i(x_v,x_{\overline{q}})`$ of the valence quark with the flavor $`i`$ and the anti-quark. Thus we have
$`\stackrel{~}{F}_h^{(n,n_v)i_1,\mathrm{},i_{n_v}}(x_{v_1},\mathrm{},x_{v_{n_v}},x_{n_{v+1}},\mathrm{},x_n)`$ (C.27)
$`={\displaystyle \underset{l=1}{\overset{n_v}{}}}\left[{\displaystyle _{x_{v_l}}^1}𝑑x_l\overline{F}_{\mathrm{part}}^{h,i_l}(x_{v_l},x_lx_{v_l})\right]{\displaystyle \underset{m=n_v+1}{\overset{n}{}}}F_{\mathrm{part}}^h(x_m)F_{\mathrm{remn}}^h\left(1{\displaystyle \underset{k=1}{\overset{n}{}}}x_k\right),`$
where $`x_l`$ is the sum of the momentum fractions of $`l^{\mathrm{th}}`$ valence quark and the corresponding anti-quark; we allow here formally any number of valence quark participants (based on the fact that multiple valence type processes give negligible contribution to the scattering amplitude). By construction the integral over $`x_{\overline{q}}`$ of the function $`\stackrel{~}{F}_h^{_{(1,1)i}}(x_v,x_{\overline{q}})`$ gives the inclusive momentum distribution of the valence quark $`i`$. Thus the function $`\overline{F}_{\mathrm{part}}^i`$ has to meet the condition
$$_0^{1x_v}𝑑x_{\overline{q}}\overline{F}_{\mathrm{part}}^{h,i}(x_v,x_{\overline{q}})F_{\mathrm{remn}}^h\left(1x_vx_{\overline{q}}\right)=q_{\mathrm{val}}^i(x_v,Q_0^2)$$
(C.28)
Assuming that the anti-quark momentum distribution behaves as $`(x_{\overline{q}})^{\alpha _R}`$, and using the above-mentioned parameterization for $`F_{\mathrm{remn}}^h`$, we get
$$\overline{F}_{\mathrm{part}}^{h,i}(x_v,x_{\overline{q}})=N^1q_{\mathrm{val}}^i(x_v,Q_0^2)(1x_v)^{\alpha _{\mathrm{I}\mathrm{R}}1\alpha _{\mathrm{remn}}}(x_{\overline{q}})^{\alpha _{\mathrm{I}\mathrm{R}}},$$
(C.29)
with the normalization factor
$$N=\frac{\mathrm{\Gamma }\left(1+\alpha _{\mathrm{remn}}\right)\mathrm{\Gamma }\left(1\alpha _{\mathrm{I}\mathrm{R}}\right)}{\mathrm{\Gamma }\left(2+\alpha _{\mathrm{remn}}\alpha _{\mathrm{I}\mathrm{R}}\right)}.$$
(C.30)
Now we can write the normalization condition (C.17) for “active” (participating in an interaction) partons as
$$_0^1𝑑xxF_{\mathrm{part}}^h\left(x\right)F_{\mathrm{remn}}^h\left(1x\right)+\underset{i}{}_0^1x_v𝑑x_0^x𝑑x_v\overline{F}_{\mathrm{part}}^{h,i}(x_v,xx_v)F_{\mathrm{remn}}^h\left(1x\right)=1,$$
(C.31)
which gives
$$\gamma _h=(1x_v)\frac{\mathrm{\Gamma }\left(3\alpha _{\mathrm{part}}+\alpha _{\mathrm{remn}}^h\right)}{\mathrm{\Gamma }\left(2\alpha _{\mathrm{part}}\right)\mathrm{\Gamma }\left(1+\alpha _{\mathrm{remn}}^h\right)},$$
(C.32)
with
$$x_v=\underset{i}{}_0^1x_v𝑑x_vq_{\mathrm{val}}^i(x_v,Q_0^2)$$
(C.33)
being the average light cone momentum fraction carried by valence quarks.
### C.3 Enhanced Diagrams
In this appendix, we demonstrate how triple Pomeron contributions appear naturally in the Gribov-Regge formalism under certain kinematical conditions, and we derive a formula for the hadron-hadron scattering amplitude in this case.
To introduce enhanced type diagrams let us come back to the process of double soft Pomeron exchange, which is a particular case of the diagram of fig 2.7. The corresponding contribution to the elastic scattering amplitude is given in eqs. (2.31), (C.11) with $`n=2`$ and with $`T_{1\mathrm{I}\mathrm{P}}`$ being replaced by $`T_{\mathrm{soft}}`$:
$`iT_{h_1h_2}^{(2)}(s,t)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{d^4k_1}{(2\pi )^4}\frac{d^4k_1^{}}{(2\pi )^4}\frac{d^4k_2}{(2\pi )^4}\frac{d^4k_2^{}}{(2\pi )^4}\frac{d^4q_1}{(2\pi )^4}\mathrm{\Theta }\left(s_1^+\right)\mathrm{\Theta }\left(s_1^{}\right)\mathrm{\Theta }\left(s_2^+\right)\mathrm{\Theta }\left(s_2^{}\right)\mathrm{\Theta }\left(s_{q_1}^+\right)\mathrm{\Theta }\left(s_{q_1}^{}\right)}`$ (C.34)
$`\times \mathrm{disc}_{s_1^+,s_2^+,s_{q_1}^+}N_{h_1}^{(2)}(p,k_1,k_2,q_1,qq_1)`$
$`\times {\displaystyle \underset{l=1}{\overset{2}{}}}\left[iT_{\mathrm{soft}}(\widehat{s}_l,q_l^2)\right]\mathrm{disc}_{s_1^{},s_2^{},s_{q_1}^{}}N_{h_2}^{(2)}(p^{},k_1^{},k_2^{},q_1,q+q_1)`$
see fig. C.2.
We are now interested in the contribution with some of the invariants
$`s_1^+`$ $`=`$ $`(pk_1)^2p^+k_1^{},`$
$`s_2^+`$ $`=`$ $`(pk_1k_2)^2p^+(k_1^{}+k_2^{}),`$ (C.35)
$`s_{q_1}^+`$ $`=`$ $`(p+q_1)^2p^+q_1^{},`$
being large, implying $`k_i^{},q_1^{}`$ to be not too small. As was shown in , one can restrict the integration region to $`s_i^+s_{q_1}^+`$, because $`s_i^+s_{q_1}^+`$, ($`k_i^{}q_1^{}`$ ) either correspond to processes with exchange of more than two Pomerons or to the Pomeron self-coupling, the latter one just renormalizing the Pomeron amplitude. Then from
$$k_i^2,k_i^2,(k_1^{}q_1)^2,(k_1+q_1)^2,(k_2q_1q)^2s_0$$
(C.36)
it follows that
$$k_i^{}\frac{s_0}{k_i^+},k_i^+\frac{s_0}{k_i^{}},q_1^+\frac{s_0}{k_1^{}},q_1^{}\frac{s_0}{k_1^+},\frac{s_0}{k_2^+}$$
(C.37)
and correspondingly $`k_1^+k_2^+`$ and the invariants $`s_1^+,s_2^+,s_{q_1}^+`$ are of the same order. The vertex $`N_h^{(2)}`$ for large $`s_{q_1}^+`$ can be described by the soft Pomeron asymptotics and we may write $`disc_{s_1^+,s_2^+,s_{q_1}^+}N_{h_1}^{(2)}(p,k_1,k_2,q_1,qq_1)`$ as a product of the Pomeron-hadron coupling $`N_h^{(1)}(p,k,q)`$, twice imaginary part of the soft Pomeron exchange amplitude $`2\mathrm{I}\mathrm{m}T_{\mathrm{soft}}((kk_{12})^2,q^2)`$ with $`k_{12}=k_1+k_2`$, and a term $`V^{3\mathrm{P}}(k_{12},k_1,k_2,q,q_1,qq_1)`$ describing the coupling of the three Pomerons<sup>1</sup><sup>1</sup>1 The triple Pomeron vertex is assumed to have nonplanar structure, corresponding to having the two lower Pomerons “in parallel”; the planar triple-Pomeron vertex would correspond to subsequent emission of these Pomerons and gives no contribution in the high energy limit . . So we get
$`{\displaystyle \frac{d^4k_1}{(2\pi )^4}\frac{d^4k_2}{(2\pi )^4}\frac{dq_1^{}}{2\pi }\mathrm{\Theta }\left(s_1^+\right)\mathrm{\Theta }\left(s_2^+\right)\mathrm{\Theta }\left(s_{q_1}^+\right)𝑑isc_{s_1^+,s_2^+,s_{q_1}^+}N_{h_1}^{(2)}(p,k_1,k_2,q_1,qq_1)}`$
$`={\displaystyle \frac{d^4k_1}{(2\pi )^4}\frac{d^4k_2}{(2\pi )^4}\frac{dq_1^{}}{2\pi }\mathrm{\Theta }\left(s_1^+\right)\mathrm{\Theta }\left(s_2^+\right)\mathrm{\Theta }\left(s_{q_1}^+\right)\frac{d^4k}{(2\pi )^4}\frac{d^4k_{12}}{(2\pi )^4}(2\pi )^4\delta (k_{12}k_1k_2)}`$
$`\times N_{h_1}^{(1)}(p,k,q)\mathrm{\hspace{0.17em}2}\mathrm{Im}T_{\mathrm{soft}}((kk_{12})^2,q^2)V^{3\mathrm{I}\mathrm{P}}(k_{12},k_1,k_2,q,q_1,qq_1)`$
$`={\displaystyle \frac{dk^+dk^2d^2k_{}}{2k^+(2\pi )^4}\mathrm{\Theta }\left(s_0^+\right)𝑑isc_{s_0^+}N_{h_1}^{(1)}(p,k,q)\frac{dk_{12}^+dk_{12}^2d^2k_{12_{}}}{2k_{12}^+(2\pi )^4}\mathrm{\hspace{0.17em}2}\mathrm{Im}T_{\mathrm{soft}}((kk_{12})^2,q^2)}`$
$`\times {\displaystyle }{\displaystyle \frac{dk_1^+dk_1^2d^2k_1_{}}{2k_1^+(2\pi )^4}}{\displaystyle \frac{d(k_{12}k_1q_1q)^2}{2\pi (k_{12}^+k_1^+)}}\mathrm{\Theta }\left(s_1^+\right)\mathrm{\Theta }\left(s_2^+\right)\mathrm{\Theta }\left(s_{q_1}^+\right)`$
$`\times V^{3\mathrm{I}\mathrm{P}}(k_{12},k_1,k_{12}k_1,q,q_1,qq_1)`$ (C.38)
see fig. C.3,
with $`s_0^+=(pk)^2p^+k^{}`$ and $`dq_1^{}=\frac{1}{k_2^+}d(k_2q_1q)^2`$. Now we can perform the integrations over $`k^2,k_{12}^2,k_1^2,(k_2q_1q)^2`$ assuming their convergence in the region $`k_i^2s_0`$, as well as over $`k_{},k_{12_{}},k_1_{}s_0`$ to transform eq. (C.38) to the form
$`{\displaystyle 𝑑x^+F_{\mathrm{part}}^{h_1}(x^+)F_{\mathrm{remn}}^{h_1}\left(1x^+\right)\mathrm{exp}\left(R_{h_1}^2q_{}^2\right)\frac{dx_{12}^+}{x_{12}^+}\frac{1}{2s^+}\mathrm{Im}T_{\mathrm{soft}}(s^+,q_{}^2)}`$ $`\times `$
$`\times {\displaystyle }{\displaystyle \frac{dz^+}{x_1^+(x_{12}^+x_1^+)p^+}}\stackrel{~}{V}^{3\mathrm{I}\mathrm{P}}(q_{}^2,q_1_{}^2,(\stackrel{}{q}_{}\stackrel{}{q}_1_{})^2)`$ (C.39)
with
$`x^+=k^+/p^+,`$ (C.40)
$`x_{12}^+=k_{12}^+/p^+,`$ (C.41)
$`x_1^+=k_1^+/p^+,`$ (C.42)
$`x_{12}^+x_1^+=(k_{12}^+k_1^+)/p^+=k_2^+/p^+,`$ (C.43)
$`z^+=k_1^+/k_{12}^+=x_1^+/x_{12}^+,`$ (C.44)
$`s^+=(kk_{12})^2k^+k_{12}^{}s_0k^+/k_{12}^+=s_0x^+/x_{12}^+.`$ (C.45)
We used (see eq. (C.11-C.18))
$$\frac{dk^2d^2k_{}}{2(2\pi )^4}\mathrm{\Theta }\left(s_0^+\right)𝑑isc_{s_0^+}N_h^{(1)}(p,k,q)=F_{\mathrm{part}}^h(x^+)F_{\mathrm{remn}}^h\left(1x^+\right)\mathrm{exp}\left(R_h^2q_{}^2\right)$$
(C.46)
and the definition
$`\stackrel{~}{V}^{3\mathrm{I}\mathrm{P}}(q_{}^2,q_1_{}^2,(\stackrel{}{q}_{}\stackrel{}{q}_1_{})^2)`$ $`=`$ $`s_0{\displaystyle \frac{dk_{12}^2d^2k_{12_{}}}{(2\pi )^4}\frac{dk_1^2d^2k_1_{}}{(2\pi )^4}\frac{d\left[(k_{12}k_1q_1q)^2\right]}{2\pi }}`$
$`\times \mathrm{\Theta }\left(s_1^+\right)\mathrm{\Theta }\left(s_2^+\right)\mathrm{\Theta }\left(s_{q_1}^+\right)V^{3\mathrm{I}\mathrm{P}}(k_{12},k_1,k_{12}k_1,q,q_1,qq_1).`$
We have also taken into account the fact that the triple Pomeron vertex $`V^{3\mathrm{I}\mathrm{P}}`$ has a scalar structure, and we therefore suppose that it can only depend on the invariants
$$k_i^2,(k_1+q_1)^2,(k_2q_1q)^2s_0,q^2q_{}^2,q_1^2q_1_{}^2,q_2^2(q_{}q_1_{})^2$$
(C.48)
and on the partition of the light cone momentum $`k_{12}^+`$ between the two lower Pomerons,
$$z^+=k_1^+/k_{12}^+=(p^{}+k_1)^2/(p^{}+k_{12})^2.$$
Furthermore, we did assume the flat distribution in $`z^+`$ in order to obtain $`k_1^+k_2^+k_{12}^+/2`$ (see the discussion above). We use the Gaussian parameterization for the $`q_i^2`$-dependence of $`\stackrel{~}{V}^{3\mathrm{I}\mathrm{P}}`$,
$$\stackrel{~}{V}^{3\mathrm{I}\mathrm{P}}(q_{}^2,q_1_{}^2,q_2_{}^2)r_{3\mathrm{I}\mathrm{P}}\mathrm{exp}\left(R_{3\mathrm{I}\mathrm{P}}^2\left[q_{}^2+q_1_{}^2+q_2_{}^2\right]\right),$$
(C.49)
where $`r_{3\mathrm{I}\mathrm{P}}`$ is the triple-Pomeron coupling and $`R_{3\mathrm{I}\mathrm{P}}^2`$ is the slope of the triple-Pomeron vertex. The slope $`R_{3\mathrm{I}\mathrm{P}}^2`$ is known to be small and will be neglected in the following.
Now, using (C.34-C.49) for the upper vertex and doing the usual treatment (C.11-C.18) of the lower one, we get for the triple Pomeron amplitude
$$iT_{h_1h_2}^{3\mathrm{I}\mathrm{P}}(s,t)=_0^1\frac{dx^+}{x^+}\frac{dx^{}}{x^{}}F_{\mathrm{remn}}^{h_1}\left(1x^+\right)F_{\mathrm{remn}}^{h_2}\left(1x^{}\right)iT_{3\mathrm{I}\mathrm{P}}^{h_1h_2}(x^+,x^{},s,t)$$
(C.50)
with
$`iT_{3\mathrm{I}\mathrm{P}+}^{h_1h_2}(x^+,x^{},s,t)=\mathrm{\hspace{0.25em}8}\pi ^2x^+x^{}s{\displaystyle \frac{r_{3\mathrm{I}\mathrm{P}}}{2}}{\displaystyle _{s_0/x^{}}^{x^+}}{\displaystyle \frac{dx_{12}^+}{x_{12}^+}}\left[{\displaystyle \frac{1}{2s^+}}\mathrm{Im}T^{h_1}(x^+,s^+,q_{}^2)\right]`$
$`\times {\displaystyle }dz^+{\displaystyle }d^2q_1_{}d^2q_2_{}{\displaystyle _0^x^{}}dx_1^{}dx_2^{}{\displaystyle \underset{l=1}{\overset{2}{}}}\left[{\displaystyle \frac{1}{8\pi ^2\widehat{s}_l}}iT^{h_2}(x_l^{},\widehat{s}_l,q_l_{}^2)\right]`$
$`\times \delta (x^{}x_1^{}x_2^{})\delta ^{(2)}\left(\stackrel{}{q}_{}\stackrel{}{q}_1_{}\stackrel{}{q}_2_{}\right),`$ (C.51)
with
$$T^h(x,\widehat{s},q_{}^2)=T_{\mathrm{soft}}^h(x,\widehat{s},q_{}^2)=T_{\mathrm{soft}}(\widehat{s},q_{}^2)F_{\mathrm{part}}^h(x)\mathrm{exp}\left(R_h^2q_{}^2\right)$$
(C.52)
and
$$\widehat{s}_1=x_{12}^+z^+x_1^{}s,\widehat{s}_2=x_{12}^+(1z^+)x_2^{}s.$$
(C.53)
The sign “$``$” in “$`3\mathrm{I}\mathrm{P}`$” refers to the Pomeron “splitting” towards the target hadron (reversed $`Y`$-diagram); the lower limit for the integral $`dx_{12}^+`$ is due to $`x_{12}^{}s_0/x_{12}^+<x^{}`$.
### C.4 Parton Generation for Triple-Pomeron Diagrams
The inclusion of the triple-Pomeron contributions results only in slight modification of the standard procedure. Now the full contribution of an elementary interaction is
$$G_{1\mathrm{I}\mathrm{P}}^{h_1h_2}(x^+,x^{},s,b)+\underset{\sigma }{}\underset{i=0}{\overset{2}{}}\widehat{\widehat{G}}_{3\mathrm{I}\mathrm{P}\sigma (i)}^{h_1h_2}(x^+,x^{},x^{\mathrm{proj}},x^{\mathrm{targ}},s,b),$$
(C.54)
where $`x^{\mathrm{proj}},x^{\mathrm{targ}}`$ are the corresponding remnant light cone momentum fractions, and with the contributions of different cuts of triple-Pomeron diagrams being obtained from eq. (5.54), together with eqs. (5.54), (5.47-5.2) as
$`\widehat{\widehat{G}}_{3\mathrm{I}\mathrm{P}(0)}^{h_1h_2}(x^+,x^{},x^{\mathrm{proj}},x^{\mathrm{targ}},s,b)`$ $`=`$ $`{\displaystyle \frac{r_{3\mathrm{I}\mathrm{P}}}{8}}{\displaystyle d^2b_1\frac{1}{x^{}}G^{h_1}(x^+,x^+x^{}s,\left|\stackrel{}{b}\stackrel{}{b}_1\right|)}`$ (C.55)
$`\times `$ $`{\displaystyle _0^{x^{\mathrm{targ}}}}𝑑\widehat{x}^{}{\displaystyle _0^1}𝑑z^+{\displaystyle _0^{\widehat{x}^{}+x^{}}}𝑑x_1^{}G^{h_2}(x_1^{},x_1^{}{\displaystyle \frac{s_0}{x^{}}}z^+s,b_1)`$
$`\times `$ $`G^{h_2}(\widehat{x}^{}+x^{}x_1^{},(\widehat{x}^{}+x^{}x_1^{}){\displaystyle \frac{s_0}{x^{}}}(1z^+)s,b_1)`$
$`\times `$ $`{\displaystyle \frac{F_{\mathrm{remn}}\left(x^{\mathrm{targ}}\widehat{x}^{}\right)}{F_{\mathrm{remn}}\left(x^{\mathrm{targ}}\right)}}`$
and
$`\widehat{\widehat{G}}_{3\mathrm{I}\mathrm{P}(1)}^{h_1h_2}(x^+,x^{},x^{\mathrm{proj}},x^{\mathrm{targ}},s,b)`$ $`=`$ $`{\displaystyle \frac{r_{3\mathrm{I}\mathrm{P}}}{2}}{\displaystyle d^2b_1_{s_0/(x^{}s)}^{x^+}\frac{dx_{12}^+}{x_{12}^+}G^{h_1}(x^+,s^+,\left|\stackrel{}{b}\stackrel{}{b}_1\right|)}`$ (C.56)
$`\times `$ $`{\displaystyle _0^1}𝑑z^+G^{h_2}(x^{},x_{12}^+z^+x^{}s,b_1)`$
$`\times `$ $`{\displaystyle _0^{x^{\mathrm{targ}}}}𝑑\widehat{x}^{}G^{h_2}(\widehat{x}^{},x_{12}^+(1z^+)\widehat{x}^{}s,b_1)`$
$`\times `$ $`{\displaystyle \frac{F_{\mathrm{remn}}\left(x^{\mathrm{targ}}\widehat{x}^{}\right)}{F_{\mathrm{remn}}\left(x^{\mathrm{targ}}\right)}}`$
and
$`\widehat{\widehat{G}}_{3\mathrm{I}\mathrm{P}(2)}^{h_1h_2}(x^+,x^{},x^{\mathrm{proj}},x^{\mathrm{targ}},s,b)`$ $`=`$ $`{\displaystyle \frac{r_{3\mathrm{I}\mathrm{P}}}{4}}{\displaystyle d^2b_1_{s_0/(sx^{})}^{x^+}\frac{dx_{12}^+}{x_{12}^+}G^{h_1}(x^+,s^+,\left|\stackrel{}{b}\stackrel{}{b}_1\right|)}`$
$`\times `$ $`{\displaystyle _0^1}𝑑z^+{\displaystyle _0^x^{}}𝑑x_1^{}G^{h_2}(x_1^{},\widehat{s}_1,b_1)G^{h_2}(x^{}x_1^{},\widehat{s}_2,b_1)`$
(similarly for $`\widehat{\widehat{G}}_{3\mathrm{I}\mathrm{P}+(i)}^{h_1h_2}`$). Thus, the elementary process splits into a single elementary scattering contribution – with the weight
$$\widehat{\widehat{G}}_{1\mathrm{I}\mathrm{P}}^{h_1h_2}=G_{1\mathrm{I}\mathrm{P}}^{h_1h_2}+\underset{\sigma }{}\underset{i=0}{\overset{1}{}}\widehat{\widehat{G}}_{3\mathrm{I}\mathrm{P}\sigma (i)}^{h_1h_2},$$
(C.58)
and the contribution with all three Pomerons being cut – with the weight
$$\widehat{\widehat{G}}_{3\mathrm{I}\mathrm{P}+(2)}^{h_1h_2}+\widehat{\widehat{G}}_{3\mathrm{I}\mathrm{P}(2)}^{h_1h_2}.$$
(C.59)
Choosing the first one, one proceeds in the usual way, with the functions $`G_{1\mathrm{I}\mathrm{P}}^{h_1h_2},`$ $`G_{\mathrm{soft}}^{h_1h_2},`$ $`G_{\mathrm{sea}\mathrm{sea}}^{h_1h_2},`$ $`G_{\mathrm{val}\mathrm{sea}}^{h_1h_2},`$ $`G_{\mathrm{sea}v\mathrm{al}}^{h_1h_2}`$ in eq. (6.4) being replaced by $`\widehat{\widehat{G}}_{1\mathrm{I}\mathrm{P}}^{h_1h_2},`$ $`\widehat{\widehat{G}}_{\mathrm{soft}}^{h_1h_2},`$ $`\widehat{\widehat{G}}_{\mathrm{sea}\mathrm{sea}}^{h_1h_2},`$ $`\widehat{\widehat{G}}_{\mathrm{val}\mathrm{sea}}^{h_1h_2},`$ $`\widehat{\widehat{G}}_{\mathrm{sea}\mathrm{val}}^{h_1h_2}`$, where we introduce the functions $`\widehat{\widehat{G_J}}^{h_1h_2}`$ via
$$\widehat{\widehat{G_J}}^{h_1h_2}=G_J^{h_1h_2}+\widehat{\widehat{G}}_{3\mathrm{I}\mathrm{P}+(0)J}^{h_1h_2}+\widehat{\widehat{G}}_{3\mathrm{I}\mathrm{P}(0)J}^{h_1h_2}+\widehat{\widehat{G}}_{3\mathrm{I}\mathrm{P}(1)J}^{h_1h_2}+\widehat{\widehat{G}}_{3\mathrm{I}\mathrm{P}(1)J}^{h_1h_2},$$
(C.60)
with $`J`$ being “soft”, “sea-sea”, “val-sea” or “sea-val”. Here $`\widehat{\widehat{G}}_{3\mathrm{I}\mathrm{P}(0)\mathrm{soft}}^{h_1h_2}`$ is given by the eq. (C.55) with $`G^{h_1}(x^+,x^+x^{}s,\left|\stackrel{}{b}\stackrel{}{b}_1\right|)`$ being replaced by its soft contribution $`G_{\mathrm{soft}}^{h_1}(x^+,x^+x^{}s,\left|\stackrel{}{b}\stackrel{}{b}_1\right|)`$ (see eq. (5.34)), whereas $`\widehat{\widehat{G}}_{3\mathrm{I}\mathrm{P}(1)\mathrm{soft}}^{h_1h_2}`$ is given in (C.56) with both $`G^{h_1}(x^+,s^+,\left|\stackrel{}{b}\stackrel{}{b}_1\right|)`$ and $`G^{h_2}(x^{},x_{12}^+z^+x^{}s,b_1)`$ being represented by the soft contributions $`G_{\mathrm{soft}}^{h_1/h_2}`$. The other functions $`\widehat{\widehat{G}}_{3\mathrm{I}\mathrm{P}+(i)J}^{h_1h_2}`$ are defined similarly to the “soft” case - considering corresponding “sea” and “valence” contributions in the cut Pomerons $`G^{h_1/h_2}`$ in eq. (C.55-C.56).
For the contribution corresponding to the cases of all three Pomerons being cut, we split the processes into three separate cut Pomeron pieces. Each piece is characterized by the function $`\overline{G}_{i\pm }^{h_1h_2}(\overline{x}^+,\overline{x}^{},s,b_1)`$, where $`b_1b/2`$ and
$`\overline{G}_1^{h_1h_2}(\overline{x}^+,\overline{x}^{},s,b_1)=G^{h_1}(\overline{x}^+,\overline{x}^+\overline{x}^{}s,b_1),\overline{x}^+=x^+,\overline{x}^{}=s_0/(x_{12}^+s),`$ (C.61)
$`\overline{G}_2^{h_1h_2}(\overline{x}^+,\overline{x}^{},s,b_1)=G^{h_2}(\overline{x}^{},\overline{x}^+\overline{x}^{}s,b_1),\overline{x}^+=z^+x_{12}^+,\overline{x}^{}=x_1^{},`$ (C.62)
$`\overline{G}_3^{h_1h_2}(\overline{x}^+,\overline{x}^{},s,b_1)=G^{h_2}(\overline{x}^{},\overline{x}^+\overline{x}^{}s,b_1),\overline{x}^+=(1z^+)x_{12}^+,\overline{x}^{}=x^{}x_1^{},`$ (C.63)
where the variables $`x_{12}^+,z^+,x_1^{}`$ are generated according to the integrand of eq. (C.4) (similar for $`\overline{G}_{i+}^{h_1h_2}`$). After that each contribution $`\overline{G}_{i\pm }^{h_1h_2}(\overline{x}^+,\overline{x}^{},s,b_1)`$ is treated separately in the usual way, starting from the eq. (6.4), with the functions $`G_{1\mathrm{I}\mathrm{P}}^{h_1h_2},`$ $`G_{\mathrm{soft}}^{h_1h_2},`$ $`G_{\mathrm{sea}\mathrm{sea}}^{h_1h_2},`$ $`G_{\mathrm{val}\mathrm{sea}}^{h_1h_2},`$ $`G_{\mathrm{sea}v\mathrm{al}}^{h_1h_2}`$ being replaced by $`\overline{G}_{i\pm }^{h_1h_2}(\overline{x}^+,\overline{x}^{},s,b_1)`$ and by the corresponding partial contributions of soft $`G_{\mathrm{soft}}^h`$, “sea-sea”-type $`G_{\mathrm{sea}\mathrm{sea}}^h`$, and “valence-sea”-type parton scattering $`G_{\mathrm{val}\mathrm{sea}}^h`$ (see eq. (5.34-5.36)), but without “valence-valence” contribution.
## Appendix D Calculation of $`\mathrm{\Phi }`$ and $`H`$
### D.1 Calculation of $`\mathrm{\Phi }_{pp}`$
Here, we present the detailed calculation of $`\mathrm{\Phi }`$ for proton-proton collisions. As shown earlier, $`\mathrm{\Phi }_{pp}`$ may be written as
$`\mathrm{\Phi }_{pp}(x^+,x^{}s,b)`$ $`=`$ $`{\displaystyle \underset{r_1=0}{\overset{\mathrm{}}{}}}\mathrm{}{\displaystyle \underset{r_N=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{r_1!}}\mathrm{}{\displaystyle \frac{1}{r_N!}}{\displaystyle \underset{\lambda =1}{\overset{r_1+\mathrm{}+r_N}{}}dx_\lambda ^+dx_\lambda ^{}}`$ (D.1)
$`\times `$ $`{\displaystyle \underset{\rho _1=1}{\overset{r_1}{}}}G_{1\rho _1}\mathrm{}{\displaystyle \underset{\rho _N=r_1+\mathrm{}+r_{N1}+1}{\overset{r_1+\mathrm{}+r_N}{}}}G_{N\rho _N}`$
$`\times `$ $`F_{\mathrm{remn}}(x^+{\displaystyle \underset{\lambda }{}}x_\lambda ^+)F_{\mathrm{remn}}(x^{}{\displaystyle \underset{\lambda }{}}x_\lambda ^{}),`$
with
$`F_{\mathrm{remn}}(x)`$ $`=`$ $`x^{\alpha _{\mathrm{remn}}}\mathrm{\Theta }(x)\mathrm{\Theta }(1x),`$ (D.2)
and with $`G_{i\lambda }`$ being of the form
$$G_{i\lambda }(x_\lambda ^+,x_\lambda ^{},s,b)=\alpha _i(x_\lambda ^+x_\lambda ^{})^{\beta _i},$$
(D.3)
with $`\alpha _i`$ and $`\beta _i`$ being $`s`$\- and $`b`$-dependent parameters. We obtain
$`\mathrm{\Phi }_{pp}(x^+,x^{},s,b)`$ $`=`$ $`{\displaystyle \underset{r_1=0}{\overset{\mathrm{}}{}}}\mathrm{}{\displaystyle \underset{r_N=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(\alpha _1)^{r_1}}{r_1!}}\mathrm{}{\displaystyle \frac{(\alpha _N)^{r_N}}{r_N!}}`$ (D.4)
$`\times `$ $`I_{r_1,\mathrm{},r_N}(x^+)I_{r_1,\mathrm{},r_N}(x^{})`$
with
$$I_{r_1,\mathrm{},r_N}(x)=\underset{\lambda =1}{\overset{r_1+\mathrm{}+r_N}{}}dx_\lambda \underset{\rho _1=1}{\overset{r_1}{}}x_{\rho _1}^{\beta _1}\mathrm{}\underset{\rho _N=r_1+\mathrm{}+r_{N1}+1}{\overset{r_1+\mathrm{}+r_N}{}}x_{\rho _N}^{\beta _N}F_{\mathrm{remn}}(x\underset{\lambda }{}x_\lambda ),$$
(D.5)
which amounts to
$$I_{r_1,\mathrm{},r_N}(x)=\underset{\lambda =1}{\overset{r_1+\mathrm{}+r_N}{}}[dx_\lambda x_\lambda ^{ϵ_\lambda }](x\underset{\lambda }{}x_\lambda )^{\alpha _{\mathrm{remn}}}\theta (x\underset{\lambda }{}x_\lambda )\theta \left(1(x\underset{\lambda }{}x_\lambda )\right),$$
(D.6)
with
$$ϵ_\lambda =\{\begin{array}{ccc}\beta _1\hfill & \mathrm{for}& \lambda r_1\hfill \\ \beta _2\hfill & \mathrm{for}& r_1<\lambda r_1+r_2\hfill \\ \mathrm{}\hfill & & \\ \beta _N\hfill & \mathrm{for}& r_1+\mathrm{}+r_{N1}<\lambda r_1+\mathrm{}+r_N\hfill \end{array}.$$
(D.7)
We define new variables,
$$\{\begin{array}{ccc}\hfill u_\lambda & =& \frac{x_\lambda }{xx_1\mathrm{}x_{\lambda 1}}\hfill \\ & & \\ \hfill du_\lambda & =& \frac{dx_\lambda }{xx_1\mathrm{}x_{\lambda 1}}\hfill \end{array},$$
(D.8)
which have the following property,
$$\underset{\alpha =1}{\overset{\lambda 1}{}}(1u_\alpha )=\underset{\alpha =1}{\overset{\lambda 1}{}}\frac{x\mathrm{}x_\alpha }{x\mathrm{}x_{\alpha 1}}=\frac{x\mathrm{}x_{\lambda 1}}{x},$$
(D.9)
and therefore
$$\{\begin{array}{ccc}\hfill x_\lambda & =& xu_\lambda _{\alpha =1}^{\lambda 1}(1u_\alpha )\hfill \\ & & \\ \hfill dx_\lambda & =& xdu_\lambda _{\alpha =1}^{\lambda 1}(1u_\alpha )\hfill \end{array}.$$
(D.10)
This leads to
$$I_{r_1,\mathrm{},r_N}(x)=\underset{\lambda =1}{\overset{r_1+\mathrm{}+r_N}{}}\left\{du_\lambda u_\lambda ^{ϵ_\lambda }x^{1+ϵ_\lambda }\underset{\alpha =1}{\overset{\lambda 1}{}}(1u_\alpha )^{1+ϵ_\lambda }\right\}\left[x\underset{\lambda =1}{\overset{r_1+\mathrm{}+r_N}{}}(1u_\lambda )\right]^{\alpha _{\mathrm{remn}}}$$
(D.11)
Defining
$$\alpha =\alpha _{\mathrm{remn}}+\underset{\lambda =1}{\overset{r_1+\mathrm{}+r_N}{}}\stackrel{~}{ϵ}_\lambda =\alpha _{\mathrm{remn}}+r_1\stackrel{~}{\beta }_1+\mathrm{}+r_N\stackrel{~}{\beta }_N$$
(D.12)
and
$$\gamma _\lambda =\alpha _{\mathrm{remn}}+\underset{\nu =\lambda +1}{\overset{r_1+\mathrm{}+r_N}{}}\stackrel{~}{ϵ}_\nu =\{\begin{array}{ccc}\alpha _{\mathrm{remn}}+(r_1\lambda )\stackrel{~}{\beta }_1+r_2\stackrel{~}{\beta }2+\mathrm{}+r_N\stackrel{~}{\beta }_N\hfill & \mathrm{if}& \lambda r_1\hfill \\ \alpha _{\mathrm{remn}}+(r_1+r_2\lambda )\stackrel{~}{\beta }_2+r_3\stackrel{~}{\beta }_3+\mathrm{}+r_N\stackrel{~}{\beta }_N\hfill & \mathrm{if}& r_1<\lambda r_1+r_2\hfill \\ \mathrm{}\hfill & & \\ \alpha _{\mathrm{remn}}+(r_1+\mathrm{}+r_N\lambda )\stackrel{~}{\beta }_N\hfill & \mathrm{if}& \lambda >r_1+\mathrm{}+r_{N1}\hfill \end{array}$$
(D.13)
with
$`\stackrel{~}{\beta }`$ $`=`$ $`\beta +1,`$ (D.14)
$`\stackrel{~}{ϵ}`$ $`=`$ $`ϵ+1,`$ (D.15)
we find
$$I_{r_1,\mathrm{},r_N}(x)=x^\alpha \underset{\lambda =1}{\overset{r_1+\mathrm{}+r_N}{}}du_\lambda u_\lambda ^{ϵ_\lambda }(1u_\lambda )^{\gamma _\lambda }$$
(D.16)
The $`u`$-integration can be done,
$$_0^1𝑑uu^ϵ(1u)^\gamma =\frac{\mathrm{\Gamma }(1+ϵ)\mathrm{\Gamma }(1+\gamma )}{\mathrm{\Gamma }(2+ϵ+\gamma )},$$
(D.17)
and we get
$$I_{r_1,\mathrm{},r_N}(x)=x^\alpha \underset{\lambda =1}{\overset{r_1+\mathrm{}+r_N}{}}\frac{\mathrm{\Gamma }(1+ϵ_\lambda )\mathrm{\Gamma }(1+\gamma _\lambda )}{\mathrm{\Gamma }(2+ϵ_\lambda +\gamma _\lambda )}$$
(D.18)
Using the relation $`1+ϵ_\lambda +\gamma _\lambda =\gamma _{\lambda 1}`$, we get
$`I_{r_1,\mathrm{},r_N}(x)`$ $`=`$ $`x^\alpha \mathrm{\Gamma }(1+\beta _1)^{r_1}\mathrm{}\mathrm{\Gamma }(1+\beta _N)^{r_N}{\displaystyle \underset{\lambda =1}{\overset{r_1+\mathrm{}+r_N}{}}}{\displaystyle \frac{\mathrm{\Gamma }(1+\gamma _\lambda )}{\mathrm{\Gamma }(1+\gamma _{\lambda 1})}}`$ (D.19)
$`=`$ $`x^\alpha \mathrm{\Gamma }(1+\beta _1)^{r_1}\mathrm{}\mathrm{\Gamma }(1+\beta _N)^{r_N}{\displaystyle \frac{\mathrm{\Gamma }\left(1+\gamma _{r_1+\mathrm{}+r_N}\right)}{\mathrm{\Gamma }(1+\gamma _0)}},`$ (D.20)
or
$$I_{r_1,\mathrm{},r_N}(x)=x^{\alpha _{\mathrm{remn}}+r_1\stackrel{~}{\beta }_1+\mathrm{}+r_N\stackrel{~}{\beta }_N}\mathrm{\Gamma }(\stackrel{~}{\beta }_1)^{r_1}\mathrm{}\mathrm{\Gamma }(\stackrel{~}{\beta }_N)^{r_N}\frac{\mathrm{\Gamma }(\stackrel{~}{\alpha }_{\mathrm{remn}})}{\mathrm{\Gamma }(\stackrel{~}{\alpha }_{\mathrm{remn}}+r_1\stackrel{~}{\beta }_1+\mathrm{}+r_N\stackrel{~}{\beta }_N)}.$$
(D.21)
The final expression for $`\mathrm{\Phi }_{pp}`$ is therefore
$`\mathrm{\Phi }_{pp}(x^+,x^{},s,b)`$ $`=`$ $`x^{\alpha _{\mathrm{remn}}}{\displaystyle \underset{r_1=0}{\overset{\mathrm{}}{}}}\mathrm{}{\displaystyle \underset{r_N=0}{\overset{\mathrm{}}{}}}\left\{{\displaystyle \frac{\mathrm{\Gamma }(1+\alpha _{\mathrm{remn}})}{\mathrm{\Gamma }(1+\alpha _{\mathrm{remn}}+r_1\stackrel{~}{\beta }_1+\mathrm{}+r_N\stackrel{~}{\beta }_N)}}\right\}^2`$ (D.22)
$`\times `$ $`{\displaystyle \frac{(\alpha _1x^{\stackrel{~}{\beta }_1}\mathrm{\Gamma }^2(\stackrel{~}{\beta }_1))^{r_1}}{r_1!}}\mathrm{}{\displaystyle \frac{(\alpha _Nx^{\stackrel{~}{\beta }_N}\mathrm{\Gamma }^2(\stackrel{~}{\beta }_N))^{r_N}}{r_N!}}`$
with $`x=x^+x^{}`$. This is the expression shown in eq. (3.15).
### D.2 Calculation of H
The function $`H`$ is defined as
$`H(z^+,z^{})`$ $`=`$ $`{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{\mu =1}{\overset{m}{}}dx_\mu ^+dx_\mu ^{}\frac{1}{m!}\underset{\mu =1}{\overset{m}{}}G(x_\mu ^+,x_\mu ^{},s,b)}`$ (D.23)
$`\times `$ $`\delta (1z^+{\displaystyle \underset{\mu =1}{\overset{m}{}}}x_\mu ^+)\delta (1z^{}{\displaystyle \underset{\mu =1}{\overset{m}{}}}x_\mu ^{}).`$
Using the expression
$`G(x_\mu ^+,x_\mu ^{},s,b)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}\alpha _i(x_\mu ^+x_\mu ^{})^{\beta _i},`$ (D.24)
with the $`\alpha _i`$ and $`\beta _i`$ are functions of the impact parameter $`b`$ and the energy squared $`s`$,
$`\alpha _i`$ $`=`$ $`\left(\alpha _{D_i}+\alpha _{D_i}^{}\right)s^{\beta _{D_i}+\gamma _{D_i}b^2}e^{\frac{b^2}{\delta _{D_i}}},`$ (D.25)
$`\beta _i`$ $`=`$ $`\beta _{D_i}+\gamma _{D_i}b^2+\beta _{D_i}^{}\alpha _{\mathrm{part}},`$ (D.26)
with $`\alpha _{D_i}^{}0`$ and $`\beta _{D_i}^{}0`$ only if $`\alpha _{D_i}=0`$, and using the same method as for the calculation of $`\mathrm{\Phi }_{pp}(x^+,x^{})`$, one may write
$`H(z^+,z^{})`$ $`=`$ $`\underset{r_1+\mathrm{}+r_N=0}{\overset{\mathrm{}}{\underset{}{{\displaystyle \underset{r_1=0}{\overset{\mathrm{}}{}}}\mathrm{}{\displaystyle \underset{r_N=0}{\overset{\mathrm{}}{}}}}}}{\displaystyle \frac{(\alpha _1)^{r_1}}{r_1!}}\mathrm{}{\displaystyle \frac{(\alpha _N)^{r_N}}{r_N!}}J_{r_1,\mathrm{},r_N}(z^+)J_{r_1,\mathrm{},r_N}(z^{})`$ (D.27)
with
$$J_{r_1,\mathrm{},r_N}(z)=\underset{\lambda =1}{\overset{r_1+\mathrm{}+r_N}{}}[dx_\lambda x_\lambda ^{ϵ_\lambda }]\delta \left(1z\underset{\lambda =1}{\overset{r_1+\mathrm{}+r_N}{}}x_\lambda \right)$$
(D.28)
and
$$ϵ_\lambda =\{\begin{array}{ccc}\beta _1\hfill & \mathrm{for}& \lambda r_1\hfill \\ \beta _2\hfill & \mathrm{for}& r_1<\lambda r_1+r_2\hfill \\ \mathrm{}\hfill & & \\ \beta _N\hfill & \mathrm{for}& r_1+\mathrm{}+r_{N1}<\lambda r_1+\mathrm{}+r_N\hfill \end{array}.$$
(D.29)
One may use the $`\delta `$ function to obtain
$$J_{r_1,\mathrm{},r_N}(z)=\underset{\lambda =2}{\overset{r_1+\mathrm{}+r_N}{}}[dx_\lambda x_\lambda ^{ϵ_\lambda }]\left((1z)\underset{\lambda =2}{\overset{r_1+\mathrm{}+r_N}{}}x_\lambda \right)^{ϵ_1}.$$
(D.30)
Introducing $`\stackrel{~}{x}_\alpha =x_{\alpha +1}`$, we define new variables,
$$\{\begin{array}{ccc}\hfill u_\lambda ^{}& =& \frac{\stackrel{~}{x}_\lambda ^{}}{1z\stackrel{~}{x}_1\mathrm{}\stackrel{~}{x}_{\lambda ^{}1}}\hfill \\ & & \\ \hfill du_\lambda ^{}& =& \frac{d\stackrel{~}{x}_\lambda ^{}}{1z\stackrel{~}{x}_1\mathrm{}\stackrel{~}{x}_{\lambda ^{}1}}\hfill \end{array},$$
(D.31)
which have the following property,
$$\underset{\alpha =1}{\overset{\lambda ^{}1}{}}(1u_\alpha )=\underset{\alpha =1}{\overset{\lambda ^{}1}{}}\frac{1z\mathrm{}\stackrel{~}{x}_\alpha }{1z\mathrm{}\stackrel{~}{x}_{\alpha 1}}=\frac{1z\mathrm{}\stackrel{~}{x}_{\lambda ^{}1}}{1z},$$
(D.32)
and therefore
$$\{\begin{array}{ccc}\hfill \stackrel{~}{x}_\lambda ^{}& =& (1z)u_\lambda ^{}_{\alpha =1}^{\lambda ^{}1}(1u_\alpha )\hfill \\ & & \\ \hfill d\stackrel{~}{x}_\lambda ^{}& =& (1z)du_\lambda ^{}_{\alpha =1}^{\lambda ^{}1}(1u_\alpha )\hfill \end{array}.$$
(D.33)
This leads to
$`J_{r,s,t}(z)`$ $`=`$ $`{\displaystyle \underset{\lambda ^{}=1}{\overset{r_1+\mathrm{}+r_N1}{}}\left\{du_\lambda ^{}u_\lambda ^{}^{ϵ_{\lambda ^{}+1}}(1z)^{1+ϵ_{\lambda ^{}+1}}\underset{\alpha =1}{\overset{\lambda ^{}1}{}}(1u_\alpha )^{1+ϵ_{\lambda ^{}+1}}\right\}}`$ (D.34)
$`\times `$ $`\left[(1z){\displaystyle \underset{\lambda ^{}=1}{\overset{r_1+\mathrm{}+r_N1}{}}}(1u_\lambda ^{})\right]^{ϵ_1}`$
Defining
$$\alpha ^{}=ϵ_1+\underset{\lambda ^{}=1}{\overset{r_1+\mathrm{}+r_N1}{}}\stackrel{~}{ϵ}_{\lambda ^{}+1}=r_1\stackrel{~}{\beta }_1+\mathrm{}+r_N\stackrel{~}{\beta }_N1$$
(D.35)
and
$`\gamma _\lambda ^{}`$ $`=`$ $`ϵ_1+{\displaystyle \underset{\nu =\lambda ^{}+1}{\overset{r_1+\mathrm{}+r_N1}{}}}\stackrel{~}{ϵ}_{\nu +1}`$ (D.36)
$`=`$ $`\{\begin{array}{ccc}ϵ_1+(r_11\lambda ^{})\stackrel{~}{\beta }_1+r_2\stackrel{~}{\beta }_2+\mathrm{}+r_N\stackrel{~}{\beta }_N\hfill & \mathrm{if}& \lambda ^{}r_11\hfill \\ ϵ_1+(r_1+r_21\lambda ^{})\stackrel{~}{\beta }_2+r_3\stackrel{~}{\beta }_3+\mathrm{}+r_N\stackrel{~}{\beta }_N\hfill & \mathrm{if}& r_11<\lambda ^{}r_1+r_21\hfill \\ \mathrm{}\hfill & & \\ ϵ_1+(r_1+\mathrm{}+r_N1\lambda ^{})\stackrel{~}{\beta }_N\hfill & \mathrm{if}& \lambda ^{}>r_1+\mathrm{}+r_{N1}1\hfill \end{array},`$ (D.41)
with
$`\stackrel{~}{\beta }`$ $`=`$ $`\beta +1,`$ (D.42)
$`\stackrel{~}{ϵ}`$ $`=`$ $`ϵ+1,`$ (D.43)
we find
$$J_{r_1,\mathrm{},r_N}(z)=(1z)^\alpha ^{}\underset{\lambda ^{}=1}{\overset{r_1+\mathrm{}+r_N1}{}}du_\lambda ^{}u_\lambda ^{}^{ϵ_{\lambda ^{}+1}}(1u_\lambda ^{})^{\gamma _\lambda ^{}}$$
(D.44)
The $`u`$-integration can be done,
$$_0^1𝑑uu^ϵ(1u)^\gamma =\frac{\mathrm{\Gamma }(1+ϵ)\mathrm{\Gamma }(1+\gamma )}{\mathrm{\Gamma }(2+ϵ+\gamma )}$$
(D.45)
and we get
$$J_{r_1,\mathrm{},r_N}(z)=(1z)^\alpha ^{}\underset{\lambda ^{}=1}{\overset{r_1+\mathrm{}+r_N1}{}}\frac{\mathrm{\Gamma }(1+ϵ_{\lambda ^{}+1})\mathrm{\Gamma }(1+\gamma _\lambda ^{})}{\mathrm{\Gamma }(2+ϵ_{\lambda ^{}+1}+\gamma _\lambda ^{})}$$
(D.46)
Using the relation $`1+ϵ_{\lambda ^{}+1}+\gamma _\lambda ^{}=\gamma _{\lambda ^{}1}`$, we get, if $`r_10`$,
$`J_{r_1,\mathrm{},r_N}(z)`$ $`=`$ $`(1z)^\alpha ^{}\mathrm{\Gamma }(\stackrel{~}{\beta }_1)^{r_11}\mathrm{\Gamma }(\stackrel{~}{\beta }_2)^{r_2}\mathrm{}\mathrm{\Gamma }(\stackrel{~}{\beta }_N)^{r_N}{\displaystyle \underset{\lambda ^{}=1}{\overset{r_1+\mathrm{}+r_N1}{}}}{\displaystyle \frac{\mathrm{\Gamma }(1+\gamma _\lambda ^{})}{\mathrm{\Gamma }(1+\gamma _{\lambda ^{}1})}}`$ (D.47)
$`=`$ $`(1z)^\alpha ^{}\mathrm{\Gamma }(\stackrel{~}{\beta }_1)^{r_11}\mathrm{\Gamma }(\stackrel{~}{\beta }_2)^{r_2}\mathrm{}\mathrm{\Gamma }(\stackrel{~}{\beta }_N)^{r_N}{\displaystyle \frac{\mathrm{\Gamma }(1+\gamma _{r_1+\mathrm{}+r_N1})}{\mathrm{\Gamma }(1+\gamma _0)}}`$ (D.48)
$`=`$ $`(1z)^\alpha ^{}\mathrm{\Gamma }(\stackrel{~}{\beta }_1)^{r_11}\mathrm{\Gamma }(\stackrel{~}{\beta }_2)^{r_2}\mathrm{}\mathrm{\Gamma }(\stackrel{~}{\beta }_N)^{r_N}{\displaystyle \frac{\mathrm{\Gamma }(\stackrel{~}{\beta }_1)}{\mathrm{\Gamma }(r\stackrel{~}{\beta }_1+\mathrm{}+r_N\stackrel{~}{\beta }_N)}}.`$ (D.49)
If $`r_1=0`$ and $`r_20`$, we get
$$J_{r_1,\mathrm{},r_N}(z)=(1z)^\alpha ^{}\mathrm{\Gamma }(\stackrel{~}{\beta }_2)^{r_21}\mathrm{\Gamma }(\stackrel{~}{\beta }_3)^{r_3}\mathrm{}\mathrm{\Gamma }(\stackrel{~}{\beta }_N)^{r_N}\frac{\mathrm{\Gamma }(\stackrel{~}{\beta }_2)}{\mathrm{\Gamma }(r_2\stackrel{~}{\beta }_2+\mathrm{}+r_N\stackrel{~}{\beta }_N)},$$
(D.50)
and so on, which corresponds finally to
$$J_{r_1,\mathrm{},r_N}(z)=(1z)^{r_1\stackrel{~}{\beta }_1+\mathrm{}+r_N\stackrel{~}{\beta }_N1}\frac{\mathrm{\Gamma }(\stackrel{~}{\beta }_1)^{r_1}\mathrm{}\mathrm{\Gamma }(\stackrel{~}{\beta }_N)^{r_N}}{\mathrm{\Gamma }(r_1\stackrel{~}{\beta }_1+\mathrm{}+r_N\stackrel{~}{\beta }_N)}.$$
(D.51)
The final expression for $`H`$ is therefore
$`H(z^+,z^{})`$ $`=`$ $`\underset{r_1+\mathrm{}+r_N0}{\underset{}{{\displaystyle \underset{r_1=0}{\overset{\mathrm{}}{}}}\mathrm{}{\displaystyle \underset{r_N=0}{\overset{\mathrm{}}{}}}}}{\displaystyle \frac{\left[(1z^+)(1z^{})\right]^{r_1\stackrel{~}{\beta }_1+\mathrm{}+r_N\stackrel{~}{\beta }_N1}}{\mathrm{\Gamma }(r_1\stackrel{~}{\beta }_1+\mathrm{}+r_N\stackrel{~}{\beta }_N)^2}}`$ (D.52)
$`\times {\displaystyle \frac{(\alpha _1\mathrm{\Gamma }(\stackrel{~}{\beta }_1)^2)^{r_1}}{r_1!}}\mathrm{}{\displaystyle \frac{(\alpha _N\mathrm{\Gamma }(\stackrel{~}{\beta }_N)^2)^{r_N}}{r_N!}}`$
### D.3 Calculation of $`\mathrm{\Phi }_{AB}`$
The expression for the virtual emissions in case of nucleus-nucleus collisions is given as
$`\mathrm{\Phi }_{AB}(X^+,X^{},s,b)`$ $`=`$ $`{\displaystyle \underset{l_1=0}{\overset{\mathrm{}}{}}}\mathrm{}{\displaystyle \underset{l_{AB}=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{k=1}{\overset{AB}{}}\left\{\underset{\lambda =1}{\overset{l_k}{}}d\stackrel{~}{x}_{k,\lambda }^+d\stackrel{~}{x}_{k,\lambda }^{}\right\}\underset{k=1}{\overset{AB}{}}\left\{\frac{1}{l_k!}\underset{\lambda =1}{\overset{l_k}{}}G(\stackrel{~}{x}_{k,\lambda }^+,\stackrel{~}{x}_{k,\lambda }^{},s,b)\right\}}`$ (D.53)
$`\times `$ $`{\displaystyle \underset{i=1}{\overset{A}{}}}F_{\mathrm{remn}}\left(x_i^+{\displaystyle \underset{\pi (k)=i}{}}\stackrel{~}{x}_{k,\lambda }^+\right){\displaystyle \underset{j=1}{\overset{B}{}}}F_{\mathrm{remn}}\left(x_j^{}{\displaystyle \underset{\tau (k)=j}{}}\stackrel{~}{x}_{k,\lambda }^{}\right),`$
where $`X^+=\left\{x_1^+\mathrm{}x_A^+\right\}`$, $`X^{}=\left\{x_1^{}\mathrm{}x_B^{}\right\}`$ and $`\pi (k)`$ and $`\tau (k)`$ represent the projectile or target nucleon linked to pair $`k`$. Using the expression
$`G(\stackrel{~}{x}_{k,\lambda }^+,\stackrel{~}{x}_{k,\lambda }^{},s,b)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}\underset{G_{i,k,\lambda }}{\underset{}{\alpha _i(\stackrel{~}{x}_{k,\lambda }^+\stackrel{~}{x}_{k,\lambda }^{})^{\beta _i}}},`$ (D.54)
with the $`\alpha _i`$ and $`\beta _i`$ are functions of the impact parameter $`b`$ and the energy squared $`s`$,
$`\alpha _i`$ $`=`$ $`\left(\alpha _{D_i}+\alpha _{D_i}^{}\right)s^{\beta _{D_i}+\gamma _{D_i}b^2}e^{\frac{b^2}{\delta _{D_i}}},`$ (D.55)
$`\beta _i`$ $`=`$ $`\beta _{D_i}+\gamma _{D_i}b^2+\beta _{D_i}^{}\alpha _{\mathrm{part}},`$ (D.56)
with $`\alpha _{D_i}^{}0`$ and $`\beta _{D_i}^{}0`$ only if $`\alpha _{D_i}=0`$. The remnant function $`F_{\mathrm{remn}}`$ is given as
$$F_{\mathrm{remn}}(x)=x^{\alpha _{\mathrm{remn}}}\mathrm{\Theta }(x)\mathrm{\Theta }(1x).$$
(D.57)
We have (see calculation of $`\mathrm{\Phi }_{pp}`$)
$`\mathrm{}{\displaystyle \underset{l_k=0}{\overset{\mathrm{}}{}}}\mathrm{}{\displaystyle \frac{1}{l_k!}}{\displaystyle \underset{\lambda =1}{\overset{l_k}{}}}G(s,\stackrel{~}{x}_{k,\lambda }^+,\stackrel{~}{x}_{k,\lambda }^{},b)`$ $`=`$ $`\mathrm{}{\displaystyle \underset{l_k=0}{\overset{\mathrm{}}{}}}\mathrm{}{\displaystyle \frac{1}{l_k!}}{\displaystyle \underset{\lambda =1}{\overset{l_k}{}}}\left(G_{1,k,\lambda }\mathrm{}G_{N,k,\lambda }\right)`$ (D.58)
$`=`$ $`\mathrm{}{\displaystyle \underset{r_{1,k}=0}{\overset{\mathrm{}}{}}}\mathrm{}{\displaystyle \underset{r_{N,k}=0}{\overset{\mathrm{}}{}}}\mathrm{}{\displaystyle \frac{1}{r_{1,k}!\mathrm{}r_{N,k}!}}`$
$`\times {\displaystyle \underset{\rho _1=1}{\overset{r_{1,k}}{}}}G_{1,k,\rho _1}\mathrm{}{\displaystyle \underset{\rho _N=r_{1,k}+\mathrm{}+r_{N1,k}+1}{\overset{r_{1,k}+\mathrm{}+r_{N,k}}{}}}G_{N,k,\rho _N}.`$
So eq. (D.53) can be written as
$`\mathrm{\Phi }_{AB}(X^+,X^{},s,b)`$ $`=`$ $`{\displaystyle \underset{r_{1,1}\mathrm{}r_{N,1}}{}}\mathrm{}{\displaystyle \underset{r_{1,AB}\mathrm{}r_{N,AB}}{}}{\displaystyle \underset{k=1}{\overset{AB}{}}}{\displaystyle \frac{1}{r_{1,k}!\mathrm{}r_{N,k}!}}{\displaystyle \underset{k=1}{\overset{AB}{}}\left\{\underset{\lambda =1}{\overset{r_{1,k}+\mathrm{}+r_{N,k}}{}}d\stackrel{~}{x}_{k,\lambda }^+d\stackrel{~}{x}_{k,\lambda }^{}\right\}}`$ (D.60)
$`\times `$ $`{\displaystyle \underset{k=1}{\overset{AB}{}}}\left\{{\displaystyle \underset{\rho _1=1}{\overset{r_{1,k}}{}}}\alpha _1(\stackrel{~}{x}_{k,\rho _1}^+\stackrel{~}{x}_{k,\rho _1}^{})^{\beta _1}\mathrm{}{\displaystyle \underset{\rho _N=r_{1,k}+\mathrm{}+r_{N1,k}+1}{\overset{r_{1,k}+\mathrm{}+r_{N,k}}{}}}\alpha _N(\stackrel{~}{x}_{k,\rho _N}^+\stackrel{~}{x}_{k,\rho _N}^{})^{\beta _N}\right\}`$
$`\times `$ $`{\displaystyle \underset{i=1}{\overset{A}{}}}F_{\mathrm{remn}}\left(x_i^+{\displaystyle \underset{\pi (k)=i}{}}\stackrel{~}{x}_{k,\lambda }^+\right){\displaystyle \underset{j=1}{\overset{B}{}}}F_{\mathrm{remn}}\left(x_j^{}{\displaystyle \underset{\tau (k)=j}{}}\stackrel{~}{x}_{k,\lambda }^{}\right),`$
which leads to
$`\mathrm{\Phi }_{AB}(X^+,X^{},s,b)`$ $`=`$ $`{\displaystyle \underset{r_{1,1}\mathrm{}r_{N,1}}{}}\mathrm{}{\displaystyle \underset{r_{1,AB}\mathrm{}r_{N,AB}}{}}{\displaystyle \underset{k=1}{\overset{AB}{}}}{\displaystyle \frac{\left(\alpha _1\right)^{r_{1,k}}}{r_{1,k}!}}\mathrm{}{\displaystyle \frac{\left(\alpha _N\right)^{r_{N,k}}}{r_{N,k}!}}`$ (D.61)
$`\times I_R^+(X^+)I_R^{}(X^{}),`$
where
$`I_R^\sigma (X)`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{AB}{}}\underset{\lambda =1}{\overset{r_{1,k}+\mathrm{}+r_{N,k}}{}}d\stackrel{~}{x}_{k,\lambda }\left(\stackrel{~}{x}_{k,\lambda }\right)^{ϵ_{k,\lambda }}\underset{i=1}{\overset{\mathrm{P}^\sigma }{}}F_{\mathrm{remn}}\left(x_i\underset{\kappa ^\sigma (k)=i}{}\stackrel{~}{x}_{k,\lambda }\right)},`$ (D.62)
with $`R=\{r_{j,k}\}`$ , $`\mathrm{P}^+=A`$, $`\mathrm{P}^{}=B`$, $`\kappa ^+(k)=\pi (k)`$, and $`\kappa ^{}(k)=\tau (k)`$. Using the property
$$\underset{k=1}{\overset{AB}{}}=\underset{i=1}{\overset{\mathrm{P}^\sigma }{}}\underset{\kappa ^\sigma (k)=i}{},$$
(D.63)
one can write
$`I_R^\sigma (X)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{\mathrm{P}^\sigma }{}}}\{{\displaystyle }{\displaystyle \underset{\kappa ^\sigma (k)=i}{}}{\displaystyle \underset{\lambda =1}{\overset{r_{1,k}+\mathrm{}+r_{N,k}}{}}}d\stackrel{~}{x}_{k,\lambda }\left(\stackrel{~}{x}_{k,\lambda }\right)^{ϵ_{k,\lambda }}`$ (D.64)
$`\times F_{\mathrm{remn}}(x_i{\displaystyle \underset{\kappa ^\sigma (k)=i}{}}\stackrel{~}{x}_{k,\lambda })\}.`$
Let us rename the $`\stackrel{~}{x}_{k,\lambda }`$ linked to the remnant $`i`$ as $`\stackrel{~}{x}_1,\stackrel{~}{x}_2,\mathrm{},\stackrel{~}{x}_{r_{1,i}^\sigma +\mathrm{}+r_{N,i}^\sigma }`$, where $`r_{p,i}^\sigma `$ is per definition the number of Pomerons of type $`p`$ linked to remnant $`i`$. So we get for the term in brackets
$$\underset{\nu =1}{\overset{r_{1,i}^\sigma +\mathrm{}+r_{N,i}^\sigma }{}}d\stackrel{~}{x}_\nu \stackrel{~}{x}_\nu ^{ϵ_\nu }F_{\mathrm{remn}}\left(x_i\underset{\nu =1}{\overset{r_{1,i}^\sigma +\mathrm{}+r_{N,i}^\sigma }{}}\stackrel{~}{x}_\nu \right).$$
(D.65)
This is exactly the corresponding $`I`$ for proton-proton scattering. So we have
$$I_R^\sigma (X)=\underset{i=1}{\overset{\mathrm{P}^\sigma }{}}x_i^{\alpha _{\mathrm{remn}}+r_{1,i}^\sigma \stackrel{~}{\beta }_1+\mathrm{}+r_{N,i}^\sigma \stackrel{~}{\beta }_N}\mathrm{\Gamma }(\stackrel{~}{\beta }_1)^{r_{1,i}^\sigma }\mathrm{}\mathrm{\Gamma }(\stackrel{~}{\beta }_N)^{r_{N,i}^\sigma }g\left(r_{1,i}^\sigma \stackrel{~}{\beta }_1+\mathrm{}+r_{N,i}^\sigma \stackrel{~}{\beta }_N\right),$$
(D.66)
with
$$g(z)=\frac{\mathrm{\Gamma }(1+\alpha _{\mathrm{remn}})}{\mathrm{\Gamma }(1+\alpha _{\mathrm{remn}}+z)},$$
(D.67)
and
$$\stackrel{~}{\beta }=\beta +1.$$
(D.68)
Since we have
$$r_{p,i}^\sigma =\underset{\kappa ^\sigma (k)=i}{}r_{p,k},$$
(D.69)
we find finally
$`\mathrm{\Phi }_{AB}(X^+,X^{},s,b)=`$
$`{\displaystyle \underset{r_{1,1}\mathrm{}r_{N,1}}{}}\mathrm{}{\displaystyle \underset{r_{1,AB}\mathrm{}r_{N,AB}}{}}{\displaystyle \underset{k=1}{\overset{AB}{}}}{\displaystyle \frac{\left(\alpha _1\right)^{r_{1,k}}}{r_{1,k}!}}\mathrm{}{\displaystyle \frac{\left(\alpha _N\right)^{r_{N,k}}}{r_{N,k}!}}`$ (D.70)
$`{\displaystyle \underset{i=1}{\overset{A}{}}}(x_i^+)^{\alpha _{\mathrm{remn}}}{\displaystyle \underset{\pi (k)=i}{}}\left(\mathrm{\Gamma }(\stackrel{~}{\beta }_1)(x_i^+)^{\stackrel{~}{\beta }_1}\right)^{r_{1,k}}\mathrm{}\left(\mathrm{\Gamma }(\stackrel{~}{\beta }_N)(x_i^+)^{\stackrel{~}{\beta }_N}\right)^{r_{N,k}}g\left({\displaystyle \underset{\pi (k)=i}{}}r_{1,k}\stackrel{~}{\beta }_1+\mathrm{}+r_{N,k}\stackrel{~}{\beta }_N\right)`$
$`{\displaystyle \underset{j=1}{\overset{B}{}}}(x_j^{})^{\alpha _{\mathrm{remn}}}{\displaystyle \underset{\tau (k)=j}{}}\left(\mathrm{\Gamma }(\stackrel{~}{\beta }_1)(x_j^{})^{\stackrel{~}{\beta }_1}\right)^{r_{1,k}}\mathrm{}\left(\mathrm{\Gamma }(\stackrel{~}{\beta }_N)(x_j^{})^{\stackrel{~}{\beta }_N}\right)^{r_{N,k}}g\left({\displaystyle \underset{\tau (k)=j}{}}r_{1,k}\stackrel{~}{\beta }_1+\mathrm{}+r_{N,k}\stackrel{~}{\beta }_N\right)`$
### D.4 Exponentiation of $`\mathrm{\Phi }_{AB}`$
We first replace in eq. (D.70) the function $`g(z)`$ by the “exponentiated” function $`g_\mathrm{e}(z)`$,
$`g_\mathrm{e}\left({\displaystyle \underset{\kappa ^\sigma (k)=i}{}}r_{1,k}\stackrel{~}{\beta }_1+\mathrm{}+r_{N,k}\stackrel{~}{\beta }_N\right)`$ $`=`$ $`\mathrm{exp}\left\{ϵ_\mathrm{e}\left({\displaystyle \underset{\kappa ^\sigma (k)=i}{}}r_{1,k}\stackrel{~}{\beta }_1+\mathrm{}+r_{N,k}\stackrel{~}{\beta }_N\right)\right\}`$ (D.71)
$`=`$ $`{\displaystyle \underset{\kappa ^\sigma (k)=i}{}}\left(e^{ϵ_\mathrm{e}\stackrel{~}{\beta }_1}\right)^{r_{1,k}}\mathrm{}\left(e^{ϵ_\mathrm{e}\stackrel{~}{\beta }_N}\right)^{r_{N,k}},`$
and we obtain
$`\mathrm{\Phi }_{\mathrm{e}_{AB}}(X^+,X^{},s,b)=`$
$`{\displaystyle \underset{i=1}{\overset{A}{}}}(x_i^+)^{\alpha _{\mathrm{remn}}}{\displaystyle \underset{j=1}{\overset{B}{}}}(x_j^{})^{\alpha _{\mathrm{remn}}}`$ (D.72)
$`{\displaystyle \underset{r_{1,1}\mathrm{}r_{1,AB}}{}}{\displaystyle \underset{k=1}{\overset{AB}{}}}{\displaystyle \frac{\left(\alpha _1\right)^{r_{1,k}}}{r_{1,k}!}}\left\{{\displaystyle \underset{i=1}{\overset{A}{}}}{\displaystyle \underset{\pi (k)=i}{}}\left(D_1(x_i^+)^{\stackrel{~}{\beta }_1}\right)^{r_{1,k}}{\displaystyle \underset{j=1}{\overset{B}{}}}{\displaystyle \underset{\tau (k)=j}{}}\left(D_1(x_j^{})^{\stackrel{~}{\beta }_1}\right)^{r_{1,k}}\right\}`$
$`\mathrm{}`$
$`{\displaystyle \underset{r_{N,1}\mathrm{}r_{N,AB}}{}}{\displaystyle \underset{k=1}{\overset{AB}{}}}{\displaystyle \frac{\left(\alpha _N\right)^{r_{N,k}}}{r_{N,k}!}}\left\{{\displaystyle \underset{i=1}{\overset{A}{}}}{\displaystyle \underset{\pi (k)=i}{}}\left(D_N(x_i^+)^{\stackrel{~}{\beta }_N}\right)^{r_{N,k}}{\displaystyle \underset{j=1}{\overset{B}{}}}{\displaystyle \underset{\tau (k)=j}{}}\left(D_N(x_j^{})^{\stackrel{~}{\beta }_N}\right)^{r_{N,k}}\right\}`$
with
$$D_j=\mathrm{\Gamma }(\stackrel{~}{\beta }_j)e^{ϵ_\mathrm{e}\stackrel{~}{\beta }_j}.$$
(D.73)
And using again the property (D.63), we can write
$`\mathrm{\Phi }_{e_{AB}}(X^+,X^{},s,b)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{A}{}}}(x_i^+)^{\alpha _{\mathrm{remn}}}{\displaystyle \underset{j=1}{\overset{B}{}}}(x_j^{})^{\alpha _{\mathrm{remn}}}`$
$`\times `$ $`{\displaystyle \underset{r_{1,1}\mathrm{}r_{1,AB}}{}}{\displaystyle \underset{k=1}{\overset{AB}{}}}{\displaystyle \frac{\left(\alpha _1D_1^2(x_{\pi (k)}^+x_{\tau (k)}^{})^{\stackrel{~}{\beta }_1}\right)^{r_{1,k}}}{r_{1,k}!}}`$
$`\times `$ $`\mathrm{}`$
$`\times `$ $`{\displaystyle \underset{r_{N,1}\mathrm{}r_{N,AB}}{}}{\displaystyle \underset{k=1}{\overset{AB}{}}}{\displaystyle \frac{\left(\alpha _ND_N^2(x_{\pi (k)}^+x_{\tau (k)}^{})^{\stackrel{~}{\beta }_N}\right)^{r_{N,k}}}{r_{N,k}!}}.`$
Now the sum can be performed and we get the final expression for the “exponentiated” $`\mathrm{\Phi }`$,
$`\mathrm{\Phi }_{\mathrm{e}_{AB}}(X^+,X^{},s,b)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{A}{}}}(x_i^+)^{\alpha _{\mathrm{remn}}}{\displaystyle \underset{j=1}{\overset{B}{}}}(x_j^{})^{\alpha _{\mathrm{remn}}}{\displaystyle \underset{k=1}{\overset{AB}{}}}e^{\stackrel{~}{G}\left(x_{\pi (k)}^+x_{\tau (k)}^{}\right)},`$
with
$$\stackrel{~}{G}(x)=\underset{i=1}{\overset{N}{}}\stackrel{~}{\alpha }_ix^{\stackrel{~}{\beta }_i},$$
(D.75)
where
$`\stackrel{~}{\alpha }_i`$ $`=`$ $`\alpha _i\mathrm{\Gamma }(\stackrel{~}{\beta }_i)^2e^{2ϵ_\mathrm{e}\stackrel{~}{\beta }_i},`$ (D.76)
$`\stackrel{~}{\beta }_i`$ $`=`$ $`\beta _i+1,`$ (D.77)
and
$`\alpha _i`$ $`=`$ $`\left(\alpha _{D_i}+\alpha _{D_i}^{}\right)s^{\beta _{D_i}+\gamma _{D_i}b^2}e^{\frac{b^2}{\delta _{D_i}}},`$ (D.78)
$`\beta _i`$ $`=`$ $`\beta _{D_i}+\gamma _{D_i}b^2+\beta _{D_i}^{}\alpha _{\mathrm{part}},`$ (D.79)
with $`\alpha _{D_i}^{}0`$ and $`\beta _{D_i}^{}0`$ only if $`\alpha _{D_i}=0`$ .
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# 1 Introduction
## 1 Introduction
The simplest and best known example of a self-dual system is electrodynamics in vacuum. The set of Maxwell’s equations is invariant under the simultaneous replacements $`\stackrel{}{E}\stackrel{}{B},\stackrel{}{B}\stackrel{}{E}`$. While being a symmetry of the Hamiltonian $`H=\stackrel{}{E}^2+\stackrel{}{B}^2`$, the Lagrangian does transform: $`L=\stackrel{}{E}^2\stackrel{}{B}^2L`$. The generalization to a $`(p1)`$-form potential $`C`$ in $`d=2p`$ dimensions with action $`S=\mathrm{d}C\mathrm{d}C`$ is immediate.
These theories are free systems with linear equations of motion. The interesting question is whether one can construct interacting self-dual systems. The main goal of these notes is to discuss the conditions (self-duality equations) which have to be satisfied by the action of a dynamical system in order to be self-dual, in the sense to be specified below. Apparently Schrödinger was the first to discuss nonlinear self-duality. In he reformulated the Born-Infeld (BI) theory in such a way that it was manifestly invariant under U(1) duality rotations. We will mainly be interested in four-dimensional nonlinear systems of gauge fields coupled to matter. For non-supersymmetric systems the results have been obtained, as a generalization of patterns of duality in extended supergravity (see also ), in and reviewed and extended in . Our special emphasis is on manifestly $`𝒩=1,2`$ supersymmetric generalizations.
As will be discussed below, self-dual theories possess quite remarkable properties. Our main concern, however, in pursuing the study of such systems lies in the fact that self-duality turns out to be intimately connected with spontaneous breaking of supersymmetry (for still not completely understood reasons). Recently several models for partial breaking of $`𝒩=2`$ supersymmetry to $`𝒩=1`$ in four dimensions have been constructed. Two most prominent models – described by the Goldstone-Maxwell multiplet and by the tensor Goldstone multiplet – are self-dual $`𝒩=1`$ supersymmetric theories; the other Goldstone multiplets are dual superfield version of the tensor one (as we will describe, self-duality may be consistent with the existence of dual formulations). In our opinion, this cannot be accidental.
It may look curious but the fact that the nonlinear superfield constraint, which underlies the Goldstone-Maxwell construction of , has turned out to be fruitful for nontrivial generalizations. This constraint was used in to derive nonlinear U$`(n)`$ duality invariant models, both in non-supersymmetric and supersymmetric cases. In the present paper, we apply the nonlinear constraint, which is at the heart of the tensor Goldstone construction of , to derive new self-dual systems.
These notes are organized as follows. In sect. 2 we review nonlinear electrodynamics: we define the notion of self-duality and state the self-duality equation which has to be satisfied by the action. The derivation can be found in Appendix A. We also discuss various properties of self-dual nonlinear electrodynamics, e.g. when coupled to a complex scalar field. We then proceed with a description of the general structure of self-dual Lagrangians, of which the Born-Infeld action is but a particular example, with very special properties, though. In sect. 3 we present, following Refs. , the generalization to a collection of U(1) vector-fields, coupled to an arbitrary number of scalar fields. Sect. 4, which is based on Ref. , is the $`𝒩=1`$ supersymmetric version of sect. 2. In sect. 5 we discuss properties of the supersymmetric Born-Infeld action and make contact with the work of Bagger and Galperin , where this action was obtained as a model of partial $`𝒩=2𝒩=1`$ supersymmetry breaking. In the next section we supersymmetrize the analysis of sect. 4. In sect. 7 we discuss self-dual models with tensor multiplets. In sect. 8, we temporarily leave supersymmetry and derive the self-duality equations and determine the maximal duality group of a $`d`$-dimensional system with $`n`$ Abelian $`(p1)`$-form potentials and $`m`$ Abelian $`(dp1)`$-form potentials, with and without coupling to scalar fields. In sect. 9 we turn to $`𝒩=2`$ supersymmetric models. We find the duality equation and demonstrate that the $`𝒩=2`$ Born-Infeld action proposed in Ref. is indeed self-dual. This action correctly reduces to the $`𝒩=1`$ Born-Infeld action when the $`(0,1/2)`$ part of the $`𝒩=2`$ vector multiplet is switched off. However, there are in fact infinitely many manifestly $`𝒩=2`$ generalization of the $`𝒩=1`$ Born-Infeld action with this property . Within the context of the D3-brane world-volume action, one has to impose additional properties (beyond self-duality), in particular the action should be invariant under translations in the transverse directions in the embedding space, or, in other words, it should contain only derivatives of the scalar fields. We show that even when allowing for nonlinear field redefinitions, the action of Ref. does not satisfy this property. It is therefore not the correct model for partial $`𝒩=4𝒩=2`$ supersymmetry breaking, based on the $`𝒩=2`$ Goldstone-Maxwell multiplet. We should mention that we know of no à priori reason why such a theory should be automatically self-dual. However this is the case for partial breaking of $`𝒩=2`$ supersymmetry to $`𝒩=1`$. In any case, the manifestly $`𝒩=2`$ supersymmetric world-volume action of a D3 brane in $`d=6`$ is still unknown (as well as the manifestly $`(1,0)`$ supersymmetric BI action in $`d=6`$, from which it might be derived via dimensional reduction).
As already mentioned, Appendix A contains the derivation of the self-duality equation in the simplest context, namely of pure nonlinear electrodynamics.
At the end of the introduction we want to mention that all our considerations are classical. The systems we study should be considered as effective theories. That they are relevant is demonstrated by the appearance of the Born-Infeld action as the world-volume action of D-branes . However the study of nonlinear self-dual systems might also be interesting in its own right.
Any nonlinear theory must possess a dimensionful parameter. Within the context of (open) string theory this is the string scale $`\alpha ^{}`$. We will always set this parameter to unity.
## 2 Self-duality in nonlinear electrodynamics
We begin with a review of self-dual models of a single U(1) gauge field with field strength $`F_{ab}=_aA_b_bA_a`$. The dynamics of such a model is determined by a nonlinear Lagrangian $`L(F_{ab})=\frac{1}{4}F^{ab}F_{ab}+𝒪(F^4)`$. With the definition<sup>1</sup><sup>1</sup>1We are working in $`d=4`$ Minkowski space, where $`\stackrel{~}{\stackrel{~}{F}}=F`$, and often use the notation $`FG=F^{ab}G_{ab}`$ implying $`F\stackrel{~}{G}=\stackrel{~}{F}G`$ and $`\stackrel{~}{F}\stackrel{~}{G}=FG`$.
$$\stackrel{~}{G}_{ab}(F)\frac{1}{2}\epsilon _{abcd}G^{cd}(F)=2\frac{L(F)}{F^{ab}},G(F)=\stackrel{~}{F}+𝒪(F^3),$$
(2.1)
the Bianchi identity and the equation of motion read
$$^b\stackrel{~}{F}_{ab}=0,^b\stackrel{~}{G}_{ab}=0.$$
(2.2)
Since these differential equations, satisfied by $`F`$, have the same form, one may consider duality transformations<sup>2</sup><sup>2</sup>2In the case of Maxwell’s electrodynamics, the field strength transforms into its Hodge dual $`\stackrel{~}{F}`$, hence the name ‘duality transformations’.
$`\left(\begin{array}{c}G^{}(F^{})\\ F^{}\end{array}\right)=\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)\left(\begin{array}{c}G(F)\\ F\end{array}\right),\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)\mathrm{GL}(2,),`$ (2.11)
such that the transformed quantities $`F^{}`$ and $`G^{}`$ also satisfy the equations (2.2). For $`G^{}`$ one should require
$$\stackrel{~}{G^{}}_{ab}(F^{})=2\frac{L^{}(F^{})}{F^{ab}},$$
(2.12)
and the transformed Lagrangian, $`L^{}(F)`$, exists (in general, $`L^{}(F)\frac{1}{4}FF+𝒪(F^4)`$) and can be determined for any $`\mathrm{GL}(2,)`$-matrix entering the transformation (2.11). In particular, for an infinitesimal duality transformation<sup>3</sup><sup>3</sup>3Throughout this paper, small Latin letters from the beginning of the alphabet denote finite duality transformation parameters, capital letters are used for infinitesimal transformations.
$`\delta \left(\begin{array}{c}G\\ F\end{array}\right)=\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)\left(\begin{array}{c}G\\ F\end{array}\right),\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)\mathrm{gl}(2,)`$ (2.21)
one finds
$$\mathrm{\Delta }L=L^{}(F)L(F)=(A+D)L(F)\frac{1}{2}D\stackrel{~}{G}F+\frac{1}{4}BF\stackrel{~}{F}\frac{1}{4}CG\stackrel{~}{G};$$
c.f. also sect. 3, eq. (3.37).
The above considerations become nontrivial if one requires the model to be self-dual, i.e.
$$L^{}(F)=L(F).$$
(2.22)
The requirement of self-duality implies:
(i) only U(1) duality rotations can be consistently defined in the nonlinear case, although Maxwell’s case is somewhat special (see sect. 3 for details)
$`\left(\begin{array}{c}G^{}(F^{})\\ F^{}\end{array}\right)=\left(\begin{array}{cc}\mathrm{cos}\lambda & \hfill \mathrm{sin}\lambda \\ \mathrm{sin}\lambda & \hfill \mathrm{cos}\lambda \end{array}\right)\left(\begin{array}{c}G(F)\\ F\end{array}\right);`$ (2.29)
(ii) the Lagrangian solves the self-duality equation
$$G^{ab}\stackrel{~}{G}_{ab}+F^{ab}\stackrel{~}{F}_{ab}=0.$$
(2.30)
A derivation of the self-duality equation is presented in Appendix A.
Due to the definition of $`G(F)`$, the self-duality equation severely constrains the possible functional form of $`L(F)`$. Any solution of the self-duality equation defines a self-dual model.
Self-dual theories possess several remarkable properties:
I. Duality-invariance of the energy-momentum tensor
Given an invariant parameter $`g`$ in the self-dual theory, the observable $`L(F,g)/g`$ is duality invariant . Indeed, using eq. (A.6) and the duality invariance of $`g`$, one gets
$$\delta \frac{}{g}L=\frac{}{g}\delta L=\frac{1}{2}\lambda \frac{}{g}\left(\stackrel{~}{G}G\right)=\frac{1}{2}\lambda \frac{}{g}\left(\stackrel{~}{G}G+\stackrel{~}{F}F\right)=0,$$
(2.31)
since $`F`$ is $`g`$-independent. Any self-dual theory can be minimally coupled to the gravitational field $`g_{mn}`$ such that the duality invariance remains intact, and $`g_{mn}`$ does not change under the curved-space duality transformations. Therefore, the energy-momentum tensor is duality invariant.
II. SL$`(2,)`$ duality invariance in the presence of dilaton and axion
Given a self-dual model $`L(F)`$, its compact U(1) duality group can be enlarged to the non-compact SL$`(2,)`$, by suitably coupling the electromagnetic field to the dilaton $`\phi `$ and axion $`a`$,
$$𝒮=𝒮_1+\mathrm{i}𝒮_2=a+\mathrm{i}\mathrm{e}^\phi .$$
(2.32)
Non-compact duality transformations read
$`\left(\begin{array}{c}G^{}\\ F^{}\end{array}\right)=\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)\left(\begin{array}{c}G\\ F\end{array}\right),𝒮^{}={\displaystyle \frac{a𝒮+b}{c𝒮+d}},\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)\mathrm{SL}(2,),`$ (2.41)
and the duality invariant Lagrangian is
$$L(F,𝒮,𝒮)=L(𝒮,𝒮)+L(\sqrt{𝒮_2}F)+\frac{1}{4}𝒮_1F\stackrel{~}{F}.$$
(2.42)
with $`L(𝒮,𝒮)`$ the SL$`(2,)`$ invariant Lagrangian for the scalar fields,
$$L(𝒮,𝒮)=\frac{\overline{𝒮}𝒮}{(𝒮\overline{𝒮})^2}.$$
(2.43)
A derivation of the self-dual model (2.42) will be described in sect. 3.
III. Self-duality under Legendre transformation
What is usually meant by ‘duality transformations’ in field theory, more precisely for models of gauge differential forms of which electrodynamics is one example, are Legendre transformations. We now show that any system which solves the self-duality equation is automatically invariant under Legendre transformation.
Let us recall the definition of Legendre transformation in the case of a generic model of nonlinear electrodynamics specified by $`L(F)`$. One associates with $`L(F)`$ an auxiliary model $`L(F,F_\mathrm{D})`$ defined by
$$L(F,F_\mathrm{D})=L(F)\frac{1}{2}F\stackrel{~}{F}_\mathrm{D},F_\mathrm{D}{}_{}{}^{ab}=^aA_\mathrm{D}{}_{}{}^{b}^bA_\mathrm{D}{}_{}{}^{a}.$$
(2.44)
$`F`$ is now an unconstrained antisymmetric tensor field, $`A_\mathrm{D}`$ a Lagrange multiplier field and $`F_\mathrm{D}`$ the dual electromagnetic field. This model is equivalent to the original one. Indeed, the equation of motion for $`A_\mathrm{D}`$ implies $`_b\stackrel{~}{F}{}_{}{}^{ab}=0`$ and therefore the second term in $`L(F,F_\mathrm{D})`$ is a total derivative, that is $`L(F,F_\mathrm{D})`$ reduces to $`L(F)`$. On the other hand, one can first consider the equation of motion for $`F`$:
$$G(F)=F_\mathrm{D}.$$
(2.45)
It is solved by expressing $`F`$ as a function of the dual field strength, $`F=F(F_\mathrm{D})`$. Inserting this solution into $`L(F,F_\mathrm{D})`$, one gets the dual model
$$L_\mathrm{D}(F_\mathrm{D})\left(L(F)\frac{1}{2}F\stackrel{~}{F}_\mathrm{D}\right)|_{F=F(F_\mathrm{D})}.$$
(2.46)
It remains to show that for any solution $`L`$ of the self-duality equation, its Legendre transform $`L_\mathrm{D}`$ satisfies:
$$L_\mathrm{D}(F)=L(F).$$
(2.47)
It follows from the results of Appendix A that the combination $`L\frac{1}{4}F\stackrel{~}{G}`$ is invariant under arbitrary duality rotations, i.e.
$$L(F)\frac{1}{4}F\stackrel{~}{G}(F)=L(F^{})\frac{1}{4}F^{}\stackrel{~}{G}^{}(F^{}).$$
(2.48)
For a finite U(1) duality rotation (2.29) by $`\lambda =\pi /2`$ this relation reads
$$L(F)\frac{1}{2}F\stackrel{~}{F}_\mathrm{D}=L(F_\mathrm{D}),F_\mathrm{D}G(F).$$
(2.49)
Comparing with (2.46) this proves (2.47).
Let us turn to a more detailed discussion of the self-duality equation (2.30). Since in four dimensions the electromagnetic field has only two independent invariants
$$\alpha =\frac{1}{4}F^{ab}F_{ab},\beta =\frac{1}{4}F^{ab}\stackrel{~}{F}_{ab},$$
(2.50)
its Lagrangian $`L(F_{ab})`$ can be considered as a real function of one complex variable
$$L(F_{ab})=L(\omega ,\overline{\omega }),\omega =\alpha +\mathrm{i}\beta .$$
(2.51)
The theory is parity invariant iff $`L(\omega ,\overline{\omega })=L(\overline{\omega },\omega )`$.
One calculates $`\stackrel{~}{G}`$ (2.1) to be
$$\stackrel{~}{G}_{ab}=\left(F_{ab}+\mathrm{i}\stackrel{~}{F}_{ab}\right)\frac{L}{\omega }+\left(F_{ab}\mathrm{i}\stackrel{~}{F}_{ab}\right)\frac{L}{\overline{\omega }},$$
(2.52)
and the self-duality equation (2.30) takes the form
$$\mathrm{Im}\left\{\omega 4\omega \left(\frac{L}{\omega }\right)^2\right\}=0.$$
(2.53)
In the literature one finds alternative forms of the self-duality equation but it is eq. (2.53) which turns out to be most convenient for supersymmetric generalizations. If one splits $`L`$ into the sum of Maxwell’s part and an interaction,
$$L=\frac{1}{2}\left(\omega +\overline{\omega }\right)+L_{\mathrm{int}},L_{\mathrm{int}}=𝒪(\omega ^2),$$
(2.54)
(2.53) becomes a condition on $`L_{\mathrm{int}}`$:
$$\mathrm{Im}\left\{\omega \frac{L_{\mathrm{int}}}{\omega }\omega \left(\frac{L_{\mathrm{int}}}{\omega }\right)^2\right\}=0.$$
(2.55)
We restrict $`L_{\mathrm{int}}`$ to a real analytic function of $`\omega `$ and $`\overline{\omega }`$. Then, every solution of eq. (2.55) is of the form<sup>4</sup><sup>4</sup>4In the Euclidean formulation of self-dual theories, it is the form (2.56) of $`L_{\mathrm{int}}`$ which allows for (anti)self-dual solutions $`\stackrel{~}{F}=\pm F`$ .
$$L_{\mathrm{int}}(\omega ,\overline{\omega })=\omega \overline{\omega }\mathrm{\Lambda }(\omega ,\overline{\omega }),\mathrm{\Lambda }=\mathrm{const}+𝒪(\omega ),$$
(2.56)
where $`\mathrm{\Lambda }`$ satisfies
$$\mathrm{Im}\left\{\frac{(\omega \mathrm{\Lambda })}{\omega }\overline{\omega }\left(\frac{(\omega \mathrm{\Lambda })}{\omega }\right)^2\right\}=0.$$
(2.57)
Note that for any solution $`L_{\mathrm{int}}(\omega ,\overline{\omega })`$ of (2.55), or any solution $`\mathrm{\Lambda }(\omega ,\overline{\omega })`$ of (2.57), the functions
$$\widehat{L}_{\mathrm{int}}(\omega ,\overline{\omega })=\frac{1}{g^2}L_{\mathrm{int}}(g^2\omega ,g^2\overline{\omega }),\widehat{\mathrm{\Lambda }}(\omega ,\overline{\omega })=g^2\mathrm{\Lambda }(g^2\omega ,g^2\overline{\omega })$$
(2.58)
are also solutions for arbitrary real parameter $`g^2`$.
In perturbation theory one looks for a parity invariant solution of the self-duality equation by considering the Ansatz
$$\mathrm{\Lambda }(\omega ,\overline{\omega })=\underset{n=0}{\overset{\mathrm{}}{}}\underset{p+q=n}{}C_{p,q}\omega ^p\overline{\omega }^q,C_{p,q}=C_{q,p},$$
(2.59)
where $`n=p+q`$ is the level of the coefficient $`C_{p,q}`$. It turns out that for odd level the self-duality equation uniquely expresses all coefficients recursively. If, however, the level is even, the self-duality equation uniquely fixes the level-$`n`$ coefficients $`C_{p,q}`$ with $`pq`$ through those at lower levels, while $`C_{r,r}`$ remain undetermined. This means that a general solution of the self-duality equation involves an arbitrary real analytic function of one real argument, $`f(\omega \overline{\omega })`$.
There are a few exact solutions of the self-duality equation known, the most prominent one being the BI Lagrangian
$`L_{\mathrm{BI}}`$ $`=`$ $`{\displaystyle \frac{1}{g^2}}\left\{\mathrm{\hspace{0.33em}1}\sqrt{det(\eta _{ab}+gF_{ab})}\right\}`$
$`=`$ $`{\displaystyle \frac{1}{g^2}}\left\{\mathrm{\hspace{0.33em}1}\sqrt{1+g^2(\omega +\overline{\omega })+{\displaystyle \frac{1}{4}}g^4(\omega \overline{\omega })^2}\right\},`$
$`\mathrm{\Lambda }_{\mathrm{BI}}`$ $`=`$ $`{\displaystyle \frac{g^2}{1+\frac{1}{2}g^2(\omega +\overline{\omega })+\sqrt{1+g^2(\omega +\overline{\omega })+\frac{1}{4}g^4(\omega \overline{\omega })^2}}},`$ (2.60)
with $`g`$ the coupling constant. In the limit $`g0`$, $`L_{\mathrm{BI}}`$ reduces to the Maxwell Lagrangian. Some other exact solutions of the self-duality equation were constructed in Ref. .
It is worth noting that the BI Lagrangian can be given in the form
$$L_{\mathrm{BI}}=\frac{1}{2}(\chi +\overline{\chi }),$$
(2.61)
where the complex field $`\chi `$ is a functions of $`\omega `$ and $`\overline{\omega }`$ which satisfies the nonlinear constraint
$$\chi +\frac{1}{2}g^2\chi \overline{\chi }\omega =0.$$
(2.62)
As will be discussed below, this form of the BI Lagrangian admits nontrivial generalizations .
We close this section with a comment. While we have limited our discussion to Lagrangians which depend on $`F`$ but not on its derivatives, the latter case can also be treated easily if one considers the action rather than the Lagrangian and if one defines
$$\stackrel{~}{G}[F]=2\frac{\delta S[F]}{\delta F},$$
etc.. This procedure is mandatory when we treat supersymmetric models.
## 3 Theory of duality invariance I: non-supersymmetric models
This section has mainly review character. We discuss the theory of duality invariance of non-supersymmetric models with Abelian gauge fields , coupled to scalar and antisymmetric tensor fields. Supersymmetric models will be treated in sects. 4-6.
### 3.1 Fundamentals
We consider a theory of $`n`$ Abelian gauge fields coupled to matter fields $`\varphi ^\mu `$. The gauge fields enter the Lagrangian only via their field strengths $`F_{ab}^i`$, where $`i=1,2,\mathrm{},n`$,
$$L=L(F_{ab}^i,\varphi ^\mu ,_a\varphi ^\mu )L(\phi ).$$
(3.1)
As in sect. 2, we introduce the dual fields
$$\stackrel{~}{G}_{ab}^i(\phi )2\frac{L(\phi )}{F^{iab}}$$
(3.2)
which arise in the equations of motion $`^b\stackrel{~}{G}_{ab}^i=0`$ for the gauge fields.
Our aim is to analyze the general conditions for the equations of motion (including the Bianchi identities) of the theory to be invariant under infinitesimal duality transformations
$`\delta \left(\begin{array}{c}G\\ F\end{array}\right)`$ $`=`$ $`\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)\left(\begin{array}{c}G\\ F\end{array}\right),\delta \varphi ^\mu =\xi ^\mu (\varphi ).`$ (3.9)
Here $`A,B,C`$ and $`D`$ are real constant $`n\times n`$ matrices, and $`\xi ^\mu `$ are some unspecified functions of the matter fields. The variation $`\delta G`$ is understood as follows
$$\delta G=G^{}(\phi ^{})G(\phi ),\stackrel{~}{G}^{}(\phi ^{})=2\frac{L(\phi ^{})}{F^{}}=2\frac{L(\phi )}{F^{}}+2\frac{}{F}\delta L,$$
(3.10)
where
$$\delta L=L(\phi ^{})L(\phi ).$$
(3.11)
Using the definitions $`F^{}=F+CG+DF`$ and $`\varphi ^{}=\varphi +\xi (\varphi )`$ of the transformed fields, one can express the derivative $`/F^{}`$ in (3.10) in terms of those w.r.t. the original fields. This gives
$$\delta \stackrel{~}{G}_{ab}^i=2\frac{}{F^{iab}}\delta LC^{jk}\stackrel{~}{G}^j\frac{G^k}{F^{iab}}D^{ji}\stackrel{~}{G}_{ab}^j,$$
(3.12)
where we have used the definition (3.2). The latter variation should coincide with $`\delta \stackrel{~}{G}`$ that follows from (3.9) and their consistency is equivalent to the relation
$`{\displaystyle \frac{}{F^{iab}}}\left[2\delta L{\displaystyle \frac{1}{2}}B^{jk}F^j\stackrel{~}{F}^k{\displaystyle \frac{1}{2}}C^{jk}G^j\stackrel{~}{G}^k\right]`$
$`=2\left(A^{ij}+D^{ji}\right){\displaystyle \frac{L}{F^{jab}}}+{\displaystyle \frac{1}{2}}\left(B^{ij}B^{ji}\right)\stackrel{~}{F}_{ab}^j+{\displaystyle \frac{1}{2}}\left(C^{kj}C^{kj}\right)\stackrel{~}{G}^j{\displaystyle \frac{G^k}{F^{iab}}}.`$ (3.13)
Here the left-hand side is a partial derivative of some function with respect to $`F`$. The right-hand side satisfies the same property iff
$$D+A^\mathrm{T}=\kappa \mathbf{\hspace{0.17em}1},B^\mathrm{T}=B,C^\mathrm{T}=C,$$
(3.14)
for some real $`\kappa `$. As a result, we find
$$\frac{}{F^i}\left[\delta L\frac{1}{4}B^{jk}F^j\stackrel{~}{F}^k\frac{1}{4}C^{jk}G^j\stackrel{~}{G}^k\kappa L\right]=0.$$
(3.15)
This relation expresses the fact that the Bianchi identities and equations of motion of the gauge fields are invariant under the duality transformation (3.9), (3.14).
Now let us turn to the transformation of the matter equation of motion:
$$E_\mu =\frac{\delta }{\delta \varphi ^\mu }S[F,\varphi ]=\left(\frac{}{\varphi ^\mu }_a\frac{}{(_a\varphi ^\mu )}\right)L.$$
(3.16)
By definition, its variation reads (it is simpler to work with the action)
$`\delta E`$ $`=`$ $`{\displaystyle \frac{\delta }{\delta \varphi ^{}}}S[F^{},\varphi ^{}]{\displaystyle \frac{\delta }{\delta \varphi }}S[F,\varphi ]`$ (3.17)
$`=`$ $`{\displaystyle \frac{\delta }{\delta \varphi ^{}}}S[F,\varphi ]+{\displaystyle \frac{\delta }{\delta \varphi }}\delta S.`$
Using $`F^{}=F+CG+DF`$ and $`\varphi ^{}=\varphi +\xi (\varphi )`$ one can express the derivative $`\delta /\delta \varphi ^{}`$ in the second line in terms of those w.r.t. the original fields. This leads to
$$\delta E_\mu =\frac{\delta }{\delta \varphi ^\mu }\left[\delta S\frac{1}{4}\mathrm{d}^4xC^{ij}\stackrel{~}{G}^iG^j\right]\frac{\xi ^\nu }{\varphi ^\mu }E_\nu .$$
(3.18)
From here it is clear that $`E_\mu `$ will transform covariantly under duality transformations,
$$\delta E_\mu =\frac{\xi ^\nu }{\varphi ^\mu }E_\nu ,$$
(3.19)
if we require
$$\frac{\delta }{\delta \varphi ^\mu }\left[\delta S\frac{1}{4}\mathrm{d}^4xC^{jk}G^j\stackrel{~}{G}^k\right]=0.$$
(3.20)
The relations (3.15) and (3.20) are compatible with each other provided $`\kappa =0`$ and hence
$$\delta L=\frac{1}{4}B^{ij}F^i\stackrel{~}{F}^j+\frac{1}{4}C^{ij}G^i\stackrel{~}{G}^j.$$
(3.21)
It is easy to check that the combination (the ‘interaction Hamiltonian’) $`L\frac{1}{4}F^i\stackrel{~}{G}^i`$ is duality invariant,
$$\delta \left(L\frac{1}{4}F^i\stackrel{~}{G}^i\right)=0.$$
(3.22)
Eq. (3.21) can be rewritten in an equivalent, but more useful, form if one directly varies $`L`$ as a function of its arguments. This leads to the self-duality equation
$$\delta _\varphi L=\frac{1}{4}B^{ij}F^i\stackrel{~}{F}^j\frac{1}{4}C^{ij}G^i\stackrel{~}{G}^j+\frac{1}{2}A^{ij}F^i\stackrel{~}{G}^j,$$
(3.23)
where
$$\delta _\varphi L=\left(\xi ^\mu \frac{}{\varphi ^\mu }+(_a\varphi ^\nu )\frac{\xi ^\mu }{\varphi ^\nu }\frac{}{(_a\varphi ^\mu )}\right)L.$$
(3.24)
Since $`\kappa =0`$, the condition (3.14) on the matrix parameters in (3.9) can be rewritten in matrix notation as
$$X^\mathrm{T}\mathrm{\Omega }+\mathrm{\Omega }X=0,$$
(3.25)
where
$`X`$ $`=`$ $`\left(\begin{array}{cc}A& B\\ C& D\end{array}\right),\mathrm{\Omega }=\left(\begin{array}{cc}0& \hfill \mathrm{𝟏}\\ \mathrm{𝟏}& \hfill 0\end{array}\right).`$ (3.30)
We conclude that Sp($`2n,`$) is the maximal group of duality transformations, although in specific models the duality group $`G`$ may actually be smaller. It should be pointed out that Sp($`2n,`$) or its non-compact subgroup $`G`$ may appear as the group of duality symmetries if the set of matter fields $`\varphi ^\mu `$ include scalar fields parameterizing the coset space $`G/H`$, with $`H`$ the maximal compact subgroup of $`G`$ (see for a more detailed discussion). Any self-dual theory without matter, $`L(F)`$, can be understood as a self-dual model with matter, $`L(F,\varphi ,\varphi )`$, with the matter fields frozen, $`\varphi (x)=\varphi _0G/H`$. The duality transformations preserving this background must thus be a subgroup of U$`(n)`$, the maximal compact subgroup of Sp($`2n,`$). If one treats the matter fields $`\varphi ^\mu `$ as coupling constants, then non-compact duality transformations relate models with different coupling constants. It is worth recalling that for the maximal compact subgroup of Sp($`2n,`$) the relations (3.25) and (3.30) should be supplemented by $`X^\mathrm{T}=X`$ and hence
$$D=A,C=B,A^\mathrm{T}=A,B^\mathrm{T}=B(B+\mathrm{i}A)^{}=(B+\mathrm{i}A).$$
(3.31)
### 3.2 U(n) duality invariant models
Let us analyze the conditions of self-duality for pure gauge theories with maximal duality group U($`n`$). Because of (3.31) and since $`\delta _\varphi L=0`$ in the absence of matter, the self-duality equation (3.23) reduces to
$$B^{ij}(F^i\stackrel{~}{F}^j+G^i\stackrel{~}{G}^j)+2A^{ij}F^i\stackrel{~}{G}^j=0.$$
Since the matrices $`A`$ and $`B`$ satisfy eq. (3.31) and otherwise arbitrary, the latter relation leads to the self-duality equations
$`G^i\stackrel{~}{G}^j+F^i\stackrel{~}{F}^j`$ $`=`$ $`0,`$ (3.32)
$`(F^i{\displaystyle \frac{}{F^j}}F^j{\displaystyle \frac{}{F^i}})L`$ $`=`$ $`0.`$ (3.33)
The first equation is a natural generalization of the self-duality equation (2.30). The second equation requires manifest SO($`n`$) invariance of the Lagrangian when $`F^i`$ transforms in the fundamental representation of SO($`n`$).
The U$`(n)`$ duality invariant models possess quite remarkable properties. In particular, they are self-dual under a Legendre transformation which acts on a single Abelian gauge field while keeping the other $`n1`$ fields invariant. The proof is similar to that given in sect. 2. Another property is that any U$`(n)`$ duality invariant model can be lifted to a model with the maximal non-compact duality symmetry Sp($`2n,`$) by coupling the gauge fields to scalar fields $`\varphi ^\mu `$ parameterizing the quotient space Sp($`2n,)/\mathrm{U}(n)`$ . The case $`n=1`$ will be discussed in the next subsection.
Nonlinear U$`(n)`$ duality invariant models with $`n>1`$ were first constructed in as a generalization of the special algebraic representation for the BI action reviewed in sect. 2. The Lagrangian reads
$$L=\frac{1}{2}\mathrm{tr}(\chi +\overline{\chi }),$$
(3.34)
where the complex $`n\times n`$ matrix $`\chi `$ is a function of $`F^i`$ which satisfies the nonlinear constraint
$$\chi ^{ij}+\frac{1}{2}\chi ^{ik}\overline{\chi }^{jk}=\omega ^{ij},\omega ^{ij}=\frac{1}{4}(F^iF^j+\mathrm{i}F^i\stackrel{~}{F}^j).$$
(3.35)
We refer the reader to for the proof of self-duality. The explicit solution of above constraint on $`\chi `$ was provided in Ref. .
One might feel uneasy with above derivation of the self-duality equations (3.32) and (3.33) in pure gauge theory $`L(F)`$ as it was essentially based on the relation (3.21) which is valid in the presence of matter. Without using the matter consistency condition (3.20) we could not have set $`\kappa =0`$ and, therefore, the variation of $`L`$ should be
$$\delta L=\frac{1}{4}B^{ij}F^i\stackrel{~}{F}^j+\frac{1}{4}C^{ij}G^i\stackrel{~}{G}^j+\kappa L.$$
(3.36)
However, practically all conclusions turn out to remain unchanged if we make use of additional physical requirements (the use of matter fields in the previous consideration simply allows to streamline the derivation). Let us consider for simplicility the case of a single gauge field, $`n=1`$. Then eq. (3.36) implies ($`\kappa =A+D`$) (c.f. eq. (2) with $`\mathrm{\Delta }L=0`$)
$$\frac{1}{4}BF\stackrel{~}{F}\frac{1}{4}CG\stackrel{~}{G}=D\frac{L}{F}F(A+D)L.$$
(3.37)
Assuming that $`L`$ is parity even, the expressions on both sides have different parities and should vanish separately
$`BF\stackrel{~}{F}CG\stackrel{~}{G}=0,`$ (3.38)
$`D{\displaystyle \frac{L}{F}}F=(A+D)L.`$ (3.39)
Let us also assume that $`L`$ reduces to Maxwell’s Lagrangian in the weak field limit, $`L=\frac{1}{4}FF+𝒪(F^4)`$, hence $`G=\stackrel{~}{F}+𝒪(F^3)`$, $`\stackrel{~}{G}=F+𝒪(F^3)`$, and therefore eq. (3.38) means
$$(B+C)F\stackrel{~}{F}+𝒪(F^4)=0.$$
(3.40)
To the lowest order, this is satisfied iff $`B=C`$. Eq. (3.39) means that $`L(F)`$ is a homogeneous function provided $`D0`$. This equation requires $`D=A`$ if $`L=\frac{1}{4}FF`$ and $`D=A=0`$ otherwise. We see that only U(1) duality rotations are possible in nonlinear electrodynamics, while in Maxwell’s theory one can also allow scale transformations. The latter are however forbidden if one requires invariance of the energy-momentum tensor under duality transformations.
### 3.3 Coupling to dilaton and axion
We are going to prove that any U(1) duality invariant model $`L(F)`$ can be uniquely coupled to the dilaton and axion such that the resulting model $`L(F,𝒮)`$ is invariant under SL$`(2,)\mathrm{Sp}(2,)`$ duality transformations . This property was stated in sect. 2.
Following the notation of subsect. 3.1, the case under consideration corresponds to $`n=1`$ and $`\varphi ^\mu =(𝒮,\overline{𝒮})`$. In accordance with eq. (2.41), the infinitesimal transformation of $`𝒮`$ reads
$$\delta 𝒮=B+2A𝒮C𝒮^2.$$
(3.41)
To describe the interaction of the dilaton and axion with the gauge field, we assume that the total Lagrangian is of the form $`L(𝒮,𝒮)+L(F,𝒮)`$ where the duality invariant kinetic term was given in (2.43). The self-duality equation (3.23) is now equivalent to the following three equations on $`L(F,𝒮)`$:
$`2𝒮{\displaystyle \frac{L}{𝒮}}+2\overline{𝒮}{\displaystyle \frac{L}{\overline{𝒮}}}`$ $`=`$ $`F{\displaystyle \frac{L}{F}},`$
$`{\displaystyle \frac{L}{𝒮}}+{\displaystyle \frac{L}{\overline{𝒮}}}`$ $`=`$ $`{\displaystyle \frac{1}{4}}F\stackrel{~}{F},`$
$`𝒮^2{\displaystyle \frac{L}{𝒮}}+\overline{𝒮}^2{\displaystyle \frac{L}{\overline{𝒮}}}`$ $`=`$ $`{\displaystyle \frac{1}{4}}G\stackrel{~}{G}.`$ (3.42)
Inspection of these equations shows that $`L(F,𝒮)`$ is
$$L(F,𝒮)=L_(\sqrt{𝒮_2}F)+\frac{1}{4}𝒮_1F\stackrel{~}{F},$$
(3.43)
where $`L(F)`$ solves the self-duality equation (2.30). Since $`L(F,𝒮)`$ is self-dual, the combination $`L\frac{1}{4}F\stackrel{~}{G}`$ is duality invariant. Its invariance under a finite duality rotation by $`\pi /2`$ is equivalent to the fact that the Legendre transform of the Lagrangian is
$$L(F,𝒮)\frac{1}{2}F\stackrel{~}{F}_\mathrm{D}=L(F_\mathrm{D},\frac{1}{𝒮}),F_\mathrm{D}G(F),$$
(3.44)
c.f. eq. (2.49).
### 3.4 Coupling to NS $`B`$-field and RR fields
Within the context of type IIB string theory, one is interested in duality-invariant couplings of the model (3.43) to the NS and RR two-forms, $`B_{ab}`$ and $`C_{ab}`$, and the RR four-form, $`C_{abcd}`$ (which are possible bosonic background fields). E.g. the self-duality of the world-volume theory of a D3-brane is inherited from the SL($`2,`$) symmetry of type IIB supergravity (see also ). These fields transform under SL$`(2,)`$ as
$`\left(\begin{array}{c}C^{}\\ B^{}\end{array}\right)`$ $`=`$ $`\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)\left(\begin{array}{c}C\\ B\end{array}\right),`$ (3.51)
$`\stackrel{~}{C_4}^{}`$ $`=`$ $`\stackrel{~}{C_4}+{\displaystyle \frac{1}{4}}bdB\stackrel{~}{B}+{\displaystyle \frac{1}{2}}bcB\stackrel{~}{C}+{\displaystyle \frac{1}{4}}acC\stackrel{~}{C}.`$ (3.52)
The transformation of $`\stackrel{~}{C_4}`$ provides a nonlinear representation of SL$`(2,)`$.<sup>5</sup><sup>5</sup>5Note that the combination $`\stackrel{~}{C}_4\frac{1}{4}C\stackrel{~}{B}`$ is SL($`2,`$) invariant. In the presence of $`B_2`$, $`C_2`$ and $`C_4`$, the Lagrangian (3.43) is extended to
$$L(F,𝒮,B,C,\stackrel{~}{C_4})=L(\sqrt{𝒮_2})+\frac{1}{4}𝒮_1\stackrel{~}{}+\stackrel{~}{C_4}\frac{1}{2}C\stackrel{~}{},$$
(3.53)
where
$$_{ab}=F_{ab}+B_{ab}.$$
(3.54)
The theory is invariant under standard gauge transformations of the gauge forms $`B_2`$, $`C_2`$ and $`C_4`$. Moreover, the theory is indeed SL$`(2,)`$ duality invariant. Given the set of matters fields $`\varphi ^\mu =(𝒮,\overline{𝒮},B_{ab},C_{ab},\stackrel{~}{C_4})`$ it is an instructive exercise to check that the self-duality equation (3.23) is satisfied.
## 4 Self-duality in $`𝒩`$ = 1 supersymmetric nonlinear electrodynamics
Gaillard and Zumino conclude their paper by posing the following problem: “When the Lagrangian is self-dual, it is natural to ask whether its supersymmetric extension possesses a self-duality property that can be formulated in a supersymmetric way.” The problem was solved in for the case when the Lagrangian is quadratic in the U(1) field strengths coupled to supersymmetric matter. The solution in the nonlinear case was obtained in for a single vector multiplet and will be extended in the sect. 6 to any number of vector multiplets coupled to scalar multiplets. In the present section we are going to review the $`𝒩=1`$ supersymmetric results of .
Let $`S[W,\overline{W}]`$ be the action generating the dynamics of a single $`𝒩=1`$ vector multiplet. The (anti) chiral superfield strengths $`\overline{W}_{\dot{\alpha }}`$ and $`W_\alpha `$,<sup>6</sup><sup>6</sup>6Our $`𝒩=1`$ conventions are those of . In particular, $`z=(x^a,\theta ^\alpha ,\overline{\theta }_{\dot{\alpha }})`$ are the coordinates of $`𝒩=1`$ superspace, $`\mathrm{d}^8z=\mathrm{d}^4x\mathrm{d}^2\theta \mathrm{d}^2\overline{\theta }`$ is the full superspace measure, and $`\mathrm{d}^6z=\mathrm{d}^4x\mathrm{d}^2\theta `$ is the measure in the chiral subspace.
$$W_\alpha =\frac{1}{4}\overline{D}^2D_\alpha V,\overline{W}_{\dot{\alpha }}=\frac{1}{4}D^2\overline{D}_{\dot{\alpha }}V,$$
(4.1)
are defined in terms of a real unconstrained prepotential $`V`$. As a consequence, the strengths are constrained superfields, that is they satisfy the Bianchi identity
$$D^\alpha W_\alpha =\overline{D}_{\dot{\alpha }}\overline{W}^{\dot{\alpha }}.$$
(4.2)
Suppose that $`S[W,\overline{W}]S[v]`$ can be unambiguously defined<sup>7</sup><sup>7</sup>7This is always possible if $`S[W,\overline{W}]`$ does not involve the combination $`D^\alpha W_\alpha `$ as an independent variable. as a functional of unconstrained (anti) chiral superfields $`\overline{W}_{\dot{\alpha }}`$ and $`W_\alpha `$. Then, one can define (anti) chiral superfields $`\overline{M}_{\dot{\alpha }}`$ and $`M_\alpha `$ as
$$\mathrm{i}M_\alpha [v]2\frac{\delta }{\delta W^\alpha }S[v],\mathrm{i}\overline{M}^{\dot{\alpha }}[v]2\frac{\delta }{\delta \overline{W}_{\dot{\alpha }}}S[v],$$
(4.3)
with the functional derivatives defined in the standard way
$`\delta S`$ $`=`$ $`{\displaystyle \mathrm{d}^6z\delta W^\alpha \frac{\delta S}{\delta W^\alpha }}+{\displaystyle \mathrm{d}^6\overline{z}\delta \overline{W}_{\dot{\alpha }}\frac{\delta S}{\delta \overline{W}_{\dot{\alpha }}}},`$
$`{\displaystyle \frac{\delta }{\delta W^\alpha (z)}}W^\beta (z^{})`$ $`=`$ $`\delta _\alpha {}_{}{}^{\beta }({\displaystyle \frac{1}{4}}\overline{D}^2)\delta ^4(xx^{})\delta ^2(\theta \theta ^{})\delta ^2(\overline{\theta }\overline{\theta }^{}).`$ (4.4)
The vector multiplet equation of motion following from the action $`S[W,\overline{W}]`$ reads
$$D^\alpha M_\alpha =\overline{D}_{\dot{\alpha }}\overline{M}^{\dot{\alpha }}.$$
(4.5)
Since the Bianchi identity (4.2) and the equation of motion (4.5) have the same functional form, one may consider, similar to the non-supersymmetric case, U(1) duality rotations
$`\left(\begin{array}{c}M_\alpha ^{}[v^{}]\\ W_\alpha ^{}\end{array}\right)=\left(\begin{array}{cc}\mathrm{cos}\lambda & \hfill \mathrm{sin}\lambda \\ \mathrm{sin}\lambda & \hfill \mathrm{cos}\lambda \end{array}\right)\left(\begin{array}{c}M_\alpha [v]\\ W_\alpha \end{array}\right),`$ (4.12)
where $`M^{}`$ should be
$$\mathrm{i}M_\alpha ^{}[v^{}]=2\frac{\delta }{\delta W^\alpha }S[v^{}].$$
(4.13)
In order for such duality transformations to be consistently defined, the action $`S[W,\overline{W}]`$ must satisfy a generalization of the self-duality equation (2.30). Its derivation follows essentially the same steps as described in Appendix A, but with a proper replacement of partial derivatives by functional derivatives. To preserve the definition (4.3) of $`M_\alpha `$ and its conjugate, the action should transform under an infinitesimal duality rotation as
$$\delta S=S[v^{}]S[v]=\frac{\mathrm{i}}{4}\lambda \mathrm{d}^6z\left\{M^\alpha M_\alpha W^\alpha W_\alpha \right\}+\mathrm{c}.\mathrm{c}.$$
(4.14)
On the other hand, $`S`$ is a functional of $`W_\alpha `$ and $`\overline{W}_{\dot{\alpha }}`$ only, and therefore its variation is
$$\delta S=\frac{\mathrm{i}}{2}\lambda \mathrm{d}^6zM^\alpha M_\alpha +\mathrm{c}.\mathrm{c}.$$
(4.15)
Since these two variations must coincide, we arrive at the following reality condition
$$\mathrm{Im}\mathrm{d}^6z\left(W^\alpha W_\alpha +M^\alpha M_\alpha \right)=0.$$
(4.16)
In eq. (4.16), the superfield $`M_\alpha `$ was defined in (4.3), and $`W_\alpha `$ should be considered as an unconstrained chiral superfields. Eq. (4.16) is the condition for the $`𝒩=1`$ supersymmetric theory to be self-dual. We call it the $`𝒩=1`$ self-duality equation.
With proper modifications, the properties of self-dual theories, which we described in sect. 2, also hold for $`𝒩=1`$ self-dual models. In particular, the derivative of the self-dual action with respect to an invariant parameter is always duality invariant. This implies duality invariance of the $`𝒩=1`$ supercurrent, i.e. the multiplet of the energy-momentum tensor (see for a review). Duality invariant couplings to the dilaton-axion multiplet will be discussed in sect. 6. Here we would like to concentrate on self-duality under $`𝒩=1`$ Legendre transformation, defined as follows. Given a vector multiplet model $`S[W,\overline{W}]`$, we introduce the auxiliary action
$$S[W,\overline{W},W_\mathrm{D},\overline{W}_\mathrm{D}]=S[W,\overline{W}]\frac{\mathrm{i}}{2}\mathrm{d}^6zW^\alpha W_{\mathrm{D}\alpha }+\frac{\mathrm{i}}{2}\mathrm{d}^6\overline{z}\overline{W}_{\dot{\alpha }}\overline{W}_\mathrm{D}{}_{}{}^{\dot{\alpha }},$$
(4.17)
where $`W_\alpha `$ is now an unconstrained chiral spinor superfield, and $`W_{\mathrm{D}\alpha }`$ the dual field strength
$$W_{\mathrm{D}\alpha }=\frac{1}{4}\overline{D}^2D_\alpha V_\mathrm{D},\overline{W}_{\mathrm{D}\dot{\alpha }}=\frac{1}{4}D^2\overline{D}_{\dot{\alpha }}V_\mathrm{D}.$$
(4.18)
This model is equivalent to the original model, since the equation of motion for $`W_\mathrm{D}`$ implies that $`W`$ satisfies the Bianchi identity (4.2), and the action (4.17) reduces to $`S[W,\overline{W}]`$. On the other hand, the equation of motion for $`W`$ is $`M[W,\overline{W}]=W_\mathrm{D}`$, with $`M`$ defined in (4.3). Solving this equation, $`W=W[W_\mathrm{D},\overline{W}_\mathrm{D}]`$, and inserting the solution back into the action (4.17), one gets the dual model $`S_\mathrm{D}[W_\mathrm{D},\overline{W}_\mathrm{D}]`$ or, what is the same, the Legendre transform of $`S[W,\overline{W}]`$. For all $`𝒩=1`$ self-dual theories, $`S_\mathrm{D}=S`$. This follows from the fact that the combination
$$S\frac{\mathrm{i}}{4}\mathrm{d}^6zW^\alpha M_\alpha +\frac{\mathrm{i}}{4}\mathrm{d}^6\overline{z}\overline{W}_{\dot{\alpha }}\overline{M}^{\dot{\alpha }}$$
(4.19)
is invariant under arbitrary U(1) duality rotations.
We now present a family of $`𝒩=1`$ supersymmetric self-dual models with actions of the general form
$$S=\frac{1}{4}\mathrm{d}^6zW^2+\frac{1}{4}\mathrm{d}^6\overline{z}\overline{W}^2+\frac{1}{4}\mathrm{d}^8zW^2\overline{W}^2\mathrm{\Lambda }(\frac{1}{8}D^2W^2,\frac{1}{8}\overline{D}^2\overline{W}^2),$$
(4.20)
where $`\mathrm{\Lambda }(u,\overline{u})`$ is a real analytic function of the complex variable
$$u\frac{1}{8}D^2W^2.$$
(4.21)
Functionals of this type naturally appear as low-energy effective actions in quantum supersymmetric gauge theories; by ‘low-energy action’ we mean here the part of the full effective action independent of the derivatives of the U(1) field strength $`F`$. In fact, the low-energy effective actions usually have the more general form (see, for instance, ):
$$S_{\mathrm{eff}}=\frac{1}{4}\mathrm{d}^6zW^2+\frac{1}{4}\mathrm{d}^6\overline{z}\overline{W}^2+\mathrm{d}^8zW^2\overline{W}^2\mathrm{\Omega }(D^2W^2,\overline{D}^2\overline{W}^2,D^\alpha W_\alpha ).$$
(4.22)
However, the combination $`D^\alpha W_\alpha `$ is nothing but the free equation of motion of the $`𝒩=1`$ vector multiplet. Contributions to effective action, which contain factors of the classical equations of motion, are ambiguous. They are often ignored. It is worth pointing out that there is no unique way to define the action (4.22) as a functional of unconstrained chiral superfield $`W_\alpha `$ and its conjugate (what is required in the framework of our approach to supersymmetric self-dual theories) when $`\mathrm{\Omega }`$ depends on $`D^\alpha W_\alpha =\overline{D}_{\dot{\alpha }}\overline{W}^{\dot{\alpha }}`$.
Let us analyze the conditions for the model (4.20) to be self-dual. One finds
$$\mathrm{i}M_\alpha =W_\alpha \left\{\mathrm{\hspace{0.33em}1}\frac{1}{4}\overline{D}^2\left[\overline{W}^2\left(\mathrm{\Lambda }+\frac{1}{8}D^2\left(W^2\frac{\mathrm{\Lambda }}{u}\right)\right)\right]\right\}.$$
(4.23)
Then, eq. (4.16) leads to
$$\mathrm{Im}\mathrm{d}^8zW^2\overline{W}^2\left(\mathrm{\Gamma }\overline{u}\mathrm{\Gamma }^2\right)=0,$$
(4.24)
where
$$\mathrm{\Gamma }\mathrm{\Lambda }+\frac{1}{8}(D^2W^2)\frac{\mathrm{\Lambda }}{u}=\frac{(u\mathrm{\Lambda })}{u}.$$
(4.25)
In deriving eq. (4.24) we have used the following property of the $`𝒩=1`$ vector multiplet:
$$W_\alpha W_\beta W_\gamma =0.$$
(4.26)
Since the functional relation (4.24) must be satisfied for arbitrary (anti) chiral superfields $`\overline{W}_{\dot{\alpha }}`$ and $`W_\alpha `$, we arrive at the following differential equation for $`\mathrm{\Lambda }(u,\overline{u})`$:
$$\mathrm{Im}\left\{\frac{(u\mathrm{\Lambda })}{u}\overline{u}\left(\frac{(u\mathrm{\Lambda })}{u}\right)^2\right\}=0.$$
(4.27)
This equation is identical to the self-duality equation (2.57).
To obtain the component form of (4.20), one applies the reduction rules
$$\mathrm{d}^8zU=\frac{1}{16}\mathrm{d}^4xD^2\overline{D}^2U|_{\theta =0},\mathrm{d}^6zU_\mathrm{c}=\frac{1}{4}\mathrm{d}^4xD^2U_\mathrm{c}|_{\theta =0}.$$
(4.28)
We also introduce the component fields of the $`𝒩=1`$ vector multiplet, $`\{\lambda _\alpha ,\overline{\lambda }_{\dot{\alpha }},F_{ab},D\}`$, in the standard way :
$`\lambda _\alpha (x)`$ $`=`$ $`W_\alpha |_{\theta =0},`$
$`F_{\alpha \beta }(x)`$ $`=`$ $`{\displaystyle \frac{\mathrm{i}}{4}}(D_\alpha W_\beta +D_\beta W_\alpha )|_{\theta =0},`$
$`D(x)`$ $`=`$ $`{\displaystyle \frac{1}{2}}D^\alpha W_\alpha |_{\theta =0},`$ (4.29)
with
$$F_{\alpha \dot{\alpha }\beta \dot{\beta }}(\sigma ^a)_{\alpha \dot{\alpha }}(\sigma ^a)_{\beta \dot{\beta }}F_{ab}=2\epsilon _{\alpha \beta }\overline{F}_{\dot{\alpha }\dot{\beta }}+2\epsilon _{\dot{\alpha }\dot{\beta }}F_{\alpha \beta }.$$
(4.30)
Here we are interested only in the bosonic sector of the model and therefore set $`\lambda _\alpha =0`$ in what follows. Under this assumption one can readily compute the component Lagrangian
$$L(F_{ab},D)=\frac{1}{2}(𝐮+\overline{𝐮})+𝐮\overline{𝐮}\mathrm{\Lambda }(𝐮,\overline{𝐮}),𝐮\frac{1}{8}D^2W^2|_{\theta =0}=\omega \frac{1}{2}D^2,$$
(4.31)
with $`\omega `$ defined in eq. (2.51). Since only even powers of the auxiliary field $`D`$ appear in $`L`$, its equation of motion has the solution $`D=0`$. If we take this solution, the duality equation (4.27) implies that the non-supersymmetric model $`L(F)=L(F,D=0)`$ is self-dual.
We arrive at the conclusion: every non-supersymmetric self-dual model of the type considered in sect. 2 admits an $`𝒩=1`$ supersymmetric extension which is self-dual under manifestly supersymmetric duality rotations. The procedure of constructing such a supersymmetric extension is constructive: given a self-dual Lagrangian $`L(F)`$, one should first derive $`\mathrm{\Lambda }(\omega ,\overline{\omega })`$ defined by eqs. (2.54) and (2.56), and then use this function to generate the action (4.20).
## 5 Properties of the $`𝒩`$ = 1 supersymmetric BI action
We use the results of sect. 3 to obtain the unique $`𝒩=1`$ supersymmetric self-dual extension of the BI theory (2.60). With the use of $`\mathrm{\Lambda }_{\mathrm{BI}}`$ one immediately gets
$`S_{\mathrm{BI}}`$ $`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle \mathrm{d}^6zW^2}+{\displaystyle \frac{1}{4}}{\displaystyle \mathrm{d}^6\overline{z}\overline{W}^2}+{\displaystyle \frac{g^2}{4}}{\displaystyle \mathrm{d}^8z\frac{W^2\overline{W}^2}{1+\frac{1}{2}A+\sqrt{1+A+\frac{1}{4}B^2}}},`$
$`A`$ $`=`$ $`{\displaystyle \frac{g^2}{8}}\left(D^2W^2+\overline{D}^2\overline{W}^2\right),B={\displaystyle \frac{g^2}{8}}\left(D^2W^2\overline{D}^2\overline{W}^2\right).`$ (5.1)
In what follows, for convenience we fix the coupling constant to $`g^2=4`$.
The above action was first introduced in as a super extension of the BI theory. However, only much later it was realized that the theory encodes a remarkably reach structure. Bagger and Galperin , and later Roček and Tseytlin discovered that (5.1) is the action for a Goldstone multiplet associated with $`𝒩=2𝒩=1`$ partial supersymmetry breaking. Using a reformulation of (5.1) with auxiliary superfields, Brace, Morariu and Zumino demonstrated that the theory is invariant under U(1) duality rotations. The latter property has turned out to be a simple consequence of the approach developed in and reviewed in the previous section. Below we give a concise review of the results of on partial $`𝒩=2𝒩=1`$ supersymmetry breaking.
Bagger and Galperin noticed that the Cecotti-Ferrara action (5.1) can be represented in the form
$$S=\frac{1}{4}\mathrm{d}^6zX+\frac{1}{4}\mathrm{d}^6\overline{z}\overline{X},$$
(5.2)
where the chiral superfield $`X`$ is a functional of $`W`$ and $`\overline{W}`$ such that it satisfies the nonlinear constraint
$$X+\frac{1}{4}X\overline{D}^2\overline{X}=W^2.$$
(5.3)
Indeed, using the action rule
$$\mathrm{d}^8zU=\frac{1}{4}\mathrm{d}^6z\overline{D}^2U$$
(5.4)
and the constraint (5.3), one can rewrite (5.2) in the form
$$S=\frac{1}{4}\mathrm{d}^6zW^2+\frac{1}{4}\mathrm{d}^6\overline{z}\overline{W}^2+\frac{1}{2}\mathrm{d}^8zX\overline{X}.$$
(5.5)
Using the constraint (5.3) once more, we can represent $`X\overline{X}`$ as
$$X\overline{X}=\frac{W^2\overline{W}^2}{(1+\frac{1}{4}\overline{D}^2\overline{X})(1+\frac{1}{4}D^2X)}.$$
(5.6)
Since $`W^3=0`$, on the right-hand side we can safely take $`D^2X`$ in an effective form $`D^2X_{\mathrm{eff}}`$ determined by the equation
$$D^2X_{\mathrm{eff}}=\frac{D^2W^2}{1+\frac{1}{4}\overline{D}^2\overline{X}_{\mathrm{eff}}}.$$
(5.7)
Using this in (5.5) one reproduces (5.1).
The dynamical system defined by eqs. (5.2) and (5.3) is manifestly $`𝒩=1`$ supersymmetric. Remarkably, it turns out to be invariant under a second, nonlinearly realized, supersymmetry transformation
$`\delta X`$ $`=`$ $`2ϵ^\alpha W_\alpha ,`$ (5.8)
$`\delta W_\alpha `$ $`=`$ $`ϵ_\alpha +{\displaystyle \frac{1}{4}}\overline{D}^2\overline{X}ϵ_\alpha +\mathrm{i}_{\alpha \dot{\alpha }}X\overline{ϵ}^{\dot{\alpha }},`$ (5.9)
with $`ϵ_\alpha `$ a constant parameter. Such transformations commute with the first, linearly realized, supersymmetry, and altogether they generate the $`𝒩=2`$ algebra without central charge. There is a simple way to derive the supersymmetry transformations (5.8) and (5.9). One first observes that the variation (5.8) leaves the action (5.2) invariant, as a consequence of the explicit form of the field strength $`W_\alpha `$, see eq. (4.1). Due to (5.3), the variation $`\delta X`$ must be induced by a variation of $`W_\alpha `$ of the form
$$\delta W_\alpha =ϵ_\alpha +\frac{1}{4}\overline{D}^2\overline{X}ϵ_\alpha +\widehat{\delta }W_\alpha ,$$
(5.10)
where $`\widehat{\delta }W`$ should satisfy
$$W^\alpha \widehat{\delta }W_\alpha =\frac{1}{4}X\overline{D}^2\overline{W}_{\dot{\alpha }}\overline{ϵ}^{\dot{\alpha }}=\mathrm{i}X_{\alpha \dot{\alpha }}W^\alpha \overline{ϵ}^{\dot{\alpha }}.$$
(5.11)
Since
$$W^\alpha X=0,$$
(5.12)
the latter relation can be rewritten as follows
$$W^\alpha \widehat{\delta }W_\alpha =\mathrm{i}W^\alpha _{\alpha \dot{\alpha }}X\overline{ϵ}^{\dot{\alpha }},$$
(5.13)
and we thus arrive at the variation (5.9). But this is not yet the end of the story, since one still has to check that the variation (5.9) is consistent with the Bianchi identity (4.2). Indeed it is. However, in sect. 9 we will see that the above procedure cannot be directly generalized to the case of $`𝒩=2`$ supersymmetry.
In Bagger and Galperin proved that the action (5.1) is self-dual under the $`𝒩=1`$ Legendre transformation. Their proof is ingenious but rather involved. The results of sect. 4 make this property obvious. The $`𝒩=1`$ super BI theory (5.1) is invariant under U(1) duality rotations, and therefore it is automatically self-dual under the $`𝒩=1`$ Legendre transformation.
## 6 Theory of self-duality II: $`𝒩`$ = 1 supersymmetric <br>models
In this section we develop a general formalism of duality invariance for $`𝒩=1`$ supersymmetric theories of $`n`$ Abelian vector multiplets, described by chiral spinor strengths $`W_\alpha ^i`$ and their conjugates $`\overline{W}_{\dot{\alpha }}^i`$, in the presence of supersymmetric matter – chiral superfields $`\mathrm{\Phi }^\mu `$ and their conjugates $`\overline{\mathrm{\Phi }}^\mu `$. We will use the condensed notation $`S[v]=S[W^{\alpha i},\overline{W}_{\dot{\alpha }}^i,\mathrm{\Phi }^\mu ,\overline{\mathrm{\Phi }}^\mu ]`$ for the action functional and, as in sect. 4, introduce (anti) chiral superfields $`\overline{M}^{\dot{\alpha }i}`$ and $`M_\alpha ^i`$ dual to $`\overline{W}_{\dot{\alpha }}^i`$ and $`W^{\alpha i}`$:
$$\mathrm{i}M_\alpha ^i[v]2\frac{\delta }{\delta W^{\alpha i}}S[v],\mathrm{i}\overline{M}^{\dot{\alpha }i}[v]2\frac{\delta }{\delta \overline{W}_{\dot{\alpha }}^i}S[v].$$
(6.1)
To simplify notation, we introduce
$$M^iM^j=\mathrm{d}^6zM^{\alpha i}M_\alpha ^j,\overline{M}^i\overline{M}^j=\mathrm{d}^6\overline{z}\overline{M}_{\dot{\alpha }}^i\overline{M}^{\dot{\alpha }j}$$
(6.2)
and similarly for superspace contractions of (anti) chiral scalar superfields.
### 6.1 General analysis
We are interested in determining the conditions for the theory to be self-dual under chiral superfield duality transformations
$`\delta \left(\begin{array}{c}M\\ W\end{array}\right)`$ $`=`$ $`\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)\left(\begin{array}{c}M\\ W\end{array}\right),\delta \mathrm{\Phi }^\mu =\xi ^\mu (\mathrm{\Phi }^\nu ),`$ (6.9)
with $`\xi ^\mu `$ a holomorphic functions of the chiral matter fields. Here $`A,B,C`$ and $`D`$ are constant real $`n\times n`$ matrices; these matrices have to be real, since the Bianchi identities $`D^\alpha W_\alpha ^i=\overline{D}_{\dot{\alpha }}\overline{W}^{\dot{\alpha }i}`$ and the equations of motion $`D^\alpha M_\alpha ^i=\overline{D}_{\dot{\alpha }}\overline{M}^{\dot{\alpha }i}`$ are special reality conditions.
By self-duality we understand the following:
I. We require
$$\mathrm{i}M^{}[v^{}]=2\frac{\delta }{\delta W^{}}S[v^{}]=2\frac{\delta }{\delta W^{}}S[v]+2\frac{\delta }{\delta W}\delta S,$$
(6.10)
where $`\delta S=S[v^{}]S[v]`$.
II. The $`\mathrm{\Phi }`$-equation of motion
$$E_\mu [v]=\frac{\delta }{\delta \mathrm{\Phi }^\mu }S[v]$$
(6.11)
transforms covariantly under duality transformations
$$\delta E_\mu =\frac{\xi ^\nu (\mathrm{\Phi })}{\mathrm{\Phi }^\mu }E_\nu ,$$
(6.12)
where
$$\delta E=E^{}[v^{}]E[v],E^{}[v^{}]=\frac{\delta }{\delta \mathrm{\Phi }^{}}S[v^{}]=\frac{\delta }{\delta \mathrm{\Phi }^{}}S[v]+\frac{\delta }{\delta \mathrm{\Phi }}\delta S.$$
(6.13)
Analysis of the self-duality conditions is similar to the non-supersymmetric case described in sect. 3. The transformation law (6.9) and condition I are consistent provided
$`{\displaystyle \frac{\delta }{\delta W^{\alpha i}}}\left[\delta S{\displaystyle \frac{\mathrm{i}}{4}}B^{jk}\left(W^jW^k\overline{W}^j\overline{W}^k\right){\displaystyle \frac{\mathrm{i}}{4}}C^{jk}\left(M^jM^k\overline{M}^j\overline{M}^k\right)\right]`$ (6.14)
$`=`$ $`+{\displaystyle \frac{\mathrm{i}}{4}}(C^{jk}C^{kj})\left(({\displaystyle \frac{\delta }{\delta W^{\alpha i}}}M^k)M^j({\displaystyle \frac{\delta }{\delta W^{\alpha i}}}\overline{M}^k)\overline{M}^j\right)`$
$`+{\displaystyle \frac{\mathrm{i}}{4}}(B^{ij}B^{ji})W_\alpha ^j+(D^{ji}+A^{ij}){\displaystyle \frac{\delta }{\delta W^{\alpha j}}}S[v].`$
Since the left-hand side is a total variational derivative, the matrices $`A,B,C`$ and $`D`$ should be constrained as in eq. (3.14). Then, the above relation turns into
$`{\displaystyle \frac{\delta }{\delta W^{\alpha i}}}[\delta S{\displaystyle \frac{\mathrm{i}}{4}}B^{jk}(W^jW^k\overline{W}^j\overline{W}^k)`$
$`{\displaystyle \frac{\mathrm{i}}{4}}C^{jk}(M^jM^k\overline{M}^j\overline{M}^k)\kappa S[v]]=0.`$ (6.15)
Furthermore, the $`\mathrm{\Phi }`$-equation of motion can be shown to change under duality transformations as
$`\delta E_\mu ={\displaystyle \frac{\xi ^\nu }{\mathrm{\Phi }^\mu }}E_\nu `$ (6.16)
$`+{\displaystyle \frac{\delta }{\delta \mathrm{\Phi }^\mu }}\left[\delta S{\displaystyle \frac{\mathrm{i}}{4}}B^{jk}\left(W^jW^k\overline{W}^j\overline{W}^k\right){\displaystyle \frac{\mathrm{i}}{4}}C^{jk}\left(M^jM^k\overline{M}^j\overline{M}^k\right)\right].`$
Consequently, condition II is satisfied if we impose the condition
$$\frac{\delta }{\delta \mathrm{\Phi }^\mu }\left[\delta S\frac{\mathrm{i}}{4}B^{jk}\left(W^jW^k\overline{W}^j\overline{W}^k\right)\frac{\mathrm{i}}{4}C^{jk}\left(M^jM^k\overline{M}^j\overline{M}^k\right)\right]=0.$$
(6.17)
The latter is consistent with (6.15) provided $`\kappa =0`$. Therefore, Sp($`2n,`$) is the maximal duality group (see sect. 3), and the action transforms as
$`\delta S`$ $`=`$ $`{\displaystyle \frac{\mathrm{i}}{4}}\delta \left(W^iM^i\overline{M}^i\overline{W}^i\right)`$ (6.18)
$`=`$ $`{\displaystyle \frac{\mathrm{i}}{4}}B^{ij}\left(W^iW^j\overline{W}^i\overline{W}^j\right)+{\displaystyle \frac{\mathrm{i}}{4}}C^{ij}\left(M^iM^j\overline{M}^i\overline{M}^j\right).`$
Equation (6.18) contains nontrivial information. The point is that the action can be varied directly,
$`\delta S`$ $`=`$ $`S[v^{}]S[v]`$ (6.19)
$`=`$ $`{\displaystyle \frac{\mathrm{i}}{2}}\left(\delta W^iM^i\delta \overline{W}^i\overline{W}^i\right)+\delta \mathrm{\Phi }^\mu {\displaystyle \frac{\delta S}{\delta \mathrm{\Phi }^\mu }}+\delta \overline{\mathrm{\Phi }}^\mu {\displaystyle \frac{\delta S}{\delta \overline{\mathrm{\Phi }}^\mu }},`$
and the two results should coincide. This gives
$`\delta \mathrm{\Phi }^\mu {\displaystyle \frac{\delta S}{\delta \mathrm{\Phi }^\mu }}`$ $`+`$ $`\delta \overline{\mathrm{\Phi }}^\mu {\displaystyle \frac{\delta S}{\delta \overline{\mathrm{\Phi }}^\mu }}`$ (6.20)
$`=`$ $`{\displaystyle \frac{\mathrm{i}}{4}}B^{ij}\left(W^iW^j\overline{W}^i\overline{W}^j\right){\displaystyle \frac{\mathrm{i}}{4}}C^{ij}\left(M^iM^j\overline{M}^i\overline{M}^j\right)`$
$`+`$ $`{\displaystyle \frac{\mathrm{i}}{2}}A^{ij}\left(W^iM^j\overline{W}^i\overline{M}^j\right).`$
This is the self-duality equation in the presence of matter.
In the absence of matter, the maximal duality group is U($`n`$) and the transformation parameters in (6.20) are constrained by $`B=C=B^\mathrm{T}`$, $`A^\mathrm{T}=A`$. If the duality group is U($`n`$), then eq. (6.20) leads to the following self-duality equations
$`\mathrm{Im}\left(W^iW^j+M^iM^j\right)`$ $`=`$ $`0,`$ (6.21)
$`\mathrm{Im}\left(W^iM^jW^jM^i\right)`$ $`=`$ $`0.`$ (6.22)
Eq. (6.22) requires the theory to be invariant under SO$`(n)`$ which acts linearly on $`W^i`$. For $`n=1`$, eq. (6.21) reduces to (4.16).
Similar to the non-supersymmetric case , a U$`(n)`$ duality invariant theory of $`n`$ Abelian vector multiplets can be lifted to an Sp($`2n,`$) duality invariant model by coupling the vector multiplets to scalar multiplets $`\mathrm{\Phi }^\mu `$ parameterizing the quotient space Sp($`2n,)/\mathrm{U}(n)`$. Below we give a proof for $`n=1`$.
### 6.2 Coupling to the dilaton-axion multiplet
Our aim here is to couple the system (4.20), (4.27) to the dilaton-axion multiplet $`\mathrm{\Phi }`$ such that the resulting model be $`\mathrm{SL}(2,)`$ duality invariant. The $`\mathrm{SL}(2,)`$-transformation of $`\mathrm{\Phi }`$ coincides with the $`𝒮`$-transformation (2.41). Its infinitesimal form is
$$\delta \mathrm{\Phi }=B+2A\mathrm{\Phi }C\mathrm{\Phi }^2.$$
(6.23)
The self-duality equation (6.20) is now equivalent to the following requirements on the action functional $`S=S[W,\mathrm{\Phi }]`$:
$`2\mathrm{\Phi }{\displaystyle \frac{\delta S}{\delta \mathrm{\Phi }}}+2\overline{\mathrm{\Phi }}{\displaystyle \frac{\delta S}{\delta \overline{\mathrm{\Phi }}}}`$ $`=`$ $`W{\displaystyle \frac{\delta S}{\delta W}}+\overline{W}{\displaystyle \frac{\delta S}{\delta \overline{W}}},`$ (6.24)
$`{\displaystyle \frac{\delta S}{\delta \mathrm{\Phi }}}1+{\displaystyle \frac{\delta S}{\delta \overline{\mathrm{\Phi }}}}1`$ $`=`$ $`{\displaystyle \frac{\mathrm{i}}{4}}\left(WW\overline{W}\overline{W}\right),`$ (6.25)
$`\mathrm{\Phi }^2{\displaystyle \frac{\delta S}{\delta \mathrm{\Phi }}}+\overline{\mathrm{\Phi }}^2{\displaystyle \frac{\delta S}{\delta \overline{\mathrm{\Phi }}}}`$ $`=`$ $`{\displaystyle \frac{\mathrm{i}}{4}}\left(MM\overline{M}\overline{M}\right).`$ (6.26)
We are interested in a solution of these equations which for $`\mathrm{\Phi }=\mathrm{i}`$ reduces to the self-dual system given by eqs. (4.20) and (4.27). A direct analysis of the self-duality equations gives the solution
$`S[W,\mathrm{\Phi }]`$ $`=`$ $`{\displaystyle \frac{\mathrm{i}}{4}}{\displaystyle \mathrm{d}^6z\mathrm{\Phi }W^2}{\displaystyle \frac{\mathrm{i}}{4}}{\displaystyle \mathrm{d}^6\overline{z}\overline{\mathrm{\Phi }}\overline{W}^2}`$
$``$ $`{\displaystyle \frac{1}{16}}{\displaystyle \mathrm{d}^8z(\mathrm{\Phi }\overline{\mathrm{\Phi }})^2W^2\overline{W}^2\mathrm{\Lambda }(\frac{\mathrm{i}}{16}(\mathrm{\Phi }\overline{\mathrm{\Phi }})D^2W^2,\frac{\mathrm{i}}{16}(\mathrm{\Phi }\overline{\mathrm{\Phi }})\overline{D}^2\overline{W}^2)}.`$
To this action one can add the dilaton-axion kinetic term $`\mathrm{d}^8zK(\mathrm{\Phi },\overline{\mathrm{\Phi }})`$, with $`K(\mathrm{\Phi },\overline{\mathrm{\Phi }})`$ the Kähler potential of the Kähler manifold SL($`2,)/\mathrm{U}(1)`$. It is worth pointing out that the dilaton and axion (2.32) are related to $`\mathrm{\Phi }`$ by the rule $`\overline{𝒮}=\mathrm{\Phi }|_{\theta =0}`$. For the $`𝒩=1`$ super BI action (5.1), the coupling to the dilaton-axion multiplet was described in .
### 6.3 Coupling to NS and RR supermultiplets
The model (6.2) can be generalized by coupling it to supermultiplets containing the NS and RR two-forms, $`B_2`$ and $`C_2`$, and the RR four-form, $`C_4`$. The extended action is
$$S[W,\mathrm{\Phi },\beta ,\gamma ,\mathrm{\Omega }]=S[𝒲,\mathrm{\Phi }]+\{\mathrm{d}^6z(\mathrm{\Omega }+\frac{1}{2}\gamma ^\alpha 𝒲_\alpha )+\mathrm{c}.\mathrm{c}.\},$$
(6.28)
where
$$𝒲_\alpha =W_\alpha +\mathrm{i}\beta _\alpha .$$
(6.29)
is the supersymmetrization of $`F+B`$. Here $`\beta _\alpha `$, $`\gamma _\alpha `$ and $`\mathrm{\Omega }`$ are unconstrained chiral superfields which include, among their components, the fields $`B_2`$, $`C_2`$ and $`C_4`$, respectively. The action is invariant under the following gauge transformations
$`\delta \beta _\alpha =\mathrm{i}\delta W_\alpha =\mathrm{i}\overline{D}^2D_\alpha K_1,`$ (6.30)
$`\delta \gamma _\alpha =\mathrm{i}\overline{D}^2D_\alpha K_2,\delta \mathrm{\Omega }={\displaystyle \frac{\mathrm{i}}{2}}𝒲^\alpha \overline{D}^2D_\alpha K_2,`$ (6.31)
$`\delta \mathrm{\Omega }=\mathrm{i}\overline{D}^2K_3,`$ (6.32)
with $`K_i`$ real unconstrained superfields. Note that $`𝒲_\alpha `$ is invariant under (6.30). The transformations of $`\beta `$ and $`\gamma `$ imply that these superfields describe two tensor multiplets; c.f. also sect. 7. Eq. (6.32) implies that all components of $`\mathrm{\Omega }`$ but $`\mathrm{Re}D^2\mathrm{\Omega }|_{\theta =0}`$ can be algebraically gauged away; the remaining component transforms as a four-form and is identified with $`\stackrel{~}{C_4}`$.
The theory (6.28) is $`\mathrm{SL}(2,)`$ duality invariant provided the superfields $`\beta _\alpha `$, $`\gamma _\alpha `$ and $`\mathrm{\Omega }`$ transform as
$$\left(\begin{array}{c}\gamma ^{}\\ \beta ^{}\end{array}\right)=\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)\left(\begin{array}{c}\gamma \\ \beta \end{array}\right),\mathrm{\Omega }^{}=\mathrm{\Omega }\frac{\mathrm{i}}{4}bd\beta ^2\frac{\mathrm{i}}{2}bc\beta \gamma \frac{\mathrm{i}}{4}ac\gamma ^2.$$
(6.33)
One can check that the self-duality equation (6.20) is satisfied, with $`\mathrm{\Phi }^\mu =(\mathrm{\Phi },\beta _\alpha ,\gamma _\alpha ,\mathrm{\Omega })`$ the set of matter chiral superfields.
### 6.4 Example of U(n) self-dual supersymmetric theory
To conclude, we give an example of U$`(n)`$ duality invariant model describing the dynamics of $`n`$ interacting Abelian vector multiplets $`W_\alpha ^i`$. The action is
$$S=\frac{1}{4}\mathrm{d}^6z\mathrm{tr}X+\frac{1}{4}\mathrm{d}^6\overline{z}\mathrm{tr}\overline{X},$$
(6.34)
where the chiral matrix superfield $`X`$ is a functional of $`W_\alpha ^i`$ and $`\overline{W}_{\dot{\alpha }}^i`$ such that it satisfies the nonlinear constraint
$$X^{ij}+\frac{1}{4}X^{ik}\overline{D}^2\overline{X}^{jk}=W^iW^j.$$
(6.35)
The proof of self-duality of this theory can be found in . Obviously, this system is a natural generalization of the Bagger–Galperin construction for the $`𝒩=1`$ super BI action, which we discussed in sect. 4.
Since for several vector multiplets $`W^30`$, after solving constraint (6.35) the action will have a more complicated form than (5.1).
## 7 Self-dual models with $`𝒩`$ = 1 tensor multiplet
In it was shown that partial breaking of $`𝒩=2`$ supersymmetry to $`𝒩=1`$ can be described with the $`𝒩=1`$ tensor multiplet as the Goldstone multiplet. The construction of Bagger and Galperin was based on an analogy between the $`𝒩=1`$ vector and tensor multiplets. Here we will pursue the same analogy to generalize the formalism of sect. 4 to construct nonlinear self-dual models of the $`𝒩=1`$ tensor multiplet.
We start with a brief description of the $`𝒩=1`$ tensor multiplet (see for more details). The multiplet is described by a real linear superfield $`L`$
$$D^2L=\overline{D}^2L=0,L=\overline{L}.$$
(7.1)
The general solution of this constraint is
$$L=\frac{1}{2}(D^\alpha \eta _\alpha +\overline{D}_{\dot{\alpha }}\overline{\eta }^{\dot{\alpha }}),\overline{D}_{\dot{\alpha }}\eta _\alpha =0.$$
(7.2)
The chiral spinor superfield $`\eta _\alpha `$ is a gauge field defined modulo transformations
$$\delta \eta _\alpha =\mathrm{i}\overline{D}^2D_\alpha K,$$
(7.3)
with $`K`$ a real unconstrained superfield, and $`L`$ is the gauge invariant field strength. The independent components of $`L`$ are a scalar $`\phi =L|_{\theta =0}`$, a Weyl spinor $`\psi _\alpha =D_\alpha L|_{\theta =0}`$ and its conjugate, and a vector $`\stackrel{~}{V}_{\alpha \dot{\alpha }}=\frac{1}{2}[D_\alpha ,\overline{D}_{\dot{\alpha }}]L|_{\theta =0}`$ constrained by $`_a\stackrel{~}{V}^a=0`$. The constraint means that $`\stackrel{~}{V}`$ is the dual field strength of an antisymmetric tensor field, $`\stackrel{~}{V}^a=\frac{1}{2}\epsilon ^{abcd}_bB_{cd}`$.
For generic models of the tensor multiplet, the gauge invariant action is a functional of $`L`$, $`S[L]`$. Here our consideration will be restricted to those models with actions of the form $`S[\mathrm{\Psi },\overline{\mathrm{\Psi }}]`$, where
$$\mathrm{\Psi }_\alpha =D_\alpha L,D_\beta \mathrm{\Psi }_\alpha =0.$$
(7.4)
For example, for the free tensor multiplet we have
$$S_{\mathrm{free}}=\mathrm{d}^8zL^2=\frac{1}{4}\mathrm{d}^6\overline{z}\mathrm{\Psi }^2+\frac{1}{4}\mathrm{d}^6z\overline{\mathrm{\Psi }}^2.$$
(7.5)
The antichiral spinor $`\mathrm{\Psi }_\alpha `$ is a constrained superfield
$$\frac{1}{4}\overline{D}^2\mathrm{\Psi }_\alpha +\mathrm{i}_{\alpha \dot{\alpha }}\overline{\mathrm{\Psi }}^{\dot{\alpha }}=0.$$
(7.6)
This constraint can be treated as the Bianchi identity. Its general solution is (7.4). The bosonic components of $`\mathrm{\Psi }_\alpha `$ are field strengths of the zero-form and two-form, $`U_a=_a\phi `$ and $`\stackrel{~}{V}_a`$, respectively.
For the theory with action $`S[\mathrm{\Psi },\overline{\mathrm{\Psi }}]`$, we introduce antichiral $`\mathrm{{\rm Y}}_\alpha `$ and chiral $`\overline{\mathrm{{\rm Y}}}^{\dot{\alpha }}`$ superfields as follows
$$\mathrm{i}\mathrm{{\rm Y}}_\alpha 2\frac{\delta }{\delta \mathrm{\Psi }^\alpha }S,\mathrm{i}\overline{\mathrm{{\rm Y}}}^{\dot{\alpha }}2\frac{\delta }{\delta \overline{\mathrm{{\rm Y}}}_{\dot{\alpha }}}S.$$
(7.7)
Then one can check that the equation of motion reads
$$\frac{1}{4}\overline{D}^2\mathrm{{\rm Y}}_\alpha +\mathrm{i}_{\alpha \dot{\alpha }}\overline{\mathrm{{\rm Y}}}^{\dot{\alpha }}=0$$
(7.8)
which has the same form as the Bianchi identity (7.6). Therefore, in analogy with sect. 4, one may consider U(1) duality rotations
$`\left(\begin{array}{c}\mathrm{{\rm Y}}^{}\\ \mathrm{\Psi }^{}\end{array}\right)=\left(\begin{array}{cc}\mathrm{cos}\lambda & \hfill \mathrm{sin}\lambda \\ \mathrm{sin}\lambda & \hfill \mathrm{cos}\lambda \end{array}\right)\left(\begin{array}{c}\mathrm{{\rm Y}}\\ \mathrm{\Psi }\end{array}\right).`$ (7.15)
The theory proves to be duality invariant iff the self-duality equation
$$\mathrm{Im}\mathrm{d}^6\overline{z}\left(\mathrm{\Psi }^\alpha \mathrm{\Psi }_\alpha +\mathrm{{\rm Y}}^\alpha \mathrm{{\rm Y}}_\alpha \right)=0$$
(7.16)
is satisfied.
Under duality rotations, the following functional
$$S\frac{\mathrm{i}}{4}\mathrm{d}^6\overline{z}\mathrm{\Psi }^\alpha \mathrm{{\rm Y}}_\alpha +\frac{\mathrm{i}}{4}\mathrm{d}^6z\overline{\mathrm{\Psi }}_{\dot{\alpha }}\overline{\mathrm{{\rm Y}}}^{\dot{\alpha }}$$
(7.17)
remains invariant. As in sect. 4, this property implies self-duality under a superfield Legendre transformation which is defined by replacing the action $`S[\mathrm{\Psi },\overline{\mathrm{\Psi }}]`$ with
$$S[\mathrm{\Psi },\overline{\mathrm{\Psi }},\mathrm{\Psi }_\mathrm{D},\overline{\mathrm{\Psi }}_\mathrm{D}]=S[\mathrm{\Psi },\overline{\mathrm{\Psi }}]\frac{\mathrm{i}}{2}\mathrm{d}^6\overline{z}\mathrm{\Psi }^\alpha \mathrm{\Psi }_{\mathrm{D}\alpha }+\frac{\mathrm{i}}{2}\mathrm{d}^6z\overline{\mathrm{\Psi }}_{\dot{\alpha }}\overline{\mathrm{\Psi }}_\mathrm{D}{}_{}{}^{\dot{\alpha }},$$
(7.18)
where $`\mathrm{\Psi }_\alpha `$ is now an unconstrained antichiral spinor superfield, and $`\mathrm{\Psi }_{\mathrm{D}\alpha }`$ the dual field strength
$$\mathrm{\Psi }_{\mathrm{D}\alpha }=D_\alpha L_\mathrm{D},D^2L_\mathrm{D}=0,\overline{L}_\mathrm{D}=L_\mathrm{D}.$$
(7.19)
Since above considerations are very similar to those in sect. 4, one can make use of the previous results to derive nonlinear self-dual models of the tensor multiplet. This is achieved by substituting $`W^2\overline{\mathrm{\Psi }}^2`$ in the action (4.20). The results of sec. 6 can also be generalized to the case of self-dual systems with several tensor multiplets.
## 8 Self-duality and gauge field democracy
The general theory of duality invariance in four space-time dimensions, which was reviewed in sect. 3, admits a natural higher-dimensional generalization . In even dimensions $`d=2p`$, one considers theories of $`n`$ gauge $`(p1)`$-forms $`B_{a_1\mathrm{}a_{p1}}^i`$ coupled to matter fields $`\varphi ^\mu `$ such that the Lagrangian is a function of the field strengths $`F_{a_1\mathrm{}a_p}^i=p_{[a_1}B_{a_2\mathrm{}a_{p]}}^i`$, <sup>8</sup><sup>8</sup>8Our normalization is $`_{[a_1}B_{a_2\mathrm{}a_p]}=\frac{1}{p!}(_{a_1}B_{a_2\mathrm{}a_p}\pm \mathrm{})`$ matter fields and their derivatives, $`L=L(F,\varphi ,\varphi )`$. The action is invariant under the Abelian gauge symmetries $`B^iB^i+d\mathrm{\Lambda }^i`$ where $`\mathrm{\Lambda }^i`$ is any $`(p2)`$-form. In complete analogy with the four-dimensional case, one can introduce duality transformations and analyze the conditions for self-duality. The maximal duality group turns out to depend on the dimension of the space-time. For $`d=4k`$ the maximal duality group is Sp($`2n,`$), while for $`d=4k+2`$ it is O$`(n,n)`$. In the absence of matter, the maximal duality group is compact: U$`(n)`$ in $`d=4k`$ dimensions, and $`\mathrm{O}(n)\times \mathrm{O}(n)`$ for $`d=4k+2`$. The fact that the maximal duality group depends on the dimension of space-time was also discussed in .
A natural question is what happens to a self-dual system upon dimensional reduction? The answer is that one finds a self-dual system with $`(p1)`$-forms and $`(dp1)`$-forms in $`d`$ space-time dimensions, where $`d`$ is not necessarily even. We now discuss the general properties of such models. In $`d=4`$ such models also appear as the bosonic sector of the self-dual systems of the $`𝒩=1`$ tensor multiplet we discussed in sect. 7. In fact, the analysis of this section was inspired by self-duality of the tensor Goldstone multiplet .
In $`d`$ space-time dimensions, we consider a theory of $`n`$ gauge $`(p1)`$-forms $`B_{a_1\mathrm{}a_{p1}}^i`$ and $`m`$ gauge $`(dp1)`$-forms $`C_{a_1\mathrm{}a_{dp1}}^I`$ coupled to matter fields $`\varphi ^\mu `$. We introduce the gauge invariant field strengths
$$U_{a_1\mathrm{}a_p}^i=p_{[a_1}B_{a_2\mathrm{}a_{p]}}^i,V_{a_1\mathrm{}a_{dp}}^I=(dp)_{[a_1}C_{a_2\mathrm{}a_{dp]}}^I.$$
(8.1)
Without loss of generality, we assume $`p<[d/2]`$ <sup>9</sup><sup>9</sup>9$`[.]`$ denotes the integer part. The case $`p=[d/2]`$ for even $`d`$ is special and was mentioned at the beginning of this section. and then introduce the Hodge-dual of $`V`$,
$$\stackrel{~}{V}_{a_1\mathrm{}a_p}^I=\frac{1}{(dp)!}\epsilon _{a_1\mathrm{}a_pb_1\mathrm{}b_{dp}}V^{Ib_1\mathrm{}b_{dp}},$$
(8.2)
which is of lower rank than $`V`$. The Bianchi identities read
$$_{[b}U_{a_1\mathrm{}a_p]}^i=0,^b\stackrel{~}{V}_{ba_1\mathrm{}a_{p1}}^I=0.$$
(8.3)
The Lagrangian is required to be a function of the field strengths, matter fields and their derivatives
$$L=L(U,\stackrel{~}{V},\varphi ,\varphi )L(\phi ).$$
(8.4)
In terms of the dual variables
$$\stackrel{~}{G}_{a_1\mathrm{}a_p}^i(\phi )=p!\frac{L(\phi )}{U^{ia_1\mathrm{}a_p}},H_{a_1\mathrm{}a_p}^I(\phi )=p!\frac{L(\phi )}{\stackrel{~}{V}^{Ia_1\mathrm{}a_p}},$$
(8.5)
the equations of motion for the gauge fields read
$$^b\stackrel{~}{G}_{ba_1\mathrm{}a_{p1}}^i=0,_{[b}H_{a_1\mathrm{}a_p]}^I=0.$$
(8.6)
The explicit structure of the Bianchi identities and equations of motion implies that one may consider duality transformations of the form
$`\delta \left(\begin{array}{c}H\\ U\end{array}\right)`$ $`=`$ $`\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)\left(\begin{array}{c}H\\ U\end{array}\right),\delta \left(\begin{array}{c}\stackrel{~}{V}\\ \stackrel{~}{G}\end{array}\right)=\left(\begin{array}{cc}M& N\\ R& S\end{array}\right)\left(\begin{array}{c}\stackrel{~}{V}\\ \stackrel{~}{G}\end{array}\right),`$ (8.19)
$`\delta \varphi ^\mu `$ $`=`$ $`\xi ^\mu (\varphi ).`$ (8.20)
Here $`A,B,C,D`$ and $`M,N,R,S`$ are real constant matrices, and $`\xi ^\mu `$ are some unspecified functions of the matter fields. We have suppressed the indices $`i,I`$. Compatibility of the duality transformations with self-duality now imposes the conditions
$$N=C^\mathrm{T},R=B^\mathrm{T},M+A^\mathrm{T}=\kappa \mathrm{𝟏},S+D^\mathrm{T}=\kappa \mathrm{𝟏},$$
(8.21)
with $`\kappa `$ some real constant, as well as the following functional relations
$`{\displaystyle \frac{}{\stackrel{~}{V}^{I\underset{¯}{a}}}}\left[\delta L{\displaystyle \frac{1}{p!}}B^{Jj}\stackrel{~}{V}^JU^j{\displaystyle \frac{1}{p!}}C^{jJ}\stackrel{~}{G}^jH^J\kappa L\right]`$ $`=`$ $`0,`$
$`{\displaystyle \frac{}{U^{i\underset{¯}{a}}}}\left[\delta L{\displaystyle \frac{1}{p!}}B^{Jj}\stackrel{~}{V}^JU^j{\displaystyle \frac{1}{p!}}C^{jJ}\stackrel{~}{G}^jH^J\kappa L\right]`$ $`=`$ $`0,`$ (8.22)
where we have introduced the notation
$$\stackrel{~}{G}^iH^J=\stackrel{~}{G}^{ia_1\mathrm{}a_p}H_{a_1\mathrm{}a_p}^J\stackrel{~}{G}^{i\underset{¯}{a}}H_{\underset{¯}{a}}^J.$$
(8.23)
Furthermore, the matter equation of motion transforms covariantly if one requires
$$\frac{\delta }{\delta \varphi ^\mu }\left[\delta S\frac{1}{p!}\mathrm{d}^4xC^{iI}\stackrel{~}{G}^iH^I\right]=0.$$
(8.24)
Eqs. (8.22) and (8.24) are then seen to be compatible if $`\kappa =0`$ and if the Lagrangian transforms as
$`\delta L`$ $`=`$ $`{\displaystyle \frac{1}{p!}}B^{Ii}\stackrel{~}{V}^IU^i+{\displaystyle \frac{1}{p!}}C^{iI}\stackrel{~}{G}^iH^I`$ (8.25)
$`=`$ $`\delta \left({\displaystyle \frac{1}{p!}}U^i\stackrel{~}{G}^i\right)=\delta \left({\displaystyle \frac{1}{p!}}\stackrel{~}{V}^IH^I\right).`$
Since $`\kappa =0`$, eq. (8.21) means that (8.20) becomes
$`\delta \left(\begin{array}{c}H\\ U\end{array}\right)`$ $`=`$ $`\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)\left(\begin{array}{c}H\\ U\end{array}\right),\delta \left(\begin{array}{c}\stackrel{~}{V}\\ \stackrel{~}{G}\end{array}\right)=\left(\begin{array}{cc}\hfill A^\mathrm{T}& \hfill C^\mathrm{T}\\ \hfill B^\mathrm{T}& \hfill D^\mathrm{T}\end{array}\right)\left(\begin{array}{c}\stackrel{~}{V}\\ \stackrel{~}{G}\end{array}\right).`$ (8.38)
One easily shows that both variations satisfy the same algebra, namely gl$`(n+m,)`$. The maximal connected duality group is therefore $`\mathrm{GL}_0(n+m,`$). The finite form for duality transformations is
$`\left(\begin{array}{c}H^{}\\ U^{}\end{array}\right)`$ $`=`$ $`g\left(\begin{array}{c}H\\ U\end{array}\right),\left(\begin{array}{c}\stackrel{~}{V}^{}\\ \stackrel{~}{G}^{}\end{array}\right)=\left(\begin{array}{cc}\mathrm{𝟏}& \hfill 0\\ 0& \hfill \mathrm{𝟏}\end{array}\right)(g^\mathrm{T})^1\left(\begin{array}{cc}\mathrm{𝟏}& \hfill 0\\ 0& \hfill \mathrm{𝟏}\end{array}\right)\left(\begin{array}{c}\stackrel{~}{V}\\ \stackrel{~}{G}\end{array}\right),`$ (8.51)
with $`g\mathrm{GL}_0(n+m,`$).
Equation (8.25) can be rewritten in a different, more useful, form if one directly computes $`\delta L`$. This gives the self-duality equation
$`p!\delta _\varphi L`$ $`=`$ $`B^{Ii}\stackrel{~}{V}^IU^iC^{iI}\stackrel{~}{G}^iH^I`$ (8.52)
$`+`$ $`A^{IJ}\stackrel{~}{V}^IH^JD^{ij}\stackrel{~}{G}^iU^j,`$
with $`\delta _\varphi L`$ as in eq. (3.24).
In the absence of matter, the maximal connected duality group becomes SO($`n+m`$), the maximal compact subgroup of $`\mathrm{GL}_0(n+m,`$); i.e. $`A^\mathrm{T}=A,D^\mathrm{T}=D,B^\mathrm{T}=C`$. Then, the self-duality equation (8.52) is equivalent to
$`U^i\stackrel{~}{V}^I+\stackrel{~}{G}^iH^I=0,`$ (8.53)
$`\left(U^i{\displaystyle \frac{}{U^j}}U^j{\displaystyle \frac{}{U^i}}\right)L=0,\left(\stackrel{~}{V}^I{\displaystyle \frac{}{\stackrel{~}{V}^J}}\stackrel{~}{V}^J{\displaystyle \frac{}{\stackrel{~}{V}^I}}\right)L=0.`$ (8.54)
Eq. (8.54) says that the theory is manifestly SO$`(n)\times \mathrm{SO}(m)`$ invariant.
By analogy with the results of , any SO($`n+m`$) duality invariant model $`L(U^i,\stackrel{~}{V}^I)`$ can be lifted to a model with the non-compact duality symmetry $`\mathrm{GL}_0(n+m,`$) by coupling the gauge fields to scalar fields $`\varphi ^\mu `$ parameterizing the quotient space $`\mathrm{GL}_0(n+m,)/\mathrm{SO}(n+m)`$.
Any SO$`(n+m)`$ duality invariant model $`L(U^i,\stackrel{~}{V}^I)`$, where $`U_p^i=\mathrm{d}B_{p1}^i`$ and $`V_{dp}^I=\mathrm{d}C_{dp1}^I`$, with $`n0`$ and $`m0`$ , enjoys self-duality under Legendre transformation which dualizes two given forms $`B_{p1}^i`$ and $`C_{dp1}^I`$ into a $`(dp1)`$-form and a $`(p1)`$-form, respectively. This is a simple consequence of the duality invariance, see sect. 2 for more details. On the other hand, one can apply a Legendre transformation which, say, leaves the gauge $`(p1)`$-forms invariant but dualizes all gauge $`(dp1)`$-forms into $`(p1)`$-forms. One then obtains a model of $`(n+m)`$ gauge $`(p1)`$-forms. Remarkably, the SO$`(n+m)`$ duality symmetry of the original model turns into a manifest (linear) SO$`(n+m)`$ symmetry of the dualized model. This is a consequence of the self-duality equations (8.53) and (8.54) and the standard properties of Legendre transformation. Therefore, in the models that we have considered here, all fields are on the same footing, hence the title of this subsection. The SO$`(n+m)`$ duality symmetry is linearly realized if all form are of the same degree.
The self-duality equations (3.32) and (3.33) are difficult to solve. However, for (8.53) and (8.54), there exists a simple scheme to derive their general solution. One starts with an SO$`(n+m)`$ invariant model of $`(n+m)`$ gauge $`(p1)`$-forms in $`d`$ dimensions, and then simply dualize $`m`$ of the fields into gauge $`(dp1)`$-forms by applying the proper Legendre transformation. The dualized model is invariant under the duality transformations.
If $`n=m`$, there are systems (we will give examples below) which are invariant under Sp($`2n,`$) rather than the maximal duality group GL($`2n,`$). This is the case if the matrix parameterizing the infinitesimal transformation of $`\stackrel{~}{V}`$ and $`\stackrel{~}{G}`$, written in the form
$`\delta \left(\begin{array}{c}\stackrel{~}{G}\\ \stackrel{~}{V}\end{array}\right)`$ $`=`$ $`\left(\begin{array}{cc}\hfill D^\mathrm{T}& \hfill B^\mathrm{T}\\ \hfill C^\mathrm{T}& \hfill A^\mathrm{T}\end{array}\right)\left(\begin{array}{c}\stackrel{~}{G}\\ \stackrel{~}{V}\end{array}\right),`$ (8.61)
is required to coincide with the transformation of $`H`$ and $`U`$ (c.f. (8.38)). In the absence of matter, the duality group of these systems reduces to U$`(n)`$ (see sect. 3) and the self-duality equations take the form (from now on, we do not distinguish between indices $`i`$ and $`I`$)
$`U^{(i}\stackrel{~}{V}^{j)}+\stackrel{~}{G}^{(i}H^{j)}=0,`$ (8.62)
$`(U^i{\displaystyle \frac{}{U^j}}+\stackrel{~}{V}^i{\displaystyle \frac{}{\stackrel{~}{V}^j}})L(ij)=0.`$ (8.63)
Eq. (8.63) means that the Lagrangian is manifestly SO$`(n)`$ invariant. Any U$`(n)`$ duality invariant model can be made Sp$`(2n,)`$ duality invariant by coupling the gauge fields to scalars valued in $`\mathrm{Sp}(2n,)/\mathrm{U}(n)`$. For $`n=1`$ the result reads
$$L(U,\stackrel{~}{V},𝒮)=\frac{1}{p!}𝒮_1U\stackrel{~}{V}+L(\sqrt{𝒮_2}U,\sqrt{𝒮_2}\stackrel{~}{V}),$$
(8.64)
with $`𝒮`$ the dilaton-axion field (2.32) transforming by the rule (2.41) under the duality group SL($`2,`$).
In contrast to U$`(1)`$ duality invariant models of a single gauge $`(2p1)`$-form in even dimensions $`d=4p`$, U$`(1)`$ duality invariant models of a gauge $`(p1)`$-form and a gauge $`(dp1)`$-form in arbitrary dimensions $`d`$ can be considered as reducible, since they involve two independent fields. However, the latter models possess ‘self-dual’ solutions
$$U_{a_1\mathrm{}a_p}=\gamma H_{a_1\mathrm{}a_p}(U,\stackrel{~}{V}),\stackrel{~}{V}_{a_1\mathrm{}a_p}=\frac{1}{\gamma }\stackrel{~}{G}_{a_1\mathrm{}a_p}(U,\stackrel{~}{V}),$$
(8.65)
with $`\gamma `$ a constant parameter. The explicit dependence of $`\gamma `$ is dictated by the self-duality equation (8.62). Such solutions of the equations of motion describe the dynamics of a single field.
To conclude, we give an example of a U(1) duality invariant model. The Lagrangian reads
$$L=\frac{1}{p!}\frac{1}{p!}\sqrt{1+UU\stackrel{~}{V}\stackrel{~}{V}(U\stackrel{~}{V})^2}.$$
(8.66)
It is easy to check that $`L`$ solves the self-duality equation (8.62), and therefore the theory is U(1) duality invariant. The theory can be equivalently represented in the form
$$L=\frac{1}{2p!}(\chi +\overline{\chi }),$$
(8.67)
where the complex field $`\chi `$ is a functions of $`U`$ and $`\stackrel{~}{V}`$ which satisfies the nonlinear constraint
$$\chi +\frac{1}{2}\chi \overline{\chi }\psi =0,\psi =\frac{1}{2}(U+\mathrm{i}\stackrel{~}{V})^2.$$
(8.68)
This representation is analogous to that for the BI theory described in sect. 2.
The above duality invariant system has a supersymmetric origin. Let us choose $`d=4`$ and then $`p=1`$ is the only interesting choice. The dynamical fields are a scalar $`\phi `$ and an antisymmetric gauge field $`B_{ab}`$ which should enter the Lagrangian only via their field strengths $`U_a=_a\phi `$ and $`\stackrel{~}{V}^a=\frac{1}{2}\epsilon ^{abcd}_bB_{cd}`$. Then, the Lagrangian (8.66) describes the bosonic sector of a model for partial $`𝒩=2𝒩=1`$ supersymmetry breaking with the tensor multiplet as the Goldstone multiplet . The antisymmetric gauge field can be dualized into a scalar, by applying the appropriate Legendre transformation. The resulting model is manifestly U(1) invariant and it describes a 3-brane in six dimensions.
Other examples of U(1) duality invariant models of the scalar and antisymmetric tensor in four dimensions can be obtained by considering the bosonic sector of the self-dual tensor multiplet systems we discussed in sect. 7. It is worth noting that not all U(1) duality invariant models of the scalar and antisymmetric tensor admit a supersymmetric extension: the two fields have to appear in the action in the combination $`\psi `$ as defined in (8.68). This is in contrast with what we found in self-dual nonlinear electrodynamics.
Using the results of , the construction just described can be generalized to derive U$`(n)`$ duality invariant models of $`n`$ gauge $`(p1)`$-forms and $`n`$ gauge $`(dp1)`$-forms in four dimensions. The Lagrangian is
$$L=\frac{1}{2p!}\mathrm{tr}(\chi +\overline{\chi }),$$
(8.69)
where the complex $`n\times n`$ matrix $`\chi `$ is a function of $`U^i`$ and $`\stackrel{~}{V}^i`$ which satisfies the nonlinear constraint
$$\chi ^{ij}+\frac{1}{2}\chi ^{ik}\overline{\chi }^{jk}=\frac{1}{2}(U^i+\mathrm{i}\stackrel{~}{V}^i)(U^j+\mathrm{i}\stackrel{~}{V}^j).$$
(8.70)
## 9 $`𝒩`$ = 2 duality rotations
The construction of sect. 4 admits a natural generalization to $`𝒩=2`$ supersymmetry , although here much less explicit results have been obtained so far. We will discuss the case of one single Abelian gauge multiplet only, the generalization to an arbitrary number being straightforward.
We will work in $`𝒩=2`$ global superspace $`^{4|8}`$ parametrized by $`𝒵^A=(x^a,\theta _i^\alpha ,\overline{\theta }_{\dot{\alpha }}^i)`$, where $`i=\underset{¯}{1},\underset{¯}{2}`$. The flat covariant derivatives $`𝒟_A=(_a,𝒟_\alpha ^i,\overline{𝒟}_i^{\dot{\alpha }})`$ satisfy the algebra
$$\{𝒟_\alpha ^i,𝒟_\beta ^j\}=\{\overline{𝒟}_{\dot{\alpha }i},\overline{𝒟}_{\dot{\beta }j}\}=0,\{𝒟_\alpha ^i,\overline{𝒟}_{\dot{\alpha }j}\}=2\mathrm{i}\delta _j^i(\sigma ^a)_{\alpha \dot{\alpha }}_a.$$
(9.1)
Throughout this section, we will use the notation:
$`𝒟^{ij}𝒟^{\alpha (i}𝒟_\alpha ^{j)}=𝒟^{\alpha i}𝒟_\alpha ^j,`$ $`\overline{𝒟}^{ij}\overline{𝒟}_{\dot{\alpha }}^{(i}\overline{𝒟}^{j)\dot{\alpha }}=\overline{𝒟}_{\dot{\alpha }}^i\overline{𝒟}^{j\dot{\alpha }}`$
$`𝒟^4{\displaystyle \frac{1}{16}}(𝒟^{\underset{¯}{1}})^2(𝒟^{\underset{¯}{2}})^2,`$ $`\overline{𝒟}^4{\displaystyle \frac{1}{16}}(\overline{𝒟}_{\underset{¯}{1}})^2(\overline{𝒟}_{\underset{¯}{2}})^2.`$ (9.2)
An integral over the full superspace (with the measure $`\mathrm{d}^{12}𝒵=\mathrm{d}^4x\mathrm{d}^4\theta \mathrm{d}^4\overline{\theta }`$) can be reduce to one over the chiral subspace (with the measure $`\mathrm{d}^8𝒵=\mathrm{d}^4x\mathrm{d}^4\theta `$) or over the antichiral subspace ($`\mathrm{d}^8\overline{𝒵}=\mathrm{d}^4x\mathrm{d}^4\overline{\theta }`$):
$$\mathrm{d}^{12}𝒵(𝒵)=\mathrm{d}^8𝒵𝒟^4(𝒵)=\mathrm{d}^8\overline{𝒵}\overline{𝒟}^4(𝒵).$$
(9.3)
The discussion in this section is completely analogous to the one presented in the first part of sect. 4. We will thus be brief. Let $`𝒮[𝒲,\overline{𝒲}]`$ be the action describing the dynamics of a single $`𝒩=2`$ vector multiplet. The (anti) chiral superfield strengths $`\overline{𝒲}`$ and $`𝒲`$ satisfy the Bianchi identity
$$𝒟^{ij}𝒲=\overline{𝒟}^{ij}\overline{𝒲}.$$
(9.4)
The general solution of the Bianchi identity ,
$$𝒲=\overline{𝒟}^4𝒟^{ij}V_{ij},\overline{𝒲}=𝒟^4\overline{𝒟}^{ij}V_{ij}$$
(9.5)
is in terms of a real unconstrained prepotential $`V_{(ij)}`$.
Suppose that $`𝒮[𝒲,\overline{𝒲}]`$ can be unambiguously defined as a functional of unconstrained (anti) chiral superfields $`\overline{𝒲}`$ and $`𝒲`$. Then, one can define (anti) chiral superfields $`\overline{}`$ and $``$ as
$$\mathrm{i}4\frac{\delta }{\delta 𝒲}𝒮[𝒲,\overline{𝒲}],\mathrm{i}\overline{}4\frac{\delta }{\delta \overline{𝒲}}𝒮[𝒲,\overline{𝒲}]$$
(9.6)
in terms of which the equations of motion are
$$𝒟^{ij}=\overline{𝒟}^{ij}\overline{}.$$
(9.7)
Again, since the Bianchi identity (9.4) and the equation of motion (9.7) have the same functional form, one can consider infinitesimal U(1) duality transformations
$$\delta 𝒲=\lambda ,\delta =\lambda 𝒲.$$
(9.8)
The analysis of Appendix A leads to
$`\delta 𝒮`$ $`=`$ $`{\displaystyle \frac{\mathrm{i}}{8}}\lambda {\displaystyle \mathrm{d}^8𝒵\left(𝒲^2^2\right)}+{\displaystyle \frac{\mathrm{i}}{8}}\lambda {\displaystyle \mathrm{d}^8\overline{𝒵}\left(\overline{𝒲}^2\overline{}^2\right)}`$
$`=`$ $`{\displaystyle \frac{\mathrm{i}}{4}}\lambda {\displaystyle \mathrm{d}^8𝒵^2}{\displaystyle \frac{\mathrm{i}}{4}}\lambda {\displaystyle \mathrm{d}^8\overline{𝒵}\overline{}^2}`$
The theory is thus duality invariant provided the following reality condition is satisfied:
$$\mathrm{d}^8𝒵\left(𝒲^2+^2\right)=\mathrm{d}^8\overline{𝒵}\left(\overline{𝒲}^2+\overline{}^2\right).$$
(9.10)
Here $``$ and $`\overline{}`$ are defined as in (9.6), and $`𝒲`$ and $`\overline{𝒲}`$ should be considered as unconstrained chiral and antichiral superfields, respectively. Eq. (9.10) serves as our master functional equation ($`𝒩=2`$ self-duality equation) to determine duality invariant models of the $`𝒩=2`$ vector multiplet.
We remark that, as in the $`𝒩=0,1`$ cases, the action itself is not duality invariant, but
$$\delta \left(𝒮\frac{\mathrm{i}}{8}\mathrm{d}^8𝒵𝒲+\frac{\mathrm{i}}{8}\mathrm{d}^8\overline{𝒵}\overline{}\overline{𝒲}\right)=0.$$
(9.11)
The invariance of the latter functional under a finite U(1) duality rotation by $`\pi /2`$, is equivalent to the self-duality of $`𝒮`$ under Legendre transformation,
$$𝒮[𝒲,\overline{𝒲}]\frac{\mathrm{i}}{4}\mathrm{d}^8𝒵𝒲𝒲_\mathrm{D}+\frac{\mathrm{i}}{4}\mathrm{d}^8\overline{𝒵}\overline{𝒲}\overline{𝒲}_\mathrm{D}=𝒮[𝒲_\mathrm{D},\overline{𝒲}_\mathrm{D}],$$
(9.12)
where $`𝒲_\mathrm{D}`$ is the dual chiral field strength,
$$𝒲_\mathrm{D}=\overline{𝒟}^4𝒟_{ij}V_\mathrm{D}{}_{}{}^{ij},$$
(9.13)
with $`V_\mathrm{D}^{ij}`$ a real unconstrained prepotential.
Apart from the $`𝒩=2`$ Maxwell action
$$𝒮_{\mathrm{free}}=\frac{1}{8}\mathrm{d}^8𝒵𝒲^2+\frac{1}{8}\mathrm{d}^8\overline{𝒵}\overline{𝒲}^2,$$
(9.14)
only one other solution of (9.10) is known :
$$𝒮=\frac{1}{4}\mathrm{d}^8𝒵𝒳+\frac{1}{4}\mathrm{d}^8\overline{𝒵}\overline{𝒳},$$
(9.15)
where the chiral superfield $`𝒳`$ is a functional of $`𝒲`$ and $`\overline{𝒲}`$ defined via the constraint
$$𝒳=𝒳\overline{𝒟}^4\overline{𝒳}+\frac{1}{2}𝒲^2.$$
(9.16)
Following , let us prove that this system provides a solution of the self-duality equation (9.10). Under an infinitesimal variation of $`𝒲`$ only, we have
$`\delta _𝒲𝒳`$ $`=`$ $`\delta _𝒲𝒳\overline{𝒟}^4\overline{𝒳}+𝒳\overline{𝒟}^4\delta _𝒲\overline{𝒳}+𝒲\delta 𝒲,`$
$`\delta _𝒲\overline{𝒳}`$ $`=`$ $`\delta _𝒲\overline{𝒳}𝒟^4𝒳+\overline{𝒳}𝒟^4\delta _𝒲𝒳.`$ (9.17)
From these relations one gets
$$\delta _𝒲𝒳=\frac{1}{1𝒬}\left[\frac{𝒲\delta 𝒲}{1\overline{𝒟}^4\overline{𝒳}}\right],\delta _𝒲\overline{𝒳}=\frac{\overline{𝒳}}{1𝒟^4𝒳}𝒟^4\delta _𝒲𝒳,$$
(9.18)
where
$`𝒬=𝒫\overline{𝒫},`$ $`\overline{𝒬}=\overline{𝒫}𝒫,`$
$`𝒫={\displaystyle \frac{𝒳}{1\overline{𝒟}^4\overline{𝒳}}}\overline{𝒟}^4,`$ $`\overline{𝒫}={\displaystyle \frac{\overline{𝒳}}{1𝒟^4𝒳}}𝒟^4.`$ (9.19)
With these results, it is easy to compute $``$:
$$\mathrm{i}=\frac{𝒲}{1\overline{𝒟}^4\overline{𝒳}}\left\{1+\overline{𝒟}^4\overline{𝒫}\frac{1}{1𝒬}\frac{𝒳}{1\overline{𝒟}^4\overline{𝒳}}+\overline{𝒟}^4\frac{1}{1\overline{𝒬}}\frac{\overline{𝒳}}{1𝒟^4𝒳}\right\}.$$
(9.20)
Now, a short calculation gives
$$\mathrm{Im}\mathrm{d}^8𝒵\left\{^2+2\frac{1}{1𝒬}\frac{𝒳}{1\overline{𝒟}^4\overline{𝒳}}\right\}=0.$$
(9.21)
On the other hand, the constraint (9.16) implies
$$\mathrm{d}^8𝒵𝒳\mathrm{d}^8\overline{𝒵}\overline{𝒳}=\frac{1}{2}\mathrm{d}^8𝒵𝒲^2\frac{1}{2}\mathrm{d}^8\overline{𝒵}\overline{𝒲}^2,$$
(9.22)
and hence
$$\frac{\delta }{\delta 𝒲}\left\{\mathrm{d}^8𝒵𝒳\mathrm{d}^8\overline{𝒵}\overline{𝒳}\right\}=𝒲.$$
(9.23)
The latter relation can be shown to be equivalent to
$$\frac{1}{1𝒬}\frac{𝒳}{1\overline{𝒟}^4\overline{𝒳}}=𝒫\frac{1}{1\overline{𝒬}}\frac{\overline{𝒳}}{1𝒟^4𝒳}+𝒳.$$
(9.24)
Using this result in eq. (9.21), we arrive at the relation
$$\mathrm{d}^8𝒵^2\mathrm{d}^8\overline{𝒵}\overline{}^2=2\mathrm{d}^8𝒵𝒳+2\mathrm{d}^8\overline{𝒵}\overline{𝒳}$$
(9.25)
which is equivalent, due to (9.22), to (9.10).
The dynamical system (9.15), (9.16) was introduced in as the $`𝒩=2`$ supersymmetric BI action (c.f. with the similar construction for the $`𝒩=1`$ super BI action we described in sect. 3). Such an interpretation is supported in part by the fact that the theory correctly reduces to the $`𝒩=1`$ BI in a special $`𝒩=1`$ limit; we now briefly discuss this issue.
Let us introduce the $`𝒩=1`$ components of the $`𝒩=2`$ vector multiplet. Given an $`𝒩=2`$ superfield $`U`$, its $`𝒩=1`$ projection is defined to be $`U|=U(𝒵)|_{\theta _{\underset{¯}{2}}=\overline{\theta }^{\underset{¯}{2}}=0}`$. The $`𝒩=2`$ vector multiplet contains two independent chiral $`𝒩=1`$ components
$$𝒲|=\sqrt{2}\mathrm{\Phi },𝒟_\alpha ^{\underset{¯}{2}}𝒲|=2\mathrm{i}W_\alpha ,(𝒟^{\underset{¯}{2}})^2𝒲|=\sqrt{2}\overline{D}^2\overline{\mathrm{\Phi }}.$$
(9.26)
Using in addition that
$$\mathrm{d}^8𝒵=\frac{1}{4}\mathrm{d}^6z(𝒟^{\underset{¯}{2}})^2,\mathrm{d}^{12}𝒵=\frac{1}{16}\mathrm{d}^8z(𝒟^{\underset{¯}{2}})^2(\overline{𝒟}_{\underset{¯}{2}})^2,$$
(9.27)
the above definitions imply that the free $`𝒩=2`$ vector multiplet action (9.14) straightforwardly reduces to $`𝒩=1`$ superfields
$$𝒮_{\mathrm{free}}=\mathrm{d}^8z\overline{\mathrm{\Phi }}\mathrm{\Phi }+\frac{1}{4}\mathrm{d}^6zW^2+\frac{1}{4}\mathrm{d}^6\overline{z}\overline{W}^2.$$
(9.28)
If one switches off $`\mathrm{\Phi }`$,
$$\mathrm{\Phi }=0(𝒟^{\underset{¯}{2}})^2𝒲|=0,$$
(9.29)
one readily observes that the theory (9.15), (9.16) reduces to the $`𝒩=1`$ BI theory (5.2), (5.3). However, it was shown in that there exist infinitely many manifestly $`𝒩=2`$ supersymmetric models possessing this very property. Of course, the specific feature of the system (9.15), (9.16) is its invariance under U(1) duality rotations, and the requirement of self-duality severely restricts the class of possible models. But it turns out that even the latter requirement is not sufficient to uniquely fix the $`𝒩=2`$ supersymmetric BI action.
The $`𝒩=2`$ supersymmetric BI action is expected to describe a single D3-brane in six dimensions
$$L_{\mathrm{D3}\mathrm{brane}}=1\sqrt{det\left(\eta _{ab}+F_{ab}+_a\overline{\phi }_b\phi \right)}.$$
(9.30)
Here the complex transverse coordinates $`\phi `$ of the brane should, in general, be related to the scalars $`\varphi =𝒲|_{\theta =0}`$ and the other components of the $`𝒩=2`$ vector multiplet by a nonlinear field redefinition (see, e.g. ). Since $`L_{\mathrm{D3}\mathrm{brane}}`$ is manifestly invariant under constant shifts of the transverse coordinates
$$\phi (x)\phi (x)+\sigma ,$$
(9.31)
the full supersymmetric theory must also be invariant under such transformations acting on $`𝒲`$ in a nonlinear way
$$𝒲(𝒵)𝒲(𝒵)+\sigma +𝒪(𝒲,\overline{𝒲}).$$
(9.32)
Moreover, the $`𝒩=2`$ supersymmetric BI action is expected to provide a model for partial $`𝒩=4𝒩=2`$ supersymmetry breaking . It means that the action should be invariant under nonlinear transformations
$$𝒲(𝒵)𝒲(𝒵)+ϵ(\theta )+𝒪(𝒲,\overline{𝒲}),ϵ(\theta )=\sigma +ϵ_i^\alpha \theta _\alpha ^i,$$
(9.33)
with $`ϵ_i^\alpha `$ a constant spinor parameter. We now demonstrate that the system (9.15), (9.16) is not compatible even with the simpler transformations (9.32).
To start with, it is worth pointing out the following. When looking for nonlinear symmetry transformations (9.32) or (9.33), one might first try to duplicate the trick<sup>10</sup><sup>10</sup>10This course was taken up in . which successfully worked in the case of the $`𝒩=1`$ supersymmetric BI action (see sect. 5). Namely, one can introduce the transformation of $`𝒳`$
$$\delta 𝒳=ϵ(\theta )𝒲,\overline{𝒟}_{\dot{\alpha }}^iϵ(\theta )=𝒟^{ij}ϵ(\theta )=0,$$
(9.34)
which obviously leaves the action (9.15) invariant. But this variation of $`𝒳`$ must be induced by a variation of $`𝒲`$ consistent with the constraint (9.16). A direct analysis shows that the variation $`\delta 𝒲`$, that is derived in this way, does not satisfy the Bianchi identity (9.4). The difference from the $`𝒩=1`$ case is simple but crucial: the $`𝒩=2`$ vector multiplet does not possess any analogue of the property $`W^3=0`$, typical for the $`𝒩=1`$ vector multiplet.
We will use the following general Ansatz
$$\delta 𝒲=\sigma +\sigma \overline{𝒟}^4\overline{𝒴}+\overline{\sigma }\mathrm{}𝒴,\overline{𝒟}_{\dot{\alpha }}^i𝒴=0$$
(9.35)
for symmetry transformations (9.32). The variation is consistent with the Bianchi identity (9.4). The chiral superfield $`𝒴`$ is some unknown functional of $`𝒲`$ and $`\overline{𝒲}`$. The precise form of $`𝒴`$ as well as of the $`𝒩=2`$ supersymmetric BI action, $`𝒮_{\mathrm{BI}}`$, should be determined, order by order in perturbation theory, from three requirements: (i) the action is to be invariant under transformations (9.35); (ii) the action should solve the self-duality equation (9.10); (iii) to order $`𝒲^4`$, the action should have the form:
$$𝒮_{\mathrm{BI}}=𝒮_{\mathrm{free}}+𝒮_{\mathrm{int}},𝒮_{\mathrm{int}}=\frac{1}{8}\mathrm{d}^{12}𝒵𝒲^2\overline{𝒲}^2+𝒪(𝒲^6).$$
(9.36)
This reproduces the known $`F^4`$ terms in the BI action.<sup>11</sup><sup>11</sup>11This is the only known superinvariant with this property. Direct calculation gives for $`𝒴`$
$`𝒴`$ $`=`$ $`{\displaystyle \frac{1}{2}}𝒲^2\left\{1+{\displaystyle \frac{1}{2}}\overline{𝒟}^4\overline{𝒲}^2+{\displaystyle \frac{1}{8}}\overline{𝒟}^4(\overline{𝒲}^2𝒟^4𝒲^2)+{\displaystyle \frac{1}{8}}(\overline{𝒟}^4\overline{𝒲}^2)^2\right\}`$ (9.37)
$`{\displaystyle \frac{1}{36}}\overline{𝒟}^4(𝒲^3\mathrm{}\overline{𝒲}^3)+𝒪(𝒲^8),`$
while $`𝒮_{\mathrm{int}}`$ reads
$`𝒮_{\mathrm{int}}`$ $`=`$ $`{\displaystyle \frac{1}{8}}{\displaystyle }\mathrm{d}^{12}𝒵𝒲^2\overline{𝒲}^2\{1+{\displaystyle \frac{1}{2}}(𝒟^4𝒲^2+\overline{𝒟}^4\overline{𝒲}^2)`$
$`+`$ $`{\displaystyle \frac{1}{4}}((𝒟^4𝒲^2)^2+(\overline{𝒟}^4\overline{𝒲}^2)^2)+{\displaystyle \frac{3}{4}}(𝒟^4𝒲^2)(\overline{𝒟}^4\overline{𝒲}^2)\}`$
$`+`$ $`{\displaystyle \frac{1}{24}}{\displaystyle \mathrm{d}^{12}𝒵\left\{\frac{1}{3}𝒲^3\mathrm{}\overline{𝒲}^3+\frac{1}{2}(𝒲^3\mathrm{}\overline{𝒲}^3)\overline{𝒟}^4\overline{𝒲}^2+\frac{1}{2}(\overline{𝒲}^3\mathrm{}𝒲^3)𝒟^4𝒲^2+\frac{1}{48}𝒲^4\mathrm{}^2\overline{𝒲}^4\right\}}`$
$`+`$ $`𝒪(𝒲^{10}).`$
The expression in the first two lines of (9) comes from the expansion of (9.15) in powers of $`𝒲`$ and its conjugate to the order indicated. As concerns the expression in the third line of (9), it is not present in the power series expansion of (9.15), but it is required for invariance under transformations (9.35). It is also worth noting that the expression in the first line of (9.37) coincides with the decomposition of $`(𝒳)`$ (9.16) to the given order.
Our conclusion is that the system (9.15), (9.16) cannot be identified with the correct $`𝒩=2`$ supersymmetric D3-brane world-volume action, and the problem of constructing such an action is still open.
A natural possibility to look for $`𝒩=2`$ supersymmetric BI action, advocated in , is first to derive a manifestly $`(1,0)`$ supersymmetric BI action in six dimensions and, then to dimensionally reduce to four dimensions. By construction, the resulting four-dimensional model should be manifestly $`𝒩=2`$ supersymmetric and invariant under constant shift transformations $`𝒲𝒲+\sigma `$, without any nonlinear terms. However, the problem of constructing the manifestly $`(1,0)`$ supersymmetric BI action in six dimensions is not simple. In $`d=6`$ there exists an off-shell formulation for the $`(1,0)`$ vector multiplet . But super-extensions of $`F^2,F^4`$ and $`F^6`$ terms, which appear in the decomposition of the $`d=6`$ BI action, cannot be represented by integrals over $`(1,0)`$ superspace or its subspace. The super-extension of $`F^2`$ term was already derived in . As to the super-extensions of $`F^4`$ and $`F^6`$ terms, candidates were proposed in . Unfortunately, the proof of their invariance under $`(1,0)`$ supersymmetry transformations was based on the use of the identity (here we follow the $`d=6`$ notation of ) $`D_{\alpha (i}\left\{W_j^{[\beta }W_{k)}^{\gamma ]}\right\}=0`$, which holds on-shell , and not off-shell as claimed in . Therefore, the super-extensions of $`F^4`$ and $`F^6`$ terms proposed in are not invariant under $`(1,0)`$ supersymmetry transformations. Thus the problem of constructing a manifestly $`(1,0)`$ supersymmetric BI action is six dimensions remains unsolved. If such an action exists, its dimensional reduction to $`d=4`$ will be manifestly supersymmetric, but not all terms in the action can be represented as integrals over $`𝒩=2`$ superspace or its supersymmetric subspaces.
Acknowledgements We are grateful to Evgeny Ivanov, Dima Sorokin and Arkady Tseytlin for their interest in this project. We thank Paolo Aschieri for helpful discussions on duality rotations. Support from DFG-SFB-375, from GIF, the German-Israeli foundation for Scientific Research and from the EEC under TMR contract ERBFMRX-CT96-0045 is gratefully acknowledged. This work was also supported in part by the NATO collaborative research grant PST.CLG 974965, by the RFBR grant No. 99-02-16617, by the INTAS grant No. 96-0308 and by the DFG-RFBR grant No. 99-02-04022.
## Appendix A Derivation of the self-duality equation
Eq. (2.30) is derived as follows. For an infinitesimal U(1) duality rotation, we have
$$\stackrel{~}{G}_{ab}^{}(F^{})=\stackrel{~}{G}_{ab}(F)\lambda \stackrel{~}{F}_{ab}=\stackrel{~}{G}_{ab}(F)+2\frac{}{F^{ab}}\left(\frac{1}{4}\lambda F\stackrel{~}{F}\right),$$
(A.1)
where we have used the infinitesimal version of eq.(2.29). At the same time, from the definition of $`\stackrel{~}{G}^{}(F^{})`$ it follows
$$\stackrel{~}{G^{}}(F^{})=2\frac{L(F^{})}{F^{}}=2\left(\frac{}{F^{}}L(F)+\frac{}{F}\delta L\right),$$
(A.2)
where
$$\delta L=L(F^{})L(F).$$
(A.3)
Using $`F^{}=F+\lambda G`$, one can express $`/F^{}`$ on the right-hand side of (A.2) via $`/F`$ with the result
$$\stackrel{~}{G}_{ab}^{}(F^{})=\stackrel{~}{G}_{ab}(F)+2\frac{}{F^{ab}}\left(\delta L\frac{1}{4}\lambda G\stackrel{~}{G}\right).$$
(A.4)
Comparing eqs. (A.1) and (A.4) gives
$$\delta L=\frac{1}{4}\lambda (G\stackrel{~}{G}F\stackrel{~}{F}).$$
(A.5)
On the other hand, the Lagrangian can be varied directly to give
$$\delta L=\frac{L}{F^{ab}}\delta F^{ab}=\frac{1}{2}\lambda \stackrel{~}{G}G.$$
(A.6)
This is consistent with eq. (A.5) iff the self-duality equation (2.30) holds.
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# 1 Introduction
## 1 Introduction
It is well known that Bose-Einstein condensation, a first order phase transition in momentum space, can not occur for an ideal gas of free particles in two dimensions. In this short note we shall show that, by contrast, the introduction of one point-like impurity on the plane just allows the Bose-Einstein condensation to take place. Furthermore, the general pattern of the above phenomenon will be obtained from a generating Hamiltonian, which encodes a two parameters model to describe the point-like impurity in two and three spatial dimensions. The key point to be gathered is how to treat in quantum mechanics the presence of a point-like, or $`\delta `$-like, or null support impurity. Formal manipulations involving some kinds of regularized $`\delta `$-like potentials might drive, in general, to misleading and incorrect conclusions, as the latter ones are mathematically ill-defined. The correct quantum mechanical framework to treat point-like impurities is by means of the analysis of the self-adjoint extensions of the Hamiltonian operator which, in the present case, turns out to be just a symmetric operator. Consequently, its domain has to be suitably defined in order to obtain a self-adjoint Hamiltonian operator which admits a complete orthonormal set of eigenstates. This construction will be referred to in the sequel as the inclusion of contact interaction. To be definite, let us consider as a starting point of our analysis the following classical one-particle Hamiltonian in two spatial dimensions: namely,
$`H(\varphi )={\displaystyle \frac{1}{2m}}\left[𝐩{\displaystyle \frac{e}{c}}𝐀(𝐫)\right]^2,𝐩,𝐫𝐑^2,`$
(1.1)
$`A_j(x_1,x_2)={\displaystyle \frac{\varphi }{2\pi }}ϵ_{jk}{\displaystyle \frac{x_k}{r^2}},r\sqrt{x_1^2+x_2^2}.`$
Here the Aharonov-Bohm type vector potential corresponds to the presence of a $`\delta `$-like vortex of flux $`\varphi `$, which provides a good classical description of one point-like impurity, as we shall better specify below. Quantization of the classical Hamiltonian (1.1) leads to a symmetric operator and, consequently, one has to face the problem of finding all its self-adjoint extensions. The most general solution has been recently obtained in Ref. and it consists in a four parameter family. However, to our aim, we can restrict ourselves to the one parameter sub-family of the $`O(2)`$ rotational invariant self-adjoint Hamiltonians . The corresponding spectral decompositions read
$`H(\alpha ,E_0)={\displaystyle \underset{l=\mathrm{}}{\overset{+\mathrm{}}{}}}{\displaystyle _0^{\mathrm{}}}𝑑k{\displaystyle \frac{\mathrm{}^2k^2}{2m}}|l,kk,l\left|+\vartheta (E_0)\right|\psi _B\psi _B|,`$
(1.2)
$`\alpha {\displaystyle \frac{e\varphi }{hc}}]1,0],\mathrm{}E_0<+\mathrm{}`$
$`\vartheta `$ being the usual Heaviside’s step distribution, in terms of the eigenfunctions
$`r,\theta |l,k={\displaystyle \frac{\mathrm{exp}\{il\theta \}}{\sqrt{2\pi }}}\psi _l(k,r;\alpha ,E_0),k0,`$ (1.3)
where the improper eigenfunctions belonging to the continuous part of the spectrum are given by
$`\psi _l(k,r)=\sqrt{k}J_{|l+\alpha |}(kr),l𝐙\{0\},`$ (1.4)
$`\psi _0(k,r;E_0)=A(k;\alpha ,E_0)J_{|\alpha |}(kr)+B(k;\alpha ,E_0)N_{|\alpha |}(kr),`$ (1.5)
in which
$`{\displaystyle \frac{B(k;\alpha ,E_0)}{A(k;\alpha ,E_0)}}={\displaystyle \frac{\mathrm{sin}(\pi \alpha )}{\mathrm{cos}(\pi \alpha )+\mathrm{sgn}(E_0)\left(\mathrm{}^2k^2/2m|E_0|\right)^\alpha }},`$ (1.6)
whereas the normalizable bound state is provided by
$`r|\psi _B=\psi _B(\kappa ,r)={\displaystyle \frac{\kappa }{\pi }}\sqrt{{\displaystyle \frac{\mathrm{sin}(\pi \alpha )}{\alpha }}}K_\alpha (\kappa r),\mathrm{}\kappa \sqrt{2m|E_0|}.`$ (1.7)
Some remarks are now in order. First, as previously noticed, the above spectral decompositions precisely provides the correct mathematical framework to introduce and properly describe contact interaction in quantum mechanics. As a matter of fact, it turns out that the rotational invariant self-adjoint Hamiltonian operators $`H(\alpha ,E_0)`$ do represent a one parameter family, which is labeled by the energy scale $`E_0`$. In the range $`\mathrm{}<E_0<0`$, and only within this range of values, a bound state $`|\psi _B`$ exists, whose energy is just $`E_0`$. More generally, the physical meaning of the characteristic energy scale $`E_0`$ is given by the resonance energy, according to the following pattern: namely,
$`E_{\mathrm{res}}=\{\begin{array}{cc}|E_0|(\mathrm{sec}\pi \alpha )^{1/|\alpha |},\hfill & \mathrm{if}\mathrm{\hspace{0.17em}0}<|\alpha |<1/2,E_0<0;\hfill \\ |E_0|,\hfill & \mathrm{if}\mathrm{\hspace{0.17em}1}/2<|\alpha |<1,E_0<0;\hfill \end{array}`$ (1.10)
$`E_{\mathrm{res}}=E_0|\mathrm{sec}\pi \alpha |^{1/|\alpha |},\mathrm{if}\mathrm{\hspace{0.17em}1}/2<|\alpha |<1,E_00.`$ (1.11)
A further observation is that only the non-integer part of the vortex flux parameter $`\alpha `$ is actually observable, as its integer part can always be gauged away by a single valued phase transformation. To sum up, we can say that a general correct quantum mechanical description of one point-like impurity is provided by the two parameters family of self-adjoint Hamiltonians of Eq. (1.2). The existence of the contact interaction just corresponds to the presence of a specific locally square integrable singularity of the wave function at the impurity position - see Eq. (1.5). In the limit $`E_0\mathrm{}`$ contact interaction is removed, the domain of the Hamiltonian is that of the regular wave functions on the whole plane and the impurity is described in terms of a pure Aharonov-Bohm vortex of non-integer vorticity $`\alpha `$. If we further take the limit $`\alpha 0`$, i.e., also the Aharonov-Bohm interaction is turned off, the two dimensional free particle Hamiltonian is truly recovered (Friedrichs’ limit). As we shall discuss in the sequel, it is curious that just in the Friedrichs’ limit the Bose-Einstein condensation disappears in the two spatial dimensional case because, in the presence of contact interaction and no matter how weak it is, a non-vanishing critical temperature for the Bose-Einstein transition always exists.
## 2 One-particle partition function
In order to discuss Bose-Einstein condensation, it is necessary to compute the average number of particles at thermal equilibrium. To this aim, let us first evaluate the diagonal Heat-Kernel and the one-particle partition function. According to the spectral decomposition of Eq. (1.2) and after separation of the truly free particle Hamiltonian $`H_0H(0,\mathrm{})`$ contribution, it is not difficult to verify that the diagonal Heat-Kernel can be cast in the following form: namely,
$`G(\alpha ,\beta ,E_0;r)`$ $``$ $`G_{\mathrm{int}}(\alpha ,\beta ,E_0;r)+G_0(\beta )`$ (2.1)
$`=`$ $`𝐫\left|\left[\mathrm{exp}\{\beta H(\alpha ,E_0)\}\mathrm{exp}\{\beta H_0\}\right]\right|𝐫+\lambda _T^2`$
$`=`$ $`I(\alpha ;r)+I(\alpha ;r)2I(0;r)I_0(\alpha ;r)I_0(\alpha ;r)+I_0(0;r)`$
$`+\vartheta (E_0)e^{\beta E_0}\left|\psi _B(\kappa r)\right|^2+_0(\alpha ,E_0;r)+\lambda _T^2,`$
where the translation invariant free particle diagonal Heat-Kernel is nothing but the inverse square thermal wavelength $`\lambda _T(h/\sqrt{2\pi m\mathrm{k}T})`$ and we have set
$`I(\alpha ;r)={\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dk}{2\pi }}ke^{\beta \mathrm{}^2k^2/2m}{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}\left[J_{l+\alpha }(kr)\right]^2,`$ (2.2)
$`I_0(\alpha ;r)={\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dk}{2\pi }}ke^{\beta \mathrm{}^2k^2/2m}\left[J_\alpha (kr)\right]^2,`$ (2.3)
$`_0(\alpha ,E_0;r)={\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dk}{2\pi }}{\displaystyle \frac{k\mathrm{exp}\{\beta \mathrm{}^2k^2/2m\}}{1+\mathrm{tan}^2[\pi \mu (k)]+2\mathrm{tan}[\pi \mu (k)]\mathrm{cos}(\alpha \pi )}}`$
$`\times \left\{\mathrm{tan}^2[\pi \mu (k)]J_\alpha ^2(kr)+J_\alpha ^2(kr)+2\mathrm{tan}[\pi \mu (k)]J_\alpha (kr)J_\alpha (kr)\right\},`$ (2.4)
$`\mathrm{tan}[\pi \mu (k)]\mathrm{sgn}(E_0)\left[{\displaystyle \frac{2m|E_0|}{\mathrm{}^2k^2}}\right]^{|\alpha |}.`$ (2.5)
Now, it is very important to realize that the impurity interaction part of the diagonal Heat-Kernel $`G_{\mathrm{int}}(\alpha ,\beta ,E_0;r)`$ is integrable on the whole plane. This leads to the following result for the one-particle partition function
$`Z_{2\mathrm{D}}(\alpha ,\beta ,E_0)`$ $`=`$ $`{\displaystyle \frac{\mathrm{A}}{\lambda _T^2}}+{\displaystyle \frac{\alpha (\alpha +1)}{2}}+\vartheta (E_0)e^{\beta |E_0|}`$ (2.6)
$`+{\displaystyle \frac{\alpha \mathrm{sin}(\pi \alpha )}{\pi }}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dx}{x^{1+\alpha }}}{\displaystyle \frac{\mathrm{sgn}(E_0)e^{\beta |E_0|x}}{1+2\mathrm{s}\mathrm{g}\mathrm{n}(E_0)x^{|\alpha |}\mathrm{cos}(\pi \alpha )+x^{2|\alpha |}}},`$
where, as usual, we have denoted by $`A`$ the area divergence, due to the presence of the translation invariant part of the free one-particle Heat-Kernel. The above expression for the one-particle partition function can be used as a generating form which encodes different specific notable cases. In particular, the one-particle partition function in two spatial dimensions and in the presence of pure contact interaction can be obtained in the limit $`\alpha 0`$ and reads
$`Z_{2\mathrm{D}}(0,\beta ,E_B)`$ $`=`$ $`{\displaystyle \frac{\mathrm{A}}{\lambda _T^2}}+e^{\beta |E_B|}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dE}{E}}{\displaystyle \frac{e^{\beta E}}{\mathrm{ln}^2(E/E_B)+\pi ^2}}`$ (2.7)
$`=`$ $`{\displaystyle \frac{\mathrm{A}}{\lambda _T^2}}+\nu (\beta |E_B|),E_B<0,`$ (2.8)
where
$`\nu (x){\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{x^t}{\mathrm{\Gamma }(t+1)}}𝑑t.`$ (2.9)
Notice that in two spatial dimensions the bound state is always present for any $`\mathrm{}<E_B<0`$. Another distinguished case that can be read off the basic formula (2.6) is the three dimensions one-particle partition function in the presence of contact interaction. As a matter of fact, thanks to dimensional transmutation , the latter case just corresponds to the value $`\alpha =1/2`$, up to a suitable redefinition of the free part: namely,
$`Z_{3\mathrm{D}}(\beta ,E_0)={\displaystyle \frac{\mathrm{V}}{\lambda _T^3}}+\vartheta (E_0)e^{\beta |E_0|}+{\displaystyle \frac{1}{2}}\mathrm{sgn}(E_0)e^{\beta |E_0|}\mathrm{erfc}(\sqrt{\beta |E_0|}).`$ (2.10)
## 3 Results and Discussion
Now we are ready to discuss the Bose-Einstein condensation for an ideal gas of particles in the presence of one point-like impurity, as generally described by the one-particle Hamiltonian of Eq. (1.2). According to the general form (2.6) of the one-particle partition function, it turns out that the average particles density at thermal equilibrium in two spatial dimensions and in the presence of one point-like impurity is given by
$`n_{2\mathrm{D}}{\displaystyle \frac{N}{\mathrm{A}}}`$ $`=`$ $`\lambda _T^2g_1(z)+{\displaystyle \frac{z\alpha (\alpha +1)}{2\mathrm{A}(1z)}}+\vartheta (E_0){\displaystyle \frac{z}{\mathrm{A}(z_0z)}}`$ (3.1)
$`+\vartheta (E_0){\displaystyle \frac{z}{\mathrm{A}(1z)}}+\mathrm{sgn}(E_0){\displaystyle \frac{z}{\mathrm{A}}}{\displaystyle _0^{\mathrm{}}}𝑑E{\displaystyle \frac{\varrho (E;\alpha ,|E_0|)e^{\beta E}}{1z\mathrm{exp}\{\beta E\}}},`$
where we have set $`z_0\mathrm{exp}\{\beta E_0\}`$ and $`g_1(z)=\mathrm{ln}(1z)`$, whereas
$`\varrho (E;\alpha ,|E_0|)={\displaystyle \frac{\alpha \mathrm{sin}(\pi \alpha )E^{|\alpha |1}}{\pi |E_0|^\alpha \left[E^{2|\alpha |}+|E_0|^{2|\alpha |}+2\mathrm{s}\mathrm{g}\mathrm{n}(E_0)(E|E_0|)^{|\alpha |}\mathrm{cos}(\pi \alpha )\right]}}.`$ (3.2)
It is important to realize that if $`E_0<0`$ the range of the fugacity is $`0zz_0<1`$, whilst $`0z1`$ if $`E_00`$. Moreover, it is not difficult to prove that, thanks to analytic continuation, the very last term in Eq. (3.1) admits a finite limit when $`z1`$ and $`E_00`$. From the above expression (3.1) for the average particle density in two spatial dimensions, it appears to be manifest that Bose-Einstein condensation occurs only in the presence of the bound state, i.e., only for the sub-family of the self-adjoint extensions of the symmetric Hamiltonian (1.1) in the range $`\mathrm{}<E_0<0`$. In those cases, the critical temperature and/or specific area can be obtained as the unique solutions of the equation
$`\mathrm{ln}\left(1e^{\beta E_0}\right)={\displaystyle \frac{h^2\beta }{2\pi m}}n_{2\mathrm{D}}.`$ (3.3)
The three spatial dimensional case can be handled in a quite similar way, as it essentially corresponds to the specific value $`\alpha =1/2`$ in Eq. (3.1), up to terms irrelevant in the thermodynamic limit: namely,
$`n_{3\mathrm{D}}{\displaystyle \frac{N}{\mathrm{V}}}`$ $`=`$ $`\lambda _T^3g_{\frac{3}{2}}(z)+\vartheta (E_0){\displaystyle \frac{z}{\mathrm{V}(z_0z)}}+\vartheta (E_0){\displaystyle \frac{z}{\mathrm{V}(1z)}}`$ (3.4)
$`+\mathrm{sgn}(E_0){\displaystyle \frac{z}{\mathrm{V}}}{\displaystyle _0^{\mathrm{}}}𝑑E{\displaystyle \frac{\varrho (E;\alpha =\frac{1}{2},|E_0|)e^{\beta E}}{1z\mathrm{exp}\{\beta E\}}}.`$
The above equation clearly indicates that Bose-Einstein condensation always takes place in three spatial dimensions, in the presence as well as in the absence of the impurity. Nonetheless, the actual values of the critical temperature and/or density do depend upon the sign of the parameter characterizing the self-adjoint extension of the Hamiltonian. In fact, for $`E_00`$, i.e. in the absence of the bound state, the critical values are the usual ones as given by the solution of the equation $`\lambda _T^3n_{3\mathrm{D}}=\zeta (3/2)`$. At variance, when $`\mathrm{}<E_0<0`$, i.e. in the presence of the bound state, the critical values can be read off the equation
$`g_{\frac{3}{2}}(z_0)=\lambda _T^3n_{3\mathrm{D}}.`$ (3.5)
The case of pure contact interaction in two spatial dimensions can also be obtained from the basic formula (3.1) taking the limit $`\alpha 0`$ and treating separately the cases $`E_00`$ and $`\mathrm{}<E_0<0`$. As a matter of fact, the result is
$`n_{2\mathrm{D}}|_{\alpha =0}`$ $`=`$ $`\lambda _T^2g_1(z)+\vartheta (E_0){\displaystyle \frac{z}{\mathrm{A}(1z)}}+\vartheta (E_0){\displaystyle \frac{z}{\mathrm{A}(z_0z)}}`$ (3.6)
$`\vartheta (E_0){\displaystyle \frac{z}{\mathrm{A}}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dE}{E}}{\displaystyle \frac{e^{\beta E}}{1z\mathrm{exp}\{\beta E\}}}{\displaystyle \frac{1}{\mathrm{ln}^2(E/E_0)+\pi ^2}},`$
which shows that condensation does not occur when $`E_00`$, whereas it appears if $`\mathrm{}<E_0<0`$, the critical temperature and/or specific volume being always determined by Eq. (3.3) which does not depend upon $`\alpha `$. In this latter case, the very same formula can be also obtained directly from Eq. (2.8) as it does.
In conclusion, we have shown in this note that the presence of contact interaction makes it possible the occurrence of Bose-Einstein condensation in two spatial dimensions. This phenomenon is connected to the presence of a bound state in the spectrum of the self-adjoint Hamiltonian. It turns out to be remarkable that the latter circumstance is always there in the pure contact interaction case, that means without Aharonov-Bohm vortex interaction. It is in fact worthwhile to notice that, in the presence of a non-vanishing vorticity $`\alpha `$, the half-family without bound state of the self-adjoint extensions does not allow Bose-Einstein condensation, whilst the remaining half-family with bound state leads to Bose-Einstein condensation though the critical temperature is independent from $`\alpha `$. Accordingly, only when $`\mathrm{}<E_0<0`$ a non-vanishing and vorticity-independent critical temperature is allowed in the two spatial dimensional case - see Eq. (3.3)- whereas the critical temperature deviates from its conventional value in the three spatial dimensional case - see Eq. (3.5).
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# Position Measurement for a Relativistic Particle: Restricted-Path-Integral AnalysisPublished in Phys. Lett. A 208, 269-275 (1995).
## 1 Introduction
In nonrelativistic quantum mechanics, measurements are described by the von Neumann’s postulate of the state reduction (wave function collapse). In relativistic theory this is impossible since an instantaneous state reduction contradicts to causality and therefore must be modified. This problem has been considered by many authors (see for example -, but consensus was not achieved. We shall consider the measurement of the position of a relativistic particle with the help of the Restricted-Path-Integral (RPI) approach to continuous quantum measurements . The results will be shown to be in a qualitative agreement with the Hellwig and Kraus’ postulate according to which the state reduction occurs in the light cone of the of the measurement event.
Restricted-Path-Integral (RPI) approach has been proposed by R.Feynman for description of continuous (prolonged in time) measurements and technically elaborated in (see also ). The idea of the approach is that evolution of the system undergoing a continuous measurement must be described by the path integral restricted on the set of paths compatible with the measurement output.
The approach proved to be effective for different types of measurements of non-relativistic systems as well as for measurements of relativistic quantum fields (electromagnetic and gravitational, see ). Its advantage is in model-independence and generality. In the present paper the RPI approach will be applied to the problem of measurement of position of a relativistic particle.
Let the position of a relativistic particle be measured at a specified time moment and the measurement output correspond to the point (event) $`a`$ belonging to the corresponding time slice. Then the paths $`[x]`$ compatible with the output $`a`$ are those crossing this point, $`a[x]`$. Therefore, only these paths contribute the evolution of a particle subject to the measurement.
The problem is therefore reduced to calculating relativistic path integrals over the sets of paths crossing the given space-time point. Technical difficulties of this calculation will be overcome due to the specific properties of the causal propagator (the path integral over all paths).
## 2 Relativistic Path Integrals
The causal propagator (transition amplitude) for a relativistic particle can be expressed in the form of a path integral if one introduces, following Stueckelberg , the fifth parameter (besides four space-time coordinates) $`\tau `$ called proper time or historical time.
Consider for simplicity a scalar particle of the mass $`m`$. Its causal propagator is equal to the integral over the proper time,<sup>1</sup><sup>1</sup>1we shall use in the present paper the natural units $`\mathrm{}=c=1`$ everywhere but in the discussion of the results
$$K(x^{\prime \prime }x^{})=_0^{\mathrm{}}𝑑\tau \mathrm{exp}\left(i(m^2iϵ)\tau \right)K_\tau (x^{\prime \prime }x^{}),$$
(1)
of a subsidiary proper-time-dependent propagator. The latter, in turn, may be given the form of a path integral:
$$K_\tau (x^{\prime \prime }x^{})=_{x^{\prime \prime }x^{}}d[x]_\tau \mathrm{exp}\left(\frac{i}{4}_0^\tau (\dot{x},\dot{x})𝑑\tau \right).$$
(2)
Here $`(,)`$ denote the Lorentzian inner product and the usual definition of the measure is taken (see for example for elementary definitions from theory of path integrals).
As a result of these definitions, the subsidiary proper-time-dependent propagator satisfies the relativistic Schrödinger-type equation
$$\frac{d}{d\tau }K_\tau (x^{\prime \prime }x^{})=i\mathrm{}K_\tau (x^{\prime \prime }x^{})$$
(3)
and the causal propagator $`K_\tau (x^{\prime \prime }x^{})`$ is a Green function of the Klein-Gordon equation.
For the analysis of continuous measurements of a relativistic particle in the framework of the Restricted-Path-Integral (RPI) approach we have to deal with the path integrals of the type of Eqs. (1), (2) but restricted on the sets of paths compatible with the corresponding measurement outputs.
## 3 The Measurement Amplitude
We shall consider measurement of the particle position at time moment $`x^0=ct`$. First the overidealized situation of an absolutely precise measurement will be analyzed. A finite measurement error will be taken into account later on (Sect. 5).
If the position of the particle is precisely measured at time $`t`$, then the measurement outputs may be described as three-vectors $`𝐚`$ or as points $`a=(ct,𝐚)`$ of the time slice $`t=\mathrm{const}`$ of the space-time, i.e. points of the space-like surface $`𝒮=\{x|x^0=ct\}`$.
If we know that the measurement of the position (at time $`t`$) has given the result $`𝐚`$, then we know that the world line (trajectory) of the particle crossed the surface $`𝒮`$ in the point $`a`$. Therefore, instead of the integral over all paths, evolution of the particle must be described by the integral over the set $`I_a`$ of paths crossing $`𝒮`$ in the point $`a`$.
The set of paths $`I_a`$ is a “corridor” describing adequately the measurement output. In the Restricted-Path-Integral method restricting of the path integral (1), (2) onto this corridor will give an amplitude of transition under the measurement, or the measurement amplitude:
$`K^{(a)}(x^{\prime \prime },x^{})`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑\tau \mathrm{exp}\left({\displaystyle \frac{i}{\mathrm{}}}(m^2iϵ)\tau \right)K_\tau ^{(a)}(x^{\prime \prime },x^{}),`$
$`K_\tau ^{(a)}(x^{\prime \prime },x^{})`$ $`=`$ $`{\displaystyle _{x^{\prime \prime }ax^{}}}d[x]_\tau \mathrm{exp}\left({\displaystyle \frac{i}{4}}{\displaystyle _0^\tau }(\dot{x},\dot{x})𝑑\tau \right).`$ (4)
This restricted path integral is evidently connected in some way with the product of two unrestricted integrals of the type of Eqs. (1), (2), one integral from the point $`x^{}`$ to the point $`a`$, the other from $`a`$ to $`x^{\prime \prime }`$:
$$K^{(a)}(x^{\prime \prime },x^{})=K(x^{\prime \prime }a)K(ax^{}).$$
(5)
The precise definition of this product is to be found.
To find the correct definition for the product (5), we shall require that the set of the amplitudes (5) be complete. This means that summation of the amplitudes corresponding to all possible values of $`𝐚`$ should give the amplitude describing the evolution without measurement:
$$_𝒮d^3𝐚K^{(a)}(x^{\prime \prime },x^{})=K(x^{\prime \prime }x^{}).$$
(6)
This relation, with the expression (5) in the integrand, resembles the known Kolmogorov-type property of the propagator,
$$i_{x^0=ct}d^3𝐱K(x^{\prime \prime }x)\stackrel{}{_0}K(xx^{})=K(x^{\prime \prime }x^{})$$
(7)
where $`x^{\prime \prime 0}>ct>x^0`$ and it is denoted
$$f(x)\stackrel{}{_0}g(x)=f(x)\stackrel{}{\frac{}{x^0}}g(x)=f(x)\frac{g(x)}{x^0}\frac{f(x)}{x^0}g(x).$$
Therefore, the completeness of the measurement amplitudes will be provided if we define the product (5) as follows:<sup>2</sup><sup>2</sup>2This amplitude will be used only for estimating relative probabilities, so that its normalization is not important. We shall consider the normalization in Sect. 5, discussing the measurement with a finite error.
$$K^{(a)}(x^{\prime \prime },x^{})=iK(x^{\prime \prime }a)\stackrel{}{\frac{}{a^0}}K(ax^{}).$$
(8)
We found the amplitude (8) requiring that summation of such amplitudes with different $`a`$ gives the propagator of a free particle (without any measurement). This requirement is natural in the framework of the theory of free particles because of the superposition principle. After this, when the form of the amplitude (8) is found, we go over from theory of a free particle to theory of a measured particle, where the role of these amplitudes will be quite different because the superposition principle does not take place.
In quantum theory of measurements the superposition principle is restricted: the amplitudes corresponding to different measurement outputs may not be added. In our case, when the precise measurement of the position is under consideration, the amplitudes (8) with different $`a`$ correspond to different measurement outputs. Therefore, each of them must be used separately, and their summation has no sense. Amplitudes $`K^{(a)}(x^{\prime \prime },x^{})`$ with different $`a=(ct,𝐚)`$ are incoherent.<sup>3</sup><sup>3</sup>3Later on we shall consider the measurement with a finite precision. Then the amplitudes with close $`𝐚`$ (differing less than by the measurement error) are coherent and may be summed up (see Sect. 5).
The amplitudes (8) are derived for a particle which is in the space-time point $`x^{}`$ before the measurement and in the point $`x^{\prime \prime }`$ after it. The realistic situation corresponds usually to the initial and final states given by the wave functions at the corresponding time moments $`t^{}`$, $`t^{\prime \prime }`$. The measurement amplitudes are then
$$K^{(a)}(\psi ^{\prime \prime },\psi ^{})=d^3𝐱^{}d^3𝐱^{\prime \prime }\overline{\psi ^(x^{\prime \prime })}\stackrel{}{\frac{}{x^{\prime \prime 0}}}K^{(a)}(x^{\prime \prime },x^{})\stackrel{}{\frac{}{x^0}}\psi ^{}(x^{}).$$
(9)
where the bar denotes a complex conjugate.
## 4 Properties of the Amplitude
So far we talked about the amplitude (8) in such a way as if it described measurement of the position with absolute precision. In other words, the measurement described by this amplitude was supposed to consist in localizing the particle in a single point $`a=(ct,𝐚)`$ on the surface $`𝒮`$. We shall see later (Sect. 6) that the only measurements which may be described correctly are those with a finite (and not too small) error. The amplitude (8) cannot be correctly interpreted. Nevertheless, it is important to investigate the properties of this amplitude. The properties of the finite-error amplitudes will follow then straightforwardly leading to physical conclusions.
For simplicity, we shall use physical terms in the analysis of the amplitude (8) as if it could be interpreted physically. In fact, the present section is devoted to investigation of mathematical properties of the subsidiary amplitude while physical interpretation is possible only for finite-error amplitudes of Sect. 5.
Consider therefore the amplitudes (8) and (9) as those describing evolution of the particle undergoing the position measurement. Then they are transition amplitudes from the point $`x^{}`$ to the point $`x^{\prime \prime }`$ (correspondingly from the state $`\psi ^{}`$ to the state $`\psi ^{\prime \prime }`$) under the condition that the measurement carried out at time $`t`$ gave the output $`𝐚`$. The amplitudes allow one to evaluate the probability that the particle achieves a certain state given an initial state and a measurement output.
Instead of this, one may interpret the same amplitudes (8), (9) as the probability amplitudes for different measurement outputs $`𝐚`$, given the initial and final states ($`x^{}`$ and $`x^{\prime \prime }`$ or $`\psi ^{}`$ and $`\psi ^{\prime \prime }`$). Relative probabilities of different measurement outputs may be estimated as square modula of the amplitudes. A mathematically rigorous definition of probabilities must include the generalized unitarity condition (see Sect. 6). Some conclusions however may be made without this.
The first conclusion may be made directly from the form of the amplitude (8). The causal propagator $`K(x^{\prime \prime },x^{})`$ is exponentially small if the interval $`x^{\prime \prime }x^{}`$ is outside the light cone. Therefore, the amplitude (8) is small if the point $`x^{}`$ is outside the past light cone of the point $`a`$ or/and $`x^{\prime \prime }`$ is outside the future light cone of $`a`$.
This property is in the qualitative agreement with the postulate of Hellwig and Kraus that reduction of the wave function of a relativistic system (for example a field) occurs not at the moment of the measurement but in the light cone of the space-time region where the measurement takes place (see also for a critical discussion of this postulate). Now we can derive the corresponding feature of the measurement rather than postulate it.
The conclusion following from the mentioned property of the amplitude (8) may be formulated as follows:
* The probability for the measurement to give the output $`a`$ is exponentially small if $`a`$ is not in the future light cone of the support of the initial wave function $`\psi ^{}`$.
* The probability that the particle will be found in the state $`\psi ^{\prime \prime }`$ after the measurement resulting in the output $`a`$, is exponentially small if $`a`$ is not in the past light cone of the support of $`\psi ^{\prime \prime }`$.
In the above statements, ‘exponentially small’ means decreasing in $`e`$ times at distance of the order of the Compton length from the boundary of the light cone.
The fact that the propagator does not abruptly disappear but rather exponentially decreases outside the light cone seems to contradict to causality, because the particle may seemingly be discovered in the point that it cannot achieve by causal evolution. However there is actually no contradiction. If one try to demonstrate experimentally (even by a thought experiment) this violation, one may see that such a demonstration is impossible because of the uncertainty relation.
One more thing guaranteed by the same property of the propagator (its exponential decreasing but not abrupt disappearing outside the light cone) is that dependence of the amplitude (8) or (9) from the position $`𝐚`$ is smoothed on scales of the order of the Compton length. Therefore, probabilities of two measurement outputs $`𝐚_1`$ and $`𝐚_2`$ are close if these outputs differ by the value of the order of Compton length or less. Though we discuss the precise measurement of position, the information (about the initial state) given by this measurement cannot be more precise than up to the Compton length. In the limits of the Compton length, the measurement output may be arbitrary.
In a sense, the preceding argument means that the precise measurement is impossible, the measurement error cannot be less than the Compton length. In fact, we shall show below (Sect. 6) that the correct theory of the measurement including a probability distribution exists for arbitrary initial and final states only if the measurement is performed with the error larger than the Compton length.
## 5 Measurement with a Finite Error
Consider now a measurement with a finite precision. Let the measurement error be $`\mathrm{\Delta }a`$. Then the measurement output $`𝐚`$ gives the information that the actual position of the particle $`𝐱`$ differs from $`𝐚`$ not more than by the value $`\mathrm{\Delta }a`$:
$$|𝐱𝐚|<\mathrm{\Delta }a.$$
(10)
Such a measurement must be described by the amplitude
$$K^{(a,\mathrm{\Delta }a)}(x^{\prime \prime },x^{})=_{|𝐛𝐚|<\mathrm{\Delta }a}d^3𝐛K^{(b)}(x^{\prime \prime },x^{}).$$
This interpretation of the measurement information corresponds to a specific property of the measuring device. The information is of this type if the device has a rectangular characteristic, equal to unity in the region (10) and zero otherwise.
In real situations measuring devices have smooth characteristics, and the information supplied by the measurement is less definite. If the measurement gives the output $`𝐚`$, this means that an actual position of the particle with high probability is very close to $`𝐚`$, with less probability is somewhat further, and it is quite improbable that it differs from $`𝐚`$ much more than by $`\mathrm{\Delta }a`$. The amplitude describing such a measurement has the form
$$K^{(a,\mathrm{\Delta }a)}(x^{\prime \prime },x^{})=d^3𝐛\rho (|𝐛𝐚|)K^{(b)}(x^{\prime \prime },x^{})$$
(11)
with the weight function $`\rho 0`$ concentrated in the region, of the dimension $`\mathrm{\Delta }a`$, around zero.
Again, just as in the case of the precise measurement, we should use the amplitudes (11) as incoherent ones. Each of them describes evolution for a certain output of measurement. Relative probabilities of different outputs may be roughly estimated by the square modula of the corresponding amplitudes. Mathematically rigorous concept of probabilities will be discussed in Sect. 6.
The analysis of Sect. 4 may be repeated with an evident change for the finite-error measurements. Now we should speak of the light cone of the region (10) around the point $`a`$ rather than the light cone of a single point $`a`$.
## 6 Generalized Unitarity
In the RPI approach to quantum continuous measurements evolution of the system undergoing the measurement is described by a set of propagators $`U_\alpha `$ depending on measurement outputs $`\alpha `$:
$$|\psi _\alpha =U_\alpha |\psi ,\rho _\alpha =U_\alpha \rho \left(U_\alpha \right)^{}.$$
(12)
This is the evolution law for the selective measurement when the measurement output is known. If it is unknown (non-selective measurement), then the density matrix after the evolution is a sum of the density matrices corresponding to all possible outputs:
$$\rho ^{}=\underset{\alpha }{}\rho _\alpha =\underset{\alpha }{}U_\alpha \rho U_{\alpha }^{}{}_{}{}^{}.$$
(13)
Probability for the measurement output to belong to the set $`A`$ is equal to
$$\mathrm{Prob}(\alpha A)=\underset{\alpha A}{}\mathrm{Tr}\rho _\alpha .$$
Conservation of probabilities (normalization of $`\rho ^{}`$) is provided by the generalized unitarity
$$\underset{\alpha }{}U_{\alpha }^{}{}_{}{}^{}U_\alpha =\mathrm{𝟏}.$$
In the case of continuous set of the measurement outputs (typical for a continuous measurement) it is more correct to speak about integration rather than summation over different outputs. Particularly, the last formulas should be rewritten as follows:
$`\rho ^{}={\displaystyle 𝑑\mu (\alpha )\rho _\alpha },\mathrm{Prob}(\alpha A)={\displaystyle _A}𝑑\mu (\alpha )\mathrm{Tr}\rho _\alpha `$ (14)
$`{\displaystyle 𝑑\mu (\alpha )U_{\alpha }^{}{}_{}{}^{}U_\alpha }=\mathrm{𝟏}.`$ (15)
The measure on the set of all outputs has to be chosen in such a way as to provide the validity of the generalized unitarity (15).
We should now introduce the corresponding concepts (probability and generalized unitarity) in our case.
The causal propagator (1) describes the evolution of a (positive-frequency) state of a free particle:
$$\psi (x^{\prime \prime })=i_𝒮^{}d^3𝐱^{}K(x^{\prime \prime },x^{})\stackrel{}{\stackrel{}{_0}}\psi (x^{})$$
(16)
where $`𝒮^{}=\{x^{}|x^0=\mathrm{const}\}`$. In the course of this evolution the inner product
$$(\psi _1,\psi _2)=i_𝒮d^3𝐱\overline{\psi _1(x)}\stackrel{}{_0}\psi _2(x)$$
(17)
is conserved (here $`𝒮=\{x|x^0=\mathrm{const}\}`$). The evolution described by the propagator $`K(x^{\prime \prime },x^{})`$ is therefore unitary in the following sense:
$$i_𝒮d^3𝐱\overline{K(x,x^{\prime \prime })}\stackrel{}{_0}K(x,x^{})=K(x^{\prime \prime },x^{}).$$
(18)
Eq. (18) represents unitarity of the theory (conservation of probabilities) in an unusual way since there are positive- and negative-frequency wave functions in relativistic theory but we are interested only in positive-frequency states of the particle. The meaning of Eq. (18) is following. Acting by the propagator $`\overline{K(x,x^{\prime \prime })}`$ which is conjugate for $`K(x,x^{\prime \prime })`$ we describe propagation in an opposite direction, from the first argument of the propagator to the second one. Therefore, the action of $`K(x,x^{})`$ followed by the action of $`\overline{K(x,x^{\prime \prime })}`$ gives the same result as the action of $`K(x^{\prime \prime },x^{})`$.
If the arguments $`x^{}`$ and $`x^{\prime \prime }`$ belong to the same time slice, then the action of the propagator $`\overline{K(x,x^{\prime \prime })}`$ (according to the formula (16)) does not change the wave function. This brings us close to the usual form of unitarity. We may suppose for simplicity that $`x^0>x^{\prime \prime 0}>x^0`$. Then positive-frequency part of the causal propagator may be used instead of the complete propagator. The last two times $`x^{\prime \prime 0}`$ and $`x^0`$ may be arbitrarily close to each other.
Unitarity (18) of the causal propagator may be shown to lead to the generalized unitarity of the measurement amplitudes,
$$i_𝒮d^3𝐚_{\stackrel{~}{𝒮}}d^3𝐱\overline{K^{(a,\mathrm{\Delta }a)}(x,x^{\prime \prime })}\stackrel{}{_0}K^{(a,\mathrm{\Delta }a)}(x,x^{})=K(x^{\prime \prime },x^{}),$$
(19)
if the error of the measurement (dimension of the region where the function $`\rho (|𝐛𝐚|)`$ is close to the maximum) is larger than the Compton length $`\lambda _C=\mathrm{}/mc`$.
More precisely, the condition for the generalized unitarity may be formulated in terms of the Fourier expansion of the following function:
$$d^3𝐚\rho (|𝐛𝐚|)\rho (|𝐛^{}𝐚|)=d^3𝐥Q(𝐥)e^{i𝐥(𝐛^{}𝐛)}.$$
(20)
The generalized unitarity takes place if
$$|𝐥|\lambda _C=\frac{mc}{\mathrm{}}$$
(21)
for all $`𝐥`$ in the region where the function $`Q(𝐥)`$ is not negligible.
Besides this, the following normalization condition should be fulfilled:
$$d^3𝐥Q(𝐥)=1.$$
(22)
The equality (19) is a concrete form of the general relation (15) in the coordinate representation and with $`d^3𝐚`$ standing instead of $`d\mu (\alpha )`$.
The probability for the measurement output $`a`$ to belong to the set $`A𝒮`$ is expressed by the integral
$`\mathrm{Prob}(aA)=i{\displaystyle _A}d^3𝐚{\displaystyle _{𝒮^{\prime \prime }}}d^3𝐱^{\prime \prime }{\displaystyle _𝒮^{}}d^3𝐱_1^{}{\displaystyle _𝒮^{}}d^3𝐱_2^{}`$
$`\times \overline{K^{(a,\mathrm{\Delta }a)}(x^{\prime \prime },x_2^{})}\stackrel{}{{\displaystyle \frac{}{x^{\prime \prime 0}}}}K^{(a,\mathrm{\Delta }a)}(x^{\prime \prime },x_1^{})\stackrel{}{{\displaystyle \frac{}{x_{1}^{}{}_{}{}^{0}}}}\stackrel{}{{\displaystyle \frac{}{x_{2}^{}{}_{}{}^{0}}}}\rho (x_1^{},x_2^{}).`$ (23)
The generalized unitarity (19) expresses conservation of probabilities in the course of the evolution of the system undergoing the position measurement.
The generalized unitarity does not take place if the measurement of position is more precise than up to the Compton length. The reason of this is evident. Localization of the particle in the region of the dimension less than the Compton length requires larger energy than the threshold of the pair creation. In this case interaction of the particle with the localizing device has a quite different character and cannot be described as a measurement of position.
However, if we expand the relation (19) in the Fourier integral, we shall see that only low-momentum components of this relation (with $`|𝐩|<\mathrm{}/\mathrm{\Delta }a`$) are violated. Therefore, the description of the measurement is correct for high-momentum states (with $`|𝐩|\mathrm{}/\mathrm{\Delta }a`$).
## 7 Conclusion
We showed that the measurement of the position of a relativistic particle can be correctly described in the framework of the method of restricted path integrals. The results obtained are in accordance with a more formal (not dynamical) consideration of the measurements of this type (see for the first attempt and and references therein for the recent papers).
It is interesting to apply the same method to other types of measurements in relativistic systems, for example to measurement of fields or to non-local measurements of different types. This will be done in a separate paper.
ACKNOWLEDGEMENT
The authors are indebted to K.Hellwig for fruitful discussions of the problem. This work was supported in part by the Deutsche Forschungsgemeinschaft and the Russian Foundation for Fundamental Research, grant 95-01-00111a.
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# 1 Introduction
## 1 Introduction
Two dimensional integrable models represent an important laboratory for testing new ideas and developing new methods for constructing exact solutions as well as for the nonperturbative quantization of 4-D non-abelian gauge theories, gravity and string theory. Among the numerous techniques for constructing 2-D integrable models and their solutions , , the hamiltonian reduction of the WZNW model (or equivalently the gauged WZNW ) associated to a finite dimensional Lie algebra $`𝒢`$ has provided an universal and simple method for deriving the equations of motion (or action ) of 2-d integrable models. In particular, the conformal Toda (CT) models were constructed by implementing a consistent set of constraints on the WZNW currents . The method was subsequently extended to infinite dimensional affine algebras ( denoted here by $`\widehat{𝒢}`$), leading to WZNW currents satisfying the so called two loop current algebra . By further imposing an infinite number of suitable constraints the conformal affine Toda (CAT) models were constructed . The power of such method was demonstrated in constructing (multi) soliton solutions of the abelian affine Toda models and certain nonsingular nonabelian (NA) affine Toda models .
The present paper is devoted to the systematic construction of a new family of axial and vector affine NA Toda integrable models (associated to an affine Kac-Moody algebra of rank r - $`\widehat{𝒢}_r`$) having one global U(1) symmetry. They represent appropriate integrable perturbations of the conformal $`\sigma `$-models (with tachyons and dilaton included) describing strings on curved backgrounds of black hole type (see ref., ). An important feature of these integrable models is that they admit U(1)-charged topological solitons for imaginary coupling constant . With topological $`\theta `$-term added to their actions, the one soliton spectrum manifests properties quite similar to the dyons of 4-D Yang-Mills -Higgs model namely, their electric charges get contributions from the magnetic (topological) one. The conformal limits (without the tachyonic terms) of axial and vector models in consideration, are known to be T-dual to each other, having $`O(r,r|Z)`$ as T-duality group (see for example ). The natural question arises whether one can extend the critical (i.e. conformal) T-duality transformations to the corresponding integrable model and which is the off-critical T-duality group. The problem of non conformal (i.e. off critical ) T-duality was first addressed in in the context of the principal chiral $`\sigma `$-models. For more general discussion of T-duality as canonical transformations of conformal and non-conformal 2-d models see . It turns out that axial and vector models similarly to their conformal $`\sigma `$\- models counterparts are forming again T-dual pairs. The critical T-duality group is however broken to $`O(1,1|Z)`$ for the family of dyonic integrable models under investigation.
As in the conformal case the subclass of T-selfdual models is of particular interest. We first investigate the (torsionless) T-selfdual conformal $`𝒢_r`$-NA Toda models obtained from the conformal $`\sigma `$-models (associated to finite dimensional Lie algebras $`𝒢_r`$) by adding all the possible marginal operators. As it shown in Sect.4.1, T-seldual conformal NA -Toda models exist only for orthogonal algebras $`𝒢_r=B_r=SO(2r+1)`$ and for specific choice of the form of the marginal operators only. The axial and vector dyonic models can be obtained by adding certain integrable relevant operators to the conformal $`𝒢_r`$-NA Toda actions. As is well known , the algebraic structure underlying these nonconformal affine $`\widehat{𝒢}_r`$-NA Toda models is encodded in the choice of grading operator Q for $`\widehat{𝒢}_r`$ and in the form of the grade $`\pm 1`$ constant elements $`ϵ_\pm `$ (which characterizes the potential). Addressing once more the question of the T-selfdual (torsionless) members of the considered family of integrable models, we derive the Lie algebraic condition (i.e. all the possible choices of the algebraic data - $`\widehat{𝒢}_r,Q,ϵ_\pm )`$ that gives rise to T-selfdual models. It shown in sect. 4.2 and 4.3 that such subclass of models exist only for the following three affine Kac-Moody algebras:$`B_r^{(1)},A_{2r}^{(2)}`$ and $`D_{r+1}^{(2)}`$ when the grading operator Q and the $`ϵ_\pm `$’s are appropriately choosen.
An interesting byproduct of our algebraic T-selfduality condition is that the above three families of torsionless models exactly reproduce the Fateev’s integrable models . Our construction also provides a simple and systematic proof of their classical integrability. It is worthwhile to mention the following simple form of our T-selfduality condition: while the generic axial and vector models ( with one global U(1) symmetry) are characterized by the fact that their physical fields $`g_0^f`$ lie in the coset $`𝒢_0/𝒢_0^0=\frac{SL(2)U(1)^{rank𝒢1}}{U(1)}`$, the corresponding cosets for the T-selfdual models are of the particular form $`𝒢_0/𝒢_0^0=\frac{SL(2)}{U(1)}U(1)^{rank𝒢1}`$ .
This paper is organized as follows. Sect.2 contains the functional integral derivation of the effective actions for generic conformal,affine and conformal affine NA -Toda theories. In the particular case when these models manifest $`𝒢_0^0`$=U(1) gauge symmetry, the actions for the corresponding singular affine NA Toda models of axial and vector type are obtained. Two particular examples based on the affine algebras $`A_r^{(1)}`$ and $`B_r^{(1)}`$ are presented. Sect.3 is devoted to the analysis of the abelian off-critical T-duality relating the axial and vector type of models. We derive the Lie algebraic condition defining the family of T-selfdual torsionless singular affine NA Toda theories in Sect.4. In Sect.5 we present the zero curvature representations of all IM’s in consideration.
## 2 Gauged WZNW Construction of NA Toda Models
The generic conformal (or affine) $`G_r`$ (or $`\widehat{G}_r`$) - NA - Toda models are classified according to a $`𝒢_0𝒢`$ embedding induced by the grading operator $`Q`$ , which defines a specific decomposition of the corresponding finite $`𝒢_r`$ or infinite( $`\widehat{𝒢}_r`$) dimensional Lie algebras $`𝒢=_i𝒢_i`$ where $`[Q,𝒢_i]=i𝒢_i`$ and $`[𝒢_i,𝒢_j]𝒢_{i+j}`$. The group element $`g`$ can then be written in terms of the Gauss decomposition as
$$g=NBM$$
(2.1)
where $`N=\mathrm{exp}𝒢_<H_{}`$, $`B=\mathrm{exp}𝒢_0`$ and $`M=\mathrm{exp}𝒢_>H_+`$. The physical fields $`B`$ lie in the zero grade subgroup $`𝒢_0`$ and the models we seek correspond to the coset $`H_{}\backslash G/H_+`$.
For consistency with the Hamiltonian reduction formalism, the phase space of the G-invariant WZNW model is reduced by specifying the constant generators $`ϵ_\pm `$ of grade $`\pm 1`$. In order to derive an action for $`B𝒢_0`$, invariant under
$`gg^{}=\alpha _{}g\alpha _+,`$ (2.2)
where $`\alpha _\pm (z,\overline{z})H_\pm `$ we have to introduce a set of auxiliary gauge fields $`A𝒢_<`$ and $`\overline{A}𝒢_>`$ transforming as
$`AA^{}=\alpha _{}A\alpha _{}^1+\alpha _{}\alpha _{}^1,\overline{A}\overline{A}^{}=\alpha _+^1\overline{A}\alpha _++\overline{}\alpha _+^1\alpha _+.`$ (2.3)
The result is given by the gauged WZNW action (see for instance , ),
$`S_{G/H}(g,A,\overline{A})`$ $`=`$ $`S_{WZNW}(g)`$
$``$ $`{\displaystyle \frac{k}{2\pi }}{\displaystyle 𝑑z^2Tr\left(A(\overline{}gg^1ϵ_+)+\overline{A}(g^1gϵ_{})+Ag\overline{A}g^1\right)}.`$
Since the action $`S_{G/H}`$ is $`H`$-invariant, we may choose $`\alpha _{}=N^1`$ and $`\alpha _+=M^1`$. From the orthogonality of the graded subspaces, i.e. $`Tr(𝒢_i𝒢_j)=0,i+j0`$, we find
$`S_{G/H}(g,A,\overline{A})`$ $`=`$ $`S_{G/H}(B,A^{},\overline{A}^{})`$ (2.4)
$`=`$ $`S_{WZNW}(B){\displaystyle \frac{k}{2\pi }}{\displaystyle 𝑑z^2Tr[A^{}ϵ_++\overline{A}^{}ϵ_{}+A^{}B\overline{A}^{}B^1]},`$
where
$`S_{WZNW}={\displaystyle \frac{k}{4\pi }}{\displaystyle d^2zTr(g^1gg^1\overline{}g)}{\displaystyle \frac{k}{24\pi }}{\displaystyle _D}ϵ_{ijk}Tr(g^1_igg^1_jgg^1_kg),`$ (2.5)
and the topological term denotes a surface integral over a ball $`D`$ identified as space-time.
The action (2.4) describes nonsingular Toda models among which we find the conformal and the affine abelian Toda models for $`Q=_{i=1}^r\frac{2\lambda _iH}{\alpha _i^2},ϵ_\pm =_{i=1}^rc_{\pm i}E_{\pm \alpha _i}`$ and $`Q=h\widehat{d}+_{i=1}^r\frac{2\lambda _iH}{\alpha _i^2},\widehat{ϵ}_\pm =_{i=1}^rc_{\pm i}E_{\pm \alpha _i}^{(0)}+E_{\pm \alpha _0}^{(\pm 1)}`$ respectively.<sup>1</sup><sup>1</sup>1 by $`\alpha _0`$ we denote the highest root, $`\lambda _i`$ \- the fundamental weights, $`h`$ \- the coxeter number of $`𝒢`$ and $`H_i`$ are the Cartan subalgebra generators in the Cartan - Weyl basis satisfying $`Tr(H_iH_j)=\delta _{ij}`$. In both cases the zero grade subgroup $`𝒢_0=U(1)^r`$ is abelian and it coincides with the Cartan subalgebra of $`𝒢_r`$. Performing the integration over the auxiliary fields $`A`$ and $`\overline{A}`$ in the functional integral
$$Z_\pm =DAD\overline{A}\mathrm{exp}(F_\pm ),$$
(2.6)
where
$$F_\pm =\frac{k}{2\pi }\left(Tr(ABϵ_{}B^1)B(\overline{A}B^1ϵ_+B)B^1\right)d^2z$$
we derive the effective action for the conformal abelian Toda theories
$$S=S_{WZNW}(B)\frac{k}{2\pi }Tr\left(ϵ_+Bϵ_{}B^1\right)d^2z$$
(2.7)
and for the (nonconformal) affine abelian Toda IM by replacing $`ϵ_\pm `$ from the finite algebra with the $`\widehat{ϵ}_\pm `$ belonging to the affine Kac - Moody algebra $`\widehat{𝒢}`$. By construction the later (affine Toda) action describes an integrable perturbation of the $`𝒢`$-conformal abelian Toda model.
More interesting conformal (and affine) Toda models arise when the grading structure (defined by the operator Q) leads to non abelian zero grade subalgebras $`𝒢_0𝒢`$. In particular, if we supress one of the fundamental weights from $`Q`$, the zero grade subspace acquires a nonabelian structure $`sl(2)u(1)^{rank𝒢1}`$. Let us consider for instance $`Q=h^{}\widehat{d}+_{ia}^r\frac{2\lambda _iH}{\alpha _i^2}`$, where $`h^{}=0`$ or $`h^{}0`$ corresponding to the conformal or affine nonabelian (NA) - Toda models respectively. The absence of $`\lambda _a`$ in $`Q`$ prevents the contribution of the simple root step operators $`E_{\pm \alpha _a}^{(0)}`$ in constructing $`ϵ_\pm `$(or in $`\widehat{ϵ}_\pm `$), since the generators $`E_{\pm \alpha _a}^{(0)}`$ now belong to the zero grade $`𝒢_0`$. The form of the corresponding actions is as (2.7), but with $`ϵ_\pm `$ and $`Bsl(2)u(1)^{rank𝒢1}`$ as described above. They are known under the name nonsingular conformal (and affine) NA - Toda models. An important feature of these models is that they manifest an additional chiral (left and right) $`U(1)`$ \- symmetry. The algebraic origin of this fact is in the specific graded structure that allows the existence of U(1) - generator $`𝒢_0^0=YH\widehat{𝒢}_0`$ such that $`[YH,ϵ_\pm ]=0`$. The first term in the action (2.7) is invariant under local (chiral) $`𝒢_0^0`$ \- transformations by construction, while the invariance of the second (potential) term is a consequence of the defining property of $`YH`$ (i.e. to commute with the $`ϵ_\pm `$). The corresponding (chiral) conserved currents $`J_{YH}`$ and $`\overline{J}_{YH}`$ have the standard form $`J_{YH}=Tr[(YH)J]`$, and $`J=g^1g`$ and $`\overline{J}=\overline{}gg^1`$. By gauge fixing of this symmetry one obtains singular conformal (and affine ) NA - Toda models with the number of physical fields reduced by one. The elimination of the $`𝒢_0^0`$ field R in the framework of the Hamiltonian reduction consists in imposing of the nonlocal constraints $`J_{YH}=\overline{J}_{YH}=0`$. The standard method of incorporating these constraints in the functional integral (2.6) is to introduce auxiliary gauge fields $`A_0=(Y.H)a_0`$ and $`\overline{A}_0=(Y.H)\overline{a}_0`$ in the action (2.7) in such a way that the constraints appear as equations of motion of the improved action. As is well known the improvement consists in adding new $`A_0`$, $`\overline{A}_0`$\- dependent terms in the manner that the new action to be invariant under the following $`𝒢_0^0`$ -transformations
$`gg^{}=\alpha _0g\alpha _0^{},A_0A_0^{}=A_0\alpha _0^1\alpha _0,\overline{A}_0\overline{A}_0^{}=\overline{A}_0\overline{}\alpha _0^{}(\alpha _0^{})^1`$ (2.8)
There exist two inequivalent cases of gauge fixing of $`𝒢_0^0=U(1)`$ \- symmetry, namely the axial gauging where $`\alpha _0^{}=\alpha _0(z,\overline{z})𝒢_0^0`$ and the vector gauging, where $`\alpha _0^{}=\alpha _{0}^{}{}_{}{}^{1}(z,\overline{z})𝒢_0^0`$. Finaly, the improved action with all these properties has the form , :
$`S(B,A_0,\overline{A}_0)`$ $`=`$ $`S(g_0^f,A_0,\overline{A}_0)=S_{WZNW}(B)`$ (2.9)
$``$ $`{\displaystyle \frac{k}{2\pi }}{\displaystyle Tr\left(\pm A_0\overline{}BB^1+\overline{A}_0B^1B\pm A_0B\overline{A}_0B^1+A_0\overline{A}_0\right)d^2z}`$
$``$ $`{\displaystyle \frac{k}{2\pi }}{\displaystyle d^2zTr\widehat{ϵ}_+B\widehat{ϵ}_{}B^1}`$
where the $`\pm `$ signs in (2.9) correspond to axial/vector gaugings respectively. For generic affine Kac - Moody algebra $`𝒢_r`$, the zero grade element $`B𝒢_0`$ is parametrized as follow <sup>2</sup><sup>2</sup>2 All algebraic notations are as in ref.. The generators $`H_i,E_{\alpha _a}`$, when written without upper indices correspond to $`H_i^{(0)},E_{\alpha _a}^{(0)}`$ respectively. The extra generators $`\widehat{c}`$ and $`\widehat{d}`$ in the affine case with the properties $`[\widehat{d},E_{\alpha _a}^n]=nE_{\alpha _a}^n,[\widehat{d},H_{\alpha _a}^n]=nH_{\alpha _a}^n,[\widehat{c},E_{\alpha _a}^n]=[\widehat{c},\widehat{d}]=0`$ represent the center and the derivative respectively.
$`B=\mathrm{exp}(\stackrel{~}{\chi }E_{\alpha _a})\mathrm{exp}(RY^jH_j+\mathrm{\Phi }(H)+\nu \widehat{c}+\eta \widehat{d})\mathrm{exp}(\stackrel{~}{\psi }E_{\alpha _a})`$ (2.10)
where $`\mathrm{\Phi }(H)=_{j=1}^r_{i=1}^{r1}\phi _iX_{}^{j}{}_{i}{}^{}h_j`$ , $`YX_i=_{j=1}^rY^jX_{}^{j}{}_{i}{}^{}=0,i=1,\mathrm{},r1`$ and $`h_j=\frac{2\alpha _jH}{\alpha _j^2},j=1,\mathrm{},r`$, i.e. $`X_i`$ denotes r-dimensional vectors ortogonal to the particularly choosen vector Y. We next define the partition function of the reduced $`𝒢_0/𝒢_0^0=(sl(2)u(1)^{rank𝒢1})/u(1)`$ \- NA - Toda models :
$$Z_{sing}=DBDA_0D\overline{A}_0e^{S(B,A_0,\overline{A}_0)}$$
(2.11)
Integrating out the auxiliary fields $`A_0`$, $`\overline{A}_0`$ in (2.11) we derive the effective action $`S_{eff}(g_0^f)`$, where $`g_0^f𝒢_0/𝒢_0^0`$ is parametrized by the physical fields of the axial or vector singular conformal (or affine) NA - Toda models. The explicite form of the corresponding effective actions depends on the specific algebraic data - $`𝒢_r,Q,ϵ_\pm `$, i.e. on the algebra (finite or infinite) $`𝒢_r`$ of rank r, on the root $`\alpha _a`$ that is missed in Q, on the choice of $`YH=𝒢_0^0`$ (fixed by Q and the form of $`ϵ_\pm `$) and finaly on the way (axial or vector) the $`𝒢_0^0(=U(1))`$ is gauge fixed. An important property common for the entire family of conformal and affine singular (axial or vector) NA - Toda models with $`𝒢_0^0=U(1)`$ is their global U(1) Noether symmetry.<sup>3</sup><sup>3</sup>3 remember that in the case $`𝒢_0`$ is abelian an invariant subalgebra $`𝒢_0^0`$ does not exist and the corresponding affine abelian Toda models do not have any Noether symmetries. Examples of singular conformal $`𝒢_r`$ \- NA - Toda models of axial type have been constructed in refs., and . In the next two subsections we generalize the conformal constructions of refs. and to the case of (infinite ) affine algebras as well as for the case of vector gauging of $`𝒢_0^0`$. The integrable models obtained in this way represent the family of singular affine <sup>4</sup><sup>4</sup>4and conformal affine NA- Toda models when the generators $`\widehat{c}`$ and $`\widehat{d}`$ are included in the zero grade subalgebra NA - Toda models of axial and vector type with one global U(1) symmetry. Their main characteristic appears to be the fact that for imaginary coupling they admit U(1) - charged topological solitons with the electric and magnetic charges similar to the dyonic spectrum of certain 4-dimensional Yang-Mills-Higgs theory. This is the reason to call them dyonic integrable models . As we shall demonstrate in Sect.3 below the fact that the axial and vector IMs are obtained from the unique nonsingular IM with local U(1) symmetry by two different gauge fixings of this symmetry, gives rise to an interesting phenomena : pairs of integrable models related by the off-critical T-duality transformation.
### 2.1 Axial Gauging
Taking into account the fact that the action (2.9) is invariant under U(1) transformations (2.8) we can gauge away the nonlocal field $`R`$ by choosing $`\alpha _0^{}=\alpha _0=e^{\frac{1}{2}YHR}`$, that corresponds to axial gauge fixing. Then the gauge fixed element B (i.e. the factor group element $`g_0^f𝒢_0/𝒢_0^0`$) becomes
$$g_0^f=\mathrm{exp}(\chi E_{\alpha _a})\mathrm{exp}(\mathrm{\Phi }(H)+\nu \widehat{c}+\eta \widehat{d})\mathrm{exp}(\psi E_{\alpha _a})$$
(2.12)
where $`\chi =\stackrel{~}{\chi }e^{\frac{1}{2}Y\alpha _aR}`$ and $`\psi =\stackrel{~}{\psi }e^{\frac{1}{2}Y\alpha _aR}`$. With this parametrization the second ($`A_0,\overline{A}_0`$) - dependent term in the action (2.9) takes the form :
$`F_0`$ $`=`$ $`{\displaystyle \frac{k}{2\pi }}{\displaystyle Tr\left(A_0\overline{}g_0^f(g_0^f)^1+\overline{A}_0(g_0^f)^1g_0^f+A_0g_0^f\overline{A}_0(g_0^f)^1+A_0\overline{A}_0\right)d^2z}`$ (2.13)
$`=`$ $`{\displaystyle \frac{k}{2\pi }}{\displaystyle \left(a_0\overline{a}_02Y^2\mathrm{\Delta }_a(\frac{2\alpha _aY}{\alpha _a^2})(\overline{a}_0\psi \chi +a_0\chi \overline{}\psi )e^{\mathrm{\Phi }(\alpha _a)}\right)d^2z}`$
where $`\mathrm{\Delta }_a=1+\frac{(Y\alpha _a)^2}{\alpha _a^2Y^2}\psi \chi e^{\mathrm{\Phi }(\alpha _a)}`$ and $`[\mathrm{\Phi }(H),E_{\alpha _a}]=\mathrm{\Phi }(\alpha _a)E_{\alpha _a}`$. As we have mentioned the effective action is obtained by integrating over the auxiliary fields $`A_0`$ and $`\overline{A}_0`$ in the functional integral (2.11)
$$Z_{sing}=Dg_0^fDA_0D\overline{A}_0\mathrm{exp}(F_0)Dg_0^fe^{S_0S_{WZNW}(g_0^f)S_{int}}$$
(2.14)
where $`S_0=\frac{k}{2\pi }(\frac{2Y\alpha _a}{\alpha _a^2})^2d^2z\frac{\psi \chi \overline{}\psi \chi }{2Y^2\mathrm{\Delta }_a}e^{2\mathrm{\Phi }(\alpha _a)}`$ and $`S_{int}=\frac{k}{2\pi }d^2zTr[\widehat{ϵ}_+g_0^f\widehat{ϵ}_{}(g_0^f)^1]`$. The total effective action (2.9) for the axial IM is therefore given as
$`S_{eff}`$ $`=`$ $`{\displaystyle \frac{k}{4\pi }}{\displaystyle }d^2z(Tr(\mathrm{\Phi }(H)\overline{}\mathrm{\Phi }(H))+{\displaystyle \frac{4\overline{}\psi \chi e^{\mathrm{\Phi }(\alpha _a)}}{\alpha _a^2\mathrm{\Delta }_a}}`$ (2.15)
$`+`$ $`\eta \overline{}\nu +\nu \overline{}\eta 2Tr(\widehat{ϵ_+}g_0^f\widehat{ϵ_{}}(g_0^f)^1))`$
Note that the second term in (2.15) contains both symmetric and antisymmetric parts:
$`{\displaystyle \frac{e^{\mathrm{\Phi }(\alpha _a)}}{\mathrm{\Delta }_a}}\overline{}\psi \chi ={\displaystyle \frac{e^{\mathrm{\Phi }(\alpha _a)}}{\mathrm{\Delta }_a}}(g^{\mu \nu }_\mu \psi _\nu \chi +ϵ_{\mu \nu }_\mu \psi _\nu \chi )`$
where $`g_{\mu \nu }`$ is the 2-D metric of signature $`g_{\mu \nu }=diag(1,1)`$, $`=_0+_1\overline{}=_0_1`$. For $`r=1`$ ($`𝒢A_1^{(1)}`$, $`\mathrm{\Phi }(\alpha _1)`$ is zero) the antisymmetric term is a total derivative:
$`ϵ_{\mu \nu }{\displaystyle \frac{_\mu \psi _\nu \chi }{1+\psi \chi }}={\displaystyle \frac{1}{2}}ϵ_{\mu \nu }_\mu \left(\mathrm{ln}\left\{1+\psi \chi \right\}_\nu \mathrm{ln}{\displaystyle \frac{\chi }{\psi }}\right),`$
and it can be neglected. This $`A_1`$-NA-Toda model (in the conformal case), is known to describe the 2-D black hole solution for (2-D) string theory . The $`G`$-NA conformal Toda models can be used in the description of specific (r+1)-dimensional black string theories , with (r-1)-flat and 2-nonflat directions ($`g^{\mu \nu }G_{ab}(X)_\mu X^a_\nu X^b`$, $`X^a=(\psi ,\chi ,\phi _i)`$), containing axions ($`ϵ_{\mu \nu }B_{ab}(X)_\mu X^a_\nu X^b`$) and tachyons ($`\mathrm{exp}\left\{k_{ij}\phi _j\right\}`$), as well. One particular example of dyonic axial IM based on $`A_r^{(1)}`$ for $`a=1`$, i.e. $`YH=\lambda _1H`$ will be discussed in Sect.3.
### 2.2 Vector Gauging
The vector gauging is implemented by choosing $`\alpha _0^{}=\alpha _0^1=e^{\frac{1}{2}YHR^{}}`$, where $`\overline{\chi }=e^{\frac{1}{2}Y\alpha _aR^{}}\stackrel{~}{\chi }=e^{\frac{1}{2}Y\alpha _aR^{}}\stackrel{~}{\psi }`$ ,i.e. $`R^{}=\frac{1}{Y\alpha _a}ln\left(\frac{\stackrel{~}{\psi }}{\stackrel{~}{\chi }}\right)`$. Then the factor group element can be parametrized as
$`g_0^f=\mathrm{exp}(\overline{\chi }E_{\alpha _a})\mathrm{exp}(\mathrm{\Phi }(H)+\nu \widehat{c}+\eta \widehat{d})\mathrm{exp}(\overline{\chi }E_{\alpha _a})`$ (2.16)
where $`\mathrm{\Phi }(H)=YHR+_{j=1}^{r1}\phi _jX_j^ih_i`$. The second term of the action (2.9) then takes the form :
$`F_0`$ $`=`$ $`{\displaystyle \frac{k}{2\pi }}{\displaystyle }({\displaystyle \frac{2}{\alpha _a^2}}(Y\alpha _a)^2a_0\overline{a}_0\overline{\chi }^2e^{\mathrm{\Phi }(\alpha _a)}\overline{a}_0(Y^2R+{\displaystyle \frac{2Y\alpha _a}{\alpha _a^2}}\overline{\chi }\overline{\chi }e^{\mathrm{\Phi }(\alpha _a)})`$ (2.17)
$`+`$ $`a_0(Y^2\overline{}R+{\displaystyle \frac{2Y\alpha _a}{\alpha _a^2}}\overline{\chi }\overline{}\overline{\chi }e^{\mathrm{\Phi }(\alpha _a)})d^2z`$
Integrating over $`a_0`$ and $`\overline{a}_0`$ in (2.11) we derive the total effective action (2.9) for the vector gauged IM :
$`S_{eff}`$ $`=`$ $`{\displaystyle \frac{k}{4\pi }}{\displaystyle }({\displaystyle \underset{i,j=1}{\overset{r1}{}}}Tr(X_ihX_jh)\phi _i\overline{}\phi _j+Y^2R\overline{}R+\eta \overline{}\nu +\nu \overline{}\eta `$
$`+`$ $`2{\displaystyle \frac{2}{\alpha _a^2}}\rho ^2(\alpha _a){\displaystyle \frac{R\overline{}R}{\overline{\chi }^2}}e^{\mathrm{\Phi }(\alpha _a)}+2{\displaystyle \frac{2}{\alpha _a^2}}\rho (\alpha _a)(R\overline{}ln\overline{\chi }+\overline{}Rln\overline{\chi })2Tr(\widehat{ϵ_+}g_0^f\widehat{ϵ_{}}(g_0^f)^1))d^2z`$
where $`\rho (\alpha _a)=\frac{Y^2\alpha _a^2}{2(Y\alpha _a)}`$. Defining the new variables
$$E=e^{\gamma R},F=E^1(1\delta \overline{\chi }^2\mathrm{exp}\mathrm{\Phi }(\alpha _a))$$
(2.19)
the action (LABEL:actionvec) becomes
$`S_{eff}`$ $`=`$ $`{\displaystyle \frac{k}{2\pi }}{\displaystyle }({\displaystyle \frac{1}{2}}{\displaystyle \underset{i,j=1}{\overset{r1}{}}}Tr(X_ihX_jh)\phi _i\overline{}\phi _j+{\displaystyle \frac{1}{2\gamma ^2}}lnE\overline{}lnE(Y^2+4\gamma ^2\mathrm{\Gamma })`$ (2.20)
$``$ $`\mathrm{\Gamma }(lnE\overline{}\mathrm{\Phi }(\alpha _a)+\overline{}lnE\mathrm{\Phi }(\alpha _a))+{\displaystyle \frac{1}{2}}\eta \overline{}\nu +{\displaystyle \frac{1}{2}}\nu \overline{}\eta \mathrm{\Gamma }{\displaystyle \frac{(E\overline{}F+\overline{}EF)}{1EF}}`$
$``$ $`Tr(\widehat{ϵ_+}g_0^f\widehat{ϵ_{}}(g_0^f)^1))d^2z`$
where $`2\gamma \mathrm{\Gamma }=\frac{2}{\alpha _a^2}\rho (\alpha _a)`$ and $`2\delta \mathrm{\Gamma }=\frac{2}{\alpha _a^2}`$ are chosen in order to eliminate the variable $`\overline{\chi }`$. Notice that the $`E,F`$\- term in the action (2.20) is symmetric and contrary to the axial model the above action is CPT - invariant. The vector gauging, therefore provides a construction of torsionless actions T-dual to its axionic counterpart as we shall demonstrate in Sect. 3. This fact raises the question whether exist T-selfdual torsionless actions, i.e.when the axial and vector gauging leads to same action. It is worthwhile to mention that the dyonic integrable models of vector type (2.20) represent integrable perturbations of the conformal $`\sigma `$ \- models studied in in the context of the string backgrounds of black hole type .
#### 2.2.1 Example1. Torsionless $`B_r^{(1)}`$ model
We consider the particular case based on $`B_r^{(1)}`$ by taking $`Y=\frac{2\lambda _r}{\alpha _r^2}\frac{2\lambda _{r1}}{\alpha _{r1}^2}=e_r,a=r,\rho (\alpha _r)=\frac{1}{2}`$, where
$$\lambda _r=\frac{1}{2}(e_1+e_2+\mathrm{}+e_r),\lambda _{r1}=(e_1+e_2+\mathrm{}+e_{r1})$$
(2.21)
and
$`ϵ_\pm ={\displaystyle \underset{i=1}{\overset{r2}{}}}E_{\pm \alpha _i}^{(0)}+E_{\pm (\alpha _{r1}+\alpha _r)}^{(0)}+E_{(\alpha _1+2(\alpha _2+\mathrm{}\alpha _{r1}+\alpha _r))}^{(\pm 1)}`$
Parametrizing the Cartan part of $`g_0^f`$ as $`\mathrm{\Phi }(H)=_{i=1}^{r1}\phi _ie_iH+2(e_rH)lnE`$, we find
$`S_{eff}`$ $`=`$ $`{\displaystyle \frac{k}{2\pi }}{\displaystyle }({\displaystyle \frac{1}{2}}{\displaystyle \underset{i,j=1}{\overset{r1}{}}}\phi _i\overline{}\phi _i+{\displaystyle \frac{1}{2}}\eta \overline{}\nu +{\displaystyle \frac{1}{2}}\nu \overline{}\eta {\displaystyle \frac{(E\overline{}F+\overline{}EF)}{1EF}}`$ (2.22)
$``$ $`Tr(\widehat{ϵ_+}g_0^f\widehat{ϵ_{}}(g_0^f)^1))d^2z`$
where we have chosen $`\mathrm{\Gamma }=1,\gamma =\frac{1}{2},\delta =1`$. When the pair of fields $`\eta `$ and $`\nu `$ are taken to be equal to zero the corresponding conformal affine $`B_r^{(1)}`$ NA - Toda model takes the form of the nonconformal affine $`B_r^{(1)}`$ NA - Toda model.
#### 2.2.2 Example2. $`A_r^{(1)}`$ vector model
In order to find the explicit form of the singular affine NA - Toda IM based on the affine algebra $`A_r^{(1)}`$ we take $`Y=\lambda _1`$ and $`ϵ_\pm =_{i=2}^rE_{\pm \alpha _i}^{(0)}+E_{(\alpha _2+\mathrm{}\alpha _r)}^{(\pm 1)}`$. We will get simpler result if we parametrize $`\mathrm{\Phi }(H)=_{i=1}^r\varphi _ih_i`$ :
$`S_{eff}`$ $`=`$ $`{\displaystyle \frac{k}{4\pi }}{\displaystyle }({\displaystyle \underset{i,j=1}{\overset{r}{}}}\eta _{ij}\varphi _i\overline{}\varphi _j+2{\displaystyle \frac{\varphi _1\overline{}\varphi _1}{\overline{\chi }^2e^{2\varphi _1\varphi _2}}}`$ (2.23)
$`+`$ $`2(\varphi _1\overline{}ln\overline{\chi }+\overline{}\varphi _1ln\overline{\chi })2Tr(\widehat{ϵ_+}g_0^f\widehat{ϵ_{}}(g_0^f)^1))d^2z`$
With the following change of the field variables
$$E=e^{\varphi _1},F=E^1(1\overline{\chi }^2e^{2\varphi _1\varphi _2})$$
(2.24)
we find
$`S_{eff}={\displaystyle \frac{k}{2\pi }}{\displaystyle \left(\frac{1}{2}\underset{i,j=2}{\overset{r}{}}\eta _{ij}\varphi _i\overline{}\varphi _j\frac{1}{2}\frac{(E\overline{}F+\overline{}EF)}{1EF}Tr(\widehat{ϵ_+}g_0^f\widehat{ϵ_{}}(g_0^f)^1)\right)d^2z}`$ (2.25)
Introducing new fields $`c_i=e^{\varphi _{i+2}\varphi _{i+1}},i=1,\mathrm{}r1`$ we rewrite the action (2.25) in the form
$`S_{eff}`$ $`=`$ $`{\displaystyle \frac{k}{2\pi }}{\displaystyle }({\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{r1}{}}}(lnc_i\overline{}lnc_1c_{i+1}\mathrm{}c_{r1}+\overline{}lnc_ilnc_1c_{i+1}\mathrm{}c_{r1})`$ (2.26)
$``$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{(E\overline{}F+\overline{}EF)}{1EF}}Tr(\widehat{ϵ_+}g_0^f\widehat{ϵ_{}}(g_0^f)^1))d^2z`$
These are the vector IM’s studied in ref. .
## 3 Off-Critical T-Duality
T-duality is known to be an important property of the string theory. It acts as canonical transformation , in the string phase space $`𝒫=\{X^M(\sigma ),\mathrm{\Pi }_M(\sigma )=g_{MN}\dot{X}^N+b_M^NX_N^{};b,c\}`$ mapping the original conformal $`\sigma `$-model <sup>5</sup><sup>5</sup>5 under certain symmetry restrictions on the geometrical data:$`e_{MN}(X)=g_{MN}(X)+b_{MN}(X)`$ and $`\phi (X)`$. In the case of abelian T-duality $`e_{MN}(X)`$ and are $`\phi (X)`$ independent of $`dD`$ of $`X_\alpha `$, $`(\alpha =1,2,\mathrm{}dD)`$ called isometric target-space coordinates.:
$`S_\sigma ^{conf}={\displaystyle \frac{1}{4\pi \alpha ^{}}}{\displaystyle d^2z\left((g_{MN}(X)\eta ^{\mu \nu }+b_{MN}(X)ϵ^{\mu \nu })_\mu X^M_\nu X^N+\frac{\alpha ^{}}{2}R^{(2)}\phi (X)\right)}`$ (3.1)
($`\mu ,\nu =0,1;M,N=1,2,\mathrm{}D`$ and $`R^{(2)}`$ is the worldsheet curvature ) to its T-dual model $`\stackrel{~}{S}_\sigma ^{conf}(G_{MN}(\stackrel{~}{X}),B_{MN}(\stackrel{~}{X}),\varphi (\stackrel{~}{X}))`$. Curved string backgrounds with d-isometric directions provide an example of an abelian T-duality transformation,:
$`E_{\alpha \beta }`$ $`=`$ $`(e^1)_{\alpha \beta },E_{mn}=e_{mn}e_{m\alpha }(e^1)^{\alpha \beta }e_{\beta n}`$
$`E_{\alpha m}`$ $`=`$ $`(e^1)_\alpha ^\beta e_{\beta m},E_{m\alpha }=e_{m\beta }(e^1)_\alpha ^\beta `$
$`\varphi `$ $`=`$ $`\phi ln[detE_{\alpha \beta }],\alpha ,\beta =1,2,\mathrm{}d,m,n=d+1,\mathrm{}D`$ (3.2)
where $`E_{MN}=G_{MN}+B_{MN}`$. The canonical transformation $`(\mathrm{\Pi }_X,X)(\mathrm{\Pi }_{\stackrel{~}{X}},\stackrel{~}{X})`$ that generates the background maps (3.2) has the following simple form :
$`\mathrm{\Pi }_{\stackrel{~}{X}_\alpha }=X_\alpha ^{},\mathrm{\Pi }_{X_\alpha }=\stackrel{~}{X}_\alpha ^{}`$ (3.3)
and all the $`\mathrm{\Pi }_{X_m}`$ and $`X_m`$ remain unchanged. By construction both $`\sigma `$-models $`S_\sigma ^{conf}(e,\phi )`$ and $`\stackrel{~}{S}_\sigma ^{conf}(E,\varphi )`$ have coinciding energy spectrum and partition functions. The corresponding Lagrangeans are related by the generating function $``$:
$`(e,\phi )`$ $`=`$ $`(E,\varphi )+{\displaystyle \frac{d}{dt}},={\displaystyle \frac{1}{8\pi \alpha ^{}}}{\displaystyle 𝑑x\left(X\stackrel{~}{X}^{}X^{}\stackrel{~}{X}\right)},`$
$`{\displaystyle \frac{\delta }{\delta X^\alpha }}`$ $`=`$ $`\mathrm{\Pi }_{X_\alpha },{\displaystyle \frac{\delta }{\delta \stackrel{~}{X}^\alpha }}=\mathrm{\Pi }_{\stackrel{~}{X}_\alpha }`$ (3.4)
An important feature of the abelian T-duality (3.2) and (3.3) is that it maps the $`U(1)^d`$ Noether charges $`Q^\alpha =_{\mathrm{}}^{\mathrm{}}J_0^\alpha 𝑑x`$ of $`S_\sigma ^{conf}`$ into the topological charges $`\stackrel{~}{Q}_{top}^\alpha =_{\mathrm{}}^{\mathrm{}}𝑑x_x\stackrel{~}{X}^\alpha `$ of the dual model $`\stackrel{~}{S}_\sigma ^{conf}`$:
$`J_\mu ^\alpha `$ $`=`$ $`{\displaystyle \frac{1}{2}}e^{\alpha \beta }(X_n)_\mu X_\beta +{\displaystyle \frac{1}{2}}e^{\alpha m}(X_n)_\mu X_m{\displaystyle \frac{1}{2}}ϵ_{\mu \nu }^\nu \stackrel{~}{X}^\alpha `$
$`\stackrel{~}{J}_\mu ^\alpha `$ $`=`$ $`{\displaystyle \frac{1}{2}}E^{\alpha \beta }(\stackrel{~}{X}_n)_\mu \stackrel{~}{X}_\beta +{\displaystyle \frac{1}{2}}E^{\alpha m}(\stackrel{~}{X}_n)_\mu \stackrel{~}{X}_m{\displaystyle \frac{1}{2}}ϵ_{\mu \nu }^\nu X^\alpha `$ (3.5)
i.e. $`(Q^\alpha ,Q_{top}^\alpha )(\stackrel{~}{Q}_{top}^\alpha ,\stackrel{~}{Q}^\alpha )`$. It is well known , that the main part of the conformal $`\sigma `$-models representing relevant string backgrounds can be derived from the axial or vector gauged $`G/H`$-WZNW models. All the models constructed in Sect.2 with vanishing potential term $`V=\frac{m^2k}{2\pi }Tr(ϵ_+Bϵ_{}B^1)`$ (i.e. $`m=0`$ ) are of this type. They have $`d=r`$ isometric directions, i.e. $`e_{mn}`$ do not depend on $`\phi _i,(i=2,\mathrm{},r)`$ and $`\theta =\frac{1}{2}ln\frac{\psi }{\chi }`$. The T-duality group in this case is known to be $`O(r,r|Z)`$ (see for example ). Adding the potential $`V`$ with $`ϵ_\pm =_{ia}E_{\pm \alpha _i}^{(0)}`$ (see eqs. (2.15) and (2.20)) specific for the conformal NA -Toda theories with one global U(1) symmetry, we are decreasing the number of the isometric coordinates from $`d=r`$ to $`d_0=1`$. Taking $`ϵ_\pm =_{ia,b}E_{\pm \alpha _i}^{(0)}`$ one can construct NA -Toda theories with $`d_0=2`$, etc.
The problem we are addressing in this section is about the T-duality between $`d_0=1`$ axial and vector integrable (nonconformal) models of Sect.2.1 and Sect.2.2 with potential terms constructed by taking
$`\widehat{ϵ}_\pm `$ $`=`$ $`{\displaystyle \underset{i=2}{\overset{r}{}}}E_{\pm \alpha _i}^{(0)}+E_{(\alpha _2+\mathrm{}+\alpha _r)}^{(\pm 1)}`$
$`V_a`$ $`=`$ $`{\displaystyle \frac{m^2k}{2\pi }}\left({\displaystyle \underset{i=2}{\overset{r}{}}}e^{\phi _{i2}+\phi _i2\phi _{i1}}+e^{\phi _2+\phi _r}(1+\psi \chi e^{\phi _2})\right),\phi _0=\phi _{r+1}=0`$
$`V_{vec}`$ $`=`$ $`{\displaystyle \frac{m^2k}{2\pi }}\left({\displaystyle \underset{i=3}{\overset{r}{}}}e^{\varphi _{i1}+\varphi _{i+1}2\varphi _i}+Ee^{2\varphi _2+\varphi _3}+Fe^{\varphi _2+\varphi _r}\right),`$ (3.6)
for the $`A_r^{(1)}`$ model of Example 2. For the $`B_r^{(1)}`$ model (with $`\eta =0`$) of Example 1 (see also Sect.4.2) we have
$`\widehat{ϵ}_\pm ={\displaystyle \underset{i=1}{\overset{r2}{}}}E_{\pm \alpha _i}^{(0)}+E_{\pm (\alpha _{r1}+\alpha _r)}^{(0)}+E_{\pm \alpha _0}^{(\pm 1)},\alpha _0+\alpha _1+2(\alpha _2+\mathrm{}\alpha _{r1}+\alpha _r)=0`$ (3.7)
with potential given by eq.(4.7)for the axial model. Its vector counterpart is given by
$`V_{vec}^{B_r^1}={\displaystyle \underset{i=1}{\overset{r2}{}}}|c_i|^2e^{\phi _i+\phi _{i+1}}+2|c_{r1}|^2(2EF1)e^{\phi _{r1}}+|c_r|^2e^{\phi _1+\phi _2}`$
In both cases the axial IM’s isometric coordinate is $`X=\theta =\frac{1}{2}ln\frac{\chi }{\psi }`$ ($`\stackrel{~}{u}^2=\frac{\psi }{\chi }`$). For the corresponding vector IM’s we choose $`\stackrel{~}{X}_{B_r^{(1)}}=R_B=2lnE`$ and $`\stackrel{~}{X}_{A_r^{(1)}}=R_A=\frac{r+1}{r}lnE`$ as isometric coordinates. It is important to mention that in the case of $`A_r^{(1)}`$ vector model the canonical transformation (3.3) with $`d=1`$ has to be accompanied by the following point transformation:
$`\varphi _j=\varphi _j^{}{\displaystyle \frac{rj+1}{2r}}R_A,j=2,\mathrm{}r`$ (3.8)
Then performing $`d_0=1`$ T-duality transformation (3.2) (together with the $`A_r^{(1)}`$ fields transformation (3.8)) we realize that $`_{vec}`$ and $`_a`$ given by eqs. (2.20) and (2.15), with potentials (3.6) and (4.7), are related by eq.(3.4) with
$`{\displaystyle \frac{d}{dt}}=2\mathrm{\Gamma }(ln\stackrel{~}{u}\overline{}lnE\overline{}ln\stackrel{~}{u}lnE)`$
Notice that the $`B_r^{(1)}`$ vector and axial Lagrangeans have the same form, i.e. they are T-selfdual.
An alternative way to perform the T-duality transformation between $`d_0=1`$ axial and vector IM’s in consideration consists in making the following nonlocal change of the field variables:
a)$`A_r^{(1)}`$ case
$`E=e^{\frac{r}{r+1}R_A},F=e^{\frac{r}{r+1}R_A}(1+\psi \chi e^{\phi _2}),\varphi _j=\phi _j+{\displaystyle \frac{rj+1}{r}}R_A,j=2,\mathrm{}r`$ (3.9)
b)$`B_r^{(1)}`$ case
$`E=e^{\frac{1}{2}R_B},F=e^{\frac{1}{2}R_B}(1+\psi \chi ),\varphi _j=\phi _j`$ (3.10)
instead of the canonical transformation(3.3) ( resulting in (3.2)). Eqs. (3.9) and (3.10) in fact represent the integrated form of (3.3). Their derivation (see Sect.5 of ref. ) is based on the comparison of the $`g_0`$ (or $`B`$) group elements written in axial and vector parametrizations (2.12) and (2.16), i.e. imposing $`g_0^{vec}=g_0^{ax}`$. An important ingredient of this calculation are the relations (3.5) between the $`U(1)`$ currents and the topological currents $`ϵ^{\mu \nu }_\nu R`$ (and $`ϵ^{\mu \nu }_\nu \theta `$). Note that $`R=2lnE`$ is a nonlocal (nonphysical) field in the axial model, but it appears to be physical in the vector model. In the case of the $`B_r^{(1)}`$ model (2.22) the $`U(1)`$ topological-currents relations (3.5) take the following explicit form:
$`ln\stackrel{~}{u}`$ $`=`$ $`{\displaystyle \frac{v}{v1}}lnE{\displaystyle \frac{1}{2}}{\displaystyle \frac{v}{v1}}`$
$`\overline{}ln\stackrel{~}{u}`$ $`=`$ $`{\displaystyle \frac{v}{v1}}\overline{}lnE+{\displaystyle \frac{1}{2}}{\displaystyle \frac{\overline{}v}{v1}}`$ (3.11)
where $`v=EF`$ and $`\stackrel{~}{u}^2=\frac{\chi }{\psi }`$.
Although the T-duality between the vector and axial integrable models is quite similar to the conformal ”free” case (i.e. $`V=0`$ )<sup>6</sup><sup>6</sup>6 the only new feature is that one should take care about the specific ”point” transformations involving the potential $`V`$ and that the isometric coordinates are reduced from $`d=r`$ to $`d_0=1`$. the off-critical T-duality addresses few new problems specific for the integrable models. In the case of imaginary coupling constant $`\beta ^2=\frac{2\pi }{k}`$, i.e. $`\beta i\beta _0`$ and $`\phi _ki\beta _0\phi _k`$, $`\psi i\beta _0\psi `$, etc. one expects that both axial and vector IM’s possess soliton solutions. One might wonder what is the relation between the solitons (and breathers ) of the T-dual integrable models, whether their soliton spectra coincides (modulo the interchanges $`Q\stackrel{~}{Q}_{top},\stackrel{~}{Q}Q_{top}`$ ) and finally about the $`O(1,1|Z)`$ symmetry of the solitons energies and massess. Partial answer of all these questions is presented in our recent work .
## 4 No Torsion Theorem
It becomes clear from our discussion in Sect.3 that the abelian T-duality between the axial and the vector dyonic IM’s allows to single out two classes of models : (a) T-selfdual torsionless IM’s and (b) pairs of T-duals axial and vector IM’s. One may wonder what is the Lie algebraic condition defining the models of class (a). In answering this question we first establish the T-selfduality condition for the conformal $`d_0=1`$ (i.e $`𝒢_0^0=U(1)`$) singular NA Toda models. We next generalize the conformal no torsion (T-selfduality) theorem to the case of U(1) symmetric affine NA Toda IM’s of sects. 4.2 and 4.3.
### 4.1 Conformal T-selfdual NA Toda models
T-selfduality requires that axial and vector models have coinciding Hamiltonians. Since the vector models are by construction torsionless the question to be addressed is about the algebraic condition under which the axial gauging generate torsionless models as well.
Consider a finite dimensional Lie algebra $`𝒢`$ with grading operator given by $`Q_a=_{ia}^r\frac{2}{\alpha _i^2}\lambda _iH`$ and take the most general constant generators of grade $`\pm 1`$, i. e.,
$`ϵ_\pm ={\displaystyle \underset{ia}{\overset{r}{}}}c_{\pm i}E_{\pm \alpha _i}+b_\pm E_{\pm (\alpha _a+\alpha _{a+1})}+d_\pm E_{\pm (\alpha _a+\alpha _{a1})}.`$ (4.1)
It is clear that if $`c_{\pm i},b_\pm ,d_\pm 0`$, there shall be no $`𝒢_0^0`$ commuting with $`ϵ_\pm `$, since that requires an orthogonal direction to all roots appearing in $`ϵ_\pm `$. These are the generalized non-singular NA-Toda models of ref. ,. The NA-Toda models of singular metric $`G_{ij}(X)`$ correspond to the cases when $`𝒢_0^0=U(1)`$ ( or $`U(1)^r`$ in general ) and $`J_{YH}=0=\overline{J}_{YH}`$ imposed as a subsidiary constraint <sup>7</sup><sup>7</sup>7If we leave $`𝒢_0^0`$ unconstrained the resulting model belongs again to the non singular NA-Toda class of models ,. From the string theory point of view the class of singular NA Toda models are of great interest since they describe strings on specific black hole backgrounds.. Depending on the choice of the constants $`c_{\pm i},b_\pm `$ and $`d_\pm `$ we distinguish four families of singular conformal NA-Toda models: (i)$`b_\pm =d_\pm =0`$, $`𝒢_0^0=\frac{2}{\alpha _a^2}\lambda _aH`$;
(ii)$`c_{\pm (a1)}=c_{\pm (a+1)}=0`$, $`𝒢_0^0=\frac{2}{\alpha _a^2}\lambda _aH\frac{2}{\alpha _{a1}^2}\lambda _{a1}H\frac{2}{\alpha _{a+1}^2}\lambda _{a+1}H`$;
(iii)$`c_{\pm (a+1)}=d_\pm =0`$, $`𝒢_0^0=\frac{2}{\alpha _a^2}\lambda _aH\frac{2}{\alpha _{a+1}^2}\lambda _{a+1}H`$;
(iv)$`b_\pm =c_{\pm (a1)}=0`$, $`𝒢_0^0=\frac{2}{\alpha _a^2}\lambda _aH\frac{2}{\alpha _{a1}^2}\lambda _{a1}H`$.
Of course, if $`c_{\pm j}=0`$, $`ja,a\pm 1`$, we find $`𝒢_0^0=\lambda _jH`$. However, since $`[\lambda _jH,E_{\pm \alpha _a}]=0`$, there will be no singular metric present and this case shall be neglected. Cases (i) and (ii) are equivalent, since they are related by the Weyl reflection $`\sigma _{\alpha _a}(\alpha _{a\pm 1})=\alpha _a+\alpha _{a\pm 1}`$ and the corresponding fields are related by non-linear change of the variables. This case has already been discussed in refs. and , and shown to present always the antisymmetric term, originated by the presence of $`e^{k_{ai}\phi _i}`$ in $`\mathrm{\Delta }_a`$ and in the kinetic term as well. Since we are removing all dependence in $`𝒢_0^0`$, when parameterizing $`g_0^f`$, cases (iii) and (iv) may be studied together with
$`g_0^f=\mathrm{exp}(\chi E_{\alpha _a})\mathrm{exp}(\mathrm{\Phi }(H))\mathrm{exp}(\psi E_{\alpha _a})`$ (4.2)
where $`\mathrm{\Phi }(H)=_{i=1}^{a2}\phi _ih_i+\phi _{}(\chi _{}H)+\phi _+(\chi _+H)+_{i=a+2}^r\phi _ih_i`$,
$`\chi _{}^{(iii)}=\alpha _{a1}+\alpha _a,\chi _+^{(iii)}=\alpha _{a+1},,\chi _{}^{(iv)}=\alpha _{a1},\chi _+^{(iv)}=\alpha _a+\alpha _{a+1}`$ (4.3)
for cases $`(iii)`$ and $`(iv)`$ respectively, and $`𝒢_0^0=YH`$, such that $`Tr(\chi _\pm H𝒢_0^0)=0`$. Such parametrization of $`g_0^f`$ yields
$`\mathrm{\Phi }(\alpha _a)={\displaystyle \underset{i=1}{\overset{a2}{}}}k_{ai}\phi _i+(\alpha _a\chi _{})\phi _{}+(\alpha _a\chi _+)\phi _++{\displaystyle \underset{i=a+2}{\overset{r}{}}}k_{ai}\phi _i`$
Now, if we consider Lie algebras whose Dynkin diagrams connect only nearest neighbours, i. e., $`\mathrm{\Phi }(\alpha _a)=(\alpha _a\chi _{})\phi _{}+(\alpha _a\chi _+)\phi _+`$, then the “no-torsion condition” implies $`\mathrm{\Phi }(\alpha _a)=0`$. Considering case (iii), we have
$`\alpha _a\chi _{}=\alpha _a(\alpha _{a1}+\alpha _a)=0,\alpha _a\chi _+=\alpha _a(\alpha _{a+1})=0.`$ (4.4)
In this case, the only solution for both equations is to take $`a=r`$ (in such a way that $`\alpha _{r+1}=0`$) and $`𝒢=B_r`$ (so that $`\alpha _r^2=\alpha _{r1}\alpha _r=1`$). This is precisely the case proposed by Leznov and Saveliev and subsequently discussed by Gervais and Saveliev and also by Bilal , for the particular case of $`B_2`$. For case (iv), the “no-torsion condition” requires that
$`\alpha _{a1}\alpha _a=0,\alpha _a(\alpha _a+\alpha _{a+1})=0,`$
Both are satisfied for $`a=1`$ , $`𝒢=C_2`$, since $`\alpha _{a1}=0`$ and also $`\alpha _1^2=\alpha _1\alpha _2=1`$, respectively.
In general, the “no-torsion condition” (T-selfduality), i. e., $`\mathrm{\Phi }(\alpha _a)=0`$, may be expressed in terms of the structure of the co-set $`𝒢_0/𝒢_0^0=\frac{u(1)^{r1}sl(2)}{u(1)}`$. The crucial ingredient for the appearence of $`\mathrm{\Phi }(\alpha _a)`$ arises from the conjugation
$`Tr(A_0g_0^f\overline{A}_0(g_0^f)^1+A_0\overline{A}_0)=2\lambda _a^2\left(1+{\displaystyle \frac{2}{\alpha _a^2}}{\displaystyle \frac{\chi \psi \mathrm{exp}(\mathrm{\Phi }(\alpha _a))}{2\lambda _a^2}}\right).`$
Henceforth, if all generators belonging to the Cartan subalgebra parameterizing $`g_0^f`$ commute with $`E_{\pm \alpha _a}`$, then $`\mathrm{\Phi }(\alpha _a)=0`$, and therefore the structure of the co-set
$`{\displaystyle \frac{G_0}{G_0^0}}={\displaystyle \frac{u(1)^{r1}sl(2)}{u(1)}}=u(1)^{r1}{\displaystyle \frac{sl(2)}{u(1)}}`$ (4.5)
is the general condition for the absence of the antisymmetric term in the action.
Summarizing, for finite Lie algebras, it was shown that the absence of the antisymmetric terms in the action of the axial NA - Toda models can only occur for $`𝒢=B_r`$, $`a=r`$ and $`ϵ_\pm =_{i=1}^{r2}c_{\pm i}E_{\pm \alpha _i}+d_\pm E_{\pm (\alpha _r+\alpha _{r1})}`$. In such case, $`𝒢_0^0`$ is generated by $`YH=(\frac{2\lambda _r}{\alpha _r^2}\frac{2\lambda _{r1}}{\alpha _{r1}^2})H`$ and $`\mathrm{\Phi }(H)=_{i=1}^{r2}\phi _ih_i+\phi _{}(\alpha _{r1}+\alpha _r)H`$. Due to the root structure of $`B_r`$, we verify that $`\mathrm{\Phi }(\alpha _r)=\alpha _r(\alpha _{r1}+\alpha _r)\phi _{}=0`$.
In order to extend the no torsion theorem (T-selfduality) to the case of infinite affine Lie algebras (i.e for the integrable perturbations of the conformal models of sect. 4.1) let we consider $`h^{}=1_{ir}\frac{2}{\alpha _i^2}\lambda _i\alpha _0`$ , where $`\alpha _0`$ is the highest root of $`𝒢`$ such that $`\alpha _0(\frac{2\lambda _r}{\alpha _r^2}\frac{2\lambda _{r1}}{\alpha _{r1}^2})=0`$ and the gradation $`Q_a(h^{})`$ that preserves the zero grade subalgebra $`𝒢_0`$, (apart from $`\widehat{c}`$ and $`\widehat{d}`$). We choose the cooresponding (affine) grade $`\pm `$ elements in the form $`\widehat{ϵ}_+=ϵ_++E_{\alpha _0}^{(1)}`$. Since conformal and the affine models differ only by the potential term, the solution for the no torsion condition is also satisfied for infinite dimensional algebras, whose Dynkin diagram possess a $`B_r`$-“tail like”. An obvious solution is the untwisted $`B_r^{(1)}`$ model. Two new other solutions are given by the twisted affine Kac-Moody algebras $`A_{2r}^{(2)}`$ and $`D_{r+1}^{(2)}`$ as we shall describe in detail in the next two subsections.
### 4.2 The $`B_r^{(1)}`$ Torsionless Affine NA Toda models
In order to generalize the $`B_r`$ conformal models to the affine $`B_r^{(1)}`$ ones we take
$`Q=2(r1)\widehat{d}+_{i=1}^{r1}\frac{2\lambda _iH}{\alpha _i^2}`$, which decomposes $`B_r^{(1)}`$ into graded subspaces. In particular, the zero grade subspace turns out to be $`𝒢_0=SL(2)U(1)^{r1}`$, generated by ($`E_{\pm \alpha _r}^{(0)},h_1,\mathrm{},h_r`$). Following the arguments of the conformal no torsion theorem of the sect. 4.1 , we have to choose
$`\widehat{ϵ_\pm }={\displaystyle \underset{i=1}{\overset{r2}{}}}c_{\pm i}E_{\pm \alpha _i}^{(0)}+c_{\pm (r1)}E_{\pm (\alpha _{r1}+\alpha _r)}^{(0)}+c_{\pm r}E_{\pm \alpha _0}^{(\pm 1)}`$
where $`\alpha _0`$ (defined by $`\alpha _0+\alpha _1+2(\alpha _2+\mathrm{}+\alpha _{r1}+\alpha _r)=0`$ ) is the highest root of $`B_r`$ and $`𝒢_0^0`$ is generated by $`YH=(\frac{2\lambda _r}{\alpha _r^2}\frac{2\lambda _{r1}}{\alpha _{r1}^2})H`$ ( such that $`[YH,\widehat{ϵ_\pm }]=0`$). The coset $`𝒢_0/𝒢_0^0`$ (with $`\eta =0`$) is then parametrized according to (2.12) with $`\mathrm{\Phi }(H)=_{i=1}^{r1}_i\phi _i`$, where $`_i=(\alpha _r+\mathrm{}\alpha _i)H`$ so that $`Tr(_i_j)=\delta _{ij},i,j=1,\mathrm{},r1`$ and the total effective action becomes
$$S=\frac{k}{4\pi }d^2z\left(\frac{1}{4}\underset{i=1}{\overset{r1}{}}g^{\mu \nu }_\mu \phi _i_\nu \phi _i+g^{\mu \nu }\frac{_\mu \psi _\nu \chi }{1+\psi \chi }2V\right)$$
(4.6)
The “affine potential” V for $`(n>2)`$ has the form
$$V=\underset{i=1}{\overset{r2}{}}|c_i|^2e^{\phi _i+\phi _{i+1}}+2|c_{r1}|^2(1+2\psi \chi )e^{\phi _{r1}}+|c_r|^2e^{\phi _1+\phi _2}$$
(4.7)
The case $`r=2`$, i.e. $`\widehat{𝒢}=\widehat{S}O(5)`$ has to be considered separately. Choosing now the grade $`\pm 1`$ constant elements as $`\widehat{ϵ_\pm }=E_{\alpha _1+\alpha _2}^{(0)}+E_{\alpha _1\alpha _2}^{(1)}`$, denoting $`\mathrm{\Phi }(\alpha _{r1})=\phi `$ and by further changing the variables
$`\psi i\psi ;\chi i\psi ^{};\phi i\phi `$
we derive the following real action
$$S=\frac{k}{4\pi }d^2z\left(\frac{g^{\mu \nu }_\mu \psi _\nu \psi ^{}}{(1\psi \psi ^{})}+\frac{1}{4}g^{\mu \nu }_\mu \phi _\nu \phi +8(12\psi \psi ^{})cos\phi \right)$$
(4.8)
### 4.3 The twisted affine NA Toda Models
The twisted affine Kac-Moody algebras are constructed from a finite dimensional algebra possessing a nontrivial symmetry of their Dynkin diagrams (this procedure is known as folding). Such symmetry can be extended to the algebra by an outer automorphism $`\sigma `$ , as
$$\sigma (E_\alpha )=\eta _\alpha E_{\sigma (\alpha )}$$
(4.9)
where $`\eta _\alpha =\pm 1`$. For the simple roots, $`\eta _{\alpha _i}=1`$. The signs can be consistently assigned to all generators since nonsimple roots can be written as sum of two other roots.
The no torsion theorem requires a $`B_r`$-“tail like” structure which is fulfilled only by the $`A_{2r}^{(2)}`$ and $`D_{r+1}^{(2)}`$ (see appendix N of ref. ). In both cases the automorphism is of order 2 (i.e. $`\sigma ^2=1`$). Let us denote by $`\alpha `$ the roots of the untwisted algebra $`𝒢`$. For the $`A_{2r}^{(2)}`$ case, the automorphism is defined by
$$\sigma (\alpha _1)=\alpha _{2r},\sigma (\alpha _2)=\alpha _{2r1}\mathrm{},\sigma (\alpha _{r1})=\alpha _r$$
(4.10)
whilst for the $`D_{r+1}^{(2)}`$, the automorphism acts only in the “fish tail” of the Dynkin diagram of $`D_{r+1}`$, i.e.
$$\sigma (E_{\alpha _1})=E_{\alpha _1},\mathrm{},\sigma (E_{\alpha _{r1}})=E_{\alpha _{r1}},\sigma (E_{\alpha _r})=E_{\alpha _{r+1}}$$
(4.11)
The automorphism $`\sigma `$ defines a decomposition of the algebra $`𝒢=𝒢_{even}𝒢_{odd}`$. The twisted affine algebra is then constructed from $`𝒢`$ assigning an affine index $`mZ`$ to the generators in $`𝒢_{even}`$ while $`mZ+\frac{1}{2}`$ to those in $`𝒢_{odd}`$ (see appendix N of ).
The simple root step operators for $`A_{2r}^{(2)}`$ are
$$E_{\beta _i}=E_{\alpha _i}^{(0)}+E_{\alpha _{2ri+1}}^{(0)},i=1,\mathrm{},rE_{\beta _0}^{(\frac{1}{2})}=E_{\alpha _1\mathrm{}\alpha _{2r}}^{(\frac{1}{2})}$$
(4.12)
corresponding to the simple and highest roots
$$\beta _i=\frac{1}{2}(\alpha _i+\alpha _{2ri+1})i=1,\mathrm{},r,\alpha _0=\alpha _1+\mathrm{}+\alpha _{2r}=2(\beta _1+\mathrm{}\beta _r)$$
(4.13)
respectively. For $`D_{r+1}^{(2)}`$, the simple root step operators are
$`E_{\beta _i}=E_{\alpha _i}^{(0)},i=1,\mathrm{},r1,E_{\beta _r}=E_{\alpha _r}^{(0)}+E_{\alpha _{r+1}}^{(0)},`$
$`E_{\beta _0}^{(\frac{1}{2})}=E_{\alpha _1\mathrm{}\alpha _{r1}\alpha _{r+1}}^{(\frac{1}{2})}E_{\alpha _1\mathrm{}\alpha _{r1}\alpha _r}^{(\frac{1}{2})}`$ (4.14)
corresponding to the simple and highest roots
$$\beta _i=\alpha _ii=1,\mathrm{},r1,\beta _r=\frac{1}{2}(\alpha _r+\alpha _{r+1}),\alpha _0=\alpha _1+\mathrm{}\alpha _{r1}+\frac{1}{2}(\alpha _r+\alpha _{r+1})=\beta _1+\mathrm{}\beta _r$$
(4.15)
where have denoted by $`\beta `$ the roots of the twisted (folded) algebra.
The corresponding torsionless affine NA Toda models are defined by introducing the following grading operators :
$$Q_{A_{2r}^{(2)}}=2(2r1)\widehat{d}+\underset{ir,r+1}{\overset{2r}{}}\frac{2\lambda _iH}{\alpha _i^2},Q_{D_{r+1}^{(2)}}=(2r2)\widehat{d}+\underset{i=1}{\overset{r}{}}\frac{2\lambda _iH}{\alpha _i^2}$$
(4.16)
for $`A_{2r}^{(2)}`$ and $`D_{r+1}^{(2)}`$ respectively, where $`\lambda _i`$ are the fundamental weights of the untwisted algebra $`𝒢`$, i.e. $`\frac{2\lambda _i\alpha _j}{\alpha _j^2}=\delta _{ij}`$. Both models are specified by the constant grade $`\pm 1`$ operators $`\widehat{ϵ}_\pm `$
$$\widehat{ϵ}_\pm =\underset{i=1}{\overset{r2}{}}c_{\pm i}E_{\pm \beta _i}+c_{\pm (r1)}E_{\pm (\beta _{r1}+\beta _r)}+c_{\pm r}E_{\beta _0}^{(\pm \frac{1}{2})}$$
(4.17)
where $`\beta _i`$ are the simple roots of the twisted affine algebra specified in (4.13) and in (4.15).
The grading operators (4.16) determine the zero grade subalgebra in both cases to be $`𝒢_0=SL(2)U(1)^{r1}`$ generated by $`E_{\pm \beta _r}^{(0)},h_1,\mathrm{},h_r`$. Hence the zero grade subgroup is parametrized as in (4.2) with $`\eta =0=\nu `$ . The factor group $`𝒢/𝒢_0^0`$ is given in (2.12), where $`𝒢_0^0`$ is generated by $`YH=(\frac{2\mu _r}{\beta _r^2}\frac{2\mu _{r1}}{\beta _{r1}^2})H`$ and $`\mu _i`$ are the fundamental weights of the twisted algebra i.e. $`\frac{2\mu _i\beta _j}{\beta _j^2}=\delta _{ij}`$ . In order to decouple the $`\phi _i,i=1,\mathrm{},r1`$ we choose an orthonormal basis for the Cartan subalgebra( i.e. $`\mathrm{\Phi }(H)=_i\phi _i`$) where
$$_i=(\alpha _i+\mathrm{}\alpha _{2ri+1})H,YH=_r,Tr(_i_j)=2\delta _{ij},i,j=1,\mathrm{}r$$
(4.18)
and
$$_i=(\alpha _{ri+1}+\mathrm{}+\alpha _{r+1})H,YH=_r,Tr(_i_j)=\delta _{ij},i,j=1,\mathrm{}r$$
(4.19)
for $`A_{2r}^{(2)}`$ and $`D_{r+1}^{(2)}`$ respectively.
The Lagrangeans of the corresponding (axial gauged IMs), based on the twisted affine algebras in consideration, are obtained from (2.15) (with $`\eta =0`$) and the $`ϵ_\pm `$ given by (4.17). Up to a multiplicative factor$`\frac{k}{2\pi }`$ they are given by
$$_{A_{2r}^{(2)}}=\frac{\chi \overline{}\psi }{1+\frac{1}{2}\psi \chi }+\frac{1}{2}\underset{i=1}{\overset{r1}{}}\phi _i\overline{}\phi _iV_{A_{2r}^{(2)}}$$
(4.20)
and
$$_{D_{r+1}^{(2)}}=2\frac{\chi \overline{}\psi }{1+\psi \chi }+\frac{1}{2}\underset{i=1}{\overset{r1}{}}\phi _i\overline{}\phi _iV_{D_{r+1}^{(2)}}$$
(4.21)
The potentials of these twisted affine singular NA - Toda models have the form
$$V_{A_{2r}^{(2)}}=\underset{i=1}{\overset{r2}{}}|c_i|^2e^{\phi _i+\phi _{i+1}}+\frac{1}{2}|c_r|^2e^{2\phi _1}+|c_{r1}|^2e^{\phi _{r1}}(1+\psi \chi )$$
(4.22)
and
$$V_{D_{r+1}^{(2)}}=\underset{i=1}{\overset{r2}{}}|c_i|^2e^{\phi _i+\phi _{i+1}}+\frac{1}{2}|c_r|^2e^{\phi _1}+|c_{r1}|^2e^{\phi _{r1}}(1+2\psi \chi )$$
(4.23)
The T -selfdual models described by (4.6), (4.20) and (4.21) turns out to coincide with those proposed by Fateev in .
## 5 Zero Curvature
The equations of motion for the NA Toda models are known to be of the form
$$\overline{}(B^1B)+[\widehat{ϵ_{}},B^1\widehat{ϵ_+}B]=0,(\overline{}BB^1)[\widehat{ϵ_+},B\widehat{ϵ_{}}B^1]=0$$
(5.1)
The subsidiary constraint $`J_{YH}=Tr(B^1BYH)=\overline{J}_{YH}=Tr(\overline{}BB^1YH)=0`$ can be consistenly imposed since $`[YH,\widehat{ϵ_\pm }]=0`$ (as can be seen from (5.1) by taking the trace with $`Y.H`$). We next consider the axial models only. As we have mentioned in Sect.3 the vector models Lagrangeans (and equations of motion) can be obtained from the axial ones by the nonlocal change of the fields (see for example eqs.(3.9) and (3.10) for the $`A_r^{(1)}`$ and $`B_r^{(1)}`$ models). Solving those equations for the nonlocal field $`R`$ yields
$$R=(\frac{Y\alpha _r}{Y^2})\frac{\psi \chi }{\mathrm{\Delta }}e^{\mathrm{\Phi }(\alpha _r)},\overline{}R=(\frac{Y\alpha _r}{Y^2})\frac{\chi \overline{}\psi }{\mathrm{\Delta }}e^{\mathrm{\Phi }(\alpha _r)}$$
(5.2)
The equations of motion for the fields $`\psi ,\chi `$ and $`\phi _i,i=1,\mathrm{},r1`$ (obtained from (5.1) by imposing the constraints (5.2)) coincide precisely with the Euler-Lagrange equations derived from (4.20) and (4.21). Alternatively, (5.1) admits a zero curvature representation
$`\overline{A}\overline{}A+[A,\overline{A}]=0`$, where
$$A=\widehat{ϵ_{}}+B^1B,\overline{A}=B^1\widehat{ϵ_+}B$$
(5.3)
Whenever the constraints (5.2) are incorporated into $`A`$ and $`\overline{A}`$ in (5.3), equations (5.1) yields the zero curvature representation of the affine singular NA -Toda models. Such argument is valid for all conformal, affine and conformal affine NA - Toda models, in particular for the torsionless class of models discussed in the previous section.
Using the explicit parametrization of $`B`$ given in (2.10),(with $`\eta =0`$) the corresponding $`\widehat{ϵ_\pm }`$ specified in (4.17), (4.13) and (4.15) together with (5.2)(with $`Y`$ given by eqs. (4.18) and (4.19)), we obtain the flat connections $`A`$ and $`\overline{A}`$ in the following form :
(a) the $`A_{2r}^{(2)}`$ affine NA - Toda model
$`A_{A_{2r}^{(2)}}`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{r2}{}}}c_i(E_{\alpha _i}^{(0)}+E_{\alpha _{2ri+1}}^{(0)})+c_{r1}(E_{\alpha _r\alpha _{r1}}^{(0)}+E_{\alpha _{r+1}\alpha _{r+2}}^{(0)})`$ (5.4)
$`+`$ $`c_rE_{\alpha _1+\mathrm{}+\alpha _{2r}}^{(\frac{1}{2})}+\psi e^{\frac{1}{2}R}(E_{\alpha _r}^{(0)}+E_{\alpha _{r+1}}^{(0)})+{\displaystyle \underset{i=1}{\overset{r1}{}}}\phi _i_i`$
$`+`$ $`{\displaystyle \frac{\chi }{\mathrm{\Delta }}}e^{\frac{1}{2}R}(E_{\alpha _r}^{(0)}+E_{\alpha _{r+1}}^{(0)})`$
and
$`\overline{A}_{A_{2r}^{(2)}}`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{r2}{}}}c_ie^{\phi _i+\phi _{i+1}}(E_{\alpha _i}^{(0)}+E_{\alpha _{2ri+1}}^{(0)})+c_re^{2\phi _1}E_{\alpha _1\mathrm{}\alpha _{2r}}^{(\frac{1}{2})}`$ (5.5)
$`+`$ $`c_{r1}e^{\phi _{r1}}(E_{\alpha _r+\alpha _{r1}}^{(0)}+E_{\alpha _{r+1}+\alpha _{r+2}}^{(0)})`$
$`+`$ $`c_{r1}\psi e^{\frac{1}{2}R\phi _{r1}}(E_{\alpha _{r+1}+\alpha _r+\alpha _{r1}}^{(0)}E_{\alpha _r+\alpha _{r+1}+\alpha _{r+2}}^{(0)})`$
$`+`$ $`c_{r1}\chi e^{\frac{1}{2}R\phi _{r1}}(E_{\alpha _{r1}}^{(0)}E_{\alpha _{r+2}}^{(0)})+c_{r1}\psi \chi e^{\phi _{r1}}(E_{\alpha _r+\alpha _{r1}}^{(0)}E_{\alpha _{r+1}+\alpha _{r+2}}^{(0)})`$
$`+`$ $`{\displaystyle \frac{1}{2}}c_{r1}\psi ^2\chi e^{\phi _{r1}\frac{1}{2}R}(E_{\alpha _{r+1}+\alpha _r+\alpha _{r1}}^{(0)}E_{\alpha _r+\alpha _{r+1}+\alpha _{r+2}}^{(0)})`$
(b) the $`D_{r+1}^{(2)}`$ affine NA - Toda model
$`A_{D_{r+1}^{(2)}}`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{r2}{}}}c_iE_{\alpha _i}^{(0)}+c_{r1}(E_{\alpha _r\alpha _{r1}}^{(0)}+E_{\alpha _{r1}\alpha _{r+1}}^{(0)})`$ (5.6)
$`+`$ $`c_r(E_{(\alpha _1+\mathrm{}+\alpha _{r1}+\alpha _{r+1})}^{(\frac{1}{2})}E_{(\alpha _1+\mathrm{}+\alpha _r+\alpha _{r+1})}^{(\frac{1}{2})})+\psi e^{\frac{1}{2}R}(E_{\alpha _r}^{(0)}+E_{\alpha _{r+1}}^{(0)})`$
$`+`$ $`{\displaystyle \underset{i=1}{\overset{r1}{}}}\phi _i_i+{\displaystyle \frac{\chi }{\mathrm{\Delta }}}e^{\frac{1}{2}R}(E_{\alpha _r}^{(0)}+E_{\alpha _{r+1}}^{(0)})`$
and
$`\overline{A}_{D_{r+1}^{(2)}}`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{r2}{}}}c_ie^{\phi _i+\phi _{i+1}}E_{\alpha _i}^{(0)}+c_{r1}e^{\phi _{r1}}(E_{\alpha _r+\alpha _{r1}}^{(0)}+E_{\alpha _{r+1}+\alpha _{r1}}^{(0)})`$ (5.7)
$`+`$ $`2c_{r1}\psi e^{\frac{1}{2}R\phi _{r1}}E_{\alpha _{r+1}+\alpha _r+\alpha _{r1}}^{(0)}+2c_{r1}\chi e^{\frac{1}{2}R\phi _{r1}}E_{\alpha _{r1}}^{(0)}`$
$`+`$ $`2c_{r1}\psi \chi e^{\phi _{r1}}(E_{\alpha _{r+1}+\alpha _{r1}}^{(0)}+E_{\alpha _{r1}+\alpha _r}^{(0)})+c_{r1}\psi ^2\chi e^{\frac{1}{2}R\phi _{r1}}E_{\alpha _{r+1}+\alpha _r+\alpha _{r1}}^{(0)}`$
$`+`$ $`c_{r+1}e^{\phi _1}(E_{(\alpha _1+\mathrm{}+\alpha _{r1}+\alpha _{r+1})}^{(\frac{1}{2})}E_{(\alpha _1+\mathrm{}+\alpha _r+\alpha _{r+1})}^{(\frac{1}{2})})`$
For the untwisted affine $`B_r^{(1)}`$ model of the previous section the zero curvature representation is obtained from the following flat connections :
$`A_{B_r^{(1)}}`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{r2}{}}}c_iE_{\alpha _i}^{(0)}+c_{r1}E_{\alpha _r\alpha _{r1}}^{(0)}+c_rE_{\alpha _1+2(\alpha _2+\mathrm{}+\alpha _r)}^{(1)}`$ (5.8)
$`+`$ $`\psi e^{\frac{1}{2}R}E_{\alpha _r}^{(0)}+{\displaystyle \underset{i=1}{\overset{r1}{}}}\phi _i_i+{\displaystyle \frac{\chi }{\mathrm{\Delta }}}e^{\frac{1}{2}R}E_{\alpha _r}^{(0)}`$
$`\overline{A}_{B_r^{(1)}}`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{r2}{}}}c_ie^{\phi _i+\phi _{i+1}}E_{\alpha _i}^{(0)}+c_re^{\phi _1+\phi _2}E_{(\alpha _1+2(\alpha _2+\mathrm{}+\alpha _r))}^{(1)}+2\chi e^{\phi _{r1}+\frac{1}{2}R}E_{\alpha _{r1}}^{(0)}`$ (5.9)
$`+`$ $`c_{r1}(1+2\psi \chi )e^{\phi _{r1}}E_{\alpha _{r1}+\alpha _r}^{(0)}2c_{r1}e^{\phi _{r1}\frac{1}{2}R}\psi (1+\psi \chi )E_{\alpha _{r1}+2\alpha _r}^{(0)}`$
The zero curvature representation of the subclass of torsionless singular affine NA - Toda models shows that they are in fact classically integrable field theories. The construction of the previous sections provides a systematic affine Lie algebraic structure underlying those models, which is known to play the crucial role in the construction of their finite energy soliton solutions.
## 6 Conclusions
We have constructed a class of affine (and conformal affine) NA-Toda models from the gauged two-loop WZNW models, in which left and right symmetries are incorporated by a suitable choice of grading operator $`Q`$ and of grade $`\pm 1`$ constant generators $`ϵ_\pm `$. We have shown that for non abelian zero grade subalgebra $`𝒢_0`$, it is possible to reduce even further the phase space by constraining to zero the currents associated to generators $`YH`$, commuting with $`ϵ_\pm `$ ($`YH𝒢_0^0)`$). There exists two inequivalent manners to gauge fix $`𝒢_0^0=U(1)`$-the axial and the vector gaugings. Similarly to the T-duality transformations between the axial and vector gauged G/H - WZNW models, one can find the off-critical counterpart of the conformal T-duality, relating now the axial and vector families of IM’s constructed in Sect. 2 . We further analize the problem of deriving the Lie algebraic condition which defines a class of T-selfdual torsionless models, for the case $`𝒢_0^0=U(1)`$. The action for those models were systematicaly constructed and shown to coincide with those proposed by Fateev , describing the strong coupling limit of specific 2-d models representing the complex sine-Gordon (i.e. Lund-Regge ) interacting with Toda-like models. Their weak coupling limit appears to be the Thirring model coupled to certain affine Toda theories .
As we have mentioned in sects. 1 and 3 the conformal $`\sigma `$-model limits of the axial and vector dyonic IM’s describe critical $`D=r+1`$ strings on black hole backgrounds . One may wonder whether the nonconformal (i.e. off-critical) dyonic IM’s (representing integrable perturbations of these string models ) have some string field theory applications. As is well known the relation between 2-d conformal model and its integrable perturbations allows to describe the off-critical behaviour of the original conformal model as well as the RG flow from ultraviolet to infrared (in the case of unitary CFT’s ) see , and references therein. Hence the properties of the admissible integrable perturbations of the conformal $`\sigma `$-models are an important ingredient in the description of the space of string backgrounds (satisfying certain low energy physical requirements). As it advocated by B. Zwiebach the nonconformal versions of the conformal string backgrounds appears to be the main tool for the construction of the off-shell string field theory (SFT) . Among all the possible nonconformal backgrounds the integrable ones (admiting ”worldsheet” soliton solutions) have the advantage to offer powerfull methods for studying the string S-duality.
The dyonic integrable models studied in the present paper represent the simplest family of IM’s with one U(1) global symmetry. Translated to the string language this means that the target-space metrics ($`E_{MN}(X),\varphi (X),T(X)`$) are independent of one of the coordinates
$`X_M`$ ($`\theta =\frac{1}{2}ln\frac{\chi }{\psi }`$ in our case) and that the relevant operator does not break this isometry. Models with more isometries (generic abelian T-duality) or those admiting isotropies ( nonabelian T-duality) appears to be useful for the construction of physically intersting string models (). The problem of construction of their integrable perturbations requires to consider more general affine $`\widehat{𝒢}_r`$-NA Toda models defined by grading operators as for example $`Q_{a,b,\mathrm{}}=h_{a,b,\mathrm{}}\widehat{d}+_{ia,b\mathrm{}}^n\frac{2\lambda _iH}{\alpha _i^2}`$ and appropriately choosen $`ϵ_\pm `$. They should allow larger(than U(1)) invariant subgroup $`𝒢_0^0`$ (non-abelian in general). An important characteristic of such IM’s is that their physical fields belongs to , say $`𝒢_0/𝒢_0^0=\frac{SL(2)U(1)^{r1}}{U(1)^s},s=2,\mathrm{}r1`$ (i.e with s isometries), or $`𝒢_0/𝒢_0^0=\frac{SL(2)SL(2)U(1)^{r2}}{U(1)^s}`$, or $`𝒢_0/𝒢_0^0=\frac{SL(3)U(1)^{r2}}{U(2)}`$ (the corresponding IM’s represent string backgrounds with one isometry and SL(2) as isotropy group ) ,etc. The methods needed in the construction of such IM’s and for the investigation of their T-duality properties appears to be straightforward generalization of the methods we have developed for the simplest case of $`𝒢_0^0=U(1)`$ . An intersting open problem is the classification of the affine NA-Toda IM’s according to the number of the physical fields ( i.e. the dimension of the string target space ) and their symmetry groups $`𝒢_0^0`$ ( i.e. the symmetries of the string backgrounds). We hope that the simplest dyonic IM’s studied in the present paper and the methods we have developed might contribute to the complete description of the space of the integrable (nonconformal) string backgrounds.
Acknowledgments We are grateful to CNPq, FAPESP and UNESP for financial support.
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# U(1) Gauge Field of the Kaluza-Klein Theory in the Presence of Branes
## I Introduction
In the standard five-dimensional Kaluza-Klein (KK) approach where the full spacetime manifold is factorized as $`M_4S^1`$, the five-dimensional gravitation theory which has the reparametrization invariance on $`S^1`$ can be interpreted as a gauge theory of the Virasoro group from the four-dimensional point of view (see and references therein). After the geometric spontaneous symmetry breaking of the Virasoro invariance, the excitations of the 5D gravitational fields are split into 4D massless gravitational fields, massless gauge fields, massless scalar field, and an infinite tower of massive spin-2 fields . In particular, the U(1) gauge symmetry of the vector fields in the 4D effective action is originated from the translational isometries in the extra dimension.
Recently, there have been lots of interest in the phenomenon of localization of gravity proposed by Randall and Sundrum (RS) (for previous relevant work see references therein). RS assumed a single positive tension 3-brane and a negative bulk cosmological constant in the five-dimensional spacetime. There have been developed a large number of brane world models afterwards . The introduction of branes usually gives rise to the “warping” of the extra dimensions, resulting in non-factorizable spacetime manifolds. More importantly, the presence of branes breaks the translational isometries in the extra dimensions. Therefore, it would be very interesting to see how the conventional Kaluza-Klein picture changes in the brane world scenarios.
Dobado and Maroto recently have incorporated the effect of the presence of the brane in the Kaluza-Klein (KK) reduction by introducing Goldstone bosons (GB) fields. The GB fields parameterize the excitations of the brane and so the GB correspond to the spontaneous symmetry breaking of the translational isometries of the compactified extra dimensions by the brane. It has been shown that in the dimensional reduction of the KK unifications a sort of Higgs mechanism related to the GB gives mass to the KK graviphotons. (see also the appendix of Ref. although it has the opposite sign for the mass term.)
On the other hand, there are other approaches to confine standard model particles on the brane by allowing the fields to live in the bulk spacetime. For example, bulk gauge bosons have been considered in Ref. . It is important to derive the zero mode effective action because its zero modes (massless modes) correspond to the standard model particles localized on the brane.
In this paper, contrary to the approach mentioned above where the brane is treated as a dynamical object, we have treated the brane as fixed, and investigate the zero mode dimensional reduction of the Kaluza-Klein unifications. The brane world model considered in this paper is the RS one with a single 3-brane in the infinite fifth dimension . As expected, the breaking of isometries in the extra dimension by the brane makes the 4D effective action not being invariant under U(1) gauge transformations. Interestingly, however, the analysis of the linearized equation around the RS background shows that the KK gauge field does possess the U(1) gauge symmetry at the linear level.
In section 2, we carry out the dimensional reduction of the KK unifications in the presence of a single brane with some ansatz for the zero mode excitations. In section 3, the linearized perturbations of the 4D effective action obtained are analyzed around the RS vacuum solution. Some physical implications of our results are discussed in section 4.
## II Kaluza-Klein reduction
The RS model with a single brane can be described by the following action in 5D spacetimes
$$I=d^4x_{\mathrm{}}^{\mathrm{}}𝑑z\frac{\sqrt{\widehat{g}}}{16\pi G_5}(\widehat{R}2\mathrm{\Lambda })d^4x\sqrt{\widehat{g}_B}\sigma .$$
(1)
Here $`G_5`$ is the 5D Newton’s constant, $`\mathrm{\Lambda }`$ the bulk cosmological constant of 5D spacetime, $`\widehat{g}_B`$ the determinant of the induced metric describing the brane, and $`\sigma `$ the tension of the brane. We assume that the value of $`\sigma `$ is fine-tunned such that $`\mathrm{\Lambda }=6k^2(<0)`$ with $`k=4\pi G_5\sigma /3`$. Let us introduce the domain-wall metric in the following form
$`ds^2`$ $`=`$ $`\widehat{g}_{MN}dx^Mdx^N=H^2(z)g_{MN}dx^Mdx^N`$ (2)
$`=`$ $`H^2(z)\left[\gamma _{\mu \nu }dx^\mu dx^\nu +\varphi ^2(dz\kappa A_\mu dx^\mu )^2\right].`$ (3)
Here $`H=k|z|+1`$, $`\varphi ^2=g_{55}`$, and $`\kappa A_\mu =g_{5\mu }/g_{55}`$. The standard Kaluza-Klein decomposition of the metric is given by
$$(g_{MN})=\left(\begin{array}{cc}\gamma _{\mu \nu }+\kappa ^2\varphi ^2A_\mu A_\nu & \kappa \varphi ^2A_\mu \\ \kappa \varphi ^2A_\nu & \varphi ^2\end{array}\right),(g^{MN})=\left(\begin{array}{cc}\gamma ^{\mu \nu }& \kappa A^\mu \\ \kappa A^\nu & \varphi ^2(1+\kappa ^2\varphi ^2AA)\end{array}\right)$$
(4)
with $`A^\mu =\gamma ^{\mu \nu }A_\nu `$ and $`AA=A_\mu A^\mu `$. Here $`\kappa `$ is the gauge coupling constant.
Under the specific class of coordinate transformations such as
$$x^\mu \stackrel{~}{x}^\mu =\stackrel{~}{x}^\mu (x),z\stackrel{~}{z}=z+f(x),$$
(5)
we obtain
$$\stackrel{~}{\gamma }_{\mu \nu }=\frac{x^\alpha }{\stackrel{~}{x}^\mu }\frac{x^\beta }{\stackrel{~}{x}^\nu }\gamma _{\alpha \beta },\stackrel{~}{A}_\mu =\frac{x^\alpha }{\stackrel{~}{x}^\mu }A_\alpha +\kappa ^1\frac{f}{\stackrel{~}{x}^\mu },\stackrel{~}{\varphi }(\stackrel{~}{x},\stackrel{~}{z})=\varphi (x,z)$$
(6)
according to $`\stackrel{~}{g}_{MN}=\frac{x^P}{\stackrel{~}{x}^M}\frac{x^Q}{\stackrel{~}{x}^N}g_{PQ}`$. We see that $`\gamma _{\mu \nu }`$ transforms like a four-dimensional metric tensor, and $`\varphi `$ a scalar field under diffeomorphisms in Eq. (5). However, we point out that the 5D diffeomorphisms are split into the 4D diffeomorphisms plus the gauge transformations for the field $`A_\mu `$.
In this paper, we are mainly interested in the zero mode effective action . In general, it is a non-trivial problem to determine what the “zero mode” is if the full spacetime is not factorizable. As an ansatz for the zero mode, we assume that $`\gamma _{\mu \nu }`$, $`A_\mu `$, and $`\varphi `$ are functions of $`x`$-coordinates only, i.e., no $`z`$-coordinate dependence. If one requires the $`Z_2`$ (e.g., $`R/Z_2`$) orbifold symmetry in the brane world model, there will be no vector zero mode fluctuations. It follows because in the presence of $`Z_2`$ orbifold symmetry in Eq. (3) the vector gauge field $`A_\mu `$ should satisfy $`A_\mu (x,z)=A_\mu (x,z)`$ and so $`A_\mu (x)=0`$. In what follows, we consider general cases without having the orbifold symmetry in the theory.
The above assumption comes from the crucial observation that the graviton zero mode $`h_{\mu \nu }`$ in $`\gamma _{\mu \nu }=\eta _{\mu \nu }+\kappa h_{\mu \nu }`$ depends only on “$`x`$” even if one starts from $`h_{\mu \nu }(x,z)=H^{3/2}\psi (z)\widehat{h}_{\mu \nu }(x)`$ in the RS approach where the transverse fluctuations are fixed (e.g., $`h_{5\mu }=h_{55}=0`$). For the zero mode solution with $`m^2=0`$, we have $`\psi ^0(z)=c_hH^{3/2}`$, thus we find $`h_{\mu \nu }^0(x,z)=c_h\widehat{h}_{\mu \nu }(x)`$ with a constant $`c_h`$. Other examples are the form of the zero modes for the bulk spin-0 and spin-1 fields on the RS background . For the spin-0 field $`\mathrm{\Phi }(x,z)=H^{3/2}\chi (z)\widehat{\varphi }(x)`$, we have $`\chi =c_\mathrm{\Phi }H^{3/2}`$ for the zero mode and hence its localized zero mode is given by $`\mathrm{\Phi }^0(x,z)=c_\mathrm{\Phi }\widehat{\varphi }(x)`$. In the case of the spin-1 field $`V_\mu (x,z)=H^{3/2}\sigma (z)v_\mu (x)`$, one finds $`\sigma =c_VH^{3/2}`$ for the zero mode and hence its zero mode is given by $`V_\mu ^0(x,z)=c_Vv_\mu (x)`$. From the observations mentioned above, we may propose that the zero modes are constants with respect to “$`z`$”. Furthermore we stress that for $`h_{\mu 5}0,h_{55}0`$, it may be not a correct way to obtain the zero modes from the linearized equations. This is because their forms are too complicated to analyze the zero modes . Even if we choose the gauge-fixing, it is hard to obtain the consistent zero mode solutions. Hence the integration of Eq. (1) over $`z`$ could be a good starting point to obtain the zero modes.
Note first that
$`\sqrt{\widehat{g}}=H^5\varphi \sqrt{\gamma },`$ (7)
$`\sqrt{\widehat{g}_B}=H^4(z=0)\sqrt{\gamma }\sqrt{|\delta _\nu ^\mu +\kappa ^2\varphi ^2A^\mu A_\nu |}.`$ (8)
Using $`\widehat{g}_{MN}=H^2g_{MN}`$, one has
$$\widehat{R}=H^2\left[R(g)+8\frac{_P^PH}{H}20\frac{_PH^PH}{H^2}\right].$$
(9)
Since
$$_P^PH=H^{\prime \prime }(\varphi ^2+\kappa ^2AA)+\kappa H^{}(\varphi ^1A^\mu _\mu \varphi +\frac{1}{2}A^\mu \gamma ^{\alpha \beta }_\mu \gamma _{\alpha \beta }+_\mu A^\mu ),$$
(10)
we have
$`8{\displaystyle \frac{_P^PH}{H}}20{\displaystyle \frac{_PH^PH}{H^2}}`$ $`=`$ $`\left(8{\displaystyle \frac{H^{\prime \prime }}{H}}20({\displaystyle \frac{H^{}}{H}})^2\right)(\varphi ^2+\kappa ^2AA)`$ (12)
$`+8\kappa {\displaystyle \frac{H^{}}{H}}\left(\varphi ^1A^\mu _\mu \varphi +{\displaystyle \frac{1}{2}}A^\mu \gamma ^{\alpha \beta }_\mu \gamma _{\alpha \beta }+_\mu A^\mu \right),`$
where the prime () denotes the differentiation with respect to $`z`$.
Then, the five-dimensional action Eq. (1) is given by
$`I`$ $`=`$ $`{\displaystyle \frac{1}{16\pi G_5}}{\displaystyle }d^4x\sqrt{\gamma }\varphi [R(g){\displaystyle }dzH^3+(\varphi ^2+\kappa ^2AA){\displaystyle }dzH^3(8{\displaystyle \frac{H^{\prime \prime }}{H}}20{\displaystyle \frac{H^2}{H^2}})`$ (15)
$`+8\kappa (\varphi ^1A^\mu _\mu \varphi +{\displaystyle \frac{1}{2}}A^\mu \gamma ^{\alpha \beta }_\mu \gamma _{\alpha \beta }+_\mu A^\mu ){\displaystyle }dz{\displaystyle \frac{H^{}}{H^4}}2\mathrm{\Lambda }{\displaystyle }dzH^5]`$
$`{\displaystyle d^4x\sqrt{\gamma }\sqrt{|\delta _\nu ^\mu +\kappa ^2\varphi ^2A^\mu A_\nu |}\sigma }.`$
It should be pointed out that the 5D Ricci scalar curvature constructed from $`g_{MN}`$, $`R(g)`$, is independent of $`z`$-coordinate since the metric elements $`g_{MN}`$ are functions of $`x^\mu `$ only. Using $`H^{}=k\theta (z)`$, $`H^{\prime \prime }=2k\delta (z)`$, $`_{\mathrm{}}^{\mathrm{}}𝑑zH^3=1/k`$, and $`_{\mathrm{}}^{\mathrm{}}𝑑zH^5=1/2k`$ for the RS model with a single brane, one gets the 4D effective action
$`I_{KK}`$ $`=`$ $`{\displaystyle \frac{1}{16\pi G_5}}{\displaystyle \frac{1}{k}}{\displaystyle }d^4x\sqrt{\gamma }[\varphi R(g)\varphi \mathrm{\Lambda }+6k^2(\varphi ^1+\kappa ^2\varphi AA)`$ (17)
$`16\pi G_5k\sigma \sqrt{|\delta _\nu ^\mu +\kappa ^2\varphi ^2A^\mu A_\nu |}]`$
$`=`$ $`{\displaystyle \frac{1}{16\pi G_4}}{\displaystyle d^4x\sqrt{\gamma }\left[\varphi R(g)+6k^2\left(\varphi ^1+\varphi 2\sqrt{|\delta _\nu ^\mu +\kappa ^2\varphi ^2A^\mu A_\nu |}+\kappa ^2\varphi AA\right)\right]}.`$ (18)
Here the 4D Newton’s constant is defined as $`G_4=G_5k`$. Notice from Eq. (17) that it reproduces the ordinary KK reduction in the presence of the cosmological constant as the brane at $`z=0`$ disappears (i.e., $`\sigma `$, $`k0`$). The 5D scalar curvature $`R(g)`$ is related to the 4D Ricci scalar curvature constructed from $`\gamma _{\mu \nu }`$, $`R(\gamma )`$, as
$$d^4x\sqrt{\gamma }\varphi R(g)=d^4x\sqrt{\gamma }\varphi \left[R(\gamma )\frac{\kappa ^2}{4}\varphi ^2F^2\right]+[\mathrm{}].$$
(19)
The last term is the surface term. The field strength is defined as $`F_{\mu \nu }=_\mu A_\nu _\nu A_\mu `$ and $`F^2=F_{\mu \nu }F^{\mu \nu }`$. Using Eq. (19), one finally obtains
$`I_{KK}`$ $`=`$ $`{\displaystyle \frac{1}{16\pi G_4}}{\displaystyle }d^4x\sqrt{\gamma }[\varphi R(\gamma ){\displaystyle \frac{\kappa ^2}{4}}\varphi ^3F^2`$ (21)
$`+6k^2(\varphi ^1+\varphi 2\sqrt{|\delta _\nu ^\mu +\kappa ^2\varphi ^2A^\mu A_\nu |}+\kappa ^2\varphi AA)],`$
where we have ommitted the surface terms.
We observe that the zero mode gravitational degrees of freedom in the 5D spacetime are split into the 4D gravitational fields $`\gamma _{\mu \nu }`$, a vector field $`A_\mu `$, and a graviscalar field $`\varphi `$ as usual. However, the properties of the vector field and the scalar field are very different from those in the conventional KK reduction. The first two terms in this effective action are the same form as in the ordinary dimensional reduction of the Kaluza-Klein unifications, and they have the U(1) gauge symmetry. The difference from the conventional KK reduction is only the last term which is proportional to the brane tension squared. If one started from the KK metric decomposition with $`A_\mu =0`$ and $`\varphi =1`$ in Eq. (4), this “potential” term would disappear and one obtains the ordinary Einstein gravity on the brane with zero effective cosmological constant as well known. As can be easily seen in Eq. (17), this happens because of the fine tunning between the brane tension $`\sigma `$ and the 5D bulk cosmological constant $`\mathrm{\Lambda }`$.
The appearance of the non-linear term (e.g., $`\sqrt{|\delta _\nu ^\mu +\kappa ^2\varphi ^2A^\mu A_\nu |}`$) as well as the squared term in $`A_\mu `$ shows not only that the 4D effective action no longer has the gauge symmetry, but also that KK photons are not massless. This arises from the presence of the brane ($`k0`$) in the five-dimensional spacetime. Because Eq. (21) contains a non-linear term which generates a lot of terms, a truncated form of the effective action may be useful to understand the dynamics of the fields easily. If one expands the full action up to the order of $`\kappa ^2`$, one has
$$I_{KK}^T\frac{1}{16\pi G_4}d^4x\sqrt{\gamma }\left[\varphi R(\gamma )\frac{\kappa ^2}{4}\varphi ^3F^2+6k^2\left(\varphi ^1+\varphi 2+\kappa ^2\varphi (1\varphi )AA\right)\right].$$
(22)
Here we used $`\sqrt{|\delta _\nu ^\mu +\kappa ^2\varphi ^2A^\mu A_\nu |}1+\frac{1}{2}\kappa ^2\varphi ^2AA`$.
In order to explicitly see how the dynamical aspect of the $`\varphi `$ field comes out, let us conformally transform the metric as
$$\gamma _{\mu \nu }\overline{\gamma }_{\mu \nu }=\varphi \gamma _{\mu \nu }.$$
(23)
Then, the zero-mode effective action Eq. (21) is written by
$`I_{KK}`$ $`=`$ $`{\displaystyle \frac{1}{16\pi G_4}}{\displaystyle }d^4x\sqrt{\overline{\gamma }}[R(\overline{\gamma }){\displaystyle \frac{\kappa ^2}{4}}\varphi ^3F^2{\displaystyle \frac{3}{2}}\varphi ^2\overline{}^\mu \varphi \overline{}_\mu \varphi `$ (25)
$`+6k^2\varphi ^2(\varphi ^1+\varphi 2\sqrt{|\delta _\nu ^\mu +\kappa ^2\varphi ^3A^\mu A_\nu |}+\kappa ^2\varphi ^2AA)]`$
up to the surface terms. Here $`F^2=\overline{\gamma }^{\mu \nu }\overline{\gamma }^{\alpha \beta }F_{\mu \alpha }F_{\nu \beta }`$, $`A^\mu =\overline{\gamma }^{\mu \alpha }A_\alpha `$, and $`AA=\overline{\gamma }^{\mu \nu }A_\mu A_\nu `$. The truncated effective action in this Einstein frame is also given by
$`I_{KK}^T`$ $``$ $`{\displaystyle \frac{1}{16\pi G_4}}{\displaystyle }d^4x\sqrt{\overline{\gamma }}\{R(\overline{\gamma }){\displaystyle \frac{\kappa ^2}{4}}\varphi ^3F^2{\displaystyle \frac{3}{2}}\varphi ^2\overline{}^\mu \varphi \overline{}_\mu \varphi `$ (27)
$`+6k^2[\varphi ^2(\varphi ^1+\varphi 2)+\kappa ^2(1\varphi )AA]\}.`$
As shall be shown below, the form of these actions above also suggests that the metric tensor which is responsible for the 4D Einstein gravitation is not the $`\gamma _{\mu \nu }`$, but indeed the $`\overline{\gamma }_{\mu \nu }`$.
Now the field equations derived from the effective action (27) are given by
$`R_{\mu \nu }={\displaystyle \frac{\kappa ^2}{2}}\varphi ^3\left(F_{\mu \alpha }F_{\nu }^{}{}_{}{}^{\alpha }{\displaystyle \frac{1}{4}}\overline{\gamma }_{\mu \nu }F^2\right)6k^2\kappa ^2(1\varphi )A_\mu A_\nu `$ (28)
$`+{\displaystyle \frac{3}{2}}\varphi ^2\left[\overline{}_\mu \varphi \overline{}_\nu \varphi 2k^2\overline{\gamma }_{\mu \nu }(\varphi ^1+\varphi 2)\right],`$ (29)
$`\overline{}^\mu F_{\mu \nu }+12k^2\kappa ^2\varphi ^3(1\varphi )A_\nu =3\varphi ^1\overline{}^\mu \varphi F_{\mu \nu },`$ (30)
$`\overline{}^\mu \overline{}_\mu \varphi \varphi ^1\overline{}^\mu \varphi \overline{}_\mu \varphi +2k^2\left[\varphi ^1(4\varphi 3\varphi ^1)\kappa ^2\varphi ^2AA\right]={\displaystyle \frac{\kappa ^2}{4}}F^2.`$ (31)
Here we used the truncated effective action for simplicity. The essential properties of solutions do not change.
It is well known in the RS model that if the flat metric on the brane is replaced by any Ricci-flat 4D metric then the 5D Einstein equations with a negative cosmological constant are still satisfied . Now note that $`A_\mu =0`$ and $`\varphi =1`$ satisfies Eqs. (30) and (31). In this case, Eq. (29) becomes $`R_{\mu \nu }=0`$. Thus any 4D Ricci-flat metric $`\overline{\gamma }_{\mu \nu }`$ is a solution. Therefore, although based on the assumption for the zero-mode which is $`z`$-coordinate independent, our results reproduce this well known property in the RS model. Such somewhat general agreement may indicates the validity of our ansatz for the zero mode. One can see that $`\overline{\gamma }_{\mu \nu }=\eta _{\mu \nu }`$ corresponds to the RS solution . Since the metric for the 4D Schwartzschild black hole is Ricci-flat, the Schwartzschild black hole can be embedded in the brane world as well known. That is, assuming the spherically symmetric background $`\widehat{g}_{MN}=H^2(z)(\overline{\gamma }_{\mu \nu }^S,\varphi ^2)`$ with $`\overline{\gamma }_{\mu \nu }^S=\mathrm{diag}[12M/r,(12M/r)^1,r^2,r^2\mathrm{sin}^2\theta ]`$ and $`\varphi =1`$, we obtain the black string solution in 5D Anti de Sitter($`AdS_5`$) spacetime . Since this black string solution is unstable near the AdS horizon, but stable far from it, it is likely to end up with a ”black cigar” solution as conjectured in Ref. .
Secondly, it would be of interest to ask whether or not the Reissner-Nordström charged black hole can be embedded in the brane world. It is straightfoward to see that the Reissner-Nordström black hole solution,
$$\overline{\gamma }_{\mu \nu }^{RN}=\mathrm{diag}[12M/r+Q^2/r^2,(12M/r+Q^2/r^2)^1,r^2,r^2\mathrm{sin}^2\theta ]$$
(32)
with $`Q^{}=\kappa Q/2`$, $`F_{tr}=Q/r^2`$, and $`\varphi =1`$, satisfies Eqs. (29) and (30). In this case we have $`F_{tr}=_rA_0`$ with $`A_0=Q/r`$. However it does not satisfy Eq. (31) (i.e., $`2k^2AA=F^2/4`$). This is so because of the presence of both brane and $`\varphi `$. Thus, the 4D Reissner-Nordström black hole with $`\varphi =1`$ cannot be embedded in the brane world. In the absence of the brane (i.e., $`k=0`$), there exist charged black hole solutions with non-trivial $`\varphi `$ field. However, such charged black hole is very different from the Reissner-Nordström black hole . It does not seem that the presence of the brane changes this feature of the Kaluza-Klein theory much. On the other hand, however, if the scalar field were frozen somehow from the beginning in Eq. (4) (e.g., $`\varphi =1`$), there would be no equation like Eq. (31). Consequently, by observing Eqs. (29) and (30), one can easily find that the Reissner-Nordström black hole is a solution.
## III Linearized perturbation
In this section we consider the linearized perturbations around the RS vacuum solution ($`\overline{\gamma }_{\mu \nu }=\eta _{\mu \nu },A_\mu =0,\varphi =1`$) for the dimensionally reduced effective action in Eq. (25). Actually, at the linear level, the truncated effective action in Eq. (27) is equivalent to the non-truncated one in Eq. (25). Let us introduce the perturbations around the RS solution
$$\gamma _{\mu \nu }=\eta _{\mu \nu }+\kappa h_{\mu \nu },A_\mu =0+a_\mu ,\varphi =1+\kappa \phi .$$
(33)
Consequently,
$$\overline{\gamma }_{\mu \nu }=\eta _{\mu \nu }+\kappa \overline{h}_{\mu \nu },\overline{h}_{\mu \nu }=h_{\mu \nu }+\phi \eta _{\mu \nu }.$$
(34)
Then the bilinear action of Eq. (25) or (27) which governs the perturbative dynamics is given by
$`I_{KK}^0`$ $`=`$ $`{\displaystyle \frac{\kappa ^2}{16\pi G_4}}{\displaystyle }d^4x\{{\displaystyle \frac{1}{4}}[^\mu \overline{h}^{\alpha \beta }_\mu \overline{h}_{\alpha \beta }^\mu \overline{h}_\mu \overline{h}+2^\mu \overline{h}_{\mu \nu }^\nu \overline{h}2^\mu \overline{h}_{\mu \alpha }^\nu \overline{h}_{\nu }^{}{}_{}{}^{\alpha }]`$ (36)
$`{\displaystyle \frac{1}{4}}(_\mu a_\nu _\nu a_\mu )(^\mu a^\nu ^\nu a^\mu ){\displaystyle \frac{3}{2}}_\mu \phi ^\mu \phi +6k^2\phi ^2\},`$
where $`\overline{h}=\eta ^{\mu \nu }\overline{h}_{\mu \nu }=h+4\phi `$. Surprisingly, it turns out that the bilinear effective action is invariant under the U(1) gauge transformation. The U(1) gauge symmetry breaking term in Eq. (27) appears as higher order term than the squared order: i.e., $`6k^2\kappa ^2(1\varphi )AA6k^2\kappa ^3\phi a^\mu a_\mu `$.
In order to understand what physical states there are, let us analyze the field equations as below. From the action Eq. (36) we have the equations of motion
$`\mathrm{}\overline{h}_{\mu \nu }+_\mu _\nu \overline{h}\left(_\mu ^\alpha \overline{h}_{\alpha \nu }+_\nu ^\alpha \overline{h}_{\alpha \mu }\right)\eta _{\mu \nu }\left(\mathrm{}\overline{h}^\alpha ^\beta \overline{h}_{\alpha \beta }\right)=0,`$ (37)
$`\mathrm{}a_\mu _\mu (_\nu a^\nu )=0,`$ (38)
$`\mathrm{}\phi +4k^2\phi =0.`$ (39)
By taking the trace of Eq. (37), we have
$$\mathrm{}\overline{h}^\alpha ^\beta \overline{h}_{\alpha \beta }=0.$$
(40)
Hence Eq. (37) becomes
$$\mathrm{}\overline{h}_{\mu \nu }+_\mu _\nu \overline{h}\left(_\mu ^\alpha \overline{h}_{\alpha \nu }+_\nu ^\alpha \overline{h}_{\alpha \mu }\right)=0.$$
(41)
So far we have not chosen any gauge for $`\overline{h}_{\mu \nu }`$. Now let us choose the transverse (or harmonic) gauge in the five-dimensional spacetime. Since
$$g_{MN}=\eta _{MN}+\kappa ϵ_{MN},\left(ϵ_{MN}\right)=\left(\begin{array}{cc}h_{\mu \nu }& a_\mu \\ a_\nu & 2\phi \end{array}\right),$$
(42)
the five-dimensional harmonic gauge $`^Mϵ_{MN}=\frac{1}{2}_Nϵ`$ is equivalent to
$$^\mu \overline{h}_{\mu \nu }=\frac{1}{2}_\nu \overline{h},_\mu a^\mu =0.$$
(43)
That is, the 5D harmonic gauge is split into the harmonic gauge for the 4D gravitational field and the Lorenz gauge for the 4D KK gauge field. Using these gauge conditions above, Eq. (41) and Eq. (38) become
$$\mathrm{}\overline{h}_{\mu \nu }=0,\mathrm{}a_\mu =0,$$
(44)
respectively. Therefore, it proves that $`\overline{h}_{\mu \nu }`$ and $`a_\mu `$ indeed represent the massless spin-2 gravitons and the massless spin-1 graviphotons on the brane, respectively.
On the other hand, the spin-0 scalar field fluctuation $`\phi `$ in Eq. (39) appears to be massive. However, it has a tachyonic mass $`m_\phi ^2=4k^2`$ proportional to the brane tension squared. It seems to indicate that the scalar fluctuation of the 5D gravitational degrees of freedom corresponds to the unstable mode in the RS background.
## IV Discussion
We have investigated the KK zero mode effective action in the presence of a single brane in the extra dimension. Although the four-dimensional gravitational modes behaves as usual, the vector and scalar modes behave quite differently. In the 4D effective action, it seems that the vector field $`A_\mu `$ does not possess the U(1) gauge symmetry and that KK photons are not massless any more. The scalar field $`\varphi `$ is also no longer massless and couples to the vector field.
However, this is not all of the story. The linearized perturbation analysis around the RS background spacetime shows that the 5D massless gravitational degrees of freedom are split into spin-2, spin-1, and spin-0 modes as the standard KK model up to $`\kappa ^2`$-order. We have observed that the massive propagation of the vector mode in the 4D effective action is not revealed in the linearized perturbation. In order to observe the effect of the brane ($`k0`$), thus, one needs to study one-loops correction rather than the linearized one. For example, we expect the relevant vertex correction ($`k^2\kappa ^3\phi a^\mu a_\mu `$) from the last term of Eq. (27).
We also have observed that the spin-0 mode propagation has a tachyonic mass, indicating some instability of the RS background spacetime through the “radion” or “modulus” field. Presumably, it suggests that some stabilization mechanism for the 55-metric component is necessary in order to have the stable RS background spacetime with a single brane as in the case of the two branes. On the other hand, if one requires the $`R/Z_2`$ orbifold symmetry in the brane world model, there will be no vector zero mode propagations as mentioned above. Thus, the $`R/Z_2`$ orbifold symmetry with the “modulus” field stabilization establishes the usual localization of gravity on the brane in the RS model .
In deriving the U(1) Maxwell term from the 5D RS brane model, we use the conventional Kaluza-Klein approach. Apparently, we find a non-linear term as well as $`AA`$. This arises from a sort of brane-Higgs effect: Here the isometry of extra dimension was broken spontaneously by the presence of the brane. Hence we expect that the gauge field becomes massive. However, we have found that the massive propagation of the KK gauge field does not reveal at the linear level. Fortunately, instead we find the massless vector propagation. What will happen if we take a $`z`$-coordinate dependent or other form of ansatz for the “zero modes”? We still expect there should be non-linear or mass terms in the reduced effective action due to the broken isometry in the extra dimension. However, in order to answer whether or not the gauge field becomes massless when linearized, some further work in detail is needed. There were other attempts to achieve the 4D U(1) symmetry from the 5D U(1) bulk gauge fields .
On the other hand, the RS solution can be extended to accommodate the Schwarzschild black hole solution on the brane as a zero mode solution. This is possible because the RS solution is Ricci-flat. Hence the Ricci-flat Schwarzschild solution can be embedded into the brane world by introducing the spherically symmetric spacetime. Now it is very important to check whether or not the RS brane world allows to have the Reissner-Nordström black hole on the brane. As is shown above, the Reissner-Nordström black hole cannot be embedded in the brane world, because this case of $`A_00`$ cannot be a solution to the effective action of Eq. (21) including the non-linear term and $`AA`$. To obtain this black hole on the brane, it seems to be necessary to introduce some U(1) bulk gauge field in the five-dimensional spacetime whose dynamics is localized on the brane .
We have observed that the naive propagations of the scalar field gives rise to the tachyonic mass proportional to the tension of the brane. It may induce the instability of the RS vacuum. In the linearized gravity, by using the residual gauge freedom, one can also impose the traceless gauge in a source-free region in addition to the harmonic gauge. It follows mainly because “$`\mathrm{}\mathrm{Trace}`$” vanishes in a source-free region. What will happen if such traceless gauge is imposed in our linearized analysis? The five-dimensional trace is $`ϵ=\eta ^{MN}ϵ_{MN}=h+2\phi =\overline{h}2\phi `$. Here $`\overline{h}=h+4\phi `$ is used. By combining the trace of Eq. (44) and Eq. (39), we have
$$\mathrm{}_{(5)}ϵ=\mathrm{}\overline{h}2\mathrm{}\phi =8k^2\phi ,$$
(45)
where $`\mathrm{}_{(5)}=\mathrm{}+_5^2`$. Thus, imposing the five-dimensional traceless gauge (i.e., $`ϵ=0`$) directly results in $`\phi =0`$, that is, no graviscalar fluctuation. Since $`h=ϵ2\phi `$ and $`\overline{h}=ϵ+2\phi `$, it also means the four-dimensional traceless gauge (i.e., $`h=\overline{h}=0`$). In other words, we notice that the existence of the tachyonic graviscalar fluctuation is mutually inconsistent with imposing the traceless gauge condition. Therefore, the resolution of the instability of the RS background spacetime due to the graviscalar transforms to whether or not one can impose the traceless gauge. In the linearized gravity, the trace of metric fluctuations can be set to be zero by using remaining gauge freedom provided that there is no matter source in the region in consideration . In our analysis, however, since the graviscalar field $`\phi `$ plays like a matter source in the trace equation above, it does not seem to be plausible imposing such gauge condition in the first place. Presumably such traceless gauge condition can be imposed on $`\overline{h}`$ ($`\mathrm{}\overline{h}=0`$), but not on $`h`$ ($`\mathrm{}h=16k^2\phi `$).
As mentioned above, another possible caveat of our result is that the tachyonic graviscalar fluctuation is merely an artifact of the ansatz for the zero mode we used in this paper. For instance, instead of $`H(z)=k|z|+1`$ in Eq. (3), let us assume $`H(x,z)=k|z|\varphi (x)+1`$ for the zero mode. This ansatz is analogous to the form used in Ref. for the case of RS model with two branes. Then the 4D effective action is given by
$`I_{KK}`$ $`=`$ $`{\displaystyle \frac{1}{16\pi G_4}}{\displaystyle }d^4x\sqrt{\gamma }[R(\gamma ){\displaystyle \frac{\kappa ^2}{4}}\varphi ^2F^2+2\varphi ^2\gamma ^{\mu \nu }_\mu \varphi _\nu \varphi `$ (47)
$`+6k^2(22\sqrt{|\delta _\nu ^\mu +\kappa ^2\varphi ^2A^\mu A_\nu |}+\kappa ^2\varphi ^2AA)]`$
$``$ $`{\displaystyle \frac{1}{16\pi G_4}}{\displaystyle d^4x\sqrt{\gamma }\left[R(\gamma )\frac{\kappa ^2}{4}\varphi ^2F^2+2\varphi ^2\gamma ^{\mu \nu }_\mu \varphi _\nu \varphi \right]}.`$ (48)
Here we see that the graviscalar fluctuation as well as that of the gravivector becomes massless in the linearized perturbations. Therefore, in order to clarify the issues discussed above, further investigation is required on what the correct form of the ansatz is for the zero mode in the brane world scenarios.
Finally, it will be interesting to extend our study to various types of brane world models as well as to the RS model with two positive and negative tension branes and see how the U(1) gauge field behaves in the 4D effective action. It will be also worth investigating how graviphotons and graviscalar particles interact with the 5D bulk standard model particles in the presence of branes.
## Acknowledgments
The authors thank Hyungwon Lee for helpful discussions and the referee for crucial comments. GK thanks Seungjoon Hyun for usefull conversations. This work was supported by the Brain Korea 21 Programme, Ministry of Education, Project No. D-0025.
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# 1 Introduction
## 1 Introduction
Massless four point Feynman diagrams contribute to many important physical amplitudes. They are much more complicated than two- and three-point diagrams because depend on many parameters: the Mandelstam variables $`s`$ and $`t`$ and the values of the external momenta squared, $`p_i^2,i=1,2,3,4`$. In the most general case, when all the legs are off the mass shell, $`p_i^20`$, there exists an explicit analytical result for the master (i.e. with powers of the propagators equal to one) double box diagram (see Fig. 1)
strictly in four dimensions. Still no similar results are available for pure off shell four point diagrams with ultraviolet, infrared and/or collinear divergences.
In the opposite case, when all the end-points are on shell, i.e. for $`p_i^2=0,i=1,2,3,4`$, the problem of the analytical evaluation of such diagrams, in expansion in $`ϵ=(4d)/2`$ in the framework of dimensional regularization with the space-time dimension $`d`$ as a regularization parameter, was completely solved during last year in . Among intermediate situations, when some legs are on shell and the rest of them off shell, the case of one leg off shell, $`q^2=p_1^20`$ and three legs on shell is very important because of the relevance to the process $`e^+e^{}3`$jets (see, e.g., ). The purpose of this paper is to analytically evaluate the master double box diagram of such type, as a function of $`q^2,s`$ and $`t`$, and thereby demonstrate that the NNLO analytical calculations for this process are indeed possible.
One of the ways to evaluate the four point diagrams with one leg off shell is to expand them in the limit $`q^20`$ and compute as many terms of the resulting expansion as possible. We explain how to do this, following the strategy of regions , in the next section and present the leading power term in this expansion which provides a very non-trivial check of the subsequent analytical result.
To analytically evaluate the considered diagram we straightforwardly apply the method of ref. : we start from the alpha-representation of the double box and, after expanding some of the involved functions in Mellin–Barnes (MB) integrals, arrive at a six-fold MB integral representation with gamma functions in the integrand. Then we use a standard procedure of taking residues and shifting contours to resolve the structure of singularities in the parameter of dimensional regularization, $`ϵ`$. This procedure leads to the appearance of multiple terms where Laurent expansion in $`ϵ`$ becomes possible. Resulting integrals in all the MB parameters but the last two are evaluated explicitly in gamma functions and their derivatives. The last two-fold MB integral is evaluated by closing an initial integration contour in the complex plane to the right, with an explicit summation of the corresponding series. A final result is expressed through (generalized) polylogarithms dependent on rational combinations of $`q^2,s`$ and $`t`$ and a one-dimensional integral with a simple integrand consisting of logarithms and dilogarithms.
## 2 Expansion in the limit $`q^20`$
The dimensionally regularized master massless double box Feynman integral with one leg off shell, $`q^2=p_1^20`$, and three legs on shell, $`p_i^2=0,i=2,3,4`$, can be written as
$`F(s,t,q^2;ϵ)`$ $`=`$ $`{\displaystyle \frac{\text{d}^dk\text{d}^dl}{(k^2+2p_1k+q^2)(k^22p_2k)k^2(kl)^2}}`$ (1)
$`\times {\displaystyle \frac{1}{(l^2+2p_1l+q^2)(l^22p_2l)(l+p_1+p_3)^2}},`$
where $`s=(p_1+p_2)^2,t=(p_1+p_3)^2`$, and $`k`$ and $`l`$ are respectively loop momenta of the left and the right box. Usual prescriptions, $`k^2=k^2+i0,s=s+i0`$, etc are implied. To expand the given diagram in the limit $`q^20`$ one can apply the so-called strategy of regions based on the analysis of various regions in the space of the loop integration momenta, Taylor expanding the integrand in the parameters that are considered small in the given region and extending resulting integrations to the whole integration domain in the loop momenta. When applying this strategy all integrals without scale are by definition put to zero.
Let us choose, for convenience, the external momenta as follows:
$$p_1=\stackrel{~}{p}_1\frac{q^2}{Q^2}\stackrel{~}{p}_2,p_2=\stackrel{~}{p}_2,\stackrel{~}{p}_{1,2}=(Q/2,0,0,Q/2),$$
where $`s=Q^2`$. The given limit $`|q^2||s|,|t|`$ is closely related to the Sudakov limit so that it is reasonable to consider each loop momentum to be one of the following types:
hard (h): $`kQ\sqrt{t},`$
1-collinear (1c): $`k_+q^2/Q,k_{}Q,\underset{¯}{k}\sqrt{q^2},`$
2-collinear (2c): $`k_+Q,k_{}q^2/Q,\underset{¯}{k}\sqrt{q^2}.`$
Here $`k_\pm =k_0\pm k_3,\underset{¯}{k}=(k_1,k_2)`$. We mean by $`kQ`$, etc. that any component of $`k_\mu `$ is of order $`Q`$.
It turns out that the (h-h), (1c-h) and (1c-1c) are the only non-zero contributions to the leading power behaviour in the limit $`q^20`$. Any term originating from the (h-h) contribution is given by the expansion of the integrand in Taylor series in $`q^2`$ and expressed through on-shell double boxes in shifted dimensions and can be analytically evaluated by the algorithm presented in . The (1c-1c) contribution is obtained by expanding propagators number 2, 4 and 7 in a special way. In particular, propagators number 2 and 4 are expanded, respectively, in $`l^2`$ and $`k^2`$. (See for instructive 2-loop examples of expansions in limits of the Sudakov type.)
The (1c-h) and (1c-1c) contributions are evaluated with the help of a two-fold (respectively, one-fold) MB representation. Still this program of the evaluation of a large number of terms of the expansion looks very complicated because one needs, for phenomenological reasons, the values of $`q^2`$ greater than $`s`$ and $`t`$ so that a reliable summation of a resulting series, using Padé approximants, requires the knowledge of at least first 20–30 terms. Such a great number of terms can be hardly evaluated since a lot of irreducible structures appear. This asymptotic expansion is however very useful for comparison with the explicit result derived below.
The leading power terms of the asymptotic expansion calculated in expansion in $`ϵ`$, up to a finite part, are
$$F(s,t,q^2;ϵ)=\frac{\left(i\pi ^{d/2}\mathrm{e}^{\gamma _\mathrm{E}ϵ}\right)^2}{(s)^{2+2ϵ}(t)}\underset{i=0}{\overset{4}{}}\frac{g_i(X,Y)}{ϵ^i}+O(q^2\mathrm{ln}^3(q^2/s))+O(ϵ),$$
(2)
where $`X=q^2/s,Y=t/s`$ and
$`g_4(X,Y)`$ $`=`$ $`1,`$
$`g_3(X,Y)`$ $`=`$ $`2(\mathrm{ln}X\mathrm{ln}Y),`$
$`g_2(X,Y)`$ $`=`$ $`{\displaystyle \frac{11\pi ^2}{12}}+3\mathrm{ln}X\mathrm{ln}Y{\displaystyle \frac{3}{2}}\mathrm{ln}^2Y,`$
$`g_1(X,Y)`$ $`=`$ $`2\mathrm{ln}Y\text{Li}_2\left(Y\right)2\text{Li}_3\left(Y\right)+{\displaystyle \frac{2}{3}}\mathrm{ln}^3X{\displaystyle \frac{3}{2}}\mathrm{ln}^2X\mathrm{ln}Y{\displaystyle \frac{1}{2}}\mathrm{ln}X\mathrm{ln}^2Y`$
$`{\displaystyle \frac{1}{6}}\mathrm{ln}^3Y+\mathrm{ln}^2Y\mathrm{ln}(1+Y)+\pi ^2\left[{\displaystyle \frac{3}{2}}\mathrm{ln}X{\displaystyle \frac{19}{6}}\mathrm{ln}Y+\mathrm{ln}(1+Y)\right]+{\displaystyle \frac{49\zeta (3)}{6}},`$
$`g_0(X,Y)`$ $`=`$ $`26\text{Li}_4\left(Y\right)2S_{2,2}(Y)2(\mathrm{ln}X+6\mathrm{ln}Y+\mathrm{ln}(1+Y))\text{Li}_3\left(Y\right)`$ (3)
$`+2\mathrm{ln}Y\text{Li}_3\left({\displaystyle \frac{Y}{1+Y}}\right)+(\mathrm{ln}^2Y+2\mathrm{ln}X\mathrm{ln}Y+4\pi ^2)\text{Li}_2\left(Y\right)`$
$`{\displaystyle \frac{1}{2}}\mathrm{ln}^4X+{\displaystyle \frac{1}{2}}\mathrm{ln}^3X\mathrm{ln}Y+{\displaystyle \frac{1}{4}}\mathrm{ln}^2X\mathrm{ln}^2Y{\displaystyle \frac{1}{2}}\mathrm{ln}X\mathrm{ln}^3Y+{\displaystyle \frac{7}{8}}\mathrm{ln}^4Y`$
$`+\mathrm{ln}(1+Y)\left[\mathrm{ln}X\mathrm{ln}^2Y{\displaystyle \frac{5}{3}}\mathrm{ln}^3Y+{\displaystyle \frac{1}{2}}\mathrm{ln}^2Y\mathrm{ln}(1+Y){\displaystyle \frac{1}{3}}\mathrm{ln}Y\mathrm{ln}^2(1+Y)\right]`$
$`+\pi ^2[{\displaystyle \frac{2}{3}}\mathrm{ln}^2X{\displaystyle \frac{7}{3}}\mathrm{ln}X\mathrm{ln}Y+{\displaystyle \frac{25}{6}}\mathrm{ln}^2Y+\mathrm{ln}X\mathrm{ln}(1+Y)2\mathrm{ln}Y\mathrm{ln}(1+Y)`$
$`+{\displaystyle \frac{1}{2}}\mathrm{ln}^2(1+Y)]+\zeta (3)[{\displaystyle \frac{19}{3}}\mathrm{ln}X{\displaystyle \frac{34}{3}}\mathrm{ln}Y+2\mathrm{ln}(1+Y)]+{\displaystyle \frac{83\pi ^4}{180}}.`$
Here $`\text{Li}_a\left(z\right)`$ is the polylogarithm and
$$S_{a,b}(z)=\frac{(1)^{a+b1}}{(a1)!b!}_0^1\frac{\mathrm{ln}^{a1}(t)\mathrm{ln}^b(1zt)}{t}\text{d}t$$
(4)
the generalized polylogarithm .
## 3 From alpha parameters through MB representation to analytical result
The alpha representation of the double box looks like:
$$F(s,t,q^2;ϵ)=\mathrm{\Gamma }(3+2ϵ)\left(i\pi ^{d/2}\right)^2_0^{\mathrm{}}\text{d}\alpha _1\mathrm{}_0^{\mathrm{}}\text{d}\alpha _7\delta \left(\alpha _i1\right)D^{1+3ϵ}A^{32ϵ},$$
(5)
where
$`D`$ $`=`$ $`(\alpha _1+\alpha _2+\alpha _7)(\alpha _3+\alpha _4+\alpha _5)+\alpha _6(\alpha _1+\alpha _2+\alpha _3+\alpha _4+\alpha _5+\alpha _7),`$ (6)
$`A`$ $`=`$ $`[\alpha _1\alpha _2(\alpha _3+\alpha _4+\alpha _5)+\alpha _3\alpha _4(\alpha _1+\alpha _2+\alpha _7)+\alpha _6(\alpha _1+\alpha _3)(\alpha _2+\alpha _4)](s)`$ (7)
$`+\alpha _5\alpha _6\alpha _7(t)+\alpha _5[(\alpha _1+\alpha _3)\alpha _6+\alpha _3(\alpha _1+\alpha _2+\alpha _7)](q^2).`$
As it is well-known, one can choose a sum of an arbitrary subset of $`\alpha _i,i=1,\mathrm{},7`$ in the argument of the delta function in (5), and we use the same choice as in .
Starting from (5) we perform the same change of variables as in and apply seven times the MB representation
$$\frac{1}{(X+Y)^\nu }=\frac{1}{\mathrm{\Gamma }(\nu )}\frac{1}{2\pi i}_i\mathrm{}^{+i\mathrm{}}\text{d}w\frac{Y^w}{X^{\nu +w}}\mathrm{\Gamma }(\nu +w)\mathrm{\Gamma }(w)$$
(8)
in order to separate terms in the functions involved to make possible an explicit parametric integration. The two extra MB integrations arise form the extra term with $`q^2`$. After such integrations we are left with a 7-fold MB integral of a ratio of gamma functions. Fortunately, one of the integrations can be explicitly taken using the first Barnes lemma and we arrive at the following nice 6-fold MB integral:
$`F(s,t,q^2;ϵ)`$ $`=`$ $`{\displaystyle \frac{\left(i\pi ^{d/2}\right)^2}{\mathrm{\Gamma }(13ϵ)(s)^{3+2ϵ}}}{\displaystyle \frac{1}{(2\pi i)^6}}{\displaystyle \text{d}v\text{d}w\text{d}w_2\text{d}w_3\text{d}z\text{d}z_1\left(\frac{q^2}{s}\right)^v\left(\frac{t}{s}\right)^w}`$ (9)
$`\times \mathrm{\Gamma }(1+w)\mathrm{\Gamma }(1+v+w)\mathrm{\Gamma }(v)\mathrm{\Gamma }(w)\mathrm{\Gamma }(1w_3+v)`$
$`\times \mathrm{\Gamma }(w_2)\mathrm{\Gamma }(12ϵww_2)\mathrm{\Gamma }(w_3v)\mathrm{\Gamma }(12ϵww_3)`$
$`\times {\displaystyle \frac{\mathrm{\Gamma }(1w_2+z_1)\mathrm{\Gamma }(1w_3+z_1)\mathrm{\Gamma }(ϵ+w+w_2+w_3z_1)\mathrm{\Gamma }(z_1)}{\mathrm{\Gamma }(1+w+w_2+w_3)\mathrm{\Gamma }(14ϵww_2w_3)\mathrm{\Gamma }(1w_3)}}`$
$`\times \mathrm{\Gamma }(1ϵ+z)\mathrm{\Gamma }(2+2ϵ+w+w_2+zz_1)\mathrm{\Gamma }(2+2ϵ+w+w_3+zz_1)`$
$`\times {\displaystyle \frac{\mathrm{\Gamma }(23ϵww_2w_3+z_1z)\mathrm{\Gamma }(z_1z)}{\mathrm{\Gamma }(3+2ϵ+w+z)}}.`$
It differs from its analog for $`q^2=0`$ by the additional integration in $`v`$. This variable enters only four gamma functions in the integrand. The integral is evaluated in expansion in $`ϵ`$, up to a finite part, by resolving singularities in $`ϵ`$ absolutely by the same strategy as in the case $`q^2=0`$ . Note that the infrared and collinear poles are a little bit softer than in the pure on-shell case, the integration variable $`v`$ playing the role of an infrared regulator. The two key gamma functions that are responsible for the generation of poles in $`ϵ`$ are the same as in the previous case:
$$\mathrm{\Gamma }(ϵ+w+w_2+w_3z_1)\mathrm{\Gamma }(23ϵww_2w_3+z_1z).$$
The labeling of resulting terms is therefore the same: the initial integral is decomposed as $`J=J_{00}+J_{01}+J_{10}+J_{11}`$, etc. (Only arguments of some gamma functions are shifted by $`v`$.) The applied strategy makes it possible to perform all the integrations apart from the last two, in $`v`$ and $`w`$. We obtain four groups of terms with 26 terms in each group: the terms without MB integration, with MB integration in $`v`$ or $`w`$ and, finally, with a two-fold integration in $`v`$ and $`w`$. The one-fold integrals are explicitly evaluated by closing contour and summing up series, using formulae from .
The contribution of the resulting two-fold MB integral takes the form
$`{\displaystyle \frac{2\left(i\pi ^{d/2}\right)^2}{s^3}}{\displaystyle \frac{1}{(2\pi i)^2}}{\displaystyle \frac{\text{d}v\text{d}w}{1+w}\left(\frac{q^2}{s}\right)^v\left(\frac{t}{s}\right)^w\mathrm{\Gamma }(1+v+w)\mathrm{\Gamma }(v)\mathrm{\Gamma }(1+w)\mathrm{\Gamma }(w)^2}`$ (10)
$`\times [\mathrm{\Gamma }(1+v+w)\mathrm{\Gamma }(vw)({\displaystyle \frac{1}{ϵ}}\gamma _\mathrm{E}2\mathrm{ln}(s){\displaystyle \frac{5}{1+w}}{\displaystyle \frac{1}{1+v+w}}`$
$`+\psi (1+v)2\psi (vw)3\psi (w)+2\psi (1+w)+\psi (1+v+w))`$
$`\mathrm{\Gamma }(1+v)\mathrm{\Gamma }(v)\mathrm{\Gamma }(1+w)\mathrm{\Gamma }(w)].`$
The integration contours are straight lines along imaginary axes with $`1<`$Re$`v`$, Re$`w`$,Re$`v+w<0`$. By closing contours it is possible to convert this integral into a two-fold series where each term is identified as a derivative of the Appell function $`F_2`$ in parameters, up to the third order. The $`1/ϵ`$ part is then explicitly summed up with a result in terms of polylogarithms. (In fact, it is proportional to the $`ϵ`$ part of the master one-loop box.)
The so obtained result can be transformed into a one-dimensional integral with a simple integrand. To present the final result let us turn to the variables $`x=s/q^2`$ and $`y=t/q^2`$ keeping in mind typical phenomenological values of the involved parameters relevant to the process $`e^+e^{}3`$jets:
$$F(s,t,q^2;ϵ)=\frac{\left(i\pi ^{d/2}\mathrm{e}^{\gamma _\mathrm{E}ϵ}\right)^2}{s^2t(q^2)^{2ϵ}}\underset{i=0}{\overset{4}{}}\frac{f_i(x,y)}{ϵ^i}+O(ϵ).$$
(11)
We obtain
$`f_4(x,y)`$ $`=`$ $`1,`$ (12)
$`f_3(x,y)`$ $`=`$ $`2(\mathrm{ln}x+\mathrm{ln}y),`$ (13)
$`f_2(x,y)`$ $`=`$ $`3\text{Li}_2\left(x\right)+\text{Li}_2\left(y\right)2(\mathrm{ln}x+\mathrm{ln}y)^2`$ (14)
$`+3\mathrm{ln}(1x)\mathrm{ln}x+\mathrm{ln}(1y)\mathrm{ln}y{\displaystyle \frac{5\pi ^2}{12}},`$
$`f_1(x,y)`$ $`=`$ $`2[\text{Li}_3\left({\displaystyle \frac{x}{1xy}}\right)+\text{Li}_3\left({\displaystyle \frac{y}{1xy}}\right)\text{Li}_3\left({\displaystyle \frac{xy}{1xy}}\right)`$ (15)
$`\mathrm{ln}x\text{Li}_2\left({\displaystyle \frac{y}{1x}}\right)\mathrm{ln}y\text{Li}_2\left({\displaystyle \frac{x}{1y}}\right)]+2\mathrm{ln}(1xy)`$
$`\times \left[{\displaystyle \frac{1}{6}}\left(\mathrm{ln}^2(1xy)+\pi ^2\right)+\mathrm{ln}(1x)\mathrm{ln}x+\mathrm{ln}(1y)\mathrm{ln}y\mathrm{ln}x\mathrm{ln}y\right]`$
$`+3\text{Li}_3\left(x\right)8\text{Li}_3\left(y\right)+4\text{Li}_3\left({\displaystyle \frac{x}{1x}}\right)2\text{Li}_3\left({\displaystyle \frac{y}{1y}}\right)`$
$`(3\mathrm{ln}x+4\mathrm{ln}y)\text{Li}_2\left(x\right)+3\mathrm{ln}y\text{Li}_2\left(y\right)+{\displaystyle \frac{4}{3}}\mathrm{ln}^3x{\displaystyle \frac{2}{3}}\mathrm{ln}^3(1x)+\mathrm{ln}^2(1x)\mathrm{ln}x`$
$`{\displaystyle \frac{9}{2}}\mathrm{ln}(1x)\mathrm{ln}^2x+{\displaystyle \frac{\pi ^2}{6}}(5\mathrm{ln}x4\mathrm{ln}(1x))+{\displaystyle \frac{4}{3}}\mathrm{ln}^3y+{\displaystyle \frac{1}{3}}\mathrm{ln}^3(1y)`$
$`2\mathrm{ln}^2(1y)\mathrm{ln}y\mathrm{ln}(1y)\mathrm{ln}^2y+{\displaystyle \frac{\pi ^2}{6}}(5\mathrm{ln}y+2\mathrm{ln}(1y))`$
$`+4\mathrm{ln}x\mathrm{ln}y(\mathrm{ln}x\mathrm{ln}(1x)+\mathrm{ln}y)+{\displaystyle \frac{25\zeta (3)}{6}}.`$
The $`ϵ^0`$ part involves a one-dimensional integral:
$`f_0(x,y)`$ $`=`$ $`{\displaystyle _0^1}\text{d}z\{z^1\mathrm{ln}(1z)(4\mathrm{ln}^2(1xyz)\mathrm{ln}^2(1yxz))`$ (16)
$`{\displaystyle \frac{4y}{1xyz}}[\mathrm{ln}(1yz)(\mathrm{ln}(1z)\mathrm{ln}(1yz)2\text{Li}_2\left(z\right))`$
$`2(\mathrm{ln}(1z)\mathrm{ln}z)\text{Li}_2((1xyz)/x)]`$
$`{\displaystyle \frac{x}{1yxz}}[\mathrm{ln}(1xz)(3\mathrm{ln}^2(1z)6\mathrm{ln}(1z)\mathrm{ln}(1xz)+2\text{Li}_2\left(z\right))`$
$`+2(6\mathrm{ln}(1z)\mathrm{ln}z)\text{Li}_2((1yxz)/y)]\}`$
$`5\text{Li}_4\left(x\right)+14\text{Li}_4\left({\displaystyle \frac{x}{1y}}\right)2\text{Li}_4\left({\displaystyle \frac{x}{1x}}\right)6\text{Li}_4\left({\displaystyle \frac{xy}{(1x)(1y)}}\right)`$
$`+8\text{Li}_4\left({\displaystyle \frac{x}{1xy}}\right)+24\text{Li}_4\left(y\right)2\text{Li}_4\left(1y\right)+8\text{Li}_4\left({\displaystyle \frac{y}{1y}}\right)2\text{Li}_4\left({\displaystyle \frac{y}{1x}}\right)`$
$`8\text{Li}_4\left(1x\right)8\text{Li}_4\left({\displaystyle \frac{y}{1xy}}\right)20\text{Li}_4\left({\displaystyle \frac{1xy}{1y}}\right)+10\text{Li}_4\left({\displaystyle \frac{1xy}{1x}}\right)`$
$`3S_{2,2}(x)8S_{2,2}(y)6S_{2,2}\left({\displaystyle \frac{x}{1y}}\right)+(2\mathrm{ln}y2\mathrm{ln}x3\mathrm{ln}(1x))\text{Li}_3\left(x\right)`$
$`+2(16\mathrm{ln}(1y)11\mathrm{ln}y2\mathrm{ln}(1xy)+\mathrm{ln}x)\text{Li}_3\left({\displaystyle \frac{x}{1y}}\right)`$
$`(8\mathrm{ln}y+2\mathrm{ln}(1x)+3\mathrm{ln}x)\text{Li}_3\left({\displaystyle \frac{x}{1x}}\right)`$
$`2(4\mathrm{ln}(1y)4\mathrm{ln}y\mathrm{ln}(1x)+\mathrm{ln}x)\text{Li}_3\left({\displaystyle \frac{xy}{(1x)(1y)}}\right)`$
$`+(14\mathrm{ln}(1y)18\mathrm{ln}y+4\mathrm{ln}(1xy))\text{Li}_3\left({\displaystyle \frac{x}{1xy}}\right)+2\mathrm{ln}y\text{Li}_3\left({\displaystyle \frac{xy}{1xy}}\right)`$
$`+(7\mathrm{ln}y8\mathrm{ln}(1y))\text{Li}_3\left(y\right)+(8\mathrm{ln}(1y)+\mathrm{ln}y+2\mathrm{ln}x)\text{Li}_3\left({\displaystyle \frac{y}{1y}}\right)`$
$`4(2\mathrm{ln}y+7\mathrm{ln}(1x)+2\mathrm{ln}(1xy)8\mathrm{ln}x)\text{Li}_3\left({\displaystyle \frac{y}{1x}}\right)`$
$`2(\mathrm{ln}y+5\mathrm{ln}(1x)+8\mathrm{ln}(1xy)7\mathrm{ln}x)\text{Li}_3\left({\displaystyle \frac{y}{1xy}}\right)`$
$`{\displaystyle \frac{1}{2}}\left(\text{Li}_2\left(x\right)\right)^2\left(\text{Li}_2\left({\displaystyle \frac{x}{1y}}\right)\right)^2{\displaystyle \frac{3}{2}}\left(\text{Li}_2\left(y\right)\right)^2+4\left(\text{Li}_2\left({\displaystyle \frac{y}{1x}}\right)\right)^2`$
$`+[\mathrm{ln}^2(1x)4\mathrm{ln}^2y+2\mathrm{ln}y(4\mathrm{ln}(1x)\mathrm{ln}x)2\mathrm{ln}(1y)\mathrm{ln}x3\mathrm{ln}(1x)\mathrm{ln}x`$
$`+{\displaystyle \frac{7}{2}}\mathrm{ln}^2x+{\displaystyle \frac{5\pi ^2}{3}}\left]\text{Li}_2\right(x)+[12\mathrm{ln}^2(1y)+15\mathrm{ln}^2y+2\mathrm{ln}(1y)(\mathrm{ln}(1xy)+\mathrm{ln}x`$
$`9\mathrm{ln}y)+2\mathrm{ln}y(\mathrm{ln}(1xy)4\mathrm{ln}(1x)+4\mathrm{ln}x)`$
$`2(\mathrm{ln}^2(1xy)+\mathrm{ln}^2x)\left]\text{Li}_2\right({\displaystyle \frac{x}{1y}})+[4\mathrm{ln}^2(1y)11\mathrm{ln}^2y`$
$`+2\mathrm{ln}y(4\mathrm{ln}(1x)3\mathrm{ln}x)+\mathrm{ln}(1y)(5\mathrm{ln}y2\mathrm{ln}x)+\mathrm{ln}^2x{\displaystyle \frac{\pi ^2}{3}}\left]\text{Li}_2\right(y)`$
$`+[8\mathrm{ln}^2y8\mathrm{ln}y\mathrm{ln}(1x)10\mathrm{ln}^2(1x)+8\mathrm{ln}^2(1xy)8\mathrm{ln}(1xy)\mathrm{ln}x`$
$`8\mathrm{ln}(1x)(\mathrm{ln}(1xy)2\mathrm{ln}x)+2\mathrm{ln}(1y)\mathrm{ln}x+\mathrm{ln}^2x{\displaystyle \frac{2\pi ^2}{3}}\left]\text{Li}_2\right({\displaystyle \frac{y}{1x}})`$
$`+[\mathrm{ln}^2(1x)4\mathrm{ln}^2(1y)8\mathrm{ln}^2y+2\mathrm{ln}y(4\mathrm{ln}(1x)3\mathrm{ln}x)`$
$`+2\mathrm{ln}(1y)(4\mathrm{ln}y\mathrm{ln}x)2\mathrm{ln}(1x)\mathrm{ln}x+2\mathrm{ln}^2x\left]\text{Li}_2\right({\displaystyle \frac{xy}{(1x)(1y)}})`$
$`+2\mathrm{ln}^4(1xy)+{\displaystyle \frac{1}{3}}\mathrm{ln}^3(1xy)\left[2\mathrm{ln}y3\mathrm{ln}(1y)9\mathrm{ln}x11\mathrm{ln}(1x)\right]`$
$`+\mathrm{ln}^2(1xy)[\pi ^23\mathrm{ln}^2(1y)+3\mathrm{ln}^2y+6\mathrm{ln}^2(1x)\mathrm{ln}(1y)(\mathrm{ln}y10\mathrm{ln}x)`$
$`+4\mathrm{ln}(1x)\mathrm{ln}x2\mathrm{ln}^2x\mathrm{ln}y(5\mathrm{ln}(1x)+\mathrm{ln}x)]+{\displaystyle \frac{1}{3}}\mathrm{ln}(1xy)[7\mathrm{ln}^3(1y)`$
$`2\mathrm{ln}(1y)(5\pi ^2+12\mathrm{ln}^2y)+7\pi ^2\mathrm{ln}(1x)+\pi ^2\mathrm{ln}x+6\mathrm{ln}^2y\mathrm{ln}x4\mathrm{ln}^3(1x)`$
$`+15\mathrm{ln}^2(1y)(\mathrm{ln}y2\mathrm{ln}x)21\mathrm{ln}^2(1x)\mathrm{ln}x+3\mathrm{ln}(1x)\mathrm{ln}^2x`$
$`+\mathrm{ln}y(2\pi ^2+9\mathrm{ln}^2x+15\mathrm{ln}^2(1x)6\mathrm{ln}(1x)\mathrm{ln}x)]`$
$`{\displaystyle \frac{5}{6}}\mathrm{ln}^4(1x){\displaystyle \frac{2}{3}}\mathrm{ln}^4x+{\displaystyle \frac{23}{6}}\mathrm{ln}^3(1x)\mathrm{ln}x{\displaystyle \frac{17}{4}}\mathrm{ln}^2(1x)\mathrm{ln}^2x+{\displaystyle \frac{7}{2}}\mathrm{ln}(1x)\mathrm{ln}^3x`$
$`{\displaystyle \frac{\pi ^2}{6}}\left(3\mathrm{ln}^2(1x)10\mathrm{ln}(1x)\mathrm{ln}x+5\mathrm{ln}^2x\right)\mathrm{ln}^4(1y){\displaystyle \frac{2}{3}}\mathrm{ln}^4y`$
$`{\displaystyle \frac{19}{6}}\mathrm{ln}^3(1y)\mathrm{ln}y+5\mathrm{ln}^2(1y)\mathrm{ln}^2y+{\displaystyle \frac{2}{3}}\mathrm{ln}(1y)\mathrm{ln}^3y`$
$`+{\displaystyle \frac{\pi ^2}{6}}\left(9\mathrm{ln}^2(1y)\mathrm{ln}(1y)\mathrm{ln}y5\mathrm{ln}^2y\right)`$
$`+{\displaystyle \frac{1}{3}}\mathrm{ln}(1x)\mathrm{ln}(1y)(\mathrm{ln}^2(1x)4\mathrm{ln}^2(1y)){\displaystyle \frac{8}{3}}(\mathrm{ln}^2x+\mathrm{ln}^2y)\mathrm{ln}x\mathrm{ln}y`$
$`+3\mathrm{ln}^3(1y)\mathrm{ln}x+\mathrm{ln}^2(1y)\left[\mathrm{ln}y(4\mathrm{ln}(1x)+\mathrm{ln}x)2\mathrm{ln}(1x)\mathrm{ln}x\right]`$
$`+{\displaystyle \frac{1}{3}}\mathrm{ln}y[\mathrm{ln}^3(1x)9\mathrm{ln}^2(1x)\mathrm{ln}x12\mathrm{ln}y\mathrm{ln}^2x`$
$`+6\mathrm{ln}(1x)\mathrm{ln}x(2\mathrm{ln}y+3\mathrm{ln}x)]\mathrm{ln}(1y)[8\mathrm{ln}^2y\mathrm{ln}(1x)`$
$`+\mathrm{ln}(1x)\mathrm{ln}x(\mathrm{ln}(1x)2\mathrm{ln}x)+\mathrm{ln}y(8\mathrm{ln}^2(1x)+6\mathrm{ln}(1x)\mathrm{ln}x+\mathrm{ln}^2x)]`$
$`+{\displaystyle \frac{\pi ^2}{3}}\left[\mathrm{ln}y(4\mathrm{ln}(1x)5\mathrm{ln}x)\mathrm{ln}(1y)\mathrm{ln}x\right]`$
$`+\zeta (3)\left[12(\mathrm{ln}(1xy)\mathrm{ln}(1y))+13\mathrm{ln}(1x){\displaystyle \frac{25}{3}}(\mathrm{ln}x+\mathrm{ln}y)\right]+{\displaystyle \frac{23\pi ^4}{180}}.`$
One may hope that the one-dimensional integral that is left can also be evaluated in terms of polylogarithms. To do this it is necessary to complete the table of integrals derived in .
This result is in agreement with the leading power behaviour when $`q^20`$ (3). When performing this comparison it is reasonable to start with (10), take minus residue at $`v=0`$ (the first pole of $`\mathrm{\Gamma }(v)`$), integrate in $`w`$ by closing the contour to the right, and take into account the three other contributions (without MB integration, and with integration in $`v`$ or $`w`$) that were not presented above. Eqs. (1216) also agree with results based on numerical integration in the space of alpha parameters (where the 1% accuracy for the $`1/ϵ`$ and $`ϵ^0`$ parts is guaranteed).
Acknowledgments. I am grateful to Z. Kunszt for involving me into this problem and for kind hospitality during my visit to ETH (Zürich) in April–May 2000 where an essential part of this work was performed. I am thankful to T. Binoth and G. Heinrich for comparison of the presented result with their results based on numerical integration. Thanks to A.I. Davydychev and O.L. Veretin for useful discussions. This work was supported by the Volkswagen Foundation, contract No. I/73611, and by the Russian Foundation for Basic Research, project 98–02–16981.
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# Post-Newtonian gravitational radiation and equations of motion via direct integration of the relaxed Einstein equations. I. Foundations
## I INTRODUCTION
The motion of multiple, isolated bodies under their mutual gravitational attraction and the resulting emission of gravitational radiation is a long-standing problem that dates back to the first years following the publication of general relativity (GR). It has at times been controversial (for a thorough review see ). In 1916 Einstein calculated the gravitational radiation emitted by a laboratory-scale object using the linearized version of GR . Some of his assumptions were questionable and his answer for the energy flux was off by a factor of two (an error pointed out by Eddington ). In 1916, de Sitter derived N-body equations of motion in what later would be termed the post-Newtonian (PN) approximation. However, his equations contained an error that was discovered in the course of a disputed claim by Levi Civita that the center of mass of a binary star system would suffer a “self-acceleration”. Eddington and Clark corrected the error, and found no self-acceleration. Einstein, Infeld and Hoffman attempted to demonstrate explicitly that the Einstein equations alone imply equations of motion, by matching solutions of the vacuum equations, expanded in a weak-field, slow motion approximation, to fields representing the near-zone fields of “point” masses, working to first PN order. The result was the well-known EIH N-body equations of motion. Other highlights in this early history of the problem of motion include the development of the post-Newtonian approximation for fluid sytems by Fock and Chandrasekhar , its extension by Chandrasekhar and later workers to 2.5PN order , and the development of equations of motion for spinning bodies by Papapetrou and collaborators .
Gravitational theory presents one problem essentially identical to that of electromagnetic theory: how to mesh the natural solution of the field equations in the near zone where the bodies reside, which involves slow-motion expansions and instantaneous fields, with the solution in the far zone, which involves retarded fields. Such a meshing is needed if one is to calculate the effects of gravitational radiation reaction that result from the emission of enegy and angular momentum to infinity. One approach to resolving this problem was that of matched asymptotic expansions. Although well-rooted in applied mathematics, it was first expounded in 1970 as a powerful technique for electromagnetic and gravitational problems by Burke . Another, related approach is the “post-Minkowskian” framework, elaborated and developed most fully by Blanchet and Damour and their collaborators .
A second important problem of gravitation, which distinguishes it from electromagnetism, is the non-linearity of Einstein’s equations. Gravitation itself acts as a source of gravitation. Consequently this source extends over all space, resulting in the possibility of divergent or ill-defined integrals. In many ways, this has been the most serious difficulty to overcome. Techniques for resolving it have ranged from sweeping the difficulties under the rug, to the sophisticated analytic regularization methods of the post-Minkowskian program. A central thrust of this paper is to present a straightforward method for resolving this difficulty.
A third “problem”, which is less a problem for gravitation than it is for electromagnetism, is that of “point” sources. In electromagnetic theory, where there is a belief that fundamental charges like the electron are point-like, the singular nature of the fields at the source has led to problems of mass regularization, especially in deriving equations of electromagnetic radiation reaction; it also raises issues of the boundary between classical and quantum electrodynamics. In gravitation theory, this is less of an issue of principle, because the primary interest is in the motion of and radiation by astrophysical systems, whose members are clearly not point masses. Instead, the use of “point”, i.e. delta-function sources is meant as an efficient means of approximating the mass distribution of bodies that are nearly spherical and that are small compared to the typical separation between them, so that tidal effects, which depend on the finite size of the bodies, can be ignored. Here the issue is how to make use of a point mass approximation (which simplifies many calculations) in a way that captures all the physics without introducing spurious effects.
A fourth problem is of a technical nature: in electromagnetic theory, radiation damping in the equations of motion occurs at order $`(v/c)^3`$ beyond the simple Coulomb forces between charges, and is relatively easy to compute in a systematic approximation method, modulo the other problems listed above. By contrast, gravitational radiation damping occurs at order $`(v/c)^5`$ beyond Newtonian gravity, and requires a higher order of approximation that captures all relevant contributions. Over the years, numerous inequivalent results have been quoted for the leading gravitational radiation reaction effects. One finds published papers in which the coefficient in the relevant formula has ranged from $`21/16`$ to the correct coefficient of unity); a study by Walker and Will showed that the divergent results were all the simple consequence of missing one or more terms that contribute to the final answer.
These four “problems” were the origin of the so-called “quadrupole controversy”, which arose from a critique by Ehlers and colleagues of the foundations of the quadrupole formula for the leading-order gravitational radiation energy flux and orbital damping. This critique had the beneficial effect of spurring new research on those foundations, including a study of the systematic structure of the approximation sequence of Einstein’s equations in a slow-motion, weak-field approach; analysis of energy balance as an argument for connecting the far-zone energy flux to the near-zone damping forces, and elaboration of the post-Minkowskian approach, among others (see for a review). The work inspired by the Ehlers critique served to confirm the quadrupole formula and to strengthen its foundations. The ultimate test, of course, came in 1979 with the announcement of the measurement of orbital damping of the binary pulsar PSR 1913+16 in agreement with the quadrupole formula ; current results agree to better than 0.5 percent .
The problem of motion and radiation has received renewed interest since 1990, with the proposal for large-scale laser interferometric gravitational-wave observatories, such as the LIGO project in the US, and the realization that a leading candidate source of detectable waves would be the radiation-reaction driven inspiral of a binary system of compact objects (neutron stars or black holes) . Furthermore, it was noted that the leading method for data analysis of signals from such systems, optimal matched filtering, would require theoretical template waveforms that are accurate (primarily in the evolution of the orbital frequency or phase) well beyond the leading-order prediction of the quadrupole formula, possibly as high as corrections of order $`(v/c)^6`$.
This presented a major theoretical challenge: to calculate the motion and radiation to very high PN order, a formidable algebraic task, while addressing each of the problems listed above sufficiently well to ensure that the results were physically meaningful. This challenge was taken up by three groups of workers.
One group, headed by Blanchet, Damour and Iyer , used the post-Minkowskian (PM) approach to derive the gravitational waveform, equations of motion and energy flux explicitly to 2PN order ($`O(v/c)^4`$) and beyond. The idea is to solve the vacuum Einstein equations in the radiation zone in an expansion in powers of Newton’s constant $`G`$, and to express the asymptotic solutions in terms of a set of formal, time-dependent, symmetric and trace-free (STF) multipole moments . Then, in a near zone within one characteristic wavelength of the radiation, the equations including the material source are solved in a slow-motion approximation (expansion in powers of $`1/c`$) that yields both equations of motion for the source bodies, as well as a set of STF source multipole moments expressed as integrals over the “effective” source, including both matter and gravitational field contributions. The solutions involving the two sets of moments are then matched in an intermediate overlap zone, resulting in a connection between the formal radiative moments and the source moments. The matching also provides a natural way, using analytic continuation, to regularize integrals involving the non-compact contributions of gravitational stress-energy, that might otherwise be divergent.
The second group of Will, Wiseman and Pati use the approach described in the present paper, Direct Integration of the Relaxed Einstein Equations (DIRE), which builds upon earlier work by Epstein, Wagoner, Will and Wiseman . Like the PM approach, it involves rewriting the Einstein equations in their “relaxed” form, namely as an inhomogeneous, flat-spacetime wave equation for a field $`h^{\alpha \beta }`$, whose source consists of both the material stress-energy, and a “gravitational stress-energy” made up of all the terms non-linear in $`h^{\alpha \beta }`$. The wave equation is accompanied by a harmonic or deDonder gauge condition on $`h^{\alpha \beta }`$, which serves to specify a coordinate system, and also imposes equations of motion on the sources. Unlike the post-Minkowskian approach, a single formal solution is written down, valid everywhere in spacetime. This formal solution, based on the flat-spacetime retarded Green function, is a retarded integral equation for $`h^{\alpha \beta }`$, which is then iterated in a slow-motion ($`v/c<1`$), weak-field ($`h^{\alpha \beta }<1`$ ) approximation, that is very similar to the corresponding procedure in electromagnetism. However, because the integrand of this retarded integral is not compact by virtue of the non-linear field contributions, one quickly runs up against integrals that are not well defined, or worse, are divergent. Although at the lowest quadrupole and first PN order, various arguments were given to justify sweeping such problems under the rug , they were not very rigorous, and provided no guarantee that the divergences would not become insurmountable at higher PN orders. Indeed it is straightforward to demonstrate that at second post-Newtonian (2PN) order, the rug is indeed pulled out from under such arguments.
DIRE resolves these problems. The solution of the relaxed Einstein equation is a retarded integral, over the past null cone of the field point. The part of the integral that extends over the intersection between the past null cone and the material source and the near zone is approximated by a slow-motion expansion involving spatial integrals of moments of the source, including the non-compact gravitational contributions, just as in the post-Minkowskian and Epstein-Wagoner frameworks. But instead of extending the spatial integrals to infinity as was implicit in earlier procedures, we terminate the integrals at the boundary of the near zone, chosen to be at a radius $``$ given roughly by one wavelength of the gravitational radiation. For the integral over the rest of the past null cone exterior to the near zone (“radiation zone”), we use a change of integration variables to convert the integral into a convenient, easy-to-calculate form, that is manifestly convergent, subject only to reasonable assumptions about the past behavior of the source, that fully accounts for the retardation of the fields comprising the source stress-energy, and that does not involve an explicit slow-motion expansion. This transformation was suggested by our earlier work on a non-linear gravitational-wave phenomenon called the Christodoulou memory (it is also implicit in Appendix D of ). Not only are all integrations now explicitly finite and convergent, we can show explicitly that all contributions from the near-zone spatial integrals that depend upon the radius $``$ are actually cancelled by corresponding terms from the radiation-zone integrals, for all powers of $``$ (including $`\mathrm{ln}`$), and for any order in the PN expansion. Thus the procedure, as expected, has no dependence on the arbitrarily chosen boundary radius $``$ of the near-zone, and provides a simple practical method for regularizing integrals over non-compact sources.
The ultimate products of this work will consist of equations of motion, gravitational waveforms, and energy flux expressions, in reasonably ready-to-use forms. The equations of motion for a binary system will have the schematic form
$$d^2𝐱/dt^2=(Gm𝐱/r^3)[1+O(ϵ)+O(ϵ^2)+O(ϵ^{5/2})+O(ϵ^3)+O(ϵ^{7/2})+\mathrm{}],$$
(1)
where $`m`$ is the total mass of the binary system, $`𝐱=𝐱_1𝐱_2`$ is the separation vector and $`r=|𝐱|`$. The expansion parameter $`ϵ`$ is related to the orbital variables by $`ϵGm/rc^2(v/c)^2`$, where $`v`$ is the relative velocity. The leading term is Newtonian gravity. The next term $`O(ϵ)`$ is the first post-Newtonian correction, which gives rise to such effects as the advance of the periastron. The terms of $`O(ϵ^2)`$ and $`O(ϵ^3)`$ are non-dissipative 2PN and 3PN corrections. The $`O(ϵ^{5/2})`$ and $`O(ϵ^{7/2})`$ terms are the leading 2.5PN and post-Newtonian corrected 3.5PN gravitational radiation-reaction terms. (We do not include in this discussion contributions from spin, whose ordering in the PN hierarchy for compact bodies follows a special convention.) Explicit formulae for terms through various orders have been calculated by various authors: non-radiative terms through 2PN order , radiation reaction terms at 2.5PN and 3.5PN order , and non-radiatve 3PN terms
In order to derive equations of motion to the 3.5PN order shown, one must derive the near-zone metric $`g_{\alpha \beta }`$ as a function of spacetime and a functional of the source variables to 3.5PN order, which implies the following specific PN orders: $`g_{00}`$ through $`O(ϵ^{9/2})`$, $`g_{0i}`$ through $`O(ϵ^4)`$, $`g_{ij}`$ through $`O(ϵ^{7/2})`$. In this paper we provide the required expressions in the form of (a) Poisson-like integrals of source densities, $`_{}f(t,𝐱^{})|𝐱𝐱^{}|^pd^3x^{}`$, where $`f(t,𝐱^{})`$ could be proportional to source stress-energy densities, and thus have compact support, or could be a function of other potentials, and thus extend over the entire near-zone region of integration $``$; and (b) expressions involving time derivatives of source multipole moments $`M^{ijk\mathrm{}}`$ contracted with spatial vectors $`x^ix^jx^k\mathrm{}`$. These expressions can be simplified, iterated, and evaluated more explicitly, depending on the application envisioned (“point” mass binary system, spinning masses, perfect fluid distributions, etc.).
The second product will be expressions for the gravitational waveform, given schematically by
$$h^{ij}=\frac{G\mu }{Rc^4}\left\{v^2[1+O(ϵ^{1/2})+O(ϵ)+O(ϵ^{3/2})+O(ϵ^2)+O(ϵ^{5/2})+O(ϵ^3)\mathrm{}]\right\}_{TT},$$
(2)
where $`\mu `$ is the reduced mass, and the subscript TT denotes the “transverse-traceless” part. The leading contribution $`G\mu v^2/Rc^4G\ddot{I}^{ij}/Rc^4`$ is the standard quadrupole formula. Explicit formulae for all terms through 2.5PN order have been derived by various authors .
From the waveform, one can also derive expressions for fluxes of energy, angular momentum and linear momentum; the energy flux can be written in the schematic form
$`dE/dt=(dE/dt)_Q[1+O(ϵ)+O(ϵ^{3/2})+O(ϵ^2)+O(ϵ^{5/2})+O(ϵ^3)+\mathrm{}],`$ (3)
where $`(dE/dt)_Q`$ denotes the lowest-order quadrupole contribution,
A third approach focusses on the limit in which one body is much less massive than the other, and employs black-hole perturbation theory to derive the gravitational waveform and energy flux, for particles orbiting both rotating and non-rotating holes. This method yields both numerically accurate results as well as analytic PN expansions up to orders as high as $`(v/c)^{11}`$ . Work is currently in progress to extend these methods beyond the test-mass approximation, in an effort to compute corrections to first order in $`\mu /M`$, the ratio of the mass of the particle to that of the black hole .
This is the first in a series of papers that will treat the problem of motion and gravitational radiation systematically using the DIRE approach. This paper lays out the foundations of the method, and derives formal solutions to the near-zone fields through 3.5PN order (order $`(v/c)^7`$ beyond Newtonian gravity), in a form useful for future applications. Subsequent papers in the series will derive the explicit equations of motion and near-zone gravitational fields for binary systems of compact object through 2PN order, and deal with radiation reaction at 2.5PN and 3.5PN order.
Our conventions and notation generally follow those of . Henceforth we use units in which $`G=c=1`$. Greek indices run over four spacetime values 0, 1, 2, 3, while Latin indices run over three spatial values 1, 2, 3; commas denote partial derivatives with respect to a chosen coordinate system, while semicolons denote covariant derivatives; repeated indices are summed over; $`\eta ^{\mu \nu }=\eta _{\mu \nu }=\mathrm{diag}(1,1,1,1)`$; $`gdet(g_{\mu \nu })`$; $`a^{(ij)}(a^{ij}+a^{ji})/2`$; $`a^{[ij]}(a^{ij}a^{ji})/2`$; $`ϵ^{ijk}`$ is the totally antisymmetric Levi-Civita symbol $`(ϵ^{123}=+1)`$. We use a multi-index notation for products of vector components and partial derivatives, and for multiple spatial indices: $`x^{ij\mathrm{}k}x^ix^j\mathrm{}x^k`$, $`_{ij\mathrm{}k}_i_j\mathrm{}_k`$, with a capital letter superscript denoting an abstract product of that dimensionality: $`x^Qx^{i_1}x^{i_2}\mathrm{}x^{i_q}`$ and $`_Q_{i_1}_{i_2}\mathrm{}_{i_q}`$. Also, for a tensor of rank $`Q`$, $`f^Qf^{i_1i_2\mathrm{}i_q}`$. Angular brackets around indices denote symmetric, trace-free (STF) combinations (see Appendix A for definitions). Spatial indices are freely raised and lowered with $`\delta ^{ij}`$ and $`\delta _{ij}`$.
## II Foundations of DIRE
### A The relaxed Einstein equations
We begin by reviewing the method for recasting the Einstein Equations
$`R^{\alpha \beta }{\displaystyle \frac{1}{2}}g^{\alpha \beta }R=8\pi T^{\alpha \beta },`$ (4)
into their “relaxed” form. Here $`R^{\alpha \beta }`$ and $`R`$ are the Ricci tensor and scalar, respectively, $`g^{\alpha \beta }`$ is the spacetime metric and $`T^{\alpha \beta }`$ is the stress-energy tensor of the matter. We define the potential
$`h^{\alpha \beta }\eta ^{\alpha \beta }(g)^{1/2}g^{\alpha \beta },`$ (5)
(see e.g. ) and choose a particular coordinate system defined by the deDonder or harmonic gauge condition
$`h^{\alpha \beta },_\beta =0.`$ (6)
With these definitions the Einstein equations (4) take the form
$`\mathrm{}h^{\alpha \beta }=16\pi \tau ^{\alpha \beta },`$ (7)
where $`\mathrm{}^2/t^2+^2`$ is the flat-spacetime wave operator. The source on the right-hand side is given by the “effective” stress-energy pseudotensor
$`\tau ^{\alpha \beta }=(g)T^{\alpha \beta }+(16\pi )^1\mathrm{\Lambda }^{\alpha \beta },`$ (8)
where $`\mathrm{\Lambda }^{\alpha \beta }`$ is the non-linear “field” contribution given by
$$\mathrm{\Lambda }^{\alpha \beta }=16\pi (g)t_{LL}^{\alpha \beta }+(h^{\alpha \mu },_\nu h^{\beta \nu },_\mu h^{\alpha \beta },_{\mu \nu }h^{\mu \nu }),$$
(9)
and $`t_{LL}^{\alpha \beta }`$ is the “Landau-Lifshitz” pseudotensor, given by
$`16\pi (g)t_{LL}^{\alpha \beta }`$ $``$ $`g_{\lambda \mu }g^{\nu \rho }h_{}^{\alpha \lambda }{}_{,\nu }{}^{}h_{}^{\beta \mu }{}_{,\rho }{}^{}+{\displaystyle \frac{1}{2}}g_{\lambda \mu }g^{\alpha \beta }h_{}^{\lambda \nu }{}_{,\rho }{}^{}h_{}^{\rho \mu }{}_{,\nu }{}^{}2g_{\mu \nu }g^{\lambda (\alpha }h_{}^{\beta )\nu }{}_{,\rho }{}^{}h_{}^{\rho \mu }{}_{,\lambda }{}^{}`$ (11)
$`+{\displaystyle \frac{1}{8}}(2g^{\alpha \lambda }g^{\beta \mu }g^{\alpha \beta }g^{\lambda \mu })(2g_{\nu \rho }g_{\sigma \tau }g_{\rho \sigma }g_{\nu \tau })h_{}^{\nu \tau }{}_{,\lambda }{}^{}h_{}^{\rho \sigma }{}_{,\mu }{}^{}.`$
By virtue of the gauge condition (6), this source term satisfies the conservation law
$$\tau _{}^{\alpha \beta }{}_{,\beta }{}^{}=0,$$
(12)
which is equivalent to the equation of motion of the matter
$$T_{}^{\alpha \beta }{}_{;\beta }{}^{}=0.$$
(13)
Equation (7) is exact, and relies only on the assumption that spacetime can be covered by harmonic coordinates. It is called “relaxed” because it can be solved formally as a functional of source variables without specifying the motion of the source. Then, the harmonic gauge condition, Eq. (6) or the equations of motion, Eq. (13) are imposed to determine the metric as a function of spacetime.
Notice that the “source” in Eq. (7) contains a gravitational part that depends explicitly on $`h^{\alpha \beta }`$, the very quantity for which we are trying to solve. Also, we can expect $`\tau ^{\alpha \beta }`$, which depends on the fields $`h^{\alpha \beta }`$, to have infinite spatial extent. Indeed the very outgoing radiation that we hope to calculate, will, at some level of approximation, serve as a contribution to the source, thus generating an additional component of the radiation.
Another complication in Eq. (7) is that the second derivative term $`h^{\alpha \beta },_{\mu \nu }h^{\mu \nu }`$ in the source really “belongs” on the left-hand side with the other second derivative terms in the wave operator. This term modifies the propagation characteristics of the field from the flat-spacetime characteristics represented by the d’Alembertian operator to those of the true null cones of the curved spacetime around the source, which deviate from the flat null cones of the harmonic coordinates. Nevertheless, the DIRE technique automatically recovers the leading manifestations of this effect, commonly known as “tails”.
The material will be modeled as perfect fluid, having stress-energy tensor
$`T^{\alpha \beta }(\rho +p)u^\alpha u^\beta +pg^{\alpha \beta },`$ (14)
where $`\rho `$ and $`p`$ are the locally measured energy density and pressure, respectively, and $`u^\alpha `$ is the four-velocity of an element of fluid. Until we begin to apply our results to specific physical situations, such as compact binary stars, we will have no need to specialize $`T^{\alpha \beta }`$ further.
### B Near-zone and radiation-zone
We consider the material source to consist of a bound system of characteristic size $`𝒮`$, with a suitably defined center of mass chosen to be at the origin of coordinates, $`𝐗=0`$. The source zone then consists of the world tube $`𝒯=\{x^\alpha |r<𝒮,\mathrm{}<t<\mathrm{}\}`$. Outside $`𝒯`$, $`T^{\alpha \beta }=0`$.
The fluid is assumed to move with characteristic velocity $`v1`$. The characteristic reduced wavelength of gravitational radiation, $`\lambda {}_{}{}^{}=\lambda /2\pi 𝒮/v`$ serves to define the boundary of the near zone, defined to be the world tube $`𝒟=\{x^\alpha |r<,\mathrm{}<t<\mathrm{}`$ }. Within the near zone, the gravitational fields can be treated as almost instantaneous functions of the source variables, i.e. retardation can be ignored or treated as a small perturbation of instantaneous solutions. For physical situations of interest, up to the point where the post-Newtonian approximation breaks down, $`𝒮`$. The region exterior to the near zone is the radiation zone, $`r>`$.
The formal “solution” to Eq. (7) with an outgoing wave boundary condition can be written down in terms of the retarded, flat-space Green function:
$`h^{\alpha \beta }(t,𝐱)=`$ $`4{\displaystyle \frac{\tau ^{\alpha \beta }(t^{},𝐱^{})\delta (t^{}t+|𝐱𝐱^{}|)}{|𝐱𝐱^{}|}d^4x^{}},`$ (15)
but is really just a conversion of the differential equation (7) to an integral equation. It represents an integration of $`\tau ^{\alpha \beta }/|𝐱𝐱^{}|`$ over the past harmonic null cone $`𝒞`$ emanating from the field point $`(t,𝐱)`$ (see Figs. 1 and 2). This past null cone intersects the world tube $`𝒟`$ enclosing the near zone at the three-dimensional hypersurface $`𝒩`$. Thus the integral of Eq. (15) consists of two pieces, an integration over the hypersurface $`𝒩`$, and an integration over the rest of the past null cone $`𝒞𝒩`$. Each of these integrations will be treated differently. We will also treat differently the two cases in which (a) the field point is outside the near zone, and (b) the field point is within the near zone (Fig. 1). The former case will be relevant for calculating the gravitational-wave signal, while the latter will be important for calculating field contributions to $`\tau ^{\alpha \beta }`$ that must be integrated over the near zone, as well as for calculating fields that enter the equations of motion for the source.
### C Radiation-zone field point, inner integration
For a field point in the radiation zone, and integration over the near zone (inner integral), we first carry out the $`t^{}`$ integration in Eq. (15), to obtain
$`h_𝒩^{\alpha \beta }(t,𝐱)=`$ $`4{\displaystyle _𝒩}{\displaystyle \frac{\tau ^{\alpha \beta }(t|𝐱𝐱^{}|,𝐱^{})}{|𝐱𝐱^{}|}}d^3x^{}.`$ (16)
Within the near zone, the spatial integration variable $`𝐱^{}`$ satisfies $`|𝐱^{}|<r`$, where the distance to the field point $`r=|𝐱|`$. Expanding the $`x^{}`$-dependence in both occurrences of $`|𝐱𝐱^{}|`$ in the integrand in powers of $`|𝐱^{}|/r`$, it is straightforward to show that
$`h_𝒩^{\alpha \beta }(t,𝐱)=`$ $`4{\displaystyle \underset{q=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^q}{q!}}_Q\left({\displaystyle \frac{1}{r}}M^{\alpha \beta Q}(u)\right),`$ (17)
where
$`M^{\alpha \beta Q}(u){\displaystyle _{}}\tau ^{\alpha \beta }(u,𝐱^{})x_{}^{}{}_{}{}^{Q}d^3x^{}.`$ (18)
In Eqs. (17) and (18), the index $`Q`$ is a multi-index, such that $`_Q_{i_1}_{i_2}\mathrm{}_{i_q}`$ and the superscript $`Q`$ in $`M^{\alpha \beta Q}`$ denotes $`i_1i_2\mathrm{}i_q`$, with summation over repeated indices assumed. The integrations in Eq. (18) are now over the hypersurface $``$, which is the intersection of the near-zone world-tube with the constant-time hypersurface $`t_{}=u=tr`$. Roughly speaking, each term in the Taylor series is smaller than its predecessor by a factor of order $`v1`$, provided we restrict attention to slow-motion sources.
Note that the field and source variables appearing in the integrand $`\tau ^{\alpha \beta }`$ are evaluated at the single retarded time $`u`$; however, because the field contributions to $`\tau ^{\alpha \beta }`$ fall off as some power of $`r`$, one can expect to encounter integrals that depend on positive powers of the radius $``$ of the boundary of integration, especially in some of the higher-order moments. If this boundary were to be formally taken to $`\mathrm{}`$ (as has been the conventional approach in the past), these integrals would diverge. Instead we shall demonstrate (Sec. II I and Appendix B) that such $``$-dependent effects are precisely cancelled by contributions from the “outer” integral.
For the gravitational-wave signal, we need only to focus on the spatial components of $`h^{\alpha \beta }`$, and on the leading component in $`1/R`$, where $`R`$ is the distance to the detector. Using the fact that $`u_{,i}=\widehat{N}^i`$, where $`\widehat{𝐍}𝐱/R`$ denotes the observation direction, we obtain
$`h_𝒩^{ij}(t,𝐱)={\displaystyle \frac{4}{R}}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{m!}}{\displaystyle \frac{^m}{t^m}}{\displaystyle _{}}\tau ^{ij}(u,𝐱^{})(\widehat{𝐍}𝐱^{})^md^3x^{}+O(R^2).`$ (19)
### D Radiation-zone field point, outer integration
By making a change of integration variable from $`(r^{},\theta ^{},\varphi ^{})`$ to $`(u^{},\theta ^{},\varphi ^{})`$, where
$$tu^{}=r^{}+|𝐱𝐱^{}|,$$
(20)
we can write the integral over the rest of the past null cone $`𝒞𝒩`$ in the form
$`h_{𝒞𝒩}^{\alpha \beta }(t,𝐱)=4{\displaystyle _{\mathrm{}}^u}𝑑u^{}{\displaystyle _{𝒞𝒩}}{\displaystyle \frac{\tau ^{\alpha \beta }(u^{}+r^{},𝐱^{})}{tu^{}\widehat{𝐧}^{}𝐱}}[r^{}(u^{},\mathrm{\Omega }^{})]^2d^2\mathrm{\Omega }^{},`$ (21)
where, from Eq. (20)
$`r^{}(u^{},\mathrm{\Omega }^{})=[(tu^{})^2r^2]/[2(tu^{}\widehat{𝐧}^{}𝐱)].`$ (22)
This change of variables represents an integration first over the two-dimensional intersection of the past null cone $`𝒞`$ with the future null cone $`t^{}=u^{}+r^{}`$ emanating from the center of mass of the system at $`t_{\mathrm{CM}}=u^{}`$ (Fig. 2), followed by the $`u^{}`$-integration over all such future-directed cones, starting from the infinite past, and terminating in the cone emanating from the center of mass at time $`u`$, which is tangent to the past null cone of the observation point.
For explicit calculations, it is useful to choose the field point $`𝐱`$ to be in the z-direction, so that $`\widehat{𝐧}^{}𝐱=r\mathrm{cos}\theta ^{}`$, and to write the outer integral in the form
$`h_{𝒞𝒩}^{\alpha \beta }(t,𝐱)`$ $`=`$ $`4{\displaystyle _{u2}^u}𝑑u^{}{\displaystyle _0^{2\pi }}𝑑\varphi ^{}{\displaystyle _{1\alpha }^1}{\displaystyle \frac{\tau ^{\alpha \beta }(u^{}+r^{},𝐱^{})}{tu^{}\widehat{𝐧}^{}𝐱}}[r^{}(u^{},\mathrm{\Omega }^{})]^2d\mathrm{cos}\theta ^{}`$ (24)
$`+4{\displaystyle _{\mathrm{}}^{u2}}𝑑u^{}{\displaystyle \frac{\tau ^{\alpha \beta }(u^{}+r^{},𝐱^{})}{tu^{}\widehat{𝐧}^{}𝐱}[r^{}(u^{},\mathrm{\Omega }^{})]^2d^2\mathrm{\Omega }^{}},`$
where
$`\alpha (u^{})=(uu^{})(2r2+uu^{})/2r.`$ (25)
The incomplete angular integration in the first integral of Eq. (24) reflects the fact that for $`uu^{}u2`$, the two-dimensional intersections meet the boundary of the near zone. For $`u^{}<u2`$, the angular integration covers the full $`4\pi `$. Note that $`\tau ^{\alpha \beta }`$ contains only field contributions evaluated in the radiation zone; because they are themselves retarded, the “time dependence” $`u^{}+r^{}=t|𝐱𝐱^{}|`$ is approximately constant over each angular integration, since it follows the hypersurface $`t|𝐱|=u=`$ constant, and the dominant contribution to the fields comes from $`|𝐱^{}|<`$. This allows a kind of slow-motion, multipole expansion to be exploited in evaluating these integrals, despite their range well outside the near zone.
### E Near-zone field point, inner integration
In this case, in Eq. (15), both $`𝐱`$ and $`𝐱^{}`$ are within the near zone, hence $`|𝐱𝐱^{}|2`$. Consequently, the variation in retarded time can be treated as a small perturbation, since $`\tau ^{\alpha \beta }`$ varies on a time scale $``$. We therefore expand the retardation in powers of $`|𝐱𝐱^{}|`$, to obtain
$`h_𝒩^{\alpha \beta }(t,𝐱)=`$ $`4{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{m!}}{\displaystyle \frac{^m}{t^m}}{\displaystyle _{}}\tau ^{\alpha \beta }(t,𝐱^{})|𝐱𝐱^{}|^{m1}d^3x^{},`$ (26)
where $``$ here denotes the intersection of the hypersurface $`t=`$ constant with the near-zone world-tube. This version will be used for explicit calculations of the near-zone metric for use in the equations of motion. However an alternative formulation will be useful for studying the $``$-dependence of the inner integrals; substituting the general Taylor expansion $`|𝐱𝐱^{}|^{m1}=\mathrm{\Sigma }_{q=0}^{\mathrm{}}(1)^q(q!)^1x_{<}^{}{}_{}{}^{Q}_{>Q}(r_{>}^{}{}_{}{}^{m1})`$, where $`<(>)`$ denotes the lesser (greater) of $`|𝐱|`$ and $`|𝐱^{}|`$, we obtain
$`h_𝒩^{\alpha \beta }(t,𝐱)`$ $`=`$ $`4{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^m}{m!}}{\displaystyle \frac{^m}{t^m}}{\displaystyle \underset{q=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^q}{q!}}{\displaystyle _{}}\tau ^{\alpha \beta }(t,𝐱^{})x_{<}^{}{}_{}{}^{Q}_{>}^{}{}_{Q}{}^{}(r_{>}^{}{}_{}{}^{m1})d^3x^{}.`$ (27)
### F Near-zone field point, outer integration
The formulae from Section II D, such as (22) and (25), carry over to this case with the result,
$`h_{𝒞𝒩}^{\alpha \beta }(t,𝐱)`$ $`=`$ $`4{\displaystyle _{u2}^{u2+2r}}𝑑u^{}{\displaystyle _0^{2\pi }}𝑑\varphi ^{}{\displaystyle _{1\alpha }^1}{\displaystyle \frac{\tau ^{\alpha \beta }(u^{}+r^{},𝐱^{})}{tu^{}\widehat{𝐧}^{}𝐱}}[r^{}(u^{},\mathrm{\Omega }^{})]^2d\mathrm{cos}\theta ^{}`$ (29)
$`+4{\displaystyle _{\mathrm{}}^{u2}}𝑑u^{}{\displaystyle \frac{\tau ^{\alpha \beta }(u^{}+r^{},𝐱^{})}{tu^{}\widehat{𝐧}^{}𝐱}[r^{}(u^{},\mathrm{\Omega }^{})]^2d^2\mathrm{\Omega }^{}}.`$
Notice that the $`u^{}`$ integration ends at $`u2+2r`$ rather than $`u`$ because that corresponds to the last future null cone that intersects points in the far zone.
### G Iteration of the relaxed Einstein equations
Because the field $`h^{\alpha \beta }`$ appears in the source of the field equation, the usual method of solution is to iterate: substitute $`h^{\alpha \beta }=0`$ in the right-hand side of Eq. (15) and solve for the first-iterated $`{}_{1}{}^{}h_{}^{\alpha \beta }`$; substitute that into Eq. (15) and solve for the second-iterated $`{}_{2}{}^{}h_{}^{\alpha \beta }`$, and so on (imposing the gauge condition Eq. (6) consistently at each order). The general sequence of iterations is shown shematically in Fig. 3. The matter variables $`m_A`$ and the $`(N1)`$-iterated field $`{}_{N1}{}^{}h_{}^{\alpha \beta }`$ are used to determine $`{}_{N1}{}^{}T_{}^{\alpha \beta }`$ and $`{}_{N1}{}^{}\mathrm{\Lambda }_{}^{\alpha \beta }`$. Eq. (15) then yields $`{}_{N}{}^{}h_{}^{\alpha \beta }`$ as a function of spacetime and a functional of the matter variables. Then, if one wishes to determine the motion of the source, one substitutes $`{}_{N}{}^{}h_{}^{\alpha \beta }`$ into the matter stress-energy tensor, and obtains the equations of motion from $`{}_{N}{}^{}_{\beta }^{}({}_{N}{}^{}T_{}^{\alpha \beta })=0`$ where $`{}_{N}{}^{}_{\beta }^{}`$ denotes the covariant derivative using the $`N`$th iterated field. If one wishes to determine the $`N`$th iterated gravitational field as a function of spacetime (i.e. with the matter variables determined as functions of spacetime to a consistent order), then one only needs to solve the equations of motion $`{}_{N1}{}^{}_{\beta }^{}({}_{N1}{}^{}T_{}^{\alpha \beta })=0`$, which are equivalent to the $`N`$-th iterated gauge condition $`{}_{N}{}^{}h_{}^{\alpha \beta }{}_{,\beta }{}^{}=0`$.
### H General structure of the outer integrals
At the first iteration, the solution is simply linearized general relativity. With $`{}_{0}{}^{}h_{}^{\alpha \beta }=0`$ substituted into the right-hand-side of Eq. (15), the outer integrals vanish, and the inner integrals over the special relativistic $`T^{\alpha \beta }`$ have compact support. There is no $``$-dependence in the integrals, trivially. For field points outside the source ($`|𝐱|>|𝐱^{}|`$), within both the near and far zones , the first-iterated $`{}_{1}{}^{}h_{}^{\alpha \beta }`$ takes the form of Eq. (17). Since $`M^{\alpha \beta Q}`$ is a function only of $`u=tr`$, the spatial gradients $`_Q`$ produce only unit radial vectors $`\widehat{n}^i`$, powers of $`r`$, and retarded time derivatives of $`M^{\alpha \beta Q}`$. Products of $`\widehat{n}^i`$ can be grouped into symmetric trace-free (STF) products $`\widehat{n}^{<L>}`$, which are analogous to $`Y_{LM}`$ (see Appendix A for useful formulae related to STF products). Thus, outside the source, $`{}_{1}{}^{}h_{}^{\alpha \beta }`$ can be written as a sequence of terms of the form
$${}_{1}{}^{}h_{}^{\alpha \beta }{}_{B,L}{}^{}(t,𝐱)=f_{B,L}(u)\widehat{n}^{<L>}r^B.$$
(30)
At the second iteration, in the far zone, $`T^{\alpha \beta }=0`$, and $`{}_{1}{}^{}\mathrm{\Lambda }_{}^{\alpha \beta }(u^{}+r^{},𝐱^{})`$ consists of products of spatial and temporal derivatives of $`{}_{1}{}^{}h_{}^{\alpha \beta }(u^{}+r^{},𝐱^{})`$. It therefore can also be expressed as a sequence of terms of the form
$$\mathrm{\Lambda }^{\alpha \beta }(u^{}+r^{},𝐱^{})f_{B,L}(u^{})\widehat{n}^{<L>}r_{}^{}{}_{}{}^{B}.$$
(31)
Whenever the source at a given $`(N1)`$ iteration takes this form, it is straightforward to evaluate the general form of the outer integrals for the $`N`$th iterate. Defining the new variables $`\zeta (tu^{})/r=1+(uu^{})/r`$, $`y=\widehat{𝐧}\widehat{𝐧}^{}=\mathrm{cos}\theta `$, we find, from Eq. (22),
$$r^{}=r(\zeta ^21)/2(\zeta y).$$
(32)
Substituting Eqs. (31) and (32) into Eq. (24), and changing to integration variables $`\zeta `$, $`y`$ and $`\varphi `$, we obtain
$${}_{N}{}^{}h_{𝒞𝒩}^{\alpha \beta }{}_{B,L}{}^{}=\frac{1}{2}\left(\frac{2}{r}\right)^{B2}\widehat{n}^{<L>}_1^1P_L(y)𝑑y_{\zeta (y)}^{\mathrm{}}\frac{(\zeta y)^{B3}}{(\zeta ^21)^{B2}}f_{B,L}(ur(\zeta 1))𝑑\zeta ,$$
(33)
where $`\zeta (y)=z+\sqrt{z^22zy+1}`$, $`z=/r`$, and $`P_L(y)`$ is the Legendre polynomial.
For far-zone field points, $`z<1`$; Taylor expanding $`f_{B,L}(ur(\zeta 1))`$ about $`u`$, we obtain, for $`B>2`$,
$${}_{N}{}^{}h_{𝒞𝒩}^{\alpha \beta }{}_{B,L}{}^{}=\left(\frac{2}{r}\right)^{B2}\widehat{n}^{<L>}\underset{q=0}{\overset{\mathrm{}}{}}𝒟_{B,L}^q(z)r^q\frac{d^qf_{B,L}(u)}{du^q},$$
(34)
where the coefficients $`𝒟_{B,L}^q(z)`$ are given by
$$𝒟_{B,L}^q(z)=\frac{()^q}{q!}_1^{1+2z}\frac{(\zeta 1)^q}{(\zeta ^21)^{B2}}A_{B,L}(\zeta ,\alpha )𝑑\zeta \underset{p=0}{\overset{q}{}}k_{B,L}^{(qp+1)}(1+2z)\frac{(2z)^p}{p!},$$
(35)
where
$`A_{B,L}(\zeta ,\alpha )`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _{1\alpha }^1}P_L(y)(\zeta y)^{B3}𝑑y,`$ (37)
$`\alpha `$ $`(\zeta 1)(\zeta +12z)/2z,`$ (38)
$`dk_{B,L}^{(m)}(\zeta )/d\zeta `$ $`k_{B,L}^{(m1)}(\zeta ),m1;`$ (39)
$`k_{B,L}^{(0)}(\zeta )`$ $`A_{B,L}(\zeta ,2)/(\zeta ^21)^{B2}.`$ (40)
The case $`B=2`$ is special, and leads to the result
$${}_{N}{}^{}h_{𝒞𝒩}^{\alpha \beta }{}_{2,L}{}^{}=\frac{\widehat{n}^{<L>}}{r}_0^{\mathrm{}}f_{2,L}(us)Q_L\left(1+\frac{s}{r}\right)𝑑s+\widehat{n}^{<L>}\underset{q=0}{\overset{\mathrm{}}{}}𝒟_{2,L}^q(z)r^q\frac{d^qf_{2,L}(u)}{du^q},$$
(41)
where
$$𝒟_{2,L}^q(z)=\frac{()^{q+1}}{2q!}_1^{1+2z}(\zeta 1)^q𝑑\zeta _1^{1\alpha }\frac{P_L(y)}{(\zeta y)}𝑑y,$$
(42)
where $`Q_L(y)`$ is the Legendre function.
Notice that, for $`B2`$, the outer integral returns a result of the same generic form as the input function. The case $`B=2`$ returns terms with a logarithmic dependence on $`r`$ (via the $`Q_L`$’s); terms of this form are called “tails”.
Similarly, for field points in the near zone, $`z>1`$, we Taylor expand $`f_{B,L}(ur(\zeta 1))`$ about $`u+r=t`$, and obtain, for $`B>2`$,
$${}_{N}{}^{}h_{𝒞𝒩}^{\alpha \beta }{}_{B,L}{}^{}=\left(\frac{2}{r}\right)^{B2}\widehat{n}^{<L>}\underset{q=0}{\overset{\mathrm{}}{}}_{B,L}^q(z)r^q\frac{d^qf_{B,L}(t)}{dt^q},$$
(43)
where the coefficients $`_{B,L}^q(z)`$ are given by
$$_{B,L}^q(z)=\frac{()^q}{q!}_{2z1}^{2z+1}\frac{\zeta ^q}{(\zeta ^21)^{B2}}A_{B,L}(\zeta ,\alpha )𝑑\zeta \underset{p=0}{\overset{q}{}}k_{B,L}^{(qp+1)}(1+2z)\frac{(12z)^p}{p!},$$
(44)
and, for $`B=2`$,
$${}_{N}{}^{}h_{𝒞𝒩}^{\alpha \beta }{}_{2,L}{}^{}=\frac{\widehat{n}^{<L>}}{r}_0^{\mathrm{}}f_{2,L}(us)Q_L\left(1+\frac{s}{r}\right)𝑑s+\widehat{n}^{<L>}\underset{q=0}{\overset{\mathrm{}}{}}_{2,L}^q(z)r^q\frac{d^qf_{2,L}(t)}{dt^q},$$
(45)
where
$$_{2,L}^q(z)=\frac{()^{q+1}}{q!}\left\{_1^{2z1}\zeta ^qQ_L(\zeta )𝑑\zeta +\frac{1}{2}_{2z1}^{2z+1}\zeta ^q𝑑\zeta _1^{1\alpha }\frac{P_L(y)}{(\zeta y)}𝑑y\right\}.$$
(46)
Notice that, for near-zone field points, the functions $`f_{B,L}`$ are evaluated at the local time $`t`$, not retarded time $`u`$.
### I Cancellation of $``$ dependence
It is evident that the inner integrals and outer integrals for the field $`h^{\alpha \beta }`$ will separately depend upon the radius $``$ of the boundary between the near zone and the far zone. But since each integral was simply a rewriting of a piece of the original integral, Eq. (15), which had no $``$ dependence, it is equally evident that the separate $``$-dependences must cancel between the inner and outer integrals. In , referred to hereafter as WW, we demonstrated such a cancellation explicitly for contributions to the gravitational waveform at 2PN order that depended on positive powers of $``$. Here we demonstrate the cancellation generally, for both near-zone and far-zone field points, for arbitrary powers of $``$ (including $`\mathrm{ln}`$) and to an order of iteration sufficient for our purposes.
The proof proceeds by induction. First, as we pointed out above, the first-iterated field $`{}_{1}{}^{}h_{}^{\alpha \beta }`$ is trivially independent of $``$.
Secondly, we assume that the $`(N1)`$-iterated field does not depend on $``$, i.e. that all $``$-dependence cancels at this order of iteration. We wish to demonstrate that this implies cancellation of $``$-dependence in the $`N`$-iterated field. The proof consists of considering the limiting behavior of the inner and outer integrals for $`{}_{N}{}^{}h_{}^{\alpha \beta }`$ in the vicinity of $`|𝐱^{}|`$. Here $`T^{\alpha \beta }`$ vanishes, and we only need to consider $`{}_{N1}{}^{}\mathrm{\Lambda }_{}^{\alpha \beta }`$, which is a functional of $`{}_{N1}{}^{}h_{}^{\alpha \beta }`$. We have already seen that, in the far zone, $`{}_{N1}{}^{}\mathrm{\Lambda }_{}^{\alpha \beta }`$ can be decomposed into terms of the form $`f_{B,L}(u)\widehat{n}^{<L>}r^B`$. (We consider tail contributions with $`\mathrm{ln}r`$ dependence separately.) Since the $`N1`$ iterated field does not depend on $``$ by assumption, continuity of the fields means that $`{}_{N1}{}^{}\mathrm{\Lambda }_{}^{\alpha \beta }`$ will have this same form just inside the near zone. Thus we will calculate the limiting behavior of the inner integral of a term of this form as the integration variable $`r^{}`$ from below, and compare its $``$ dependence with that of the outer integral of the same term.
For far-zone field points, we must calculate the $``$ dependence of the moments $`M^{\alpha \beta Q}`$, and substitute into Eq. (17); after considerable algebra (see Appendix B), we obtain, for the limiting behavior of $`{}_{N}{}^{}h_{𝒩}^{\alpha \beta }`$ as the integration variable approaches $``$ from inside,
$${}_{N}{}^{}h_{𝒩}^{\alpha \beta }{}_{B,L}{}^{}\left(\frac{2}{r}\right)^{B2}\widehat{n}^{<L>}\underset{q=0}{\overset{\mathrm{}}{}}𝒟_{B,L}^{\mathrm{in},q}(z)r^q\frac{d^qf_{B,L}(u)}{du^q},$$
(47)
where
$`𝒟_{B,L}^{\mathrm{in},q}(z)`$ $`=`$ $`{\displaystyle \underset{m=0}{\overset{q}{}}}{\displaystyle \underset{j=0}{\overset{j_{max}}{}}}{\displaystyle \frac{()^m2^{2+jB}}{m!(qmj+2L+1)!}}{\displaystyle \frac{[\frac{1}{2}(qmj+2L)]!}{[\frac{1}{2}(qmj)]!}}{\displaystyle \frac{(2Lj)!}{j!(Lj)!}}`$ (51)
$`\times \{\begin{array}{cc}z^{3+LB+qj}/(3+LB+qj)\hfill & 3+LB+qj0\hfill \\ \mathrm{ln}\hfill & 3+LB+qj=0.\hfill \end{array}`$
where $`j_{max}=\mathrm{lesser}\mathrm{of}\{qm,L\}`$, and $`qmj=\mathrm{even}\mathrm{integer}0`$. Eq. (47) is of the same form as the outer integral for far-zone field points, Eq. (34). The coefficients $`𝒟_{B,L}^q(z)`$ from the outer integrals are most easily evaluated using computer algebra methods (we calculated the coefficients using independent Maple and Mathematica programs); the result is, for each $`B`$, $`L`$ and $`q`$,
$$𝒟_{B,L}^{\mathrm{in},q}(z)+\{z\mathrm{dependent}\mathrm{part}\mathrm{of}𝒟_{B,L}^q(z)\}=0.$$
(52)
Thus the $``$ dependence cancels term by term. Similar cancellation occurs for near-zone field points, as well as for the case where the integrand has $`r^B\mathrm{ln}r`$ dependence. Details are given in Appendix B.
This cancellation, while inevitable, has practical consequences, in the following sense. In calculating the inner contributions to the fields, we must integrate over a finite hypersurface, $``$, sources that extend throughout $``$. Consequently, any such integral will have terms that are independent of $``$, as well as terms that depend on $`^q`$ or $`\mathrm{ln}`$. Because we know that all terms of the latter form cancel with contributions from the outer integrals in the final expression for the field, we can drop them in any individual result. Similarly, we can drop all $``$-dependent terms that arise in any individual outer integral. This provides a kind of regularization of integrals, that cures the problem of divergent integrals that haunted earlier slow-motion methods. In fact, one can show that there is a close connection between this method of regularization and the method of analytic continuation used by Blanchet .
Thus, our procedure for determining the field is to determine separately $`h_𝒩^{\alpha \beta }`$ and $`h_{𝒞𝒩}^{\alpha \beta }`$ to a given PN order, keeping only $``$-independent terms in each expression, then sum them to obtain
$$h^{\alpha \beta }=h_𝒩^{\alpha \beta }+h_{𝒞𝒩}^{\alpha \beta }.$$
(53)
## III Weak field, slow-motion approximation
We now turn to a discussion of the number $`N`$ of iterations needed to derive equations of motion or gravitational waveforms of a desired accuracy, in a weak-field, slow-motion approximation.
We assume that, for the fluid source,
$`v^2m/𝒮p/\rho ϵ1,`$ (54)
where $`ϵ`$ will be used as an expansion parameter. But from the nature of the iteration procedure, it is evident that each iteration of the field introduces corrections of order $`m/𝒮`$. In terms of $`ϵ`$, $`m`$ and $`𝒮`$ the equations of motion (1) can be rewritten schematically as
$$d𝐯/dt(m/𝒮^2)[1+O(ϵ)+O(ϵ^2)+O(ϵ^{5/2})+O(ϵ^3)+O(ϵ^{7/2})+\mathrm{}],$$
(55)
where the terms inside the square brackets represent the Newtonian, post-Newtonian, 2PN, 2.5PN (radiation-reaction), 3PN, and 3.5PN (radiation-reaction) terms respectively. For a term of order $`ϵ^N`$, the largest number of powers of $`m/𝒮`$ that can appear in it (including one power from the $`m/𝒮^2`$ prefactor) is $`N+1`$. The radiation reaction terms of order $`ϵ^{N+1/2}`$ must contain an odd number of velocities (in order to be odd under time-reversal), thus the maximum number of powers of $`m/𝒮`$ for them is also $`N+1`$. Since one iteration gives the Newtonian potential, which yields the Newtonian equations of motion ($`N=1`$), then, to obtain the 1PN terms ($`N=2`$), one must have the second iterated field, to obtain the 2PN and 2.5PN terms ($`N=3`$), one must have the third iterated field, while to obtain the 3PN and 3.5PN terms ($`N=4`$), one must have the fourth iterated field.
Similarly, to obtain a result for the waveform accurate to the order of the quadrupole formula, $`h\ddot{}^{ij}/R(m/R)(v^2+m/𝒮)ϵ^2`$ ($`N=2`$), the second-iterated field is needed. Note that the term $`m/𝒮`$ in $`\ddot{}^{ij}`$ arises through the use of the Newtonian equation of motion. Then, to obtain the 1PN, 2PN and 3PN corrections to the quadrupole approximation, the third, fourth, and fifth-iterated fields are needed, respectively. This would be an impossible task, if it weren’t for the judicious use of the conservation law, Eq. (12). Consider for example, the source $`{}_{N1}{}^{}\tau _{}^{ij}`$ of the $`N`$th iterated gravitational-wave field $`{}_{N}{}^{}h_{𝒩}^{ij}`$, Eq. (19), specifically the leading, $`m=0`$ term. The conservation law, Eq. (12), converts $`{}_{N1}{}^{}\tau _{}^{ij}`$ into two time derivatives of $`{}_{N1}{}^{}\tau _{}^{00}x^ix^j`$ (modulo total divergences). Because of the slow-motion approximation, two time derivatives increase the order by $`ϵ`$, and thus, to sufficient accuracy, only $`{}_{N2}{}^{}\tau _{}^{00}`$ is needed in practice. An important caveat to this is that the surface terms that arise from the total divergences and the outer integrals must formally be evaluated using the $`N1`$ expressions. However, in practice, these terms contribute at sufficiently high order that they can be treated without resort to explicit $`N1`$ expressions. Effectively, the burden of accuracy has been shifted from the $`N`$th-iteration of the field, to the $`N1`$-iterated equations of motion, which enter via the two time-derivatives, and which are needed anyway to evaluate the field as a function of spacetime. Thus, for $`N=2`$, the leading quadrupole approximation, only $`{}_{0}{}^{}\tau _{}^{00}=\rho `$ is needed, together with the Newtonian equations of motion. This circumstance is responsible for the prevalent, but erroneous view that linearized gravity (one iteration) suffices to derive the quadrupole formula. The formula so derived turns out to be “correct”, but its foundation is not (see for discussion).
Thus, in WW, to evaluate the 2PN waveforms (fourth iteration), only second-iterated fields were needed in the source terms. For 3PN waveforms, only third-iterated fields will be needed.
## IV Formal structure of Near-Zone Fields
### A Metric and stress-energy pseudotensor in terms of the fields
We begin by defining a simplified notation for the field $`h^{\alpha \beta }`$:
$`N`$ $``$ $`h^{00}O(ϵ),`$ (56)
$`K^i`$ $``$ $`h^{0i}O(ϵ^{3/2}),`$ (57)
$`B^{ij}`$ $``$ $`h^{ij}O(ϵ^2),`$ (58)
$`B`$ $``$ $`h^{ii}{\displaystyle \underset{i}{}}h^{ii}O(ϵ^2),`$ (59)
where we show the leading order dependence on $`ϵ`$ in the near zone. To obtain the equations of motion to 3.5PN order, we need to determine the components of the physical metric to the following orders: $`g_{00}`$ to $`O(ϵ^{9/2})`$, $`g_{0i}`$ to $`O(ϵ^4)`$ , and $`g_{ij}`$ to $`O(ϵ^{7/2})`$. From the definition (5), one can invert to find $`g_{\alpha \beta }`$ in terms of $`h^{\alpha \beta }`$. Expanding to the required order, we find,
$`g_{00}`$ $`=`$ $`(1{\displaystyle \frac{1}{2}}N+{\displaystyle \frac{3}{8}}N^2{\displaystyle \frac{5}{16}}N^3+{\displaystyle \frac{35}{128}}N^4)+{\displaystyle \frac{1}{2}}B(1{\displaystyle \frac{1}{2}}N+{\displaystyle \frac{3}{8}}N^2)`$ (62)
$`+{\displaystyle \frac{1}{4}}(B^{ij}B^{ij}{\displaystyle \frac{1}{2}}B^2)+{\displaystyle \frac{1}{2}}K^jK^j{\displaystyle \frac{3}{4}}NK^jK^j+O(ϵ^5),`$
$`g_{0i}`$ $`=`$ $`K^i(1{\displaystyle \frac{1}{2}}N{\displaystyle \frac{1}{2}}B+{\displaystyle \frac{3}{8}}N^2)K^jB^{ij}+O(ϵ^{9/2}),`$ (63)
$`g_{ij}`$ $`=`$ $`\delta ^{ij}(1+{\displaystyle \frac{1}{2}}N{\displaystyle \frac{1}{8}}N^2+{\displaystyle \frac{1}{16}}N^3{\displaystyle \frac{1}{4}}NB+{\displaystyle \frac{1}{2}}K^kK^k)`$ (65)
$`+B^{ij}{\displaystyle \frac{1}{2}}B\delta ^{ij}K^iK^j+{\displaystyle \frac{1}{2}}NB^{ij}+O(ϵ^4),`$
$`(g)`$ $`=`$ $`1+NBNB+K^iK^i+O(ϵ^4).`$ (66)
Notice that, in order to find the metric $`g_{\alpha \beta }`$ to the desired order, we must obtain $`N`$ to $`O(ϵ^{9/2})`$, $`K^i`$ to $`O(ϵ^4)`$, $`B^{ij}`$ to $`O(ϵ^{7/2})`$ and $`B`$ to $`O(ϵ^{9/2})`$. In fact, because $`B`$ contributes linearly to $`g_{00}`$, we will treat $`B`$ differently from $`B^{ij}`$.
Using Eq. (IV A), we can express the matter stress-energy tensor $`T^{\alpha \beta }`$, Eq. (14), as a PN expansion. However, the details of such an expansion will depend on the basic variables used to characterize the matter. For example, to discuss the structure of a star in a PN expansion, it is convenient to use the mass-energy density $`\rho `$ and pressure $`p`$, together with an equation of state. However, to discuss the motion of compact bodies in an effective “point-mass” limit, it is more convenient to ignore the pressure totally, and to use the so-called “conserved”, or baryon density, $`\rho ^{}\rho \sqrt{g}u^0`$. For now, we follow the convention of Damour et al., and define the quantities
$`\sigma `$ $``$ $`T^{00}+T^{ii},`$ (67)
$`\sigma ^i`$ $``$ $`T^{0i},`$ (68)
$`\sigma ^{ij}`$ $``$ $`T^{ij}.`$ (69)
We will express various potentials formally in terms of these densities, and later make a PN expansion of them in terms of the densities most appropriate to the application.
Substituting the formulae for $`h^{\alpha \beta }`$ and $`g_{\alpha \beta }`$ into Eqs. (9) and (11) for $`\mathrm{\Lambda }^{\alpha \beta }`$, we obtain, to the required order,
$`\mathrm{\Lambda }^{00}`$ $`=`$ $`{\displaystyle \frac{7}{8}}(N)^2`$ (75)
$`+\{{\displaystyle \frac{5}{8}}\dot{N}^2\ddot{N}N2\dot{N}^{,k}K^k+{\displaystyle \frac{1}{2}}K^{i,j}(3K^{j,i}+K^{i,j})`$
$`+\dot{K}^jN^{,j}B^{ij}N^{,ij}+{\displaystyle \frac{1}{4}}NB+{\displaystyle \frac{7}{8}}N(N)^2\}`$
$`+\{K^{k,j}\dot{B}^{jk}+{\displaystyle \frac{1}{4}}B^{jk,l}(B^{jk,l}2B^{kl,j})+{\displaystyle \frac{1}{4}}\dot{N}\dot{B}{\displaystyle \frac{1}{8}}(B)^2+{\displaystyle \frac{1}{4}}\dot{N}N^{,j}K^j`$
$`+{\displaystyle \frac{7}{8}}N^{,j}N^{,k}B^{jk}{\displaystyle \frac{1}{2}}K^jN^{,k}(3K^{j,k}+4K^{k,j}){\displaystyle \frac{7}{8}}N^2(N)^2\}+O(\rho ϵ^4),`$
$`\mathrm{\Lambda }^{0i}`$ $`=`$ $`\left\{N^{,k}(K^{k,i}K^{i,k})+{\displaystyle \frac{3}{4}}\dot{N}N^{,i}\right\}`$ (79)
$`+\{\dot{N}\dot{K}^iN\ddot{K}^i2K^k\dot{K}^{i,k}B^{lm}K^{i,lm}+K^{k,l}(B^{il,k}+B^{ik,l}B^{kl,i})+N^{,k}\dot{B}^{ik}`$
$`{\displaystyle \frac{1}{4}}\dot{N}B^{,i}{\displaystyle \frac{1}{4}}N^{,i}\dot{B}NN^{,k}(K^{k,i}K^{i,k}){\displaystyle \frac{3}{4}}N\dot{N}N^{,i}+{\displaystyle \frac{1}{8}}K^i(N)^2{\displaystyle \frac{1}{4}}K^kN^{,k}N^{,i}\}`$
$`+O(\rho ϵ^{7/2}),`$
$`\mathrm{\Lambda }^{ij}`$ $`=`$ $`{\displaystyle \frac{1}{4}}\{N^{,i}N^{,j}{\displaystyle \frac{1}{2}}\delta ^{ij}(N)^2\}`$ (89)
$`+\{2K^{k,(i}K^{j),k}K^{k,i}K^{k,j}K^{i,k}K^{j,k}+2N^{,(i}\dot{K}^{j)}+{\displaystyle \frac{1}{2}}N^{,(i}B^{,j)}`$
$`{\displaystyle \frac{1}{2}}N(N^{,i}N^{,j}{\displaystyle \frac{1}{2}}\delta ^{ij}(N)^2)\delta ^{ij}(K^{l,k}K^{[k,l]}+N^{,k}\dot{K}^k+{\displaystyle \frac{3}{8}}\dot{N}^2+{\displaystyle \frac{1}{4}}NB)\}`$
$`+\{2\dot{K}^i\dot{K}^j+\dot{B}^{k(i}(K^{j),k}K^{k,j)})2\dot{B}^{ij,k}K^kN\ddot{B}^{ij}B^{ij,lm}B^{lm}`$
$`+B^{ik,l}(B^{jl,k}+B^{jk,l})2B^{kl,(i}B^{j)k,l}+{\displaystyle \frac{1}{2}}B^{kl,i}B^{kl,j}{\displaystyle \frac{1}{4}}B^{,i}B^{,j}`$
$`N(2K^{k,(i}K^{j),k}K^{k,i}K^{k,j}K^{i,k}K^{j,k})+K^kK^{k,(i}N^{,j)}2NN^{,(i}\dot{K}^{j)}{\displaystyle \frac{1}{2}}\dot{N}N^{,(i}K^{j)}`$
$`{\displaystyle \frac{1}{2}}N^{,k}N^{,(i}B^{j)k}{\displaystyle \frac{1}{2}}NN^{,(i}B^{,j)}+{\displaystyle \frac{1}{8}}(N)^2B^{ij}+{\displaystyle \frac{3}{4}}N^2(N^{,i}N^{,j}{\displaystyle \frac{1}{2}}\delta ^{ij}(N)^2)`$
$`+{\displaystyle \frac{1}{8}}\delta ^{ij}[(B)^2+2\dot{N}\dot{B}+8K^{k,l}\dot{B}^{kl}+4B^{kl,m}B^{km,l}2B^{kl,m}B^{kl,m}+3N\dot{N}^2+8NN^{,k}\dot{K}^k`$
$`+8NK^{l,k}K^{[k,l]}4K^kN^{,l}K^{k,l}+2\dot{N}K^kN^{,k}+N^{,k}N^{,l}B^{kl}+2NNB]\}`$
$`+O(\rho ϵ^4),`$
$`\mathrm{\Lambda }^{ii}`$ $`=`$ $`{\displaystyle \frac{1}{8}}(N)^2`$ (95)
$`+\left\{K^{l,k}K^{[k,l]}N^{,k}\dot{K}^k{\displaystyle \frac{1}{4}}NB{\displaystyle \frac{9}{8}}\dot{N}^2+{\displaystyle \frac{1}{4}}N(N)^2\right\}`$
$`+\{2\dot{K}^k\dot{K}^k2\dot{B}^{,k}K^k+3\dot{B}^{kl}K^{k,l}N\ddot{B}+{\displaystyle \frac{3}{4}}\dot{N}\dot{B}+{\displaystyle \frac{1}{8}}(B)^2B^{,lm}B^{lm}`$
$`+{\displaystyle \frac{3}{4}}B^{kl,m}B^{kl,m}+{\displaystyle \frac{1}{2}}B^{kl,m}B^{km,l}NK^{l,k}K^{[k,l]}{\displaystyle \frac{1}{2}}N^{,l}K^kK^{k,l}+NN^{,k}\dot{K}^k`$
$`+{\displaystyle \frac{1}{4}}\dot{N}N^{,k}K^k{\displaystyle \frac{1}{8}}N^{,k}N^{,l}B^{kl}+{\displaystyle \frac{1}{4}}NNB+{\displaystyle \frac{1}{8}}(N)^2B+{\displaystyle \frac{9}{8}}N\dot{N}^2{\displaystyle \frac{3}{8}}N^2(N)^2\}`$
$`+O(\rho ϵ^4),`$
where an overdot denotes $`/t`$. In the above expressions, terms grouped within braces make leading contributions of the same order. For example, in $`\mathrm{\Lambda }^{00}`$, the three groupings correspond to $`O(\rho ϵ)`$, $`O(\rho ϵ^2)`$, and $`O(\rho ϵ^3)`$, respectively.
### B Source moments and other integral quantities
Throughout our calculations, a number of integrals of the source stress-energy pseudotensor occur, for example, in the multipole expansions of Eq. (18). It is useful to define and collect these quantities, and to discuss their properties. All integrals are carried out over a constant time (or constant retarded time) hypersurface $``$, within the near-zone. In general, these integrals will have $``$ dependence, but, in line with the foregoing discussion, we shall consistently drop such terms. The relevant integrals are:
$`P^\mu `$ $``$ $`M^{\mu 0}={\displaystyle _{}}\tau ^{\mu 0}d^3x,`$ (97)
$`^Q`$ $``$ $`M^{00Q}={\displaystyle _{}}\tau ^{00}x^Qd^3x,`$ (98)
$`𝒥^{iQ}`$ $``$ $`ϵ^{iab}M^{0baQ}=ϵ^{iab}{\displaystyle _{}}\tau ^{0b}x^{aQ}d^3x,`$ (99)
$`𝒫^{ijabQ}`$ $``$ $`{\displaystyle _{}}x^{[a}\tau ^{i][j}x^{b]Q}d^3x.`$ (100)
By making use of the equations of motion $`\tau _{}^{\alpha \beta }{}_{,\beta }{}^{}=0`$, we can transform some of these integrals into other forms, modulo surface integrals at the boundary $``$ of the near zone. For example,
$`\dot{P}^\mu `$ $`=`$ $`{\displaystyle _{}}\tau ^{\mu j}d^2S_j,`$ (101)
$`\dot{𝒥}^i`$ $`=`$ $`ϵ^{iab}{\displaystyle _{}}\tau ^{jb}x^ad^2S_j,`$ (102)
$`\dot{}^i`$ $`=`$ $`P^i{\displaystyle _{}}\tau ^{0j}x^id^2S_j.`$ (103)
These identities express the conservation of total energy, momentum and angular momentum, and uniform center-of-mass motion, modulo a flux of gravitational radiation from the system. In calculations, the surface terms must be checked carefully to see if they make $``$-independent contributions to the order considered. For the most part, such surface terms turn out to make no contribution.
Henceforth, we shall set $`^i=\dot{}^i=0`$, which amounts to attaching the origin of coordinates to the center of mass of the system.
Other useful identities include:
$`M^{ij}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\ddot{}^{ij}+{\displaystyle \frac{1}{2}}{\displaystyle _{}}[\tau ^{lm}(x^{ij})_{,l}+\dot{\tau }^{m0}x^{ij}]d^2S_m,`$ (105)
$`M^{ijk}`$ $`=`$ $`{\displaystyle \frac{1}{6}}\ddot{}^{ijk}+{\displaystyle \frac{2}{3}}ϵ^{lk(i}\dot{𝒥}^{l|j)}`$ (107)
$`+{\displaystyle \frac{1}{6}}{\displaystyle _{}}[\tau ^{lm}(x^{ijk})_{,l}+\dot{\tau }^{m0}x^{ijk}]d^2S_m{\displaystyle \frac{2}{3}}{\displaystyle _{}}[\tau ^{l[k}x^{i]j}+\tau ^{l[k}x^{j]i}]d^2S_l,`$
$`M^{ijQ}`$ $`=`$ $`{\displaystyle \frac{1}{(q+1)(q+2)}}\ddot{}^{ijQ}+{\displaystyle \frac{2}{(q+2)}}ϵ^{mk_1(i}\dot{𝒥}^{m|j)k_2\mathrm{}k_q}(\mathrm{sym}\mathrm{k}:\mathrm{Q})+{\displaystyle \frac{8(q1)}{(q+1)}}𝒫^{ij(k_1k_2\mathrm{}k_q)}`$ (110)
$`+{\displaystyle \frac{1}{(q+1)(q+2)}}{\displaystyle _{}}[\tau ^{lm}(x^{ijQ})_{,l}+\dot{\tau }^{m0}x^{ijQ}]d^2S_m`$
$`{\displaystyle \frac{2}{(q+2)}}{\displaystyle _{}}[\tau ^{l[k_1}x^{i]jk_2\mathrm{}k_q}+\tau ^{l[k_1}x^{j]ik_2\mathrm{}k_q}]d^2S_l(\mathrm{sym}\mathrm{k}:\mathrm{Q}),`$
$`M^{0jQ}`$ $`=`$ $`{\displaystyle \frac{1}{q+1}}\dot{}^{jQ}{\displaystyle \frac{q}{q+1}}ϵ^{mj(k_1}𝒥^{m|k_2\mathrm{}k_q)}+{\displaystyle \frac{1}{q+1}}{\displaystyle _{}}\tau ^{0m}x^{jQ}d^2S_m,`$ (111)
where the notation (sym k:Q) means symmetrize on the indices $`k_1`$ through $`k_q`$, and the superscript notation <sup>m|k…)</sup> means that only the indices following the vertical line are involved in symmetrization.
### C Near-zone field expanded to 3.5 PN order
We now carry out the explicit expansion of the near-zone field through 3.5PN order, beginning with the inner integral, Eq. (26), applying the above identities where possible. Inserting powers of $`ϵ`$ to indicate the leading order of each term, we obtain the result
$`N_𝒩`$ $`=`$ $`4ϵ{\displaystyle _{}}{\displaystyle \frac{\tau ^{00}(t,𝐱^{})}{|𝐱𝐱^{}|}}d^3x^{}+2ϵ^2_t^2{\displaystyle _{}}\tau ^{00}(t,𝐱^{})|𝐱𝐱^{}|d^3x^{}{\displaystyle \frac{2}{3}}ϵ^{5/2}\stackrel{(3)}{^{kk}(t)}`$ (119)
$`+{\displaystyle \frac{1}{6}}ϵ^3_t^4{\displaystyle _{}}\tau ^{00}(t,𝐱^{})|𝐱𝐱^{}|^3d^3x^{}`$
$`{\displaystyle \frac{1}{30}}ϵ^{7/2}\{(4x^{kl}+2r^2\delta ^{kl})\stackrel{(5)}{^{kl}(t)}4x^k\stackrel{(5)}{^{kll}(t)}+\stackrel{(5)}{^{kkll}(t)}\}`$
$`+{\displaystyle \frac{1}{180}}ϵ^4_t^6{\displaystyle _{}}\tau ^{00}(t,𝐱^{})|𝐱𝐱^{}|^5d^3x^{}`$
$`{\displaystyle \frac{1}{1260}}ϵ^{9/2}\{3r^4\stackrel{(7)}{^{kk}(t)}+12r^2x^{ij}\stackrel{(7)}{^{ij}(t)}12r^2x^i\stackrel{(7)}{^{ikk}(t)}`$
$`8x^{ijk}\stackrel{(7)}{^{ijk}(t)}+3r^2\stackrel{(7)}{^{iikk}(t)}+12x^{ij}\stackrel{(7)}{^{ijkk}(t)}6x^i\stackrel{(7)}{^{ikkll}(t)}+\stackrel{(7)}{^{iikkll}(t)}\}`$
$`+N_{}+O(ϵ^5),`$
$`K_𝒩^i`$ $`=`$ $`4ϵ^{3/2}{\displaystyle _{}}{\displaystyle \frac{\tau ^{0i}(t,𝐱^{})}{|𝐱𝐱^{}|}}d^3x^{}+2ϵ^{5/2}_t^2{\displaystyle _{}}\tau ^{0i}(t,𝐱^{})|𝐱𝐱^{}|d^3x^{}`$ (124)
$`+{\displaystyle \frac{2}{9}}ϵ^3\{3x^k\stackrel{(4)}{^{ik}(t)}\stackrel{(4)}{^{ikk}(t)}+2ϵ^{mik}\stackrel{(3)}{𝒥^{mk}(t)}\}+{\displaystyle \frac{1}{6}}ϵ^{7/2}_t^4{\displaystyle _{}}\tau ^{0i}(t,𝐱^{})|𝐱𝐱^{}|^3d^3x^{}`$
$`+{\displaystyle \frac{1}{450}}ϵ^4\{30r^2x^k\stackrel{(6)}{^{ik}(t)}10r^2\stackrel{(6)}{^{ikk}(t)}20x^{kl}\stackrel{(6)}{^{ikl}(t)}+15x^k\stackrel{(6)}{^{ikll}(t)}3\stackrel{(6)}{^{ikkll}(t)}`$
$`+ϵ^{mil}[20r^2\stackrel{(5)}{𝒥^{ml}(t)}+40x^{kl}\stackrel{(5)}{𝒥^{mk}(t)}15x^l\stackrel{(5)}{𝒥^{mkk}(t)}30x^k\stackrel{(5)}{𝒥^{mkl}(t)}+12\stackrel{(5)}{𝒥^{mlkk}(t)}]\}`$
$`+K_{}^i+O(ϵ^{9/2}),`$
$`B_𝒩^{ij}`$ $`=`$ $`4ϵ^2{\displaystyle _{}}{\displaystyle \frac{\tau ^{ij}(t,𝐱^{})}{|𝐱𝐱^{}|}}d^3x^{}2ϵ^{5/2}\stackrel{(3)}{^{ij}(t)}+2ϵ^3_t^2{\displaystyle _{}}\tau ^{ij}(t,𝐱^{})|𝐱𝐱^{}|d^3x^{}`$ (130)
$`{\displaystyle \frac{1}{9}}ϵ^{7/2}\{3r^2\stackrel{(5)}{^{ij}(t)}2x^k\stackrel{(5)}{^{ijk}(t)}8x^kϵ^{mk(i}\stackrel{(4)}{𝒥^{m|j)}(t)}+6\stackrel{(3)}{M^{ijkk}(t)}\}`$
$`+{\displaystyle \frac{1}{6}}ϵ^4_t^4{\displaystyle _{}}\tau ^{ij}(t,𝐱^{})|𝐱𝐱^{}|^3d^3x^{}`$
$`{\displaystyle \frac{1}{180}}ϵ^{9/2}\{3r^4\stackrel{(7)}{^{ij}(t)}4r^2x^k\stackrel{(7)}{^{ijk}(t)}16r^2x^kϵ^{mk(i}\stackrel{(6)}{𝒥^{m|j)}(t)}`$
$`+12r^2\stackrel{(5)}{M^{ijkk}(t)}+24x^{kl}\stackrel{(5)}{M^{ijkl}(t)}24x^k\stackrel{(5)}{M^{ijkll}(t)}+6\stackrel{(5)}{M^{ijkkll}(t)}\}`$
$`+B_{}^{ij}+O(ϵ^5).`$
Explicit formulae for the boundary terms $`N_{}`$, $`K_{}^i`$ and $`B_{}^{ij}`$ are given in Appendix C. Through 3.5PN order, the terms in Eq. (IV C) divide naturally into two types: even terms, i.e. terms of integer powers in $`ϵ`$ in $`N`$ and $`B^{ij}`$ and odd-half integer powers in $`K^i`$, and odd terms, of odd-half integer powers in $`N`$ and $`B^{ij}`$ and integer powers in $`K^i`$. The even terms produce the leading Newtonian, PN, 2PN and 3PN contributions to the equations of motion, while the odd terms produce the gravitational radiation reaction forces. (Note that the even terms have odd contributions embedded within them, via contributions of the metric itself to $`\tau ^{\alpha \beta }`$). Through 3.5PN order, there is a clean division between even and odd terms, in the sense that even terms produce non-dissipative contributions to the equations of motion, while odd terms produce radiation reaction effects. At 4PN order this separation fails, because of the presence of tails – these are $`O(ϵ^{3/2})`$ modifications of the leading 2.5PN radiation-reaction terms, which result in disspative effects at 4PN order. We derive the leading contributions of these 4PN tail terms in Sec. VI C.
The outer integrals for near-zone field points turn out to contribute only beginning at 3PN order (and, as we will see, do not contribute observable effects until 4PN order). This can be seen schematically as follows: for a source term of the form $`f_{B,L}(u)\widehat{n}^{<L>}r^B`$, the outer integral has the form
$$f(tr^{})(\widehat{n}^{})^{<L>}(r^{})^B\frac{d^3x^{}}{|𝐱𝐱^{}|}_{}^{\mathrm{}}\frac{d^qf(t)}{dt^q}r_{}^{}{}_{}{}^{q+1B}𝑑r^{}\frac{d^qf(t)}{dt^q}^{q+2B},$$
(131)
where we have used the fact that $`|𝐱||𝐱^{}|`$. The only possible $``$-independent terms come from the case $`q=B2`$. Thus the outer integral gives a schematic contribution $`h_{𝒞𝒩}^{\alpha \beta }f^{(B2)}(t)`$ where the superscript $`(B2)`$ denotes $`B2`$ time derivatives. From Eq. (IV Aa), the leading contribution to the source comes from $`(N)^2`$, where, from Eq. (17), $`N`$ has the far-zone form $`N4/r+2(3\widehat{n}^{<kl>}^{kl}/r^3+3\widehat{n}^{<kl>}\dot{}^{kl}/r^2+\widehat{n}^{kl}\ddot{}^{kl}/r)+\mathrm{}`$. Taking the gradient of this expression and squaring, we get, schematically $`(N)^2^2/r^4+(^{kl}/r^6+\dot{}^{kl}/r^5+\ddot{}^{kl}/r^4+\mathrm{})`$. The first term ($`B=4`$) gives no contribution, since $``$ is constant to the order considered ($``$ varies only via gravitational radiation energy loss). The second, third and fourth terms, ($`B=6,\mathrm{\hspace{0.17em}5},\mathrm{\hspace{0.17em}4}`$) together give $`h^{kk(4)}(t)`$. Since $`^{kk}mr^2`$, we find $`h(m/r)^2v^4O(ϵ^4)`$, which is a 3PN contribution. Thus, for near zone field points, the outer integrals can be ignored until 3PN order. A similar argument for far-zone field points reveals that outer integrals begin to contribute only at 2PN order, as was found by WW .
### D Compendium of useful post-Newtonian near-zone potentials
The even terms in Eq. (IV C) have the form of ordinary Poisson-like potentials and their generalizations, sometimes called superpotentials. For a source $`f`$, we define the Poisson potential, superpotential, and superduperpotential to be
$`P(f)`$ $``$ $`{\displaystyle \frac{1}{4\pi }}{\displaystyle _{}}{\displaystyle \frac{f(t,𝐱^{})}{|𝐱𝐱^{}|}}d^3x^{},^2P(f)=f,`$ (133)
$`S(f)`$ $``$ $`{\displaystyle \frac{1}{4\pi }}{\displaystyle _{}}f(t,𝐱^{})|𝐱𝐱^{}|d^3x^{},^2S(f)=2P(f),`$ (134)
$`SD(f)`$ $``$ $`{\displaystyle \frac{1}{4\pi }}{\displaystyle _{}}f(t,𝐱^{})|𝐱𝐱^{}|^3d^3x^{},^2SD(f)=12S(f).`$ (135)
We also define potentials based on the “densities” $`\sigma `$, $`\sigma ^i`$ and $`\sigma ^{ij}`$ constructed from $`T^{\alpha \beta }`$,
$`\mathrm{\Sigma }(f)`$ $``$ $`{\displaystyle _{}}{\displaystyle \frac{\sigma (t,𝐱^{})f(t,𝐱^{})}{|𝐱𝐱^{}|}}d^3x^{}=P(4\pi \sigma f),`$ (137)
$`\mathrm{\Sigma }^i(f)`$ $``$ $`{\displaystyle _{}}{\displaystyle \frac{\sigma ^i(t,𝐱^{})f(t,𝐱^{})}{|𝐱𝐱^{}|}}d^3x^{}=P(4\pi \sigma ^if),`$ (138)
$`\mathrm{\Sigma }^{ij}(f)`$ $``$ $`{\displaystyle _{}}{\displaystyle \frac{\sigma ^{ij}(t,𝐱^{})f(t,𝐱^{})}{|𝐱𝐱^{}|}}d^3x^{}=P(4\pi \sigma ^{ij}f),`$ (139)
along with the superpotentials
$`X(f)`$ $``$ $`{\displaystyle _{}}\sigma (t,𝐱^{})f(t,𝐱^{})|𝐱𝐱^{}|d^3x^{}=S(4\pi \sigma f),`$ (141)
$`Y(f)`$ $``$ $`{\displaystyle _{}}\sigma (t,𝐱^{})f(t,𝐱^{})|𝐱𝐱^{}|^3d^3x^{}=SD(4\pi \sigma f),`$ (142)
$`Z(f)`$ $``$ $`{\displaystyle _{}}\sigma (t,𝐱^{})f(t,𝐱^{})|𝐱𝐱^{}|^5d^3x^{},`$ (143)
and their obvious counterparts $`X^i`$, $`X^{ij}`$, $`Y^i`$, $`Y^{ij}`$, and so on. A number of potentials occur sufficiently frequently in the PN expansion that is it useful to define them specifically. First and foremost is the “Newtonian” potential,
$$U_{}\frac{\sigma (t,𝐱^{})}{|𝐱𝐱^{}|}d^3x^{}=P(4\pi \sigma )=\mathrm{\Sigma }(1).$$
(144)
The potentials needed for the post-Newtonian limit are:
$`V^i\mathrm{\Sigma }^i(1),`$ $`\mathrm{\Phi }_1^{ij}\mathrm{\Sigma }^{ij}(1),`$ $`\mathrm{\Phi }_1\mathrm{\Sigma }^{ii}(1),`$ (145)
$`\mathrm{\Phi }_2\mathrm{\Sigma }(U),`$ $`XX(1).`$ (146)
Useful 2PN potentials include:
$`V_2^i\mathrm{\Sigma }^i(U),`$ $`\mathrm{\Phi }_2^i\mathrm{\Sigma }(V^i),`$ (147)
$`X^iX^i(1),`$ $`X^{ij}X^{ij}(1),`$ (148)
$`X_1X^{ii},`$ $`X_2X(U),`$ (149)
$`P_2^{ij}P(U^{,i}U^{,j}),`$ $`P_2P_2^{ii}=\mathrm{\Phi }_2{\displaystyle \frac{1}{2}}U^2,`$ (150)
$`G_1P(\dot{U}^2),`$ $`G_2P(U\ddot{U}),`$ (151)
$`G_3P(\dot{U}^{,k}V^k),`$ $`G_4P(V^{i,j}V^{j,i}),`$ (152)
$`G_5P(\dot{V}^kU^{,k}),`$ $`G_6P(U^{,ij}\mathrm{\Phi }_1^{ij}),`$ (153)
$`G_7^iP(U^{,k}V^{k,i})+{\displaystyle \frac{3}{4}}P(U^{,i}\dot{U}),`$ $`HP(U^{,ij}P_2^{ij}).`$ (154)
At 3PN order, the following potentials are useful:
$`Y_1Y^{ii},`$ $`Y_2Y(U),ZZ(1).`$ (155)
A variety of properties of these and general Poisson potentials are described in Appendix D. Note that, in evaluating Poisson potentials and superpotentials of sources that do not have compact support, our rule is to evaluate them on the finite, constant time hypersurface $``$, and to discard all terms that depend on $``$.
## V Expansion of near-zone fields to 2.5PN order
We now turn to explicit evaluation of the near-zone fields and the metric to higher PN order, in terms of Poisson potentials and multipole moments. In addition to evaluating the inner integrals shown above, we must evaluate the outer integrals consistently at each PN order, to check whether any finite, $``$-independent contributions result.
In evaluating the contributions at each order, we shall use the following notation,
$`N`$ $`=`$ $`ϵ(N_0+ϵN_1+ϵ^{3/2}N_{1.5}+ϵ^2N_2+ϵ^{5/2}N_{2.5}+ϵ^3N_3+ϵ^{7/2}N_{3.5})+O(ϵ^5),`$ (157)
$`K^i`$ $`=`$ $`ϵ^{3/2}(K_1^i+ϵK_2^i+ϵ^{3/2}K_{2.5}^i+ϵ^2K_3^i+ϵ^{5/2}K_{3.5}^i)+O(ϵ^{9/2}),`$ (158)
$`B`$ $`=`$ $`ϵ^2(B_1+ϵ^{1/2}B_{1.5}+ϵB_2+ϵ^{3/2}B_{2.5}+ϵ^2B_3+ϵ^{5/2}B_{3.5})+O(ϵ^5),`$ (159)
$`B^{ij}`$ $`=`$ $`ϵ^2(B_2^{ij}+ϵ^{1/2}B_{2.5}^{ij}+ϵB_3^{ij}+ϵ^{3/2}B_{3.5}^{ij})+O(ϵ^4),`$ (160)
where the subscript on each term indicates the level (1PN, 2PN, 2.5PN, etc.) of its leading contribution to the equations of motion. Notice that our separate treatment of $`B`$ and $`B^{ij}`$ leads to the slightly awkward notational circumstance that, for example, $`B_2^{ii}=B_1`$.
### A The Newtonian and 1.5PN solution
At lowest order in the PN expansion, we only need to evaluate $`\tau ^{00}=(g)T^{00}+O(\rho ϵ)=\sigma +O(\rho ϵ)`$ (recall that $`\sigma ^{ii}ϵ\sigma `$). Since this has compact support, the outer integrals vanish, and we find
$$N_0=4U.$$
(161)
To this order, $`(g)=1+4U+O(ϵ^2)`$.
To the next PN order, we obtain, from Eqs. (8), (IV A) and (161),
$`\tau ^{00}`$ $`=`$ $`\sigma \sigma ^{ii}+4\sigma U{\displaystyle \frac{7}{8\pi }}U^2+O(\rho ϵ^2),`$ (162)
$`\tau ^{0i}`$ $`=`$ $`\sigma ^i+O(\rho ϵ^{3/2}),`$ (163)
$`\tau ^{ii}`$ $`=`$ $`\sigma ^{ii}{\displaystyle \frac{1}{8\pi }}U^2+O(\rho ϵ^2),`$ (164)
$`\tau ^{ij}`$ $`=`$ $`O(\rho ϵ).`$ (165)
Substituting into Eqs. (IV C), and calculating terms through 1.5PN order (e.g. $`O(ϵ^{5/2})`$ in $`N`$), we obtain
$`N_1`$ $`=`$ $`7U^24\mathrm{\Phi }_1+2\mathrm{\Phi }_2+2\ddot{X},`$ (167)
$`K_1^i`$ $`=`$ $`4V^i,`$ (168)
$`B_1`$ $`=`$ $`U^2+4\mathrm{\Phi }_12\mathrm{\Phi }_2,`$ (169)
$`N_{1.5}`$ $`=`$ $`{\displaystyle \frac{2}{3}}\stackrel{(3)}{^{kk}(t)},`$ (170)
$`B_{1.5}`$ $`=`$ $`2\stackrel{(3)}{^{kk}(t)}.`$ (171)
It is straightforward in this case to show that the outer integrals and surface terms give no $``$-independent terms. It is useful to illustrate the cancellation of an $``$-dependent term in this simple case. In the far zone to Newtonian order, the field, from Eq. (17), is given by $`N4/r`$, where we focus on the monopole contribution. This contributes to $`\mathrm{\Lambda }^{00}`$ in the far-zone a term of the form $`\mathrm{\Lambda }^{00}=14^2/r^4`$. To evaluate the near-zone contribution of the outer integral of this term, we must evaluate the coefficient $`_{B,L}^q`$ in Eq. (43) with $`q=0`$ (no time derivative, since $``$ is constant, to lowest order), $`B=4`$, $`L=0`$. From Eqs. (35) and (44), this yields a contribution to $`N`$ given by $`N_{𝒞𝒩}=7^2/^2`$. However, in evaluating $`N_𝒩`$, we encounter the Poisson potential $`14P(U^2)=14P_2`$ (see Eq. (154)). Upon integrating by parts and keeping the surface term at $``$ (see Eq. (Da)), this gives a contribution $`7U^214\mathrm{\Phi }_2+7^2/^2`$, whose $``$-dependent term cancels that from the outer integral.
The physical metric to 1.5PN order is then
$`g_{00}`$ $`=`$ $`1+2U2U^2+\ddot{X}{\displaystyle \frac{4}{3}}\stackrel{(3)}{^{kk}(t)}+O(ϵ^3),`$ (173)
$`g_{0i}`$ $`=`$ $`4V^i+O(ϵ^{5/2}),`$ (174)
$`g_{ij}`$ $`=`$ $`\delta _{ij}(1+2U)+O(ϵ^2).`$ (175)
Notice that, in our formulation, the potential $`U`$ is not a retarded potential; the retardation is expressed by the PN potential $`\ddot{X}`$ and the 1.5PN term $`\frac{4}{3}^{kk(3)}(t)`$. This contrasts with the PM approach, where retarded, rather than Poisson potentials are used, and the retardation is expanded only much later in the computation. The apparently 1.5PN term $`\frac{4}{3}^{kk(3)}(t)`$ in $`g_{00}`$ actually doesn’t contribute to the equations of motion at this order because it is purely a function of time, and the leading contribution is through a spatial gradient. As a consequence, the lowest-order observable contribution to radiation reaction is at 2.5PN order. (An alternative way to treat this 1.5PN term would be to absorb it in a redefinition of the time coordinate.)
We note here the useful identity, which follows from Eqs. (165), $`\sigma =\tau ^{00}+\tau ^{ii}(1/2\pi )^2(U^2)+O(\rho ϵ^2)`$, whose consequence is,
$$_{}\sigma (t,𝐱)d^3x=+\frac{1}{2}\ddot{}^{ii}+O(ϵ^2),$$
(176)
where surface terms make no $``$-independent contribution.
### B The 2.5PN solution
At 2PN and 2.5PN order, we obtain from Eqs. (8), (IV Ab,c), (161) and (V A)
$`\tau ^{ij}`$ $`=`$ $`\sigma ^{ij}+{\displaystyle \frac{1}{4\pi }}(U^{,i}U^{,j}{\displaystyle \frac{1}{2}}\delta ^{ij}U^2)+O(\rho ϵ^2),`$ (178)
$`\tau ^{0i}`$ $`=`$ $`\sigma ^i+4\sigma ^iU+{\displaystyle \frac{2}{\pi }}U^{,j}V^{[j,i]}+{\displaystyle \frac{3}{4\pi }}\dot{U}U^{,i}+O(\rho ϵ^{5/2}).`$ (179)
Including outer integrals and boundary terms (which contribute nothing), we obtain, from Eq. (IV Cc),
$`B_2^{ij}`$ $`=`$ $`4\mathrm{\Phi }_1^{ij}+4P_2^{ij}\delta ^{ij}(2\mathrm{\Phi }_2U^2),`$ (181)
$`K_2^i`$ $`=`$ $`8V_2^i8\mathrm{\Phi }_2^i+8UV^i+16G_7^i+2\ddot{X}^i,`$ (182)
$`B_{2.5}^{ij}`$ $`=`$ $`2\stackrel{(3)}{^{ij}(t)},`$ (183)
$`K_{2.5}^i`$ $`=`$ $`{\displaystyle \frac{2}{3}}x^k\stackrel{(4)}{^{ik}(t)}{\displaystyle \frac{2}{9}}\stackrel{(4)}{^{ikk}(t)}+{\displaystyle \frac{4}{9}}ϵ^{mik}\stackrel{(3)}{𝒥^{mk}(t)}.`$ (184)
The solutions for $`B_2^{ij}`$ and $`B_{2.5}^{ij}`$, along with the earlier 1.5PN solutions must now be substituted into $`(g)T^{\alpha \beta }`$ and Eq. (IV Aa,d), with the result
$`\tau ^{00}`$ $`=`$ $`\sigma \sigma ^{ii}+4\sigma U{\displaystyle \frac{7}{8\pi }}U^2`$ (190)
$`+\sigma (7U^28\mathrm{\Phi }_1+2\mathrm{\Phi }_2+2\ddot{X})4\sigma ^{ii}U`$
$`+{\displaystyle \frac{1}{4\pi }}\{{\displaystyle \frac{5}{2}}\dot{U}^24U\ddot{U}8\dot{U}^{,k}V^k+2V^{i,j}(3V^{j,i}+V^{i,j})+4\dot{V}^jU^{,j}4U^{,ij}\mathrm{\Phi }_1^{ij}`$
$`+8U\mathrm{\Phi }_14U\mathrm{\Phi }_2{\displaystyle \frac{7}{2}}U\ddot{X}10UU^24U^{,ij}P_2^{ij}\}`$
$`+{\displaystyle \frac{4}{3}}\sigma \stackrel{(3)}{^{kk}(t)}+{\displaystyle \frac{1}{2\pi }}U^{,ij}\stackrel{(3)}{^{ij}(t)},`$
$`\tau ^{ii}`$ $`=`$ $`\sigma ^{ii}{\displaystyle \frac{1}{8\pi }}U^2`$ (192)
$`+4\sigma ^{ii}U{\displaystyle \frac{1}{4\pi }}\left\{{\displaystyle \frac{9}{2}}\dot{U}^2+4V^{i,j}V^{[i,j]}+4\dot{V}^jU^{,j}+{\displaystyle \frac{1}{2}}U\ddot{X}\right\}.`$
Substituting into Eq. (IV Ca) and (IV Cc) and evaluating terms through $`O(ϵ^{7/2})`$, and verifying that the outer integrals and surface terms make no $``$-independent contributions, we obtain,
$`N_2`$ $`=`$ $`16U\mathrm{\Phi }_1+8U\mathrm{\Phi }_2+7U\ddot{X}+{\displaystyle \frac{20}{3}}U^34V^iV^i16\mathrm{\Sigma }(\mathrm{\Phi }_1)+\mathrm{\Sigma }(\ddot{X})+8\mathrm{\Sigma }^i(V^i)`$ (195)
$`2\ddot{X}_1+\ddot{X}_2+{\displaystyle \frac{1}{6}}\stackrel{(4)}{Y}4G_116G_2+32G_3+24G_416G_516G_616H,`$
$`B_2`$ $`=`$ $`U\ddot{X}+4V^iV^i\mathrm{\Sigma }(\ddot{X})8\mathrm{\Sigma }^i(V^i)+16\mathrm{\Sigma }^{ii}(U)`$ (197)
$`+2\ddot{X}_1\ddot{X}_220G_1+8G_4+16G_5,`$
$`N_{2.5}`$ $`=`$ $`{\displaystyle \frac{1}{15}}(2x^{kl}+r^2\delta ^{kl})\stackrel{(5)}{^{kl}(t)}+{\displaystyle \frac{2}{15}}x^k\stackrel{(5)}{^{kll}(t)}{\displaystyle \frac{1}{30}}\stackrel{(5)}{^{kkll}(t)}`$ (199)
$`+{\displaystyle \frac{16}{3}}U\stackrel{(3)}{^{kk}(t)}4X^{,kl}\stackrel{(3)}{^{kl}(t)},`$
$`B_{2.5}`$ $`=`$ $`{\displaystyle \frac{1}{3}}r^2\stackrel{(5)}{^{ii}(t)}+{\displaystyle \frac{2}{9}}x^k\stackrel{(5)}{^{iik}(t)}+{\displaystyle \frac{8}{9}}x^kϵ^{mki}\stackrel{(4)}{𝒥^{mi}(t)}{\displaystyle \frac{2}{3}}\stackrel{(3)}{M^{iikk}(t)}.`$ (200)
### C The far-zone field to 1.5PN order
In anticipation of finding non-zero outer-integral contributions to the near-zone field at 3PN order, we must determine the far-zone field to an order needed for the source $`\mathrm{\Lambda }^{\alpha \beta }`$. Our foregoing discussion indicates that counting PN orders for outer integrals is different than the standard method, because the inverse radial variable $`r^1<^1v/𝒮`$; in other words, when considering contributions to the outer integrals, additional powers of $`r`$ in a term in the far-zone field can be regarded as increasing the effective order of that term by half a power in $`ϵ`$. For example, expanding $`h_𝒩^{00}=N_𝒩`$ in the far-zone, Eq. (17), we obtain
$`N_𝒩`$ $`=`$ $`4\left\{{\displaystyle \frac{}{r}}+{\displaystyle \frac{1}{2}}_{kl}\left({\displaystyle \frac{^{kl}(u)}{r}}\right){\displaystyle \frac{1}{6}}_{klm}\left({\displaystyle \frac{^{klm}(u)}{r}}\right)+\mathrm{}\right\},`$ (202)
$`ϵϵ^2ϵ^{5/2}`$
where the effective PN order of each term is indicated. In ordinary applications, the second potential in Eq. (202) would contribute a term of order $`ϵ`$ of the form $`\widehat{n}^{<kl>}^{kl}/r^3`$, which is simply the Newtonian quadrupole potential. But in the outer integral, this term contributes an $``$-independent term only through several time derivatives, and thus its effective contribution is higher order, in fact of the same order as that of the term $`\widehat{n}^{kl}\ddot{}^{kl}/r`$, which also comes from the second potential.
At this order, we must also be careful to include any outer integral and boundary contributions to the far-zone field. From the lowest-order far-zone field, we find, to the order needed, that $`\mathrm{\Lambda }^{00}=14(/r^2)^2`$, $`\mathrm{\Lambda }^{ij}=4(/r^2)^2(n^{<ij>}\delta ^{ij}/6)`$. Evaluating the coefficients $`𝒟_{4,2}^0`$ and $`𝒟_{4,0}^0`$, Eq. (35), we obtain, in the far zone, $`N_{𝒞𝒩}=7(/r)^2`$ and $`B_{𝒞𝒩}^{ij}=(/r)^2\widehat{n}^{ij}`$. Combining the multipole expansions of Eq. (17) with the outer integral contributions, we obtain in the far-zone, to the order needed,
$`N`$ $`=`$ $`4{\displaystyle \frac{}{r}}+2_{kl}\left({\displaystyle \frac{^{kl}(u)}{r}}\right){\displaystyle \frac{2}{3}}_{klm}\left({\displaystyle \frac{^{klm}(u)}{r}}\right)+7{\displaystyle \frac{^2}{r^2}}+O(ϵ^3),`$ (204)
$`K^i`$ $`=`$ $`2_k\left({\displaystyle \frac{\dot{}^{ik}(u)}{r}}\right)+2ϵ^{aib}{\displaystyle \frac{\widehat{n}^a𝒥^b}{r^2}}+{\displaystyle \frac{2}{3}}_{kl}\left({\displaystyle \frac{\dot{}^{ikl}(u)}{r}}\right)+{\displaystyle \frac{4}{3}}ϵ^{aib}_{ak}\left({\displaystyle \frac{𝒥^{bk}(u)}{r}}\right)+O(ϵ^3),`$ (205)
$`B^{ij}`$ $`=`$ $`2{\displaystyle \frac{\ddot{}^{ij}(u)}{r}}+{\displaystyle \frac{2}{3}}_k\left({\displaystyle \frac{\ddot{}^{ijk}(u)}{r}}\right)+{\displaystyle \frac{8}{3}}ϵ^{ak(i}_k\left({\displaystyle \frac{\dot{𝒥}^{a|j)}(u)}{r}}\right)+{\displaystyle \frac{^2}{r^2}}\widehat{n}^{ij}+O(ϵ^3).`$ (206)
It will turn out, however, that, despite the formal possibility of 3PN contributions from the outer integrals, the actual contributions will not begin until 4PN order (see Sec. VI C)
## VI Expansion of Near-Zone Fields to 3.5PN order
### A $`B^{ij}`$ and $`K^j`$ to 3PN and 3.5PN order
At 3PN and 3.5PN order, we obtain from Eqs. (8), (IV Ab,c), (161) and (V A)
$`\tau ^{ij}`$ $`=`$ $`\sigma ^{ij}+{\displaystyle \frac{1}{4\pi }}(U^{,i}U^{,j}{\displaystyle \frac{1}{2}}\delta ^{ij}U^2)`$ (210)
$`+4\sigma ^{ij}U+{\displaystyle \frac{1}{4\pi }}\{U^{,(i}\ddot{X}^{,j)}16V^{[i,k]}V^{[j,k]}+8U^{,(i}\dot{V}^{j)}`$
$`\delta ^{ij}({\displaystyle \frac{1}{2}}U\ddot{X}4V^{[l,k]}V^{[l,k]}+4U^{,(k}\dot{V}^{k)}+{\displaystyle \frac{3}{2}}\dot{U}^2)\}+O(\rho ϵ^3),`$
$`\tau ^{0i}`$ $`=`$ $`\sigma ^i+4\sigma ^iU+{\displaystyle \frac{2}{\pi }}U^{,j}V^{[j,i]}+{\displaystyle \frac{3}{4\pi }}\dot{U}U^{,i}`$ (218)
$`+\sigma ^i(7U^28\mathrm{\Phi }_1+2\mathrm{\Phi }_2+2\ddot{X})`$
$`+{\displaystyle \frac{1}{16\pi }}\{64U^{,k}(V_2^{[k,i]}\mathrm{\Phi }_2^{[k,i]})+32UU^{,k}V^{k,i}16UU^{,k}V^{i,k}+16U^{,i}U^{,k}V^k`$
$`24V^i(U)^2+16U^{,k}\ddot{X}^{[k,i]}+128U^{,k}G_7^{[k,i]}32\mathrm{\Phi }_1^{,k}V^{[k,i]}16\mathrm{\Phi }_2^{,k}V^{i,k}`$
$`16\ddot{X}^{,k}V^{[k,i]}16\dot{U}\mathrm{\Phi }_1^{,i}+48U\dot{U}U^{,i}+6\dot{U}\ddot{X}^{,i}+6U^{,i}\stackrel{(3)}{X}16U^{,i}\dot{\mathrm{\Phi }}_1+16\dot{U}\dot{V}^i`$
$`16U\ddot{V}^i32V^k\dot{V}^{i,k}16V^{i,kl}(\mathrm{\Phi }_1^{kl}+P_2^{kl})+16U^{,k}(\dot{\mathrm{\Phi }}_1^{ik}+\dot{P}_2^{ik})`$
$`+16V^{k,l}(\mathrm{\Phi }_1^{il,k}+\mathrm{\Phi }_1^{ik,l}\mathrm{\Phi }_1^{kl,i})+16V^{k,l}(P_2^{il,k}+P_2^{ik,l}P_2^{kl,i})\}`$
$`+{\displaystyle \frac{4}{3}}\sigma ^i\stackrel{(3)}{^{kk}(t)}+{\displaystyle \frac{1}{2\pi }}(V^{i,kl}\stackrel{(3)}{^{kl}(t)}U^{,k}\stackrel{(4)}{^{ik}(t)}),`$
where the first line in each expression is the contribution through 2PN order obtained earlier. Substituting into Eq. (IV Cb,c) and keeping contributions through $`O(ϵ^{7/2})`$, and checking that surface terms and outer integrals make no contribution to this order, we obtain,
$`B_3^{ij}`$ $`=`$ $`16\mathrm{\Sigma }^{ij}(U)+4P(U^{,(i}\ddot{X}^{,j)})64P(V^{[i,k]}V^{[j,k]})+32P(U^{,(i}\dot{V}^{j)})+2\ddot{X}^{ij}+2\ddot{S}(U^{,i}U^{,j})`$ (221)
$`+\delta ^{ij}(U\ddot{X}4V^kV^k\mathrm{\Sigma }(\ddot{X})+8\mathrm{\Sigma }^k(V^k)\ddot{X}_28G_18G_4+16G_5),`$
$`K_3^i`$ $`=`$ $`12U^2V^i+16UV_2^i16U\mathrm{\Phi }_2^i+4U\ddot{X}^i+32UG_7^i+4V^i\ddot{X}8\mathrm{\Phi }_1V^i+8\mathrm{\Phi }_2V^i`$ (230)
$`8V^k\mathrm{\Phi }_1^{ik}8V^kP_2^{ik}16\mathrm{\Sigma }(V_2^i)+16\mathrm{\Sigma }(\mathrm{\Phi }_2^i)16\mathrm{\Sigma }(UV^i)4\mathrm{\Sigma }(\ddot{X}^i)32\mathrm{\Sigma }(G_7^i)`$
$`24\mathrm{\Sigma }^i(\mathrm{\Phi }_1)+4\mathrm{\Sigma }^i(\ddot{X})+8\mathrm{\Sigma }^{kk}(V^i)+8\mathrm{\Sigma }^k(\mathrm{\Phi }_1^{ik})+8\mathrm{\Sigma }^k(P_2^{ik})+8\mathrm{\Sigma }^{ik}(V^k)`$
$`+24P(U\dot{U}U^{,i})+24P(U^{,k}U^{,i}V^k)+32P(U^{,k}V_2^{k,i})32P(U^{,k}\mathrm{\Phi }_2^{k,i})+64P(U^{,k}G_7^{k,i})`$
$`+8P(U^{,k}\ddot{X}^{k,i})+16P(U^{,k}\dot{\mathrm{\Phi }}_1^{ik})+16P(U^{,k}\dot{P}_2^{ik})16P(U^{,i}\dot{\mathrm{\Phi }}_1)+6P(U^{,i}\stackrel{(3)}{X})`$
$`16P(\dot{U}\mathrm{\Phi }_1^{,i})+6P(\dot{U}\ddot{X}^{,i})+32P(UU^{,k}V^{i,k})16P(U\ddot{V}^i)+16P(V^{k,l}\mathrm{\Phi }_1^{il,k})`$
$`16P(V^{k,l}\mathrm{\Phi }_1^{kl,i})+16P(V^{k,l}P_2^{il,k})32P(V^k\dot{V}^{i,k})16P(V^{k,i}\mathrm{\Phi }_1^{,k})+8P(V^{k,i}\ddot{X}^{,k})`$
$`16P(V^{k,l}P_2^{kl,i})16P(V^{i,lm}\mathrm{\Phi }_1^{lm})16P(V^{i,lm}P_2^{lm})`$
$`+4\ddot{X}^i(U)4\ddot{X}(V^i)+8\ddot{S}(U^{,k}V^{k,i})+6\ddot{S}(U^{,i}\dot{U})+{\displaystyle \frac{1}{6}}\stackrel{(4)}{Y^i},`$
$`B_{3.5}^{ij}`$ $`=`$ $`{\displaystyle \frac{1}{3}}r^2\stackrel{(5)}{^{ij}(t)}+{\displaystyle \frac{2}{9}}x^k\stackrel{(5)}{^{ijk}(t)}+{\displaystyle \frac{8}{9}}x^kϵ^{mk(i}\stackrel{(4)}{𝒥^{m|j)}(t)}{\displaystyle \frac{2}{3}}\stackrel{(3)}{M^{ijkk}(t)},`$ (231)
$`K_{3.5}^i`$ $`=`$ $`{\displaystyle \frac{1}{15}}r^2x^k\stackrel{(6)}{^{ik}(t)}{\displaystyle \frac{1}{45}}r^2\stackrel{(6)}{^{ikk}(t)}{\displaystyle \frac{2}{45}}x^{kl}\stackrel{(6)}{^{ikl}(t)}+{\displaystyle \frac{1}{30}}x^k\stackrel{(6)}{^{ikll}(t)}{\displaystyle \frac{1}{150}}\stackrel{(6)}{^{ikkll}(t)}`$ (234)
$`+ϵ^{mil}[{\displaystyle \frac{2}{45}}r^2\stackrel{(5)}{𝒥^{ml}(t)}+{\displaystyle \frac{4}{45}}x^{kl}\stackrel{(5)}{𝒥^{mk}(t)}{\displaystyle \frac{1}{30}}x^l\stackrel{(5)}{𝒥^{mkk}(t)}{\displaystyle \frac{1}{15}}x^k\stackrel{(5)}{𝒥^{mkl}(t)}+{\displaystyle \frac{2}{75}}\stackrel{(5)}{𝒥^{mlkk}(t)}]`$
$`+{\displaystyle \frac{16}{3}}V^i\stackrel{(3)}{^{kk}(t)}4X^{i,kl}\stackrel{(3)}{^{kl}(t)}+4X^{,k}\stackrel{(4)}{^{ik}(t)}.`$
### B $`N`$ and $`B`$ to 3PN and 3.5PN order
The expressions for $`\tau ^{00}`$ and $`\tau ^{ii}`$ to 3PN and 3.5PN order are too lengthy to be reproduced explicitly. Instead, by substituting the expansions (V) into Eqs. (8) and (IV Ab,c), and keeping terms of $`O(\rho ϵ^3)`$ and $`O(\rho ϵ^{3.5})`$, we obtain the formal contributions
$`\tau _3^{00}`$ $`=`$ $`\sigma (N_2B_2N_0B_1+K_1^iK_1^i)\sigma ^{ii}(N_1B_1)`$ (241)
$`+{\displaystyle \frac{1}{16\pi }}\{{\displaystyle \frac{7}{8}}(2N_0N_2+(N_1)^2)+{\displaystyle \frac{5}{4}}\dot{N}_0\dot{N}_1\ddot{N}_0N_1N_0\ddot{N}_1`$
$`2\dot{N}_0^{,i}K_2^i2\dot{N}_1^{,i}K_1^i+K_1^{i,j}(3K_2^{j,i}+K_2^{i,j})+N_0^{,i}\dot{K}_2^i+N_1^{,i}\dot{K}_1^iN_0^{,ij}B_3^{ij}`$
$`N_1^{,ij}B_2^{ij}+{\displaystyle \frac{1}{4}}(N_0B_2+N_1B_1)+{\displaystyle \frac{7}{8}}N_1(N_0)^2`$
$`+{\displaystyle \frac{7}{4}}N_0(N_0N_1)+K_1^{i,j}\dot{B}_2^{ij}+{\displaystyle \frac{1}{4}}B_2^{ij,k}(B_2^{ij,k}2B_2^{jk,i})+{\displaystyle \frac{1}{4}}\dot{N}_0\dot{B}_1{\displaystyle \frac{1}{8}}(B_1)^2`$
$`+{\displaystyle \frac{1}{4}}\dot{N}_0N_0^{,i}K_1^i+{\displaystyle \frac{7}{8}}N_0^{,i}N_0^{,j}B_2^{ij}{\displaystyle \frac{1}{2}}N_0^{,i}K_1^j(4K_1^{i,j}+3K_1^{j,i}){\displaystyle \frac{7}{8}}N_0^2(N_0)^2\},`$
$`\tau _{3.5}^{00}`$ $`=`$ $`\sigma (N_{2.5}B_{2.5}N_0B_{1.5})\sigma ^{ii}(N_{1.5}B_{1.5})`$ (245)
$`+{\displaystyle \frac{1}{16\pi }}\{{\displaystyle \frac{7}{4}}N_0N_{2.5}+{\displaystyle \frac{5}{4}}\dot{N}_0\dot{N}_{1.5}\ddot{N}_0N_{1.5}N_0\ddot{N}_{1.5}2\dot{N}_0^{,i}K_{2.5}^i+N_0^{,i}\dot{K}_{2.5}^i`$
$`+K_1^{i,j}(3K_{2.5}^{j,i}+K_{2.5}^{i,j})N_0^{,ij}B_{3.5}^{ij}N_1^{,ij}B_{2.5}^{ij}+{\displaystyle \frac{1}{4}}N_0B_{2.5}`$
$`+{\displaystyle \frac{7}{8}}N_{1.5}(N_0)^2+K_1^{i,j}\dot{B}_{2.5}^{ij}+{\displaystyle \frac{1}{4}}\dot{N}_0\dot{B}_{1.5}+{\displaystyle \frac{7}{8}}N_0^{,i}N_0^{,j}B_{2.5}^{ij}\},`$
$`\tau _3^{ii}`$ $`=`$ $`\sigma ^{ii}(N_1B_1)`$ (251)
$`+{\displaystyle \frac{1}{16\pi }}\{{\displaystyle \frac{1}{8}}(2N_0N_2+(N_1)^2)+2K_1^{i,j}K_2^{[j,i]}N_0^{,i}\dot{K}_2^iN_1^{,i}\dot{K}_1^i`$
$`{\displaystyle \frac{1}{4}}(N_0B_2+N_1B_1){\displaystyle \frac{9}{4}}\dot{N}_0\dot{N}_1+{\displaystyle \frac{1}{4}}N_1(N_0)^2+{\displaystyle \frac{1}{2}}N_0(N_0N_1)`$
$`+2\dot{K}_1^i\dot{K}_1^i2\dot{B}_1^{,i}K_1^i+3\dot{B}_2^{ij}K_1^{i,j}N_0\ddot{B}_1+{\displaystyle \frac{3}{4}}\dot{N}_0\dot{B}_1+{\displaystyle \frac{1}{8}}(B_1)^2B_1^{,ij}B_2^{ij}`$
$`+{\displaystyle \frac{3}{4}}B_2^{ij,k}B_2^{ij,k}+{\displaystyle \frac{1}{2}}B_2^{ij,k}B_2^{ik,j}N_0K_1^{i,j}K_1^{[j,i]}{\displaystyle \frac{1}{2}}N_0^{,i}K_1^jK_1^{j,i}+N_0N_0^{,i}\dot{K}_1^i+{\displaystyle \frac{1}{4}}\dot{N}_0N_0^{,i}K_1^i`$
$`{\displaystyle \frac{1}{8}}N_0^{,i}N_0^{,j}B_2^{ij}+{\displaystyle \frac{1}{4}}N_0N_0B_1+{\displaystyle \frac{1}{8}}(N_0)^2B_1+{\displaystyle \frac{9}{8}}N_0\dot{N}_0^2{\displaystyle \frac{3}{8}}N_0^2(N_0)^2\},`$
$`\tau _{3.5}^{ii}`$ $`=`$ $`\sigma ^{ii}(N_{1.5}B_{1.5})`$ (255)
$`+{\displaystyle \frac{1}{16\pi }}\{{\displaystyle \frac{1}{4}}N_0N_{2.5}+2K_1^{i,j}K_{2.5}^{[j,i]}N_0^{,i}\dot{K}_{2.5}^i{\displaystyle \frac{1}{4}}N_0B_{2.5}`$
$`{\displaystyle \frac{9}{4}}\dot{N}_0\dot{N}_{1.5}+{\displaystyle \frac{1}{4}}N_{1.5}(N_0)^2+3\dot{B}_{2.5}^{ij}K_1^{i,j}N_0\ddot{B}_{1.5}`$
$`+{\displaystyle \frac{3}{4}}\dot{N}_0\dot{B}_{1.5}B_1^{,ij}B_{2.5}^{ij}{\displaystyle \frac{1}{8}}N_0^{,i}N_0^{,j}B_{2.5}^{ij}+{\displaystyle \frac{1}{8}}(N_0)^2B_{1.5}\}.`$
We have simplified the expressions slightly by taking into account the fact that $`N_{1.5}`$, $`B_{1.5}`$ and $`B_{2.5}^{ij}`$ are purely functions of time, so that spatial gradients of them vanish. To obtain the full expressions, one substitutes for $`N_0`$, $`N_1`$, $`B_1`$, $`K_1^i`$, etc. from Eqs. (161), (V A), (V B), (V B), and (VI A). Substituting this into Eq. (IV Ca) and (IV Cc) (the latter contracted on indices $`ij`$), and including surface terms and outer integrals, we obtain the final 3PN and 3.5PN results for $`N`$ and $`B`$:
$`N_3`$ $`=`$ $`{\displaystyle \frac{19}{6}}U^428U^2\mathrm{\Phi }_1+14U^2\mathrm{\Phi }_2+10U^2\ddot{X}8U\ddot{X}_1+4U\ddot{X}_2+4U\mathrm{\Sigma }(\ddot{X})`$ (277)
$`4UG_156UG_2+112UG_3+80UG_464UG_556UG_656UH+{\displaystyle \frac{7}{12}}U\stackrel{(4)}{Y}`$
$`+32U\mathrm{\Sigma }^i(V^i)56U\mathrm{\Sigma }(\mathrm{\Phi }_1)8U\mathrm{\Sigma }^{ii}(U)+10\mathrm{\Phi }_1^28\mathrm{\Phi }_1\mathrm{\Phi }_2+2\mathrm{\Phi }_2^28\mathrm{\Phi }_1\ddot{X}+4\mathrm{\Phi }_2\ddot{X}`$
$`16V^iV_2^i+16V^i\mathrm{\Phi }_2^i32V^iG_7^i4V^i\ddot{X}^i+{\displaystyle \frac{7}{4}}\ddot{X}^22\mathrm{\Phi }_1^{ij}\mathrm{\Phi }_1^{ij}4\mathrm{\Phi }_1^{ij}P_2^{ij}2P_2^{ij}P_2^{ij}`$
$`8\mathrm{\Sigma }(U\mathrm{\Phi }_1)+36\mathrm{\Sigma }(G_1)8\mathrm{\Sigma }(G_2)+16\mathrm{\Sigma }(G_3)48\mathrm{\Sigma }(G_4)8\mathrm{\Sigma }(G_6)8\mathrm{\Sigma }(H)`$
$`8\mathrm{\Sigma }(\ddot{X}_1)+{\displaystyle \frac{1}{12}}\mathrm{\Sigma }(\stackrel{(4)}{Y})+16\mathrm{\Sigma }^i(V_2^i)16\mathrm{\Sigma }^i(\mathrm{\Phi }_2^i)16\mathrm{\Sigma }^i(UV^i)+32\mathrm{\Sigma }^i(G_7^i)+4\mathrm{\Sigma }^i(\ddot{X}^i)`$
$`+8\mathrm{\Sigma }^{ii}(U^2)+12\mathrm{\Sigma }^{ii}(\mathrm{\Phi }_1)+4\mathrm{\Sigma }^{ij}(\mathrm{\Phi }_1^{ij})+4\mathrm{\Sigma }^{ij}(P_2^{ij})8\mathrm{\Sigma }(\mathrm{\Sigma }(\mathrm{\Phi }_1))+64\mathrm{\Sigma }(\mathrm{\Sigma }^i(V^i))`$
$`56\mathrm{\Sigma }(\mathrm{\Sigma }^{ii}(U))32P(U^2\ddot{U})28P(U\dot{U}^2)+16P(\ddot{U}\mathrm{\Phi }_1)+16P(\dot{U}\dot{\mathrm{\Phi }}_1)`$
$`+16P(U\ddot{\mathrm{\Phi }}_1)8P(U\ddot{\mathrm{\Phi }}_2)8P(U\stackrel{(4)}{X})4P(\dot{U}\stackrel{(3)}{X})8P(\ddot{U}\ddot{X})+32P(UU^{,i}\dot{V}^i)`$
$`64P(U\dot{U}^{,i}V^i)40P(\dot{U}U^{,i}V^i)+32P(UV^{i,j}V^{i,j})+16P(UV^{i,j}V^{j,i})`$
$`+32P(U^{,i}\dot{V}_2^i)32P(U^{,i}\dot{\mathrm{\Phi }}_2^i)+64P(U^{,i}\dot{G}_7^i)+8P(U^{,i}\stackrel{(3)}{X^i})`$
$`+64P(\dot{U}^{,i}\mathrm{\Phi }_2^i)64P(\dot{U}^{,i}V_2^i)128P(\dot{U}^{,i}G_7^i)16P(\dot{U}^{,i}\ddot{X}^i)`$
$`+4P(U^{,i}U^{,j}\mathrm{\Phi }_1^{ij})+4P(U^{,i}U^{,j}P_2^{ij})64P(U^{,ij}\mathrm{\Sigma }^{ij}(U))8P(U^{,ij}\ddot{S}(U^{,i}U^{,j}))`$
$`128P(U^{,ij}P(U^{,(i}\dot{V}^{j)}))16P(U^{,ij}P(U^{,(i}\ddot{X}^{,j)}))+256P(U^{,ij}P(V^{[k,i]}V^{[k,j]}))`$
$`8P(U^{,ij}\ddot{X}^{ij})+32P(V^i\dot{\mathrm{\Phi }}_1^{,i})16P(V^i\dot{\mathrm{\Phi }}_2^{,i})16P(V^i\stackrel{(3)}{X^{,i}})+16P(V^{i,j}\dot{\mathrm{\Phi }}_1^{ij})`$
$`+16P(V^{i,j}\dot{P}_2^{ij})+96P(V^{i,j}V_2^{j,i})96P(V^{i,j}\mathrm{\Phi }_2^{j,i})+192P(V^{i,j}G_7^{j,i})`$
$`+24P(V^{i,j}\ddot{X}^{j,i})+8P(\dot{V}^i\dot{V}^i)16P(\dot{V}^i\mathrm{\Phi }_1^{,i})+8P(\dot{V}^i\ddot{X}^{,i})8P(\mathrm{\Phi }_1^{ij,k}\mathrm{\Phi }_1^{jk,i})`$
$`16P(\mathrm{\Phi }_1^{ij,k}P_2^{jk,i})8P(P_2^{ij,k}P_2^{jk,i})+16P(\mathrm{\Phi }_1^{,ij}\mathrm{\Phi }_1^{ij})+16P(\mathrm{\Phi }_1^{,ij}P_2^{ij})8P(\mathrm{\Phi }_1^{ij}\ddot{X}^{,ij})`$
$`8P(\mathrm{\Phi }_2^{,ij}\mathrm{\Phi }_1^{ij})8P(\mathrm{\Phi }_2^{,ij}P_2^{ij})8P(P_2^{ij}\ddot{X}^{,ij}){\displaystyle \frac{1}{6}}\stackrel{(4)}{Y_1}+{\displaystyle \frac{1}{12}}\stackrel{(4)}{Y_2}+{\displaystyle \frac{1}{180}}\stackrel{(6)}{Z}`$
$`8\ddot{X}(\mathrm{\Phi }_1)+{\displaystyle \frac{1}{2}}\ddot{X}(\ddot{X})+4\ddot{X}^i(V^i)2\ddot{S}(\dot{U}^2)8\ddot{S}(U\ddot{U})`$
$`8\ddot{S}(U^{,ij}\mathrm{\Phi }_1^{ij})8\ddot{S}(U^{,ij}P_2^{ij})+12\ddot{S}(V^{i,j}V^{j,i})+8\ddot{S}(U^{,i}\dot{V}^i)16\ddot{S}(V^i\dot{U}^{,i}),`$
$`B_3`$ $`=`$ $`{\displaystyle \frac{1}{2}}U^44U^2\mathrm{\Phi }_12U^2\mathrm{\Phi }_2+16UV^iV^i12UG_18UG_2+16UG_3+16UG_4`$ (291)
$`8UG_68UH+{\displaystyle \frac{1}{12}}U\stackrel{(4)}{Y}8U\mathrm{\Sigma }(\mathrm{\Phi }_1)+8U\mathrm{\Sigma }^{ii}(U)2\mathrm{\Phi }_1^2+8\mathrm{\Phi }_1\mathrm{\Phi }_2+2\mathrm{\Phi }_2^2`$
$`+16V^iV_2^i16V^i\mathrm{\Phi }_2^i+32V^iG_7^i+4V^i\ddot{X}^i+{\displaystyle \frac{1}{4}}\ddot{X}^26\mathrm{\Phi }_1^{ij}\mathrm{\Phi }_1^{ij}12\mathrm{\Phi }_1^{ij}P_2^{ij}6P_2^{ij}P_2^{ij}`$
$`+8\mathrm{\Sigma }(U\mathrm{\Phi }_1)+12\mathrm{\Sigma }(G_1)+8\mathrm{\Sigma }(G_2)16\mathrm{\Sigma }(G_3)16\mathrm{\Sigma }(G_4)+8\mathrm{\Sigma }(G_6)+8\mathrm{\Sigma }(H)`$
$`{\displaystyle \frac{1}{12}}\mathrm{\Sigma }(\stackrel{(4)}{Y})16\mathrm{\Sigma }^i(V_2^i)+16\mathrm{\Sigma }^i(\mathrm{\Phi }_2^i)48\mathrm{\Sigma }^i(UV^i)32\mathrm{\Sigma }^i(G_7^i)4\mathrm{\Sigma }^i(\ddot{X}^i)+24\mathrm{\Sigma }^{ii}(U^2)`$
$`28\mathrm{\Sigma }^{ii}(\mathrm{\Phi }_1)+8\mathrm{\Sigma }^{ii}(\ddot{X})+12\mathrm{\Sigma }^{ij}(\mathrm{\Phi }_1^{ij})+12\mathrm{\Sigma }^{ij}(P_2^{ij})+8\mathrm{\Sigma }(\mathrm{\Sigma }(\mathrm{\Phi }_1))8\mathrm{\Sigma }(\mathrm{\Sigma }^{ii}(U))`$
$`68P(U\dot{U}^2)16P(U\ddot{\mathrm{\Phi }}_1)+8P(U\ddot{\mathrm{\Phi }}_2)+48P(\dot{U}\dot{\mathrm{\Phi }}_1)20P(\dot{U}\stackrel{(3)}{X})56P(\dot{U}U^{,i}V^i)`$
$`32P(UU^{,i}\dot{V}^i)+32P(UV^{i,j}V^{i,j})16P(UV^{i,j}V^{j,i})32P(U^{,i}\dot{V}_2^i)+32P(U^{,i}\dot{\mathrm{\Phi }}_2^i)`$
$`64P(U^{,i}\dot{G}_7^i)8P(U^{,i}\stackrel{(3)}{X^i})4P(U^{,i}U^{,j}\mathrm{\Phi }_1^{ij})4P(U^{,i}U^{,j}P_2^{ij})32P(V^i\dot{\mathrm{\Phi }}_1^{,i})`$
$`+16P(V^i\dot{\mathrm{\Phi }}_2^{,i})+48P(V^{i,j}\dot{\mathrm{\Phi }}_1^{ij})+48P(V^{i,j}\dot{P}_2^{ij})+32P(V^{i,j}V_2^{j,i})32P(V^{i,j}\mathrm{\Phi }_2^{j,i})`$
$`+64P(V^{i,j}G_7^{j,i})+8P(V^{i,j}\ddot{X}^{j,i})+24P(\dot{V}^i\dot{V}^i)+16P(\dot{V}^i\mathrm{\Phi }_1^{,i})8P(\dot{V}^i\ddot{X}^{,i})`$
$`+8P(\mathrm{\Phi }_1^{ij,k}\mathrm{\Phi }_1^{jk,i})+16P(\mathrm{\Phi }_1^{ij,k}P_2^{jk,i})+8P(P_2^{ij,k}P_2^{jk,i})16P(\mathrm{\Phi }_1^{,ij}\mathrm{\Phi }_1^{ij})16P(\mathrm{\Phi }_1^{,ij}P_2^{ij})`$
$`+8P(\mathrm{\Phi }_2^{,ij}\mathrm{\Phi }_1^{ij})+8P(\mathrm{\Phi }_2^{,ij}P_2^{ij})+{\displaystyle \frac{1}{6}}\stackrel{(4)}{Y_1}{\displaystyle \frac{1}{12}}\stackrel{(4)}{Y_2}{\displaystyle \frac{1}{2}}\ddot{X}(\ddot{X})4\ddot{X}^i(V^i)+8\ddot{X}^{ii}(U)`$
$`10\ddot{S}(\dot{U}^2)+4\ddot{S}(V^{i,j}V^{j,i})8\ddot{S}(U^{,i}\dot{V}^i)+{\displaystyle \frac{4}{3}}\stackrel{(4)}{^{jj}(t)},`$
$`N_{3.5}`$ $`=`$ $`{\displaystyle \frac{1}{420}}r^4\stackrel{(7)}{^{jj}(t)}{\displaystyle \frac{1}{105}}r^2x^{ij}\stackrel{(7)}{^{ij}(t)}+{\displaystyle \frac{1}{105}}r^2x^i\stackrel{(7)}{^{ijj}(t)}+{\displaystyle \frac{2}{315}}x^{ijk}\stackrel{(7)}{^{ijk}(t)}`$ (300)
$`{\displaystyle \frac{1}{420}}r^2\stackrel{(7)}{^{iijj}(t)}{\displaystyle \frac{1}{105}}x^{ij}\stackrel{(7)}{^{ijkk}(t)}+{\displaystyle \frac{1}{210}}x^i\stackrel{(7)}{^{ijjkk}(t)}{\displaystyle \frac{1}{1260}}\stackrel{(7)}{^{iijjkk}(t)}`$
$`\left({\displaystyle \frac{8}{15}}x^{ij}U+{\displaystyle \frac{6}{5}}x^iX^{,j}+{\displaystyle \frac{2}{3}}r^2X^{,ij}{\displaystyle \frac{2}{9}}x^kY^{,ijk}+{\displaystyle \frac{11}{45}}Y^{,ij}\right)\stackrel{(5)}{^{ij}(t)}`$
$`+\left({\displaystyle \frac{16}{15}}r^2U{\displaystyle \frac{34}{15}}x^iX^{,i}+{\displaystyle \frac{16}{5}}X\right)\stackrel{(5)}{^{jj}(t)}{\displaystyle \frac{1}{45}}(16x^iU52X^{,i})\stackrel{(5)}{^{ijj}(t)}`$
$`+\left({\displaystyle \frac{4}{9}}x^kX^{,ij}{\displaystyle \frac{2}{27}}Y^{,ijk}\right)\stackrel{(5)}{^{ijk}(t)}{\displaystyle \frac{2}{15}}U\stackrel{(5)}{^{iijj}(t)}`$
$`\left({\displaystyle \frac{4}{3}}X^{i,j}{\displaystyle \frac{8}{3}}x^i\dot{X}^{,j}+{\displaystyle \frac{10}{9}}\dot{Y}^{,ij}\right)\stackrel{(4)}{^{ij}(t)}+8\dot{X}\stackrel{(4)}{^{jj}(t)}{\displaystyle \frac{8}{9}}\dot{X}^{,i}\stackrel{(4)}{^{ijj}(t)}`$
$`\left(14UX^{,ij}+2\mathrm{\Sigma }(X^{,ij})4X_1^{,ij}+2X_2^{,ij}+{\displaystyle \frac{2}{3}}\ddot{Y}^{,ij}\right)\stackrel{(3)}{^{ij}(t)}`$
$`+{\displaystyle \frac{1}{3}}(70U^216\mathrm{\Phi }_1+20\mathrm{\Phi }_2+4\ddot{X})\stackrel{(3)}{^{jj}(t)}+{\displaystyle \frac{8}{3}}U\stackrel{(3)}{M^{jjkk}(t)}{\displaystyle \frac{4}{3}}X^{,ij}\stackrel{(3)}{M^{ijkk}(t)}`$
$`+{\displaystyle \frac{16}{9}}x^kX^{,ij}ϵ^{mk(i}\stackrel{(4)}{𝒥^{m|j)}(t)}{\displaystyle \frac{4}{9}}(8x^kU5X^{,k})ϵ^{mkj}\stackrel{(4)}{𝒥^{mj}(t)}+{\displaystyle \frac{16}{9}}\dot{X}^{,k}ϵ^{mkj}\stackrel{(3)}{𝒥^{mj}(t)},`$
$`B_{3.5}`$ $`=`$ $`{\displaystyle \frac{1}{60}}r^4\stackrel{(7)}{^{jj}(t)}+{\displaystyle \frac{1}{45}}r^2x^i\stackrel{(7)}{^{ijj}(t)}{\displaystyle \frac{1}{5}}(Y^{,ij}6x^iX^{,j})\stackrel{(5)}{^{ij}(t)}{\displaystyle \frac{1}{5}}(16X+2x^iX^{,i})\stackrel{(5)}{^{jj}(t)}`$ (305)
$`{\displaystyle \frac{4}{15}}X^{,i}\stackrel{(5)}{^{ijj}(t)}+12X^{i,j}\stackrel{(4)}{^{ij}(t)}+\left(6U^2+{\displaystyle \frac{16}{3}}\mathrm{\Phi }_112\mathrm{\Phi }_2\right)\stackrel{(3)}{^{jj}(t)}`$
$`(2UX^{,ij}2\mathrm{\Sigma }(X^{,ij})+4X_1^{,ij}2X_2^{,ij}8P_2^{ij})\stackrel{(3)}{^{ij}(t)}`$
$`+{\displaystyle \frac{4}{45}}r^2x^kϵ^{mkj}\stackrel{(6)}{𝒥^{mj}(t)}+{\displaystyle \frac{4}{3}}X^{,k}ϵ^{mkj}\stackrel{(4)}{𝒥^{mj}(t)}`$
$`{\displaystyle \frac{1}{15}}r^2\stackrel{(5)}{M^{jjkk}(t)}{\displaystyle \frac{2}{15}}x^{jk}\stackrel{(5)}{M^{iijk}(t)}+{\displaystyle \frac{2}{15}}x^j\stackrel{(5)}{M^{iijkk}(t)}{\displaystyle \frac{1}{30}}\stackrel{(5)}{M^{iijjkk}(t)}.`$
The final term in the expression for $`B_3`$ is purely a function of time, and as such does not affect the equations of motion through 3.5PN order. It comes in part from the surface terms Eqs. (C), and in part from various integrations by parts of Poisson potentials to achieve the expressions shown. In $`N_3`$, all such terms cancel. Similarly, purely time-dependent terms which appear in $`N_{3.5}`$ and $`B_{3.5}`$ do not contribute to the equations of motion.
As expected, the outer integrals make their first formal contribution to the field at 3PN order, however, the observable contribution vanishes to this order, so we have not shown any such contributions explicitly in Eqs. (VI B). In the next subsection, we study the contributions of the outer integrals in more detail, and show that through 3.5PN order, all contributions from the outer integrals are pure gauge terms.
### C Outer integrals and the contributions of “tails”
Our earlier qualitative discussion suggested that terms involving products of the monopole moment $``$ and the quadrupole moment $`^{ij}`$ of the far-zone fields would contribute via the outer integrals at 3PN order. Because higher multipole moments involve higher powers of $`1/r`$ or higher time derivatives, they would be expected to contribute at even higher PN order. Thus working through 3.5PN order, we might expect at most that products of $``$ with quadrupole $`^{ij}`$, octupole $`^{ijk}`$ or current quadrupole $`𝒥^{ij}`$ moments would contribute. Other terms, such as products of $``$ with higher-order moments, or products of higher-order moments, such as terms quadratic in $`^{ij}`$, will be 4PN order or higher. In studying the contribution of the outer integrals to the fields at 3.5PN order, therefore, it suffices to employ the far-zone field given in Eq. (V C). However, to illustrate the first non-trivial “tail” contribution, we will evaluate certain pieces of the outer integrals through 4PN order.
We substitute Eqs. (V C) into Eqs. (IV A) using the “quick and dirty” rule expressed by Eq. (131) to determine which terms to keep, and obtain, in the far zone,
$`\mathrm{\Lambda }^{00}`$ $`=`$ $`14r^2\widehat{n}^i_{ijk}(^{jk}/r)8r^1_{ij}(\ddot{}^{ij}/r)+8r^2\widehat{n}^i_j(\ddot{}^{ij}/r)24r^4\widehat{n}^{<ij>}\ddot{}^{ij}`$ (311)
$`2r^2\widehat{n}^i_i(\ddot{}^{jj}/r){\displaystyle \frac{14}{3}}r^2\widehat{n}^i_{ijkl}(^{jkl}/r)+{\displaystyle \frac{8}{3}}r^1_{ijk}(\ddot{}^{ijk}/r)`$
$`{\displaystyle \frac{8}{3}}r^2\widehat{n}^i_{jk}(\ddot{}^{ijk}/r)8r^3\widehat{n}^{<ij>}_k(\ddot{}^{ijk}/r){\displaystyle \frac{2}{3}}r^2\widehat{n}^i_{ij}(\ddot{}^{jkk}/r)`$
$`+{\displaystyle \frac{16}{3}}r^2\widehat{n}^iϵ^{aib}_{ak}(\dot{𝒥}^{bk}/r)32r^3\widehat{n}^{<ij>}ϵ^{aki}_k(\dot{𝒥}^{aj}/r){\displaystyle \frac{8}{3}}r^2\widehat{n}^iϵ^{akj}_{ik}(\dot{𝒥}^{aj}/r)`$
$`+O(\rho ϵ^4),`$
$`\mathrm{\Lambda }^{0i}`$ $`=`$ $`8r^2\widehat{n}^j_{ik}(\dot{}^{jk}/r)8r^2\widehat{n}^j_{jk}(\dot{}^{ik}/r)6r^2\widehat{n}^i_{jk}(\dot{}^{jk}/r)`$ (314)
$`+8r^1_j(\stackrel{(3)}{^{ij}}/r){\displaystyle \frac{8}{3}}r^1_{jk}(\stackrel{(3)}{^{ijk}}/r)+{\displaystyle \frac{16}{3}}r^1ϵ^{iab}_{ak}(\ddot{𝒥}^{bk}/r)`$
$`8r^2\widehat{n}^j(\stackrel{(3)}{^{ij}}/r)+2r^2\widehat{n}^i(\stackrel{(3)}{^{jj}}/r)+O(\rho ϵ^{7/2}),`$
$`\mathrm{\Lambda }^{ii}`$ $`=`$ $`2r^2\widehat{n}^i_{ijk}(^{jk}/r)8r^2\widehat{n}^i_j(\ddot{}^{ij}/r)+2r^2\widehat{n}^i_i(\ddot{}^{jj}/r)`$ (318)
$`8r^2\stackrel{(4)}{^{ii}}{\displaystyle \frac{2}{3}}r^2\widehat{n}^i_{ijkl}(^{jkl}/r)+{\displaystyle \frac{8}{3}}r^2\widehat{n}^i_{jk}(\ddot{}^{ijk}/r)`$
$`+{\displaystyle \frac{2}{3}}r^2\widehat{n}^i_{ij}(\ddot{}^{jkk}/r){\displaystyle \frac{8}{3}}r^1_i(\stackrel{(4)}{^{ijj}}/r){\displaystyle \frac{16}{3}}r^2\widehat{n}^iϵ^{aib}_{ak}(\dot{𝒥}^{bk}/r)`$
$`+{\displaystyle \frac{8}{3}}r^2\widehat{n}^iϵ^{akj}_{ik}(\dot{𝒥}^{aj}/r){\displaystyle \frac{32}{3}}r^1ϵ^{aki}_k(\stackrel{(3)}{𝒥^{ai}}/r)+O(\rho ϵ^4),`$
$`\mathrm{\Lambda }^{ij}`$ $`=`$ $`8r^2\stackrel{(4)}{^{ij}}+O(\rho ϵ^3).`$ (319)
All moments $`^{ij}`$, $`^{ijk}`$, and $`𝒥^{ij}`$ in these expression are functions of retarded time $`tr`$. Notice that the term kept in $`\mathrm{\Lambda }^{ij}`$ is actually of $`O(\rho ϵ^3)`$ (4PN order) according to our scheme, however because it has $`1/r^2`$ dependence, it will yield a 4PN tail contribution of a form which we wish to keep.
We expand the derivatives and evaluate the coefficients $`_{B,L}^q`$ and $`_{2,L}^q`$ \[Eqs. (44) and (46)\] for each term, throwing away all $``$-dependent terms. Terms with $`1/r^2`$ fall-off yield integrals over Legendre functions $`Q_L`$, as in Eq. (45). The result, through 3.5PN order (and keeping all formally 4PN terms involving integrals over $`Q_L`$), is
$`(N_3)_{𝒞𝒩}`$ $`=`$ $`\{8\widehat{n}^{<ij>}{\displaystyle _1^{\mathrm{}}}\stackrel{(4)}{^{ij}}(tr\zeta )Q_2(\zeta )d\zeta {\displaystyle \frac{8}{3}}{\displaystyle _1^{\mathrm{}}}\stackrel{(4)}{^{jj}}(tr\zeta )Q_0(\zeta )d\zeta `$ (322)
$`+{\displaystyle \frac{4}{3}}(\widehat{n}^{<ij>}2\delta ^{ij}2\delta ^{ij}\mathrm{ln}r)\stackrel{(4)}{^{ij}}(t)\},`$
$`(N_{3.5})_{𝒞𝒩}`$ $`=`$ $`\{{\displaystyle \frac{8}{3}}\widehat{n}^{<ijk>}{\displaystyle _1^{\mathrm{}}}\stackrel{(5)}{^{ijk}}(tr\zeta )Q_3(\zeta )d\zeta {\displaystyle \frac{8}{5}}\widehat{n}^i{\displaystyle _1^{\mathrm{}}}\stackrel{(5)}{^{ijj}}(tr\zeta )Q_1(\zeta )d\zeta `$ (324)
$`{\displaystyle \frac{2}{3}}r(3\widehat{n}^{<ij>}2\delta ^{ij})\stackrel{(5)}{^{ij}}(t)+{\displaystyle \frac{2}{45}}(5\widehat{n}^{<ijk>}+18\widehat{n}^i\delta ^{jk})\stackrel{(5)}{^{ijk}}(t)\},`$
$`(K_{3.5})_{𝒞𝒩}^i`$ $`=`$ $`\left\{8\widehat{n}^j{\displaystyle _1^{\mathrm{}}}\stackrel{(4)}{^{ij}}(tr\zeta )Q_1(\zeta )d\zeta +4\widehat{n}^j\stackrel{(4)}{^{ij}}(t)\right\},`$ (325)
$`(K_4)_{𝒞𝒩}^i`$ $`=`$ $`\{{\displaystyle \frac{8}{3}}\widehat{n}^{<jk>}{\displaystyle _1^{\mathrm{}}}\stackrel{(5)}{^{ijk}}(tr\zeta )Q_2(\zeta )d\zeta {\displaystyle \frac{8}{9}}{\displaystyle _1^{\mathrm{}}}\stackrel{(5)}{^{ijj}}(tr\zeta )Q_0(\zeta )d\zeta `$ (328)
$`+{\displaystyle \frac{16}{3}}\widehat{n}^{<ak>}ϵ^{iaj}{\displaystyle _1^{\mathrm{}}}\stackrel{(4)}{𝒥^{jk}}(tr\zeta )Q_2(\zeta )d\zeta `$
$`+{\displaystyle \frac{16}{9}}ϵ^{ikj}{\displaystyle _1^{\mathrm{}}}\stackrel{(4)}{𝒥^{jk}}(tr\zeta )Q_0(\zeta )d\zeta \},`$
$`(B_3)_{𝒞𝒩}`$ $`=`$ $`\left\{8{\displaystyle _1^{\mathrm{}}}\stackrel{(4)}{^{ii}}(tr\zeta )Q_0(\zeta )d\zeta +8(1\mathrm{ln}r)\stackrel{(4)}{^{ii}}(t)\right\},`$ (329)
$`(B_{3.5})_{𝒞𝒩}`$ $`=`$ $`\{+{\displaystyle \frac{8}{3}}\widehat{n}^i{\displaystyle _1^{\mathrm{}}}\stackrel{(5)}{^{ijj}}(tr\zeta )Q_1(\zeta )d\zeta +{\displaystyle \frac{32}{3}}ϵ^{aij}\widehat{n}^i{\displaystyle _1^{\mathrm{}}}\stackrel{(4)}{𝒥^{aj}}(tr\zeta )Q_1(\zeta )d\zeta `$ (331)
$`+4r\stackrel{(5)}{^{ii}}(t){\displaystyle \frac{4}{3}}\widehat{n}^i\stackrel{(5)}{^{ijj}}(t){\displaystyle \frac{16}{3}}ϵ^{aij}\widehat{n}^i\stackrel{(4)}{𝒥^{aj}}(t)\},`$
$`(B_4)_{𝒞𝒩}^{ij}`$ $`=`$ $`8{\displaystyle _1^{\mathrm{}}}\stackrel{(4)}{^{ij}}(tr\zeta )Q_0(\zeta )d\zeta .`$ (332)
Using the recursion relations satisfied by Legendre functions, we can establish the general formulae:
$`{\displaystyle _1^{\mathrm{}}}X(tr\zeta )Q_L(\zeta )𝑑\zeta `$ $`=`$ $`{\displaystyle \frac{1}{L(L+1)}}X(tr)`$ (334)
$`{\displaystyle \frac{1}{2L+1}}{\displaystyle _1^{\mathrm{}}}X^{}(tr\zeta )(Q_{L+1}(\zeta )Q_{L1}(\zeta ))𝑑\zeta ,`$
$`{\displaystyle _1^{\mathrm{}}}X(tr\zeta )Q_0(\zeta )𝑑\zeta `$ $`=`$ $`X(tr){\displaystyle _1^{\mathrm{}}}X^{}(tr\zeta )(Q_1(\zeta )+Q_0(\zeta ))𝑑\zeta `$ (336)
$`+{\displaystyle _0^{\mathrm{}}}\dot{X}(trs)\mathrm{ln}(s/2r)𝑑s,`$
where prime denotes $`/\zeta `$, $`s=r(\zeta 1)`$, $`X`$ represents one of the multipole moments of the system ($`^{ij}`$ and higher), and we assume that, in the distant past the system becomes sufficiently “stationary” that as $`s\mathrm{}`$, $`X(trs)\mathrm{ln}s0`$. Since for a binary system that becomes unbound ($`rv_0s`$) in the infinite past (because of gravitational-radiation anti-damping, looking backwards), $`X`$ in the worst case is proportional to $`(d/dt)^4^{ij}mv^4/r^2mv_0^2/s^2`$, then this boundary condition is satisfied (see for a detailed discussion of the past behavior of binary systems whose evolution includes gravitational radiation reaction). Repeated use of these identities allows us to convert many of the integrals in Eqs. (VI C) into integrals of higher time-derivatives of the expressions, which are thus of higher PN order, plus residual terms that cancel many of the non-integral terms in Eqs. (VI C). It is also useful to expand the retarded time $`trs`$ about $`ts`$, and to separate the $`\mathrm{ln}r`$ terms from the $`\mathrm{ln}(s/2)`$ terms in the integrals, leaving only terms proportional to $`X^{(n)}(t)`$ and $`_0^{\mathrm{}}X^{(n)}(ts)\mathrm{ln}(s/2)𝑑s`$. In the end, the only terms that remain at 3PN and 3.5PN order are
$`N_{𝒞𝒩}`$ $`=`$ $`\left\{{\displaystyle \frac{16}{3}}\stackrel{(4)}{^{ii}}(t){\displaystyle \frac{8}{3}}{\displaystyle _0^{\mathrm{}}}\stackrel{(5)}{^{ii}}(ts)\mathrm{ln}(s/2)ds\right\}+O(ϵ^5),`$ (338)
$`K_{𝒞𝒩}^i`$ $`=`$ $`O(ϵ^{9/2}),`$ (339)
$`B_{𝒞𝒩}^{ij}`$ $`=`$ $`O(ϵ^4),`$ (340)
$`B_{𝒞𝒩}`$ $`=`$ $`8{\displaystyle _0^{\mathrm{}}}\stackrel{(5)}{^{ii}}(ts)\mathrm{ln}(s/2)ds+O(ϵ^5).`$ (341)
As these are purely functions of time, they do not contribute to the equations of motion through 3.5PN order. Alternatively, one can show that the terms in Eqs. (VI C) turn out to be purely gauge terms through 3.5PN order. In fact, by making the gauge transformation $`h^{\mu \nu }h^{\mu \nu }\xi ^{\mu ,\nu }\xi ^{\nu ,\mu }+\eta ^{\mu \nu }\xi _{,\alpha }^\alpha `$ (the linear transformation suffices to this order), with
$`\xi ^0`$ $`=`$ $`\{{\displaystyle \frac{4}{3}}\stackrel{(3)}{^{ii}}(t)+{\displaystyle \frac{8}{3}}{\displaystyle _0^{\mathrm{}}}\stackrel{(4)}{^{ii}}(ts)\mathrm{ln}(s/2)ds{\displaystyle \frac{2}{3}}x^{ij}{\displaystyle _0^{\mathrm{}}}\stackrel{(6)}{^{ij}}(ts)\mathrm{ln}(s/2)ds`$ (344)
$`+{\displaystyle \frac{2}{3}}r^2{\displaystyle _0^{\mathrm{}}}\stackrel{(6)}{^{ii}}(ts)\mathrm{ln}(s/2)ds+{\displaystyle \frac{4}{45}}x^i{\displaystyle _0^{\mathrm{}}}\stackrel{(6)}{^{ijj}}(ts)\mathrm{ln}(s/2)ds`$
$`+{\displaystyle \frac{8}{9}}x^iϵ^{ikj}{\displaystyle _0^{\mathrm{}}}\stackrel{(5)}{𝒥^{jk}}(ts)\mathrm{ln}(s/2)ds\},`$
$`\xi ^i`$ $`=`$ $`\{4x^j{\displaystyle _0^{\mathrm{}}}\stackrel{(5)}{^{ij}}(ts)\mathrm{ln}(s/2)ds+{\displaystyle \frac{4}{3}}x^i{\displaystyle _0^{\mathrm{}}}\stackrel{(5)}{^{jj}}(ts)\mathrm{ln}(s/2)ds`$ (346)
$`+{\displaystyle \frac{44}{45}}{\displaystyle _0^{\mathrm{}}}\stackrel{(5)}{^{ijj}}(ts)\mathrm{ln}(s/2)ds{\displaystyle \frac{8}{9}}ϵ^{ikj}{\displaystyle _0^{\mathrm{}}}\stackrel{(4)}{𝒥^{jk}}(ts)\mathrm{ln}(s/2)ds\},`$
we can convert the outer integral contributions to $`h^{\alpha \beta }`$ in Eq. (VI C) to a form consisting of nothing but a 4PN tail term:
$`(N+B)_{𝒞𝒩}`$ $`=`$ $`{\displaystyle \frac{16}{5}}x^{ij}{\displaystyle _0^{\mathrm{}}}\stackrel{(7)}{^{<ij>}}(ts)\mathrm{ln}(s/2)ds+O(ϵ^5),`$ (347)
$`(K^i)_{𝒞𝒩}`$ $`=`$ $`O(ϵ^{9/2}),`$ (348)
$`(B^{ij})_{𝒞𝒩}`$ $`=`$ $`O(ϵ^4).`$ (349)
Note that, to this order, $`N+B=2g_{00}`$, and only the gradient of the term in Eq. (349) contributes to the acceleration, hence this term can be thought of as a 4PN tail modification of the Newtonian gravitational potential, or as a 1.5PN modification due to tails of the 2.5PN radiation-reaction potentials. This result is in complete agreement with the near-zone tail contribution derived by Blanchet and Damour using matched asymptotic expansions within the post-Minkowskian formalism.
## VII Discussion
We have presented a method for direct integration of the relaxed Einstein equations in a post-Newtonian expansion, applicable to equations of motion and gravitational radiation from isolated gravitating systems. As a foundation for future work, we presented a solution for the near-zone gravitational field through 3.5 post-Newtonian order in terms of Poisson potentials, together with a prescription for ensuring that no divergent or undefined integrals occur. In subsequent work, we will apply the near-zone results to the derivation of equations of motion for binary systems of compact objects through 2.5 PN order, and including 3.5 PN radiation reaction terms. Work on the 3PN contributions to the equations of motion is in progress.
The results presented here can also be applied to the gravitational radiation wave-form and energy flux from binary systems to as high as 3PN order beyond the quadrupole approximation. It can also be used to discuss equations of motion and radiation damping of systems containing spinning bodies, as well as the structure and evolution of fluid bodies. These will be the subject of future work.
###### Acknowledgements.
We gratefully acknowledge useful discussions with Luc Blanchet. This research was supported in part by the National Science Foundation under Grant No. PHY 96-00049. The initial phase of this work was also supported by Fellowships to CW from the Fulbright Foundation and the J. S. Guggenheim Foundation. CW is grateful to the Observatoire de Paris, Meudon and the Centre National de la Recherche Scientifique in France, and to the Hebrew University of Jerusalem, for hospitality and support during a sabbatical year where this work began in earnest.
## A STF Tensors and their properties
Throughout this series of papers, we shall make frequent use of the properties of symmetric, trace-free (STF) products of unit vectors. The general formula for such STF products is
$`\widehat{n}^{<L>}{\displaystyle \underset{p=0}{\overset{[l/2]}{}}}(1)^p{\displaystyle \frac{(2ll2p)!!}{(2l1)!!}}\left[\widehat{n}^{L2P}\delta ^P+\mathrm{sym}(\mathrm{q})\right],`$ (A1)
where $`[l/2]`$ denotes the integer just less than or equal to $`l/2`$, the capitalized superscripts denote the dimensionality, $`l2p`$ or $`p`$, of products of $`\widehat{n}^i`$ or $`\delta ^{ij}`$ respectively, and “sym(q)” denotes all distinct terms arising from permutations of indices, where $`q=l!/[(2^pp!(l2p)!]`$ is the total number of such terms (see for compendia of formulae). For convenience, we display the first several examples explicitly
$`\widehat{n}^{<ij>}`$ $`=`$ $`\widehat{n}^{ij}{\displaystyle \frac{1}{3}}\delta ^{ij},`$ (A3)
$`\widehat{n}^{<ijk>}`$ $`=`$ $`\widehat{n}^{ijk}{\displaystyle \frac{1}{5}}(\widehat{n}^i\delta ^{jk}+\widehat{n}^j\delta ^{ik}+\widehat{n}^k\delta ^{ij}),`$ (A4)
$`\widehat{n}^{<ijkl>}`$ $`=`$ $`\widehat{n}^{ijkl}{\displaystyle \frac{1}{7}}(\widehat{n}^{ij}\delta ^{kl}+\mathrm{sym}(6))+{\displaystyle \frac{1}{35}}(\delta ^{ij}\delta ^{kl}+\delta ^{ik}\delta ^{jl}+\delta ^{il}\delta ^{jk}),`$ (A5)
$`\widehat{n}^{<ijklm>}`$ $`=`$ $`\widehat{n}^{ijklm}{\displaystyle \frac{1}{9}}(\widehat{n}^{ijk}\delta ^{kl}+\mathrm{sym}(10))+{\displaystyle \frac{1}{63}}(\widehat{n}^i\delta ^{jk}\delta ^{lm}+\mathrm{sym}(15)),`$ (A6)
$`\widehat{n}^{<ijklmn>}`$ $`=`$ $`\widehat{n}^{ijklmn}{\displaystyle \frac{1}{11}}(\widehat{n}^{ijkl}\delta ^{mn}+\mathrm{sym}(15))+{\displaystyle \frac{1}{99}}(\widehat{n}^{ij}\delta ^{kl}\delta ^{mn}+\mathrm{sym}(45))`$ (A8)
$`{\displaystyle \frac{1}{693}}(\delta ^{ij}\delta ^{kl}\delta ^{mn}+\mathrm{sym}(15)).`$
There is a close connection between these STF tensors and spherical harmonics. For example, it is straightforward to show that, for any unit vector $`\widehat{𝐍}`$, the contraction of $`\widehat{N}^L`$ with $`\widehat{n}^{<L>}`$ is given by
$`\widehat{N}^L\widehat{n}^{<L>}={\displaystyle \frac{l!}{(2l1)!!}}P_l(\widehat{𝐍}\widehat{𝐧}),`$ (A9)
where $`P_l`$ is a Legendre polynomial.
## B Cancellation of $``$-dependence between inner and outer integrals
Here we demonstrate explicitly the cancellation of $``$-dependent terms between the inner and outer integrals. We assume that, at each iteration step, from just inside the boundary of the near zone out into the far zone, the source stress-energy tensor $`{}_{N1}{}^{}\mathrm{\Lambda }_{}^{\alpha \beta }`$ can be decomposed into terms of the form $`f_{B,L}(u)\widehat{n}^{<L>}r^B`$, where $`u=tr`$ is retarded time, and $`\widehat{n}^{<L>}`$ is a STF product of unit radial vectors. We calculate the behavior of the inner integral of such a term as the integration variable approaches $``$ from below with the result obtained from the outer integral of the same term. We consider far-zone and near-zone field points separately.
### 1 Far-zone field points
The inner integral is given by Eq. (17), with the multipole moment given by Eq. (18). We want to examine the behavior of the moment, as $`|𝐱^{}|`$, that is
$`M^{\alpha \beta \overline{Q}}(u)`$ $``$ $`{\displaystyle \frac{1}{16\pi }}{\displaystyle ^{}}f_{B,L}(ur^{}){\displaystyle \frac{\widehat{n}^{<L>}}{r_{}^{}{}_{}{}^{B}}}x_{}^{}{}_{}{}^{\overline{Q}}r_{}^{}{}_{}{}^{2}𝑑r^{}𝑑\mathrm{\Omega }^{}`$ (B1)
$`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^m}{m!}}f_{B,L}^{(m)}(u)G_{B,L,\overline{Q}}^m()\mathrm{\Delta }^{L,\overline{Q}},`$ (B2)
where the superscript $`(m)`$ denotes $`m`$ retarded time derivatives, and where
$`\mathrm{\Delta }^{L,\overline{Q}}`$ $`=`$ $`{\displaystyle \frac{1}{4\pi }}{\displaystyle \widehat{n}^{<L>}\widehat{n}^{\overline{Q}}𝑑\mathrm{\Omega }},`$ (B4)
$`G_{B,L,\overline{Q}}^m()`$ $`=`$ $`{\displaystyle ^{}}r_{}^{}{}_{}{}^{2+\overline{q}B+m}𝑑r^{}`$ (B5)
$`=`$ $`\{\begin{array}{cc}^{3+\overline{q}B+m}/(3+\overline{q}B+m)\hfill & 3+\overline{q}B+m0\hfill \\ \mathrm{ln}\hfill & 3+\overline{q}B+m=0\hfill \end{array}.`$ (B8)
Then, from inside $``$,
$$h_{𝒩}^{\alpha \beta }{}_{B,L}{}^{}\underset{\overline{q}=0}{\overset{\mathrm{}}{}}\frac{(1)^{\overline{q}}}{\overline{q}!}\underset{m=0}{\overset{\mathrm{}}{}}\frac{(1)^m}{m!}_{\overline{Q}}\left(\frac{1}{r}f_{B,L}^{(m)}(u)\right)G_{B,L,\overline{Q}}^m()\mathrm{\Delta }^{L,\overline{Q}}.$$
(B9)
It is straightforward to show that the contraction of $`_{\overline{Q}}`$ with $`\mathrm{\Delta }^{L,\overline{Q}}`$ is given by
$$\mathrm{\Delta }^{L,\overline{Q}}_{\overline{Q}}=\{\begin{array}{cc}0\hfill & \overline{q}<L\hfill \\ 0\hfill & L+\overline{q}=\mathrm{odd}\hfill \\ \frac{2^L\overline{q}!((\overline{q}+L)/2)!}{(\overline{q}+L+1)!((\overline{q}L)/2)!}|^2|^{(\overline{q}L)/2}_{<L>}\hfill & \overline{q}L\hfill \end{array}.$$
(B10)
Using the fact that
$`^2\left({\displaystyle \frac{f(u)}{r}}\right)`$ $`=`$ $`{\displaystyle \frac{\ddot{f}}{r}},`$ (B12)
$`_{<L>}\left({\displaystyle \frac{f(u)}{r}}\right)`$ $`=`$ $`(1)^L\widehat{n}^{<L>}{\displaystyle \underset{k=0}{\overset{L}{}}}{\displaystyle \frac{(L+k)!}{2^kk!(Lk)!}}{\displaystyle \frac{f^{(Lk)}(u)}{r^{k+1}}},`$ (B13)
(see eg. ) and redefining summation variables, $`q=m+\overline{q}k`$, $`j=Lk`$, we obtain Eqs. (47) and (51). Evaluating the outer integral for the same term yields $`z`$-dependent or $`\mathrm{ln}`$-dependent terms that are precisely equal and opposite those of Eq. (51).
### 2 Near-zone field point
In the near zone, for $`|𝐱^{}|>|𝐱|`$, Eq. (27) together with the specific decomposition of $`\mathrm{\Lambda }^{\alpha \beta }`$ gives
$${}_{N}{}^{}h_{𝒩}^{\alpha \beta }{}_{B,L}{}^{}\frac{1}{4\pi }\underset{\overline{q}=0}{\overset{\mathrm{}}{}}\frac{(1)^{\overline{q}}}{\overline{q}!}\underset{m=0}{\overset{\mathrm{}}{}}\frac{(1)^m}{m!}x^{\overline{Q}}_t^n^{}f_{B,L}(tr^{})\frac{\widehat{n}^{<L>}}{r_{}^{}{}_{}{}^{B}}_{\overline{Q}}^{}(r_{}^{}{}_{}{}^{m1})r_{}^{}{}_{}{}^{2}dr^{}d\mathrm{\Omega }^{}.$$
(B14)
We use the fact that
$`_{\overline{Q}}^{}r_{}^{}{}_{}{}^{m1}`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{k_m}{}}}{\displaystyle \frac{(2\overline{q}4k+1)!!}{(2\overline{q}2k+1)!!}}{\displaystyle \frac{m!}{(m2k)!}}\left[{\displaystyle \frac{(m2k1)!!}{(m2\overline{q}+2k1)!!}}\right]`$ (B16)
$`\times {\displaystyle \frac{\overline{q}!}{2^kk!(\overline{q}2k)!}}\delta ^K\widehat{n}^{<\overline{Q}2K>}r_{}^{}{}_{}{}^{m\overline{q}1},`$
where $`k_m=\mathrm{lesser}\mathrm{of}\{[\overline{q}/2],[m/2]\}`$, $`\delta ^K`$ denotes a product of $`K`$ Kronecker deltas, the quantity in square brackets can be evaluated for negative or positive values of the arguments, and the expression $`\delta ^K\widehat{n}^{<\overline{Q}2K>}`$ is to be symmetrized on all indices (since the expression ultimately is to be contracted on $`x^{\overline{Q}}`$ no explicit symmetrization is needed). It can then be shown that
$$\widehat{n}^{\overline{Q}}\frac{1}{4\pi }\delta ^K\widehat{n}^{<\overline{Q}2K>}\widehat{n}^{<L>}𝑑\mathrm{\Omega }^{}=\frac{L!}{(2L+1)!!}\delta _{L,\overline{q}2k}\widehat{n}^{<L>}.$$
(B17)
We then expand $`f(tr^{})=\mathrm{\Sigma }_{n=0}^{\mathrm{}}(1)^nf^{(n)}(t)r_{}^{}{}_{}{}^{n}/n!`$, integrate over $`r^{}`$ toward $``$, rearrange the summations, and define $`r=(\overline{q}L)/2`$, and $`q=m+n`$, and obtain
$${}_{N}{}^{}h_{𝒩}^{\alpha \beta }{}_{B,L}{}^{}\left(\frac{2}{r}\right)^{B2}\widehat{n}^{<L>}\underset{q=0}{\overset{\mathrm{}}{}}_{B,L}^{\mathrm{in},q}(z)r^q\frac{d^qf_{B,L}(t)}{dt^q},$$
(B18)
with
$`_{B,L}^{\mathrm{in},q}(z)`$ $`=`$ $`{\displaystyle \underset{r=0}{\overset{[q/2]}{}}}{\displaystyle \underset{m=2r}{\overset{q}{}}}{\displaystyle \frac{(1)^{L+q}(2)^{2+LB}(L+r)!}{(qm)!(2L+2r+1)!(m2r)!r!}}\left[{\displaystyle \frac{(m12r)!!}{(m12r2L)!!}}\right]`$ (B22)
$`\times \{\begin{array}{cc}z^{qL2rB+2}/(qL2rB+2)\hfill & qL2rB+20\hfill \\ \mathrm{ln}\hfill & qL2rB+2=0.\hfill \end{array}`$
Here too, evaluating the outer integral for the same term, for each $`B`$, $`L`$ and $`q`$ yields $`z`$-dependent or $`\mathrm{ln}`$-dependent terms that are precisely equal and opposite those of Eq. (B22).
### 3 Source terms with $`\mathrm{ln}r`$ dependence
Until now we have assumed that the stress-energy source $`\mathrm{\Lambda }^{\alpha \beta }`$ can be decomposed into terms of the form $`f_{B,L}(u)\widehat{n}^{<L>}r^B`$. At sufficiently high PN order, tail contributions to the fields will arise, leading to the possibility of $`\mathrm{ln}r`$ dependence in $`\mathrm{\Lambda }^{\alpha \beta }`$. To illustrate that cancellation of $``$ dependence occurs in this event also, we consider source terms of the form $`f_{B,L}(u^{})\widehat{n}^{<L>}r_{}^{}{}_{}{}^{B}\mathrm{ln}r^{}`$. Noting that, from Eq. (32), $`\mathrm{ln}r^{}=\mathrm{ln}[2(\zeta y)/r(\zeta ^21)]`$, and incorporating this logarithmic term into the outer integral, Eq. (33), we obtain
$`{}_{N}{}^{}h_{𝒞𝒩}^{\alpha \beta (\mathrm{ln})}{}_{B,L}{}^{}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\widehat{n}^{<L>}{\displaystyle _1^1}P_L(y)𝑑y{\displaystyle _{\zeta (y)}^{\mathrm{}}}\left({\displaystyle \frac{2(\zeta y)}{r(\zeta ^21)}}\right)^{B2}\mathrm{ln}\left({\displaystyle \frac{2(\zeta y)}{r(\zeta ^21)}}\right)`$ (B24)
$`\times f_{B,L}(ur(\zeta 1)){\displaystyle \frac{d\zeta }{\zeta y}}`$
$`=`$ $`{\displaystyle \frac{}{B}}{}_{N}{}^{}h_{𝒞𝒩}^{\alpha \beta }{}_{B,L}{}^{}.`$ (B25)
For the inner integral, the only difference which the logarithmic term makes is in the radial integral, now given by
$`G_{B,L,\overline{Q}}^m()^{(\mathrm{ln})}`$ $`=`$ $`{\displaystyle ^{}}r_{}^{}{}_{}{}^{2+\overline{q}B+m}\mathrm{ln}r^{}dr^{}`$ (B26)
$`=`$ $`{\displaystyle \frac{}{B}}G_{B,L,\overline{Q}}^m().`$ (B27)
Thus, if the original coefficients cancel for all $`B`$ (and if we can treat $`B`$ formally as a continuous parameter), then the coefficients generated by $`\mathrm{ln}r`$ terms cancel.
An alternative method is to show directly from the definitions (eg. Eqs. (B25) and (B27)) that, for both the inner and outer integrals and for $`z<1`$ and $`z>1`$,
$${}_{N}{}^{}h_{B,L}^{\alpha \beta (\mathrm{ln})}=\mathrm{ln}_Nh_{B,L}^{\alpha \beta }_1^z{}_{N}{}^{}h_{B,L}^{\alpha \beta }𝑑\overline{z}/\overline{z},$$
(B28)
modulo $`z`$\- or $``$-independent terms. Then, if the $`z`$-dependent parts of $`{}_{N}{}^{}h_{B,L}^{\alpha \beta }`$ cancel between outer and inner integrals, so too do the $`z`$-dependent parts of $`{}_{N}{}^{}h_{B,L}^{\alpha \beta (\mathrm{ln})}`$.
## C Boundary Terms
The boundary terms in $`h_𝒩^{\alpha \beta }`$ that arise from integrating by parts various integrals over $``$ are given by
$`N_{}`$ $`=`$ $`4{\displaystyle _{}}\tau ^{0j}(t,𝐱^{})d^2S_j^{}+{\displaystyle \frac{2}{3}}r^2_t^2{\displaystyle _{}}\tau ^{0j}(t,𝐱^{})d^2S_j^{}{\displaystyle \frac{4}{3}}x^i_t{\displaystyle _{}}\tau ^{ij}(t,𝐱^{})d^2S_j^{}`$ (C4)
$`{\displaystyle \frac{4}{3}}x^i_t^2{\displaystyle _{}}\tau ^{0j}(t,𝐱^{})x_{}^{}{}_{}{}^{i}d^2S_j^{}+{\displaystyle \frac{1}{30}}r^4_t^4{\displaystyle _{}}\tau ^{0j}(t,𝐱^{})d^2S_j^{}`$
$`{\displaystyle \frac{2}{15}}r^2x^i_t^4{\displaystyle _{}}\tau ^{0j}(t,𝐱^{})x_{}^{}{}_{}{}^{i}d^2S_j^{}{\displaystyle \frac{2}{15}}r^2x^i_t^3{\displaystyle _{}}\tau ^{ij}(t,𝐱^{})d^2S_j^{},`$
$`K_{}^i`$ $`=`$ $`4{\displaystyle _{}}\tau ^{ij}(t,𝐱^{})d^2S_j^{}+{\displaystyle \frac{2}{3}}r^2_t^2{\displaystyle _{}}\tau ^{ij}(t,𝐱^{})d^2S_j^{}+{\displaystyle \frac{2}{3}}x^k_t^3{\displaystyle _{}}\tau ^{0j}(t,𝐱^{})x_{}^{}{}_{}{}^{ik}d^2S_j^{}`$ (C6)
$`{\displaystyle \frac{4}{3}}x^k_t^2{\displaystyle _{}}\tau ^{j[i}(t,𝐱^{})x_{}^{}{}_{}{}^{k]}d^2S_j^{}{\displaystyle \frac{2}{9}}_t^3{\displaystyle _{}}\tau ^{0j}(t,𝐱^{})r_{}^{}{}_{}{}^{2}x_{}^{}{}_{}{}^{i}d^2S_j^{},`$
$`B_{}^{ij}`$ $`=`$ $`4_t{\displaystyle _{}}\tau ^{k(i}(t,𝐱^{})x_{}^{}{}_{}{}^{j)}d^2S_k^{}2_t^2{\displaystyle _{}}\tau ^{0k}(t,𝐱^{})x_{}^{}{}_{}{}^{ij}d^2S_k^{}`$ (C15)
$`{\displaystyle \frac{2}{3}}r^2_t^3{\displaystyle _{}}\tau ^{k(i}(t,𝐱^{})x_{}^{}{}_{}{}^{j)}d^2S_k^{}{\displaystyle \frac{1}{3}}r^2_t^4{\displaystyle _{}}\tau ^{0k}(t,𝐱^{})x_{}^{}{}_{}{}^{ij}d^2S_k^{}`$
$`+{\displaystyle \frac{2}{3}}x^l_t^3{\displaystyle _{}}\tau ^{k(i}(t,𝐱^{})x_{}^{}{}_{}{}^{jl)}d^2S_k^{}+{\displaystyle \frac{2}{9}}x^l_t^4{\displaystyle _{}}\tau ^{0k}(t,𝐱^{})x_{}^{}{}_{}{}^{ijl}d^2S_k^{}`$
$`+{\displaystyle \frac{8}{9}}x^l_t^3{\displaystyle _{}}(\tau ^{k[i}(t,𝐱^{})x_{}^{}{}_{}{}^{l]j}+\tau ^{k[j}(t,𝐱^{})x_{}^{}{}_{}{}^{l]i})d^2S_k^{}`$
$`{\displaystyle \frac{1}{18}}_t^3{\displaystyle _{}}[\tau ^{lk}(t,𝐱^{})(r_{}^{}{}_{}{}^{2}x_{}^{}{}_{}{}^{ij})_{,l}+\dot{\tau }^{0k}(t,𝐱^{})r_{}^{}{}_{}{}^{2}x_{}^{}{}_{}{}^{ij}]d^2S_k^{}`$
$`+{\displaystyle \frac{1}{3}}_t^3{\displaystyle _{}}(\tau ^{k[l}(t,𝐱^{})x_{}^{}{}_{}{}^{i]jl}+\tau ^{k[l}(t,𝐱^{})x_{}^{}{}_{}{}^{j]il})d^2S_k^{}`$
$`{\displaystyle \frac{1}{30}}r^4_t^5{\displaystyle _{}}\tau ^{k(i}(t,𝐱^{})x_{}^{}{}_{}{}^{j)}d^2S_k^{}{\displaystyle \frac{1}{60}}r^4_t^6{\displaystyle _{}}\tau ^{0k}(t,𝐱^{})x_{}^{}{}_{}{}^{ij}d^2S_k^{}`$
$`+{\displaystyle \frac{1}{15}}r^2x^l_t^5{\displaystyle _{}}\tau ^{k(i}(t,𝐱^{})x_{}^{}{}_{}{}^{jl)}d^2S_k^{}+{\displaystyle \frac{1}{45}}r^2x^l_t^6{\displaystyle _{}}\tau ^{0k}(t,𝐱^{})x_{}^{}{}_{}{}^{ijl}d^2S_k^{}`$
$`+{\displaystyle \frac{4}{45}}r^2x^l_t^5{\displaystyle _{}}(\tau ^{k[i}(t,𝐱^{})x_{}^{}{}_{}{}^{l]j}+\tau ^{k[j}(t,𝐱^{})x_{}^{}{}_{}{}^{l]i})d^2S_k^{}.`$
## D Properties of Poisson Potentials
Here we list some useful properties of Poisson potentials and superpotentials, given by Eqs. (IV D). These rely upon the general result, which can be obtained by integration by parts,
$$P(^2g)=g+_P(g),$$
(D1)
where $`_P(g)`$ denotes the boundary term, given by
$$_P(g)\frac{1}{4\pi }_{}\left[\frac{g(t,𝐱^{})}{|𝐱𝐱^{}|}_r^{}\mathrm{ln}(g(t,𝐱^{})|𝐱𝐱^{}|)\right]_{r^{}=}^2𝑑\mathrm{\Omega }^{}.$$
(D2)
The boundary terms must be carefully evaluated case by case to determine if any $``$-independent terms survive. All $``$-dependent terms can be discarded. Some useful formulae that result from this include:
$`P(|g|^2)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\{g^2+2P(g^2g)_P(g^2)\},`$ (D4)
$`P(gf)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\{fg+P(f^2g)+P(g^2f)_P(fg)\},`$ (D5)
$`P(f|U|^2)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\{fU^2+P(U^2^2f)2\mathrm{\Sigma }(fU)+4P(UUf)_P(fU^2)\}.`$ (D6)
In many specific cases, the boundary terms can be dropped:
$`P(U)`$ $`=`$ $`{\displaystyle \frac{1}{2}}X,`$ (D8)
$`P(X)`$ $`=`$ $`{\displaystyle \frac{1}{12}}Y,`$ (D9)
$`P(|U|^2)`$ $`=`$ $`{\displaystyle \frac{1}{2}}U^2+\mathrm{\Phi }_2,`$ (D10)
$`P(x^iU^{,jk\mathrm{}})`$ $`=`$ $`{\displaystyle \frac{1}{2}}x^iX^{,jk\mathrm{}}+{\displaystyle \frac{1}{12}}Y^{,ijk\mathrm{}},`$ (D11)
$`P(r^2U^{,ij})`$ $`=`$ $`{\displaystyle \frac{1}{2}}r^2X^{,ij}{\displaystyle \frac{1}{12}}Y^{,ij}+{\displaystyle \frac{1}{6}}x^kY^{,ijk},`$ (D12)
while in others, there are contributions from the boundary terms. For example, in the 2PN potential $`P(U\ddot{X})`$, the boundary term yields the term $`\frac{1}{2}_{}\sigma (t,𝐱)d^3x_t^2_{}\sigma (t,𝐲)d^3y`$. Using Eq. (176), we obtain, to the necessary order,
$$P(U\ddot{X})=\frac{1}{2}\{U\ddot{X}\mathrm{\Sigma }(\ddot{X})+2G_2\frac{1}{2}\stackrel{(4)}{^{ii}(t)}\}+O(ϵ^5).$$
(D13)
Similarly, we find for the 3PN potential,
$$P(U\stackrel{(4)}{Y})=\frac{1}{2}\{U\stackrel{(4)}{Y}\mathrm{\Sigma }(\stackrel{(4)}{Y})+12P(U\stackrel{(4)}{X})2\stackrel{(4)}{^{ii}(t)}\}+O(ϵ^5).$$
(D14)
For the Poisson Superpotential $`S(f)`$, we have
$$S(^2g)=2P(g)+_S(g),$$
(D15)
where
$$_S(g)\frac{1}{4\pi }_{}\left[g(t,𝐱^{})|𝐱𝐱^{}|_r^{}\mathrm{ln}\left(\frac{g(t,𝐱^{})}{|𝐱𝐱^{}|}\right)\right]_{r^{}=}^2𝑑\mathrm{\Omega }^{}.$$
(D16)
Thus, for example, in the superpotential $`(/t)^2_{}\tau ^{00}|𝐱𝐱^{}|d^3x^{}`$, we find the term
$`\ddot{S}(^2U^2)`$ $`=`$ $`2\ddot{P}(U^2)3(d/dt)^2({\displaystyle _{}}\sigma d^3x)^2+O(ϵ^5)`$ (D17)
$`=`$ $`4G_1+4G_23\stackrel{(4)}{^{ii}(t)}+O(ϵ^5).`$ (D18)
Other useful identities include
$`\mathrm{\Sigma }(x^i)`$ $`=`$ $`x^iUX^{,i},`$ (D20)
$`\mathrm{\Sigma }(x^{ij})`$ $`=`$ $`{\displaystyle \frac{1}{3}}Y^{,ij}\delta ^{ij}X+x^{ij}U2x^{(i}X^{,j)}.`$ (D21)
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# Microscopic theory for quantum mirages in quantum corrals
## I Introduction
Scanning Tunneling Microscopy (STM) allows the manipulation of single atoms on top of a surface as well as the construction of quantum structures of arbitrary shape. Additionally, the differential conductance, $`G(V)dI/dV`$, is proportional to the local density of states (LDOS) of the surface spot below the tip . Hence, STM can be used to modify and to measure the LDOS.
A STM was used by Crommie et al. to build a quantum corral, i.e., a 71 $`\AA `$ radius circle made with 48 atoms of iron on top of a surface of copper . The free motion of the electrons along the surface changed in the presence of the Fe atoms so that quasi-bound states appeared inside the corral. The measured LDOS was quite similar to that of a gas of noninteracting electrons inside a circular confining potential.
More recently, STM has permitted to study the problem of a single magnetic impurity embedded in the two dimensional electron gas formed on a metallic surface . This is the famous Kondo problem. Below the Kondo temperature, $`T_K`$, a many electron singlet state forms so that the spin of the magnetic impurity is screened by the conduction electrons. As a consequence, the impurity density of states develops a resonance at the Fermi energy (the Abrikosov-Suhl resonance ). When the STM tip is placed on top of the magnetic impurity, $`G(V)`$ displays a narrow dip around the Fermi level . The dip (instead of the resonance) is due to a Fano type interference between the direct tunneling from the tip to the surface and the additional chanels that appear due to the presence of the impurity . The depth of the dip decreases gradually as the lateral distance between the tip and the impurity is increased. This permits to image the magnetic atom. The dip vanishes when lateral tip-magnetic impurity distance is bigger than $`10\AA `$, which is twice $`k_F^1`$, the inverse of the Fermi vector.
This situation is dramatically changed when the magnetic impurity is placed at the focus of an elliptical corral of size smaller than 150 $`\AA `$ , built on the Cu(111) surface. In this configuration, the Kondo dip is observed not only on the focus where the magnetic impurity is located but also on the empty focus, which can be as far as 110 $`\AA `$ away from the impurity. Remarkably, the phantom dip is not observed if either the tip or the impurity are not at the foci. The phenomenon of the phantom dip is referred to as the quantum mirage . In this paper we provide a quantum mechanical theory for this phenomenon. In particular we want to address the issue of under which conditions the quantum mirage can be observed and whether an elliptical corral is necessary to obtain the mirage. We show that the elliptical geometry is convenient but not necessary and we show that there is no need to invoke semiclassical arguments to explain the mirage.
The structure of this paper is the following: In section II we review the theoretical framework adequate to study the quantum mirage. First, we present the Hamiltonian of a surface with both a magnetic impurity and a quantum corral. Then we give a formal expression for the relation between $`G(V)`$ and the surface LDOS. Our original contribution starts in section III, where we give a qualitative explanation for the quantum mirage. In section IV we present quantitative results for elliptical quantum corrals and in section V we discuss our results as well as the limitations of our theory.
## II Theoretical Framework
### A The Hamiltonian of the surface
The Hamiltonian of the surface is an extension of the well known Anderson model to the case in which the electrons feel the potential produced by the atoms creating the corral:
$`H_{\mathrm{surf}}={\displaystyle \underset{n,\sigma }{}}ϵ_nc_{n,\sigma }^{}c_{n,\sigma }+ϵ_d{\displaystyle \underset{\sigma }{}}d_\sigma ^{}d_\sigma `$ (1)
$`+Ud_{}^{}d_{}d_{}^{}d_{}+V_h{\displaystyle \underset{n,\sigma }{}}\psi _n^{}(\stackrel{}{R}_I)c_{n,\sigma }^{}d_\sigma +h.c..`$ (2)
$`ϵ_n`$ and $`\psi _n(\stackrel{}{R})`$ are the eigenvalues and eigenfunctions of the surface corral Hamiltonian. $`c_{n,\sigma }^{}`$ and $`d_\sigma ^{}`$ create an electron in the state $`n`$ of the corral, and in the magnetic impurity, respectively. In this work we only consider the states from the metallic surface band, which seem to give the main contribution to the LDOS measured by the STM . The first term in (2) describes the Fermi sea formed by filling these states. The second term in (2) is the impurity single particle energy. The third term is the on-site repulsion felt whenever two electrons are at the impurity site. The last term describes the hopping between the surface and the impurity states. In the Anderson model, this coupling is localized at the impurity site, $`\stackrel{}{R}_I`$. From the formal point of view, the presence of the corral is accounted for by replacing the plane waves, which diagonalize the free surface electron Hamiltonian, by the corral states. Throughout the paper we neglect the magnetic moment of the corral atoms. This is justified because the mirage appears also when the corral atoms are non-magnetic .
It must also be noted that Hamiltonian (2) does not contain any scattering from the surface states to the bulk states, a process which could occur due to the presence of both the impurity and the corral atoms. These physical processes should be considered in order to have a more quantitative theory of this system, something beyond the scope of this paper.
### B $`G(V,\stackrel{}{R})`$ vs. LDOS
We now review the link between the quantity measured in the experiments, $`G(V,\stackrel{}{R})=dI/dV(\stackrel{}{R})`$, the differential conductance measured when the tip is at position $`\stackrel{}{R}`$ on the surface, and the surface Green’s function, $`𝒢_S(\stackrel{}{R},ϵ^+)`$. The Hamiltonian of the whole system, tip and surface, can be written as the sum of three terms, $`H=H_{\mathrm{tip}}+H_{\mathrm{surf}}+H_{\mathrm{tun}}`$. The first is the Hamiltonian of the tip. The second, given in equation (2), corresponds to the Hamiltonian of the surface, including the corral and the magnetic impurity. The third is the tunneling (Bardeen) Hamiltonian, which describes processes in which an electron is transferred between the tip and the surface :
$$H_{\mathrm{tun}}=\underset{\sigma }{}A_\sigma ^{}\left(t_c\mathrm{\Psi }_\sigma (\stackrel{}{R})+t_dd_\sigma \right)+h.c.,$$
(3)
where
$$\mathrm{\Psi }_\sigma ^{}(\stackrel{}{R})=\underset{n}{}\psi _n^{}(\stackrel{}{R})c_{n,\sigma }^{},$$
(4)
creates a surface electron in the spin state $`\sigma `$ at the position $`\stackrel{}{R}`$ of the surface and $`A_\sigma ^{}`$ creates an electron in the tip. $`t_c`$ is the tunneling amplitude to the surface states and $`t_d`$ is the amplitude for tunneling directly to the magnetic impurity. $`t_d`$ has to be taken into account only when the tip is located very near the magnetic adatom ($`\stackrel{}{R}\stackrel{}{R}_I`$). We assume the knowledge of the eigenstates of the tip and the surface Hamiltonians and treat the tunneling term as a perturbation. To lowest order in the tunneling Hamiltonian and low enough temperatures, linear response predicts :
$$\frac{dI}{dV}(\stackrel{}{R})G(V,\stackrel{}{R})=\frac{4e^2}{\pi \mathrm{}}\rho _T\rho _S(ϵ_F+eV,\stackrel{}{R}),$$
(5)
where $`eV`$ is the voltage drop and $`\rho _T`$ is the density of states of the tip (assumed to be energy independent in the vicinity of $`ϵ_F`$). We follow the convention that positive $`eV`$ means electrons flowing towards the surface. Finally, the local density of states of the surface, $`\rho _S`$, is related to the retarded surface Green’s function through the relation:
$$\rho _S(\stackrel{}{R},ϵ_F+eV)=\frac{1}{\pi }Im\left[𝒢_S(\stackrel{}{R},ϵ_F+eV)\right].$$
(6)
$`𝒢_𝒮`$ is the retarded Green’s function corresponding to the operator $`t_c\mathrm{\Psi }_\sigma (\stackrel{}{R})+t_dd_\sigma `$, and is given by :
$`𝒢_S(\stackrel{}{R},ϵ^+)=t_c^2𝒢_c(\stackrel{}{R},\stackrel{}{R},ϵ^+)+𝒢_d(ϵ^+)\times `$ (7)
$`\left(t_d+t_cV_h𝒢_c(\stackrel{}{R},\stackrel{}{R}_I,ϵ^+)\right)\left(t_d+t_cV_h𝒢_c(\stackrel{}{R}_I,\stackrel{}{R},ϵ^+)\right),`$ (8)
where $`ϵ^+ϵ+i\eta `$. In the surface Green’s function (8), two different propagators appear. The first is the impurity free ($`U=0`$, $`V_h=0`$) surface Green’s function:
$$𝒢_c(\stackrel{}{R}_1,\stackrel{}{R}_2,ϵ^+)=\underset{n}{}\frac{\psi _n^{}(\stackrel{}{R}_1)\psi _n(\stackrel{}{R}_2)}{ϵ^+ϵ_n}.$$
(9)
The second is the Green’s function at the impurity site, whose evaluation is the difficult part of the many body problem . For temperatures much lower than $`T_K`$, $`𝒢_d`$ can be approximated by the Green’s function of an effective resonant level with a broadening $`T_K`$:
$$𝒢_d(ϵ^+)=\frac{Z_K}{ϵϵ_F+ik_BT_K},$$
(10)
where $`Z_K`$ is chosen so that the impurity propagator fulfills the Friedel sum rule :
$$Z_K\frac{T_K}{\pi V_h^2\rho },$$
(11)
where $`\rho =\frac{1}{\pi }Im\left[𝒢_c(\stackrel{}{R}_I,\stackrel{}{R}_I,ϵ_F)\right]`$ is the impurity free surface LDOS at the impurity site and at $`ϵ_F`$. A necessary condition for the appearance of the Kondo resonance is that the conduction band is formed by a quasi-continuum of states, with energy spacing $`\mathrm{\Delta }<T_K`$ . In the case of the quantum corrals that we study below $`\mathrm{\Delta }>T_K`$. However, the broadening, $`\delta `$, of these states, fulfills $`\delta >T_K`$, so that the density of states (in the absence of the magnetic impurity) is almost flat close to $`ϵ_F`$ and we can use equation (10).
The surface Green’s function can be expressed now as:
$`𝒢_S(\stackrel{}{R},ϵ^+)=t_c^2(𝒢_c(\stackrel{}{R},\stackrel{}{R},ϵ^+)+{\displaystyle \frac{T_K/\pi \rho }{ϵϵ_F+ik_BT_k}}\times `$ (12)
$`({\displaystyle \frac{t_d}{t_cV_h}}+𝒢_c(\stackrel{}{R},\stackrel{}{R}_I,ϵ^+))({\displaystyle \frac{t_d}{t_cV_h}}+𝒢_c(\stackrel{}{R}_I,\stackrel{}{R},ϵ^+))).`$ (13)
For the case $`t_d=0`$ (tip located far from the magnetic impurity) we can eliminate the parameter $`V_h`$ from our problem, due to the Friedel sum rule. When the tip is placed exactly at the magnetic impurity site ($`\stackrel{}{R}=\stackrel{}{R}_I`$), $`t_d`$ is no longer zero and we need to estimate the parameter $`t_d/(t_cV_h)`$. To do that, we proceed as follows. In the absence of corral atoms, we can approximate the impurity free surface Green’s function by $`𝒢_c(\stackrel{}{R}_I,\stackrel{}{R}_I,ϵ^+)i\rho _0`$ and one obtains the well known Fano function for the differential conductance through the tip :
$$G(V,\stackrel{}{R}_I)=\frac{4e^2}{\pi \mathrm{}}\rho _T\rho _0t_c^2\frac{(q+ϵ^{})^2}{1+ϵ_{}^{}{}_{}{}^{2}},$$
(14)
where $`ϵ^{}=(eVϵ_F)/T_K`$, $`\pi \rho _0q=t_d/(t_cV_h)`$ and $`\rho _0`$ is the LDOS at the Fermi Level for the surface states in the absence of quantum corral. $`q`$ is the Fano parameter which determines the shape of the $`G(V,\stackrel{}{R}_I)`$ curves. It can can take values between 0 (symmetric dip) and $`\mathrm{}`$ (Breit-Wigner). We obtain $`q`$, and therefore $`t_d/(t_cV_h)`$, by fitting the $`G(V,\stackrel{}{R}_I)`$ curve to the Fano lineshape in the case of tunneling through the magnetic adatom in the absence of corral.
## III The Quantum Mirage: Qualitative Explanation
In this section we give a qualitative explanation of the mirage, based on the general formalism of the previous section. We need to do several plausible hypothesis:
We suppose that the mirage is produced by quasi-bound states of the corral (an assumption that is consistent with the experiments ). Hence, we approximate the conduction Green’s function by:
$$𝒢_c(\stackrel{}{R}_1,\stackrel{}{R}_2,ϵ^+)\underset{n}{}\frac{\psi _n^{}(\stackrel{}{R}_1)\psi _n(\stackrel{}{R}_2)}{ϵϵ_n+i\delta },$$
(15)
where $`\delta `$, the broadening of the quasi-bound corral states, is roughly 20 meV .
An additional approximation can be done if any of the two following statements holds:
* The level spacing between the energies of the quasi-bound states is much bigger than $`\delta `$.
* The level spacing is lower than $`\delta `$ but, due to the geometry of the quantum corral, only a few of the bound wavefunctions take a non-negligible value at the magnetic impurity site, $`\stackrel{}{R}_I`$. If the energy separation of these states is bigger than $`\delta `$, then only one of these states will transmit the quantum mirage, as it is evident from equation (15). This condition is fulfilled in the case of the elliptic corral, as we will show below.
In any of these two situations, whenever there is a quasi-bound state which simultaneously has an energy near $`ϵ_F`$ and a non-negligible density at $`\stackrel{}{R}_I`$, we can replace equation (15) in (13) by:
$$𝒢_c(\stackrel{}{R},\stackrel{}{R}_I,ϵ^+)\frac{\psi _{ϵ_F}^{}(\stackrel{}{R})\psi _{ϵ_F}(\stackrel{}{R}_I)}{ϵϵ_F+i\delta }.$$
(16)
In next section we shall use the complete expression (15) for our calculations.
Our last approximation is to assume $`t_d<<t_c`$. A finite $`t_d`$ is considered in next section.
When we put together all these approximations, the change in $`G(V,\stackrel{}{R})`$ due to the presence of the impurity in $`\stackrel{}{R}_I`$ reads:
$`\delta G(V,\stackrel{}{R}){\displaystyle \frac{4e^2V_h^2t_c^2}{\pi ^2\mathrm{}}}\rho _T|\psi _{ϵ_F}(\stackrel{}{R})|^2|\psi _{ϵ_F}(\stackrel{}{R}_I)|^2\times `$ (17)
$`Im\left({\displaystyle \frac{1}{(eV+i\delta )^2}}{\displaystyle \frac{1}{eV+ik_BT_K}}\right).`$ (18)
For $`eV<<\delta `$ we can write:
$`\delta G(V,\stackrel{}{R})|\psi _{ϵ_F}(\stackrel{}{R})|^2|\psi _{ϵ_F}(\stackrel{}{R}_I)|^2{\displaystyle \frac{k_BT_K}{(eV)^2+(k_BT_K)^2}}.`$ (19)
Equation (19) is the most important result of this section. We want to highlight several points:
i) The spectral change in $`G(V)`$ is a dip of width $`k_BT_K`$ centered around $`eV=0`$, as observed in the experiments .
ii) According to equation (19), the dip is projected to any point $`\stackrel{}{R}`$ of the corral with an strength given by $`|\psi _{ϵ_F}(\stackrel{}{R})|^2|\psi _{ϵ_F}(\stackrel{}{R}_I)|^2`$. Therefore, the projection is magnified when both the impurity and the tip are at points where the Fermi level corral wavefunction peaks. The projection disappears when either the impurity or the observation point, are located in a minimum. As we show in the next section, for the eccentricity of the experiment the wavefunction of the elliptical corral at the Fermi level has its maxima close, but not at the foci. This result is in agreement with the experimental observation, but reduces the importance of the role played by the foci.
iii) The wavefunction at the Fermi level of any quantum corral has several maxima so that we predict that the mirage can be observed in other geometries. A possible candidate is the stadium corral shown in figure 3 of reference . Therefore, an elliptical geometry is not needed to observe the mirage.
To conclude this section, we compare equation (19), valid for a confined geometry, with the case of an impurity in a translationally invariant surface. In both cases the surface Green’s function, $`𝒢_s`$, is the sum of two contributions, the impurity free contribution, $`𝒢_c`$, and the scattering contribution (see equation (8)). The first accounts for the paths in which the electron does not interact with the impurity and the second accounts for the paths in which the electron does indeed interact with the impurity. Hence, the local density of states in any point of the surface contains information about the impurity.
In the case of the free surface (without corral), a continuum of quantum states with different $`\stackrel{}{k}`$ carries that information so that destructive interference takes place at distances of the order of $`2k_F^1`$, the inverse of the Fermi vector . In contrast, when the electrons interact with the corral atoms, the information is carried, essentially, by a few quantum states, so that the destructive interference is less efficient. Equation (19) is derived assuming that a single quantum state is carrying the information so that there is no interference at all.
## IV The mirage in the ellipse
In this section we study the mirage in an elliptical corral. Following the ideas of the previous section, we model the Green’s function of the surface states by that of the electrons confined in an hard wall elliptical corral. In order to compare with experiment , we consider the case in which the corral is built on a Cu(111) surface. We replace the real eigenvalues of the corral, $`ϵ_n`$, by $`ϵ_ni\delta `$, in order to model the inelastic processes, such as scattering to the bulk states. It turns out that the problem of a quantum particle confined in an ellipse can be solved analytically . To do that, we write the Schrödinger equation in elliptical coordinates:
$`x`$ $`=`$ $`aeCos[\theta ]Cosh[\eta ]`$ (20)
$`y`$ $`=`$ $`aeSin[\theta ]Sinh[\eta ],`$ (21)
where $`a`$ and $`e`$ are the semimajor axis and eccentricity, respectively. The Helmholtz equation in this coordinate system is separable, so that the eigenstates of the problem can be written as:
$$\psi (\theta ,\eta )=\mathrm{\Theta }(\theta )\mathrm{\Lambda }(\eta ).$$
(22)
The Schrödinger equation is written as
$`{\displaystyle \frac{d^2\mathrm{\Lambda }(\eta )}{d\eta ^2}}`$ $``$ $`(\alpha 2kCosh[2\eta ])\mathrm{\Lambda }(\eta )=0`$ (23)
$`{\displaystyle \frac{d^2\mathrm{\Theta }(\theta )}{d\theta ^2}}`$ $`+`$ $`(\alpha 2kCos[2\theta ])\mathrm{\Theta }(\theta )=0`$ (24)
$`k`$ $`=`$ $`{\displaystyle \frac{m^{}(ea)^2ϵ}{2\mathrm{}^2}},`$ (25)
where $`\alpha `$ is the separation constant, $`ϵ`$ is the particle energy and $`m^{}`$ is the electron effective mass which, in the Cu(111) surface band is $`0.38`$ $`m_e`$ . For a given $`k`$, only a discrete set of $`\alpha _r(k)`$ meet the requirement $`\mathrm{\Theta }(\theta )=\mathrm{\Theta }(\theta +2pi)`$. The elliptical hard wall condition reads $`\mathrm{\Lambda }(\eta _0)=0`$. It is clear that $`\eta =\eta _0`$ defines an ellipse of eccentricity $`e=(Cosh[\eta _0])^{(1)}`$. For each $`\alpha _r(k)`$ there is a discrete number of $`k_n`$ compatible with the hard wall boundary condition. With all this in mind, we find 2 types of physically possible solutions for the particle inside the hard wall ellipse:
$`\psi _{n,c}(\theta ,\eta )`$ $`=`$ $`ce_r(k_n^c,\theta )Ce_r(k_n^c,\eta )`$ (26)
$`\psi _{n,s}(\theta ,\eta )`$ $`=`$ $`se_r(k_n^s,\theta )Se_r(k_n^s,\eta ),`$ (27)
where $`ce`$, $`se`$, $`Se`$ and $`Ce`$ are the Mathieu functions . Of course, we have $`Se_r(k_n^s,\eta _0)=0`$ and $`Ce_r(k_n^c,\eta _0)=0`$. These equations permit to find the spectrum. In figure 1 we plot a part of the spectrum of an ellipse with $`e=0.5`$ and $`a=71.3\AA `$.
In figure 2 we plot the LDOS at the focus in the absence of a magnetic impurity. It is clear that only a few states of figure 1 contribute significantly to the LDOS at the focus. The energy separation between these levels is much larger than $`\delta =20meV`$. There is one of these quasi-bound wavefunctions that has an energy of 447.5 meV, very near $`ϵ_F`$ (which, for the Cu(111) surface band is 450 meV). We can thus explain the experimental observation of a quantum mirage in this quantum corral using the results of section III.
In the left panel of figure 3 we show a contour plot of the wavefunction at the Fermi level for this ellipse. It must be stressed that the Fermi wavefunction maxima are located at a distance of 3.28 $`\AA `$ of the closest focus. The lattice constant of Cu(111) is 2.55 $`\AA `$. Hence, experimentally it is very difficult to distinguish between the foci and the maxima.
The knowledge of the corral spectrum and wavefunctions permits to calculate $`G(V,\stackrel{}{R})`$ for the elliptical corral via equations (5), (6), (9), (13). In figure 3 (right panel), we plot the difference between the $`G(V,\stackrel{}{R})`$ map with and without the impurity, for $`eV=0`$ . In our calculations we take the value $`k_BT_K=4.6meV`$ ($`T_K=50K`$), as observed in . In this experiment $`T=4K`$, so that condition $`T<<T_K`$ is fulfilled. The change in the differential conductance occurs not only at the focus where the impurity is located but also at the empty focus, located 71 $`\AA `$ away from the impurity. The fingerprint of the Kondo effect is thus dominantly located around the impurity and around the empty focus. The similitude between the left and the right panel in figure 3 supports our claim that the wavefunction of the corral at the Fermi level projects the Kondo dip from the impurity to the other focus. Our figure 3 should be compared with figures 3-c and 3-d of reference . In the case of surface without corral, the Kondo signature would be localized around the impurity, being negligible at a distance larger than 2$`k_F^1`$ .
In figure 4 (left panel) we plot $`G(V)`$ when the tip is on top of the focus where the impurity is located (compare to figure 4(a) of reference ). For this calculation we use $`q0.2`$, the value used to fit $`G(V)`$ without corral , and consider a nonzero value for $`t_d`$, following the method outlined in section III. In the right panel of our figure 4 we plot $`G(V)`$ measured on top of the empty focus. Hence, our theory is in agreement with the main experimental result: the existence of a Kondo resonance at the empty focus, more than 80 $`\AA `$ away from the magnetic impurity. It must be stressed, however, that in our model the dip observed on top of the magnetic impurity and the one observed on top of the empty focus have different lineshapes, and there exists a factor of 2 between their intensities. In the experiment the attenuation factor is approximately $`8`$ and both the original dip and the ghost are more symmetric. In order to remove this discrepancy a less phenomenological theory for inelastic processes, like scattering from surface states to bulk states, would be necessary. We predict also that combinations of surface and adatoms for which inelastic scattering is smaller than for Co and Cu would increase the size of the mirage. In figure 5 we also show $`G(V)`$ when the tip is not at a maximum of the Fermi corral eigenstate. In those situations the mirage is not present, in agreement with the experiments .
In figure 6 we plot the intensity of the mirage (the dip amplitude) as a function of $`a`$, keeping $`e=0.5`$. In reference an oscillatory dependence of the mirage effect as a function of $`a`$ (for fixed $`e`$) was mentioned, the oscillation period being $`\lambda _F/4`$. Our calculation is consistent with that claim. However, we obtain a curve with more structure. The Fourier transform of the intensity of the mirage shows several peaks, the largest of which is located at $`\lambda _F/4`$, in agreement with the experiment. In figure 6 we also plot the number of occupied states inside the corral, as a function of $`a`$, and keeping $`ϵ_F`$ constant at 450 meV. We see that most of the changes in the occupation number do not lead to large changes in the mirage strength. The mirage is only enhanced when a particular kind of states, whose wavefunction is heavily peaked very close to the foci, is occupied. This rule was also observed in the experiments .
For the ellipse with $`e=0.786`$ in , we have been able to reproduce all the results obtained for the ellipse with $`e=0.5`$, assuming that the Fermi Level is somewhat below 450 meV. This indicates that the position of the resonances given by the hard wall ellipse might not coincide with the experimental results.
Since the maxima of the Fermi wavefunction are not exactly located at the foci, it is our contention that the important issue is to place the impurity at the maximum of the Fermi wavefunction. Therefore, geometrical or semiclassical interpretations of the mirage might not be adequate to address this phenomenon. To check this, we have studied a square corral, obtaining the mirage effect. Elliptical corrals are very convenient because some states with high quantum numbers (such as the state at the Fermi level for the $`e=0.5`$ ellipse with the adequate $`a`$) have two main maxima located close to the foci of the ellipse. In contrast, all the maxima in a square corral have the same height so that the projection effect is less pronounced than in the ellipse.
## V Discussion and conclusions
We now comment on some of the limitations of our theory. The first has to do with the approximation of the eigenstates inside the quantum corral as quasi-bound states broadened in energy. For a quantitative description of the corral energy spectrum a more detailed calculation is needed, taking into account the role of the corral atoms as tunneling centers to the bulk states . The second is the use of the Friedel sum rule in a resonant level model. A more realistic calculation of the impurity Green’s function would imply to take into account the real wavefunctions inside the corral and the possibility of tunneling from the magnetic impurity to the bulk states. The quantitative discrepancy with the experiments, in what concerns the attenuation of the mirage, should be solved including these effects. A more complete theory for the STM through magnetic impurities in metallic surfaces without quantum corrals has been developed in .
The emphasis of this paper is placed on the qualitative understanding of the mirage rather than on a detailed description of the experiments. Our main results are the following: i) The LDOS evaluated at an arbitrary surface point, $`\stackrel{}{R}`$, in the Anderson model, contains information about the LDOS at the impurity site, $`\stackrel{}{R}_I`$. A mirage will appear in a remote point, $`\stackrel{}{R}`$, whenever there is a single quantum state at the Fermi level whose amplitude $`\psi _{ϵ_F}`$ peaks both at the impurity ($`\stackrel{}{R}_I`$) and at $`\stackrel{}{R}`$. In order to avoid destructive interference between different states it is necessary that the energy spacing between states with a non-negligible amplitude at the impurity site is bigger than the energy broadening, $`\delta `$. ii) The mirage can be obtained in corrals with shapes other than elliptical. However, the elliptical shape is quite convenient because some of the corral eigenstates peak strongly at two points very close to the foci. iii) Our theory predicts that the intensity of the mirage in an elliptical corral oscillates, as a function of the semimajor axis length, keeping fixed the eccentricity, with a dominant period of $`\lambda _F/4`$, in agreement with reference .
We want to acknowledge C. Piermarocchi and H. Manoharan for fruitful discussions. JFR and DP acknowledge Spanish Ministry of Education (MEC) for postdoctoral research fellowship and grant FPU/AP98. Work supported in part by MEC under contract PB96-0085.
During the completion of this manuscript we became aware of a theoretical work addressing the problem of the mirage in an elliptical quantum corral . In that work the states of the ellipse are described by a more detailed method, assuming that the wall atoms are magnetic, and the issue of the existence of the mirage in different geometries is not addressed. Our theory can be applied to the general case of a quantum corral formed by non-magnetic scatterers (in which quantum mirages have also been observed ).
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# Critical behavior and conservation in directed sandpiles
## I Introduction
Sandpile cellular automata are the most famous example of self-organized critical (SOC) behavior . Under an external drive consisting of a slow addition of sand (energy) grains and the action of dissipation through the loss of energy on the lattice boundaries, these models reach a stationary steady state. In the limit of infinitesimal driving and dissipation (this last achieved in the thermodynamic limit), the stationary state of sandpile models exhibits diverging response functions associated to a characteristic avalanche dynamics. This is the hallmark of a critical behavior that has attracted an enormous amount of interest as a plausible explanation of the avalanche-like critical behavior empirically observed in many natural systems .
Sandpile models have been at the center of an intense research activity made of both analytical studies and numerical simulations. Despite the simple definition of these automata, it turns out that their full analytical understanding is a very problematic task . As a further complication, also the numerical inspection of these models results to be particularly difficult. For example, the precise identification of universality classes has resisted for many years even the most careful numerical analysis, and only recent results have partially settled this problem . On the other hand, these refined analyses have pointed out that several sandpile models do not follow the simple finite size scaling (FSS) form usually adopted in the description of critical behavior . For instance, the more sophisticated multiscaling approach seems to be required for a full description of the scaling properties of the original Bak, Tang, and Wiesenfeld (BTW) model .
Many sandpile features have been underlined as the possible origin of these scaling anomalies. The deterministic dynamical rules of the BTW model induce nonergodic effects , that are certainly missing in stochastic models, such as the Manna model , which shows a perfect FSS behavior, even for moderate system sizes. A further complication of sandpile automata stems from the peculiar role of the boundary dissipation, that makes the lattice size scaling entangled with the system dynamics. In such cases, the thermodynamic limit is essential for the dissipative dynamics of large avalanches. A clear understanding of the interplay between dissipation and size scaling has not yet been achieved and it has been recently the subject of several studies .
In this paper we address some of the aforementioned problems in the case of directed sandpile models . In this case Dhar and Ramaswamy obtained an exact solution for the Abelian deterministic directed sandpile (DDS), that can be used as a milestone to check the numerical simulation analysis. Directed sandpiles thus become an interesting test field to study how the critical behavior is affected by the introduction of stochastic elements and dissipation. We perform large scale numerical simulations of two directed sandpile automata: the deterministic directed sandpile model and the stochastic directed sandpile model . We study both models in the case of boundary and bulk dissipation . We find, in agreement with the results in Ref. , that the models define two different universality classes. In addition we show that the universality class of the models does not depend on the way in which dissipation is implemented. Finally we analyze the properties of anisotropic models in which the dynamics is not fully directed . In this case we observe that on large scales the critical behavior is the same of that of fully directed models. Results for the stochastic model are compared with a recent theoretical approach by Paczuski and Bassler, that provides values for the critical exponents in perfect agreement with numerical simulations. These results are also recovered in Ref. .
The numerical analysis also points out that some critical exponent values, such as the correlation length exponents or the affinity exponent (to be defined later on), are independent of the particular universality classes and common to all models considered. In order to explain this numerical evidence, we provide a phenomenological characterization of directed sandpiles based on the basic symmetries introduced by the conserved dynamics of these automata. Following balance of energy arguments inspired in Refs. , we derive a series of results and predictions on the value of critical exponents which are a straightforward consequence of conservation. These general results can be considered as super-universal, because they characterize the critical behavior of all directed sandpiles with local dynamical rules, independently on the specific universality class. The results presented here provide a general picture of directed models and the role of boundary and bulk dissipation in the process of self-organization.
The paper is arranged as follows. In Sec. II we introduce and define the various directed models considered. Secs. III and IV present and discuss, from the standpoint of universality, the numerical results for directed models with boundary and bulk dissipation. In Sec. V we introduce anisotropic models, and present the numerical results obtained, in comparison with those of directed models. Sec. VI is devoted to an analytical approach based on the conservation of energy. Finally, in Sec. VII we draw our conclusions and perspectives.
## II Directed Models
Sandpile models are usually defined on a $`d`$-dimensional hyper-cubic lattice of size $`L`$. To each node of the lattice is assigned an integer variable $`z_i`$, called “energy”. Energy is added to the system uniformly at randomly chosen sites ($`z_iz_i+1`$). When a site becomes active, that is, when its energy becomes larger than or equal to a certain threshold $`z_c`$, it topples. A toppling site loses an energy $`z_c`$, that is distributed among its neighbors according to a certain set of rules. The neighbors that receive energy can become active and topple on their turn, thus generating an avalanche. The slow driving condition is effectively imposed by stopping the random energy addition during the avalanche spreading. This means that the driving time scale is infinitely large with respect to the toppling characteristic time scale.
The models we consider in this Section are directed, in the sense that the energy is always transported along a preferred fixed direction. We denote this preferred direction by the coordinate $`x_{}`$, whose positive direction is usually defined as “downwards”. The transverse direction (subspace of dimension $`d1`$ perpendicular to $`x_{}`$) will be denoted by $`\stackrel{}{x}_{}`$.
The toppling rules of the models define two main classes:
(i) Deterministic directed sandpile (DDS): In $`d`$ dimensions, the threshold is set to $`z_c=2d+1`$. When a site in a given hyper-plane $`x_{}`$ topples, it sends deterministically one grain of energy to each one of its nearest and next-nearest neighbors on the hyper-plane $`x_{}+1`$ (see Fig. 1a)). Our definition is somewhat different from the original model of Dhar and Ramaswamy , in both the driving and the orientation of the lattice. Both models, however, are expected to share the same universality class, being deterministic and directed. Numerical simulations confirm indeed this point.
(ii) Stochastic directed sandpile (SDS): In this case, the threshold is $`z_c=2`$, independently of the dimensionality of the lattice. When a site in the hyper-plane $`x_{}`$ topples, it sends two grains of energy to two sites, randomly chosen among its $`2d1`$ nearest and next-nearest neighbors on the hyper-plane $`x_{}+1`$. The toppling rules of this model can be defined exclusive if the two energy grains are always distributed on different sites, Fig. 1b). On the other hand, the model can be defined non-exclusive if the dynamics allows the transfer of two energy grains onto the same site, Fig. 1c). We therefore report simulations on the exclusive stochastic directed sandpile (ESDS) and on the non-exclusive stochastic directed sandpile (NESDS). In spite of the stochastic nature of these models, we must bear in mind that they are nevertheless Abelian . The discussion therefore focuses on the difference between stochastic and deterministic models.
Once the toppling rules have been determined, the models are finally defined by specifying the dissipation mechanism. For systems with boundary dissipation, we impose periodic boundary conditions in the transverse directions $`\stackrel{}{x}_{}`$ and open at the hyper-plane $`x_{}=L`$. In this way, the models are locally conserved; energy can only leave the system at the bottom of the lattice. In models with bulk dissipation, we impose periodic boundary conditions in both the $`x_{}`$ and $`\stackrel{}{x}_{}`$ directions. Dissipation is implemented by allowing a toppling site to loose an energy $`z_c`$ without transferring it with probability $`p`$. This means that, on average, an energy $`ϵ=z_cp`$ is dissipated in each toppling. In the limit $`ϵ0`$, the system shows critical behavior .
In the stationary state we can define the probability that the addition of a single energy grain is followed by an avalanche of toppling events. Avalanches are then characterized by the total number of topplings $`s`$ and the time duration $`t`$. In the limit of infinitesimal driving (slow driving condition) the system shows scaling behavior and the probability distributions of these quantities follow the finite-size scaling (FSS) forms:
$`P(s)`$ $`=`$ $`s^{\tau _s}𝒢(s/s_c),`$ (1)
$`P(t)`$ $`=`$ $`t^{\tau _t}(t/t_c),`$ (2)
where $`s_c`$ and $`t_c`$ are the characteristic size and time, respectively. The exponents $`\tau _s`$ and $`\tau _t`$ characterize the critical behavior and define the universality classes to which the models belong. In the critical region the characteristic time and size are determined only by the system size $`L`$ or the dissipation $`ϵ`$, in the case of boundary and bulk dissipation, respectively. In directed models, the affinity exponent $`\zeta `$ is of particular importance; it relates the avalanche characteristic lengths in the perpendicular direction, $`\xi _{}`$, and in the parallel direction, $`\xi _{}`$, through the relation $`\xi _{}\xi _{}^\zeta `$. This exponent characterizes the degree of anisotropy due to the preferential direction present in the transport of the energy. In other words, it expresses the self-affine properties in the scaling of avalanches. A general result concerns the average avalanche size $`s`$, that also scales linearly with $`L`$ : a new injected grain of energy has to travel, on average, a distance of order $`L`$ before reaching the boundary. In the stationary state, to each energy grain drop must correspond, on average, an energy grain flowing out of the system. This implies that the average avalanche size corresponds to the average number of topplings needed for a grain to reach the boundary; i.e., $`sL`$. The same result can be exactly obtained by inspecting the conservation symmetry of the model as we shall see in Sec. VI.
For the DDS, the exact analytical solution in $`d=2`$ yields the exponents $`\tau _s=4/3`$ and $`\tau _t=D=3/2`$ . The upper critical dimension is found to be $`d_c=3`$, and it is also possible to find exactly the logarithmic corrections to scaling . The introduction of stochastic ingredients in the toppling dynamics of directed sandpiles has been studied recently in a model that randomly stores energy on each toppling . This model is strictly related to directed percolation and defines a universality class “per se”. In our case stochasticity affects only the partition of energy during topplings, and there is a priori no obvious relations between the critical behavior of these models.
## III Numerical simulations with boundary dissipation
In this Section we report results from computer simulations of deterministic and stochastic directed sandpiles, performed with boundary dissipation. The system sizes considered range from $`L=100`$ to $`L=6400`$. The statistical distribution functions have been computed averaging over $`10^7`$ nonzero avalanches.
In the case of boundary dissipation, the lattice size $`L`$ is the only characteristic length present in the system. Approaching the thermodynamic limit ($`L\mathrm{}`$), the avalanche characteristic size and time in Eqs. 1 and 2 diverge as $`s_cL^D`$ and $`t_cL^z`$, respectively. The exponent $`D`$ defines the fractal dimension of the avalanche cluster and $`z`$ is the usual dynamic critical exponent. The directed nature of the model introduces a drastic simplification, since it imposes $`z=1`$. In order to compute the different exponents characterizing the dynamics of the avalanches, we have performed the moment analysis of the distributions, in analogy to the method developed by De Menech et al. . We define the $`q`$-th moment of the avalanche size distribution on a lattice of size $`L`$ as $`s^q_L=𝑑ss^qP(s)`$. If the FSS hypothesis (1) is valid in the asymptotic limit of large $`s`$, then the $`q`$-th moment has the following dependence on system size:
$$s^q_L=L^{D(q+1\tau _s)}𝑑yy^{(q\tau _s)}𝒢(y)L^{\sigma _s(q)}.$$
(3)
The exponent $`\sigma _s(q)=D(q+1\tau _s)`$ is computed as the slope of the log-log plot of $`s^q_L`$ as a function of $`L`$. For large enough values of $`q`$ (i.e., away from the region where the integral in (3) is dominated by its lower cut-off), one can compute the fractal dimension $`D`$ as the slope of $`\sigma _s(q)`$ as a function of $`q`$: $`D=\sigma _s(q)/q`$. On the other hand, as we have argued in the previous Section, the first moment must scale linearly with $`L`$, which imposes $`\sigma _s(1)=1`$. Once $`D`$ is known we can estimate $`\tau _s`$ using the relation $`\sigma _s(1)=D(2\tau _s)`$
Along the same lines we can obtain the moments of the avalanche time distribution. In this case, $`t^q_LL^{\sigma _t(q)}`$, with $`\sigma _t(q)/q=z`$. Analogous considerations for small $`q`$ apply also for the time moment analysis. Here, an estimate of the asymptotic convergence of the numerical results is provided by the constraint $`z=1`$, that must hold for large enough sizes. Then, the $`\tau _t`$ exponent can be found using the scaling relation $`(2\tau _t)=\sigma _t(1)`$.
Once the exponents have been estimated numerically, we can check the accuracy of the moment analysis’ predictions using the FSS hypothesis. If the FSS hypothesis of Eq.s (1,2) is correct, then the plots of the distributions, under the rescaling $`ss/L^D`$ and $`P(s)P(s)L^{D\tau _s}`$ and correspondingly $`tt/L^z`$ and $`P(t)P(t)L^{z\tau _t}`$, should collapse onto the same universal function, for different values of $`L`$.
In Table I we report the exponents found for the DDS, ESDS, and NESDS models in $`d=2`$. Figure 2 shows the moments $`\sigma _s(q)`$ and $`\sigma _t(q)`$. Figures 3 and 4 plot the FSS data collapse for sizes and times, respectively. The exponents obtained for the DDS are in perfect agreement with the expected analytical results. This fact supports the idea that the system sizes used in the present work allow to recover the correct asymptotic behavior. Results for the ESDS and NESDS are identical within the error bars, pointing out that these two models are in the same universality class. On the other hand, the obtained exponents prove beyond any doubts that deterministic and stochastic directed sandpile models do not belong to the same universality class.
We have also directly computed the characteristic lengths in the parallel and transversal directions, $`\xi _{}`$ and $`\xi _{}`$, as a function of the system size. The anisotropy of the system is reflected in the different definitions of both characteristic lengths. In this sense, we define them with the same spirit as in directed percolation .
Consider a given avalanche, labeled $`\alpha `$, that has started at the site $`(x_{}^{(0)},\stackrel{}{x}_{}^{(0)})`$, and has affected the set of different sites $`\{(x_{}^{(i)},\stackrel{}{x}_{}^{(i)})\}`$, for $`i=0\mathrm{}a1`$ (i.e., it has covered an area $`a`$). Let us define the quantities
$$R_{}(\alpha )=\frac{1}{a}\underset{i=1}{\overset{a1}{}}|x_{}^{(0)}x_{}^{(i)}|$$
(4)
and
$$R_{}^2(\alpha )=\frac{1}{a}\underset{i=1}{\overset{a1}{}}(\stackrel{}{x}_{}^{(0)}\stackrel{}{x}_{}^{(i)})^2.$$
(5)
Furthermore, let us define $`R_{}(a)`$ and $`R_{}^2(a)`$ as the averages of the previous quantities, over all avalanches of the same fixed area $`a`$. Let be $`P(a)`$ the probability of observing an avalanche of area $`a`$. We define the correlation lengths by
$$\xi _{}=\frac{_aR_{}(a)aP(a)}{_aaP(a)},\xi _{}^2=\frac{_aR_{}^2(a)aP(a)}{_aaP(a)}.$$
(6)
The different definitions (4) and (5) are obviously due to the different nature of the avalanche spreading in the directions $`x_{}`$ and $`x_{}`$. In the former case, the spreading is isotropic, and thus the second moment of the relative distance distribution is needed to define a meaningful correlation length. In the latter case, on the other hand, the spreading is always in the direction of growing $`x_{}`$, and therefore the first moment is sufficient.
The system being critical, both correlation lengths should scale with the system size, defining the exponents $`\nu _{}`$ and $`\nu _{}`$ by
$$\xi _{}L^\nu _{},\xi _{}L^\nu _{}.$$
(7)
The affinity exponent, defined by
$$\xi _{}\xi _{}^\zeta $$
(8)
is thus given by $`\zeta =\nu _{}/\nu _{}`$.
We have calculated the correlations lengths in the models DDS, ESDS, and NESDS, given by the definition (6). The results, plotted in Fig. 5, give the following dependence of the correlation lengths with system size for all models:
$$\xi _{}L,\xi _{}L^{1/2}.$$
(9)
These relations define the exponents $`\nu _{}=1`$ and $`\nu _{}=1/2`$, and an affinity exponent $`\zeta =1/2`$. It is interesting to note that this exponent is independent of the universality class of the model, defining a sort of super-universal property of directed models.
As pointed out in Ref., the stochastic dynamics of SDS models introduces multiple toppling events on the same site, which are by definition absent in the deterministic case. This gives rise to very a different avalanche structure, eventually reflected in the different asymptotic critical behavior. It is worth remarking that the universality class of SDS appears robust to modifications of the stochastic microscopic dynamics as pointed out in Ref., where it is shown that modifications of SDS models with stochastic toppling threshold still belong to the same universality class. Recently, Paczuski and Bassler , have proposed a theoretical approach that allows the calculation of critical exponents in directed models with multiple topplings. The analysis goes through the mapping of the avalanche evolution into the dynamics of an interface moving in a random medium, as also proposed in . This theoretical result gives the exponents $`\tau _s=10/7`$ and $`\tau _t=D=7/4`$, in perfect agreement with the values obtained by numerical simulations, Table I. The same exponent values are also found in the approach of Ref. .
## IV Numerical simulations with bulk dissipation
In this Section we report results from computer simulations of deterministic and stochastic sandpiles, performed with bulk dissipation. In this case, dissipation is implemented as described in Sec. II. That is, in a system with periodic boundary conditions, each toppling site has a probability $`ϵ/z_c`$ of losing an energy $`z_c`$, and a probability $`1ϵ/z_c`$ of transferring it to its neighbors. The dissipation rates range from $`ϵ=0.0016`$ to $`0.0512`$, and the (fixed) system size considered is $`L=6400`$. Statistical distribution functions have been computed averaging over $`10^7`$ nonzero avalanches.
In the presence of bulk dissipation the characteristic sizes are determined by the dissipation rate $`ϵ`$, which defines the only characteristic length in the system. Approaching the limit $`ϵ0`$, the avalanche characteristic size and time diverge as $`s_cϵ^{\mathrm{\Delta }_s}`$ and $`t_cϵ^{\mathrm{\Delta }_t}`$, respectively. It is also very easy to relate the mean avalanche size to the dissipation rate $`ϵ`$. On average, each added grain must be dissipated in the evolution of the avalanche, resulting in $`ϵs=1`$. This readily yields $`s=ϵ^1`$. In this case it is extremely important that the characteristic length of the avalanche $`\xi _{}`$ is always smaller than the size of the lattice used. This allows us to study only finite size effects introduced by the dissipation probability, without spurious effects due to the finite lattice size.
The moment analysis can be straightforwardly generalized to systems with bulk dissipation. In this case the role of the system size $`L`$ as scaling parameter is played by the dissipation $`ϵ`$. If the FSS hypothesis holds, the $`q`$-th moment for, say the size distribution, has an explicit dependence on the dissipation rate that reads:
$$s^q_ϵϵ^{\mathrm{\Delta }_s(q+1\tau _s)}=ϵ^{\rho _s(q)}.$$
(10)
The new moment $`\rho _s(q)=\mathrm{\Delta }_s(q+1\tau _s)`$ can be estimated by linear regression in a log-log plot of $`s^q_ϵ`$ as a function of $`ϵ^1`$. Once this moment is computed, the exponent $`\mathrm{\Delta }_s`$ is given by $`\mathrm{\Delta }_s=\rho _s(q)/q`$. The relation $`s=ϵ^1`$ imposes $`\rho _s(1)=1`$, and from here, once known $`\mathrm{\Delta }_s`$, we compute $`\tau _s`$ using the relation $`\rho _s(1)=\mathrm{\Delta }_s(2\tau _s)`$. Analogous considerations allow to compute the exponents of the time distribution, $`\mathrm{\Delta }_t`$ and $`\tau _t`$. Finally, to check the exponents with the data collapse technique, one must plot the rescaled functions $`P(s)ϵ^{\mathrm{\Delta }_s\tau _s}`$ as a function of $`s/ϵ^{\mathrm{\Delta }_s}`$ and $`P(t)ϵ^{\mathrm{\Delta }_t\tau _t}`$ as a function of $`t/ϵ^{\mathrm{\Delta }_t}`$, respectively.
In Table II we report the exponents computed in $`d=2`$ for the directed models DDS, ESDS, and NESDS with bulk dissipation. The corresponding moments $`\rho _s(q)`$ and $`\rho _t(q)`$ are shown in Figures 6, while Figs. 7 and 8 plot the data collapse for sizes and times, respectively.
To conclude our analysis of directed sandpiles with bulk dissipation, we have proceeded to compute the correlation length of the models. In this case, the scaling of the correlation lengths with vanishing dissipation define the scaling exponents
$$\xi _{}ϵ^\nu _{}^{},\xi _{}ϵ^\nu _{}^{}.$$
(11)
and an affinity exponent $`\zeta =\nu _{}^{}/\nu _{}^{}`$. Using an analogous definition as in the case of boundary dissipation, we compute the exponents $`\nu _{}^{}=1`$, $`\nu _{}^{}=1/2`$, and $`\zeta =1/2`$, as shown in Fig. 9. That is, the correlation length exponents are identical for both boundary and bulk dissipation. These results again imply an affinity exponent $`\zeta =1/2`$ in all the models studied so far.
In view of these results, we have confirmation that the stochastic models belong to a different universality class than the deterministic directed sandpile. These results also point out in a very clear way that the critical behavior of models with boundary or bulk dissipation is identical. In fact, all critical exponents $`\tau _s,\tau _t,z,`$ and $`\zeta `$ are identical in both cases. This further confirms the complete equivalence of both points of view with respect to sandpiles and shows that, at least in the directed case, the open boundary conditions usually implemented in simulations do not affect the scaling behavior in a peculiar way. Of course, the open boundary conditions breaks the translational invariance of the system, but in the thermodynamic limit this effect is negligible for the asymptotic critical behavior. Finally, these results validate theoretical approaches in which it is assumed a homogeneous dissipation that is much easier to treat analytically.
As a last observation it is worth remarking that also in this case, a series of exponents such as $`\zeta `$ and $`\nu _{}^{}`$ assume values independently of the universality class of the model under study. This sort of super-universality can be explained in terms of energy conservation as we shall see in Sec. VI.
## V Numerical simulations of anisotropic models
An important question to study in directed sandpile models is the effect on the scaling properties of any amount of diffusion along the preferred direction of transport $`x_{}`$. One would expect that the broken symmetry introduced by the preferential direction should prevail on large scales, so that the dynamical scaling in directed and simply anisotropic sandpiles become indistinguishable in the thermodynamic limit. This fact hints towards the possibility of a unique universality class for both directed and anisotropic sandpiles. This universality class is determined uniquely by the lack of symmetry along the $`x_{}`$ direction, and the presence or absence of stochastic elements in the definition of the models.
In order to test this conjecture, we have performed numerical simulations of an anisotropic stochastic sandpile model, defined according to the following rules: on a hyper-cubic lattice of size $`L`$, we consider a model with threshold $`z_c=2`$. When a site topples, it sends two grains of energy to two sites, randomly selected among the $`2d+1`$ nearest and next-nearest neighbors on the hyper-plane $`x_{}+1`$, and the nearest neighbor on the hyper-plane $`x_{}1`$, see Fig. 10. The rules in this model are defined non-exclusive, in such a way that the same site can receive the two sand grain expelled by an active site. The model is clearly anisotropic, because the probability to transfer energy in the downwards direction is three times larger that in the upwards direction. It would thus correspond to a non-exclusive stochastic anisotropic sandpile (NESAS). We consider only the case of boundary dissipation, performing simulations for sizes ranging from $`L=100`$ up to $`6400`$, and averaging over $`10^7`$ nonzero avalanches.
In Fig. 11 we plot the correlation length $`\xi _{}`$ and $`\xi _{}`$, measured according to the rules given in Eqs. (6). We confirm the expectation that anisotropic models have the same scaling properties, as regards the scaling of the correlation lengths, as directed models with the same deterministic or stochastic ingredients. We have also measured the exponents $`\tau _s,\tau _t,D`$, and $`z`$ for this model, using the moment analysis technique. The values found are $`\tau _s=1.43(1)`$, $`D=1.75(1)`$, $`\tau _t=1.72(2)`$, $`z=0.98(2)`$. These results, compared with Tables I and II, show that this anisotropic models belongs to the same universality class of the ESDS and NESDS directed models, confirming the irrelevance of the diffusion along the preferred direction $`x_{}`$.
## VI The role of conservation in sandpile models
We have seen in the preceding Sections that a subset of critical exponents characterizing the critical behavior of directed and anisotropic models have an interesting super-universal property; i.e. they are independent of the universality class of the models. In order to understand this feature we perform a theoretical analysis based on the conservation of energy, that is the basic symmetry in standard sandpile automata. We shall see in the following that the super-universal character of some critical exponents is dictated by simple energy conservation considerations. The use of this approach allows also to establish a relation between boundary and bulk dissipation models by introducing an effective dissipation that depends on the system size.
The avalanche dynamics in sandpile models is implicitly due to the imposed infinite time scale separation between driving and dissipation . In order to devise a theory that can take into account the symmetry introduced by the energy conservation, one must first regularize the rules of the models in such a way that a single time scale is ruling the dynamics. One way to do so is to introduce a nonzero driving rate, defined as the probability per unit time $`h`$ of a site to receive a grain of energy . This driving rate plays the role of an external field and leads to the SOC behavior in the limit $`h0^+`$. On the other hand, given that the toppling rules are conserved, energy can leave the system only at the boundaries. Boundary dissipation is a natural choice in computer simulations. However, it introduces undesirable complications due to its singular character in a local theory. It is therefore convenient to use an homogeneous effective dissipation $`ϵ`$, defined as the average energy lost in each toppling event. As observed in previous Sections, one can define models with periodic boundary conditions and built-in bulk dissipation. When constructing the local theory for models with open boundary conditions, the bulk dissipation $`ϵ`$ amounts to an effective parameter that is to be related to the system size $`L`$.
With all these ingredients, we are ready to formulate conservation of energy as a continuous equation. In sandpiles, we define the order parameter $`\rho _a`$ as the density of active sites (i.e., whose height $`zz_c`$). The only dynamics in the model is obviously due to the field $`\rho _a(\stackrel{}{x},t)`$, which is coupled to the local energy density, $`E(\stackrel{}{x},t)`$ (i.e. the local density of sandgrains), which enhances or suppresses the generation of new active sites. A Langevin description for sandpile automata is possible by considering the dynamics of the local order-parameter field $`\rho _a(\stackrel{}{x},t)`$ in a coarse-grained picture, bearing in mind that the energy density $`E(\stackrel{}{x},t)`$ is a conserved field. In Refs. , in analogy with absorbing-state phase transitions , a pair of coupled dynamical equations for the fields $`\rho _a(\stackrel{}{x},t)`$ and $`E(\stackrel{}{x},t)`$ were proposed. In the following we elucidate the consequences of energy conservation and we focus only on the latter equation. The interested reader can find the full set of equations in Ref. . In the next subsections we shall consider separately directed and anisotropic models.
### A Directed sandpiles
We seek a continuous equation for the coarse-grained local density of energy, $`E(\stackrel{}{x},t)`$. In the limit of zero driving and dissipation, energy is conserved. Therefore, the evolution equation fulfilled by the local field $`E`$ is:
$$\frac{E(\stackrel{}{x},t)}{t}=\stackrel{}{}\stackrel{}{J}_Eϵ\rho _a(\stackrel{}{x},t)+h(\stackrel{}{x},t)+\eta _E(\stackrel{}{x},t).$$
(12)
The first term simply represents the diffusion of energy; the second term accounts for the dissipation that is associated with every toppling event; the third term represents the external driving. Finally, the last term is a source of stochastic noise, that accounts for the randomness in the flow of energy. The noise term can be generated by the toppling rules in a stochastic model, or by the initial conditions plus the random driving in a deterministic model. We will require the noise to have zero average:
$$\eta _E(\stackrel{}{x},t)=0.$$
(13)
The noise correlator $`\eta _E(\stackrel{}{x},t)\eta _E(\stackrel{}{x}^{},t^{}`$ is of fundamental importance for the determination of universality classes and the critical behavior of the order parameter. However, for our present purposes we do not need precise knowledge of its analytical form (for a detailed discussion, see Refs. ).
The current can be constructed by appealing to the symmetries of the model. The transport of energy is due to topplings. These are isotropic along the transversal direction $`\stackrel{}{x}_{}`$, therefore the current along this direction will be proportional to the gradient of the density of active sites. In the preferred direction, on the other hand, all the energy is transferred downwards; therefore, the current in this direction must be proportional to the density of active sites. The final form of the current is then
$$\stackrel{}{J}_E(\stackrel{}{x},t)=D_{}\stackrel{}{}_{}\rho _a(\stackrel{}{x},t)+2\lambda \rho _a(\stackrel{}{x},t)\stackrel{}{e}_{}.$$
(14)
Plugging this expression into the equation for the energy, we have the final result:
$$\frac{E(\stackrel{}{x},t)}{t}=D_{}_{}^2\rho _a(\stackrel{}{x},t)2\lambda _{}\rho _a(\stackrel{}{x},t)ϵ\rho _a(\stackrel{}{x},t)+h(\stackrel{}{x},t)+\eta _E(\stackrel{}{x},t),$$
(15)
where the symbol $`_{}`$ stands for the partial derivative $`/x_{}`$. This is the general conservation equation for any directed sandpile model. It is worth remarking at this point that the energy field is a static field, in the sense that energy diffuses only if active sites are present in the system. This is intuitively understood in sandpile models, where energy (sand) grains diffuse only from toppling sites.
To analyze the consequences of Eq. (15), it proves useful to define the susceptibility $`\chi (\stackrel{}{x},t)`$ :
$$\chi (\stackrel{}{x}\stackrel{}{x}^{},tt^{})=\frac{\delta \rho _a(\stackrel{}{x},t)}{\delta h(\stackrel{}{x}^{},t^{})}_\eta ,$$
(16)
where the symbol $`_\eta `$ denotes an average over the noise distribution. By definition, the susceptibility measures the average increase in the number of active sites due to an impulsive perturbation, that is, to the addition of a single energy grain. Since we measure the size of the avalanches by the total number of topplings, the average avalanche size is given by
$$s=d^dx𝑑t\chi (\stackrel{}{x},t).$$
(17)
Taking the functional derivative of Eq. (15) and averaging over time and noise, we obtain, in the limit $`t\mathrm{}`$, in which the sandpile is in a stationary state with constant average energy, the following equation for the static susceptibility:
$$D_{}_{}^2\chi (\stackrel{}{x})2\lambda _{}\chi (\stackrel{}{x})ϵ\chi (\stackrel{}{x})=\delta ^{(d)}(\stackrel{}{x}).$$
(18)
This equation can be easily solved in Fourier space. Defining the transformation
$$\chi (x_{},\stackrel{}{x}_{})=\frac{1}{(2\pi )^d}d^{d1}k𝑑q\chi (q,\stackrel{}{k})e^{i\stackrel{}{k}\stackrel{}{x}_{}}e^{iqx_{}}$$
(19)
and substituting into Eq. (18), we obtain the solution
$$\chi (q,\stackrel{}{k})=\frac{1}{D_{}k^2+2i\lambda q+ϵ},$$
(20)
which yields the susceptibility in real space
$$\chi (x_{},\stackrel{}{x}_{})=\frac{1}{(2\pi )^d}d^{d1}ke^{i\stackrel{}{k}\stackrel{}{x}_{}}_{\mathrm{}}^{\mathrm{}}𝑑q\frac{e^{iqx_{}}}{D_{}k^2+2i\lambda q+ϵ}.$$
(21)
The integral in $`q`$ is easily performed by residues. The remaining $`d1`$ integrals in $`\stackrel{}{k}`$ become then decoupled Gaussian integrals that yield the result, setting $`D_{}=1`$:
$$\chi (x_{},\stackrel{}{x}_{})=\frac{1}{2\lambda }\left(\frac{\lambda }{2\pi }\right)^{(d1)/2}x_{}^{(1d)/2}e^{x_{}ϵ/2\lambda }e^{\lambda x_{}^2/2x_{}}.$$
(22)
Eq. (22) can be conveniently rewritten into the scaling form
$$\chi (x_{},\stackrel{}{x}_{})=x_{}^{(1d)/2}\mathrm{\Gamma }(\frac{x_{}}{\xi _{}},\frac{x_{}}{\xi _{}}),$$
(23)
where $`\mathrm{\Gamma }`$ is a cut-off function that decreases exponentially in both its arguments. Comparing this last expression with Eq. (22), we can identify the parallel and transversal correlation lengths:
$$\xi _{}ϵ^1,\xi _{}ϵ^{1/2}.$$
(24)
In more general terms, if we define the exponents $`\nu _{}^{}`$ and $`\nu _{}^{}`$ by Eqs. (11), then we have for directed sandpiles $`\nu _{}^{}=1`$ and $`\nu _{}^{}=1/2`$. From these last expressions, we can read off a first exact result for directed sandpiles: the avalanches produced in those models are elongated, with characteristic length in the parallel and transversal directions related by an affinity exponent $`\zeta =1/2`$. It is very important to stress that these results are independent of the particular model consider and of the dimensionality $`d`$ of the system, dictated only by the energy balance in the stationary state.
We can use the result (24) to relate the effective bulk dissipation with the system size in a model with open boundary conditions. To sustain a steady state with constant average energy, avalanches must reach the bottom boundary in order to be able to dissipate. This means that the characteristic length of the avalanches in the parallel direction must be proportional to the system size $`\xi _{}L`$. We have therefore that in boundary dissipation models we can define an effective dissipation rate $`ϵ`$ that is related with the system size by:
$$ϵL^1.$$
(25)
From this relations we easily find that $`\mathrm{\Delta }_s=D`$ and $`\mathrm{\Delta }_t=z`$. These identities are recovered in numerical simulations (see Tables I and II). Finally, from Eq. (20), we can recover the well-known result linking the system size and the average avalanche size, $`s=\chi (q=0,\stackrel{}{k}=0)ϵ^1L`$ . That is, for any directed model, the average avalanche size must scale as the inverse of the dissipation, in the case of model with bulk dissipation, or as the size of the system along the preferred direction of transport for sandpiles defined with boundary dissipation.
### B Anisotropic sandpiles
Having completed the analysis of directed sandpiles, we turn our attention to the more complex case of anisotropic sandpiles. In this kind of model, the transport of energy is not strictly directed in the parallel direction, but is simply stronger in the direction $`+x_{}`$ than in the opposite direction $`x_{}`$. The presence of backwards flow allows the possibility of diffusion in the preferred direction, and thus the equation for the conservation of energy becomes in this case
$$\frac{E(\stackrel{}{x},t)}{t}=D_{}_{}^2\rho _a(\stackrel{}{x},t)+D_{}_{}^2\rho _a(\stackrel{}{x},t)2\lambda _{}\rho _a(\stackrel{}{x},t)ϵ\rho _a(\stackrel{}{x},t)+h(\stackrel{}{x},t)+\eta _E(\stackrel{}{x},t).$$
(26)
Even though it is straightforward to obtain Eq. (26) by symmetry arguments, it is instructive to derive it also starting from the anisotropy of the microscopic rules. To this end, let us consider a one dimensional model in which active sites, when toppling, send a fraction of energy $`p`$ to the left neighbor, and a fraction $`(1p)`$ to the right neighbor, with $`0p1`$. Let us coarse-grain the line into cells of a certain (small) size. The variation of energy of the $`i`$-th cell after a parallel updating will be given by
$$\mathrm{\Delta }E^{(i)}\rho _a^{(i)}+(1ϵ)p\rho _a^{(i+1)}+(1ϵ)(1p)\rho _a^{(i1)},$$
(27)
where $`\rho _a^{(i)}`$ is the density of active sites in the $`i`$-th cell. The proportionality stems from the fact that the actual contribution comes only from the boundary sites in the cell. This last equation can be rewritten
$$\mathrm{\Delta }E^{(i)}ϵ\rho _a^{(i)}(1ϵ)p\left[\rho _a^{(i)}\rho _a^{(i+1)}\right](1ϵ)(1p)\left[\rho _a^{(i)}\rho _a^{(i1)}\right].$$
(28)
For $`p=0`$ or $`p=1`$ (for a strictly directed model), only one term remains and, after passing to a continuous formulation, we recover Eq. (15), restricted to the preferred direction. For $`p0,1`$, and after performing some algebraic manipulations, we obtain
$`\mathrm{\Delta }E^{(i)}`$ $``$ $`ϵ\rho _a^{(i)}+(1ϵ)\left[\rho _a^{(i+1)}+\rho _a^{(i1)}2\rho _a^{(i)}\right]`$
$``$ $`(1ϵ)p\left[\rho _a^{(i)}\rho _a^{(i1)}\right]+(1ϵ)(1p)\left[\rho _a^{(i)}\rho _a^{(i+1)}\right]`$
$``$ $`ϵ\rho _a^{(i)}+(1ϵ)p_x^2\rho _a^{(i)}+(1ϵ)p_x\rho _a^{(i)}+(1ϵ)(1p)(_x\rho _a^{(i)})`$
$`=`$ $`ϵ\rho _a^{(i)}+(1ϵ)p_x^2\rho _a^{(i)}(1ϵ)(12p)_x\rho _a^{(i)},`$
where in the third line we have passed to the continuum limit, defining the discretized first derivative $`_x\rho _a^{(i)}=\rho _a^{(i)}\rho _a^{(i+1)}`$ and second derivative $`_x^2\rho _a^{(i)}=\rho _a^{(i+1)}+\rho _a^{(i1)}2\rho _a^{(i)}`$. For $`p=1/2`$ we recover, as expected, an isotropic equation. For all other values of $`p`$, the equation for the energy will be effectively given by (26), with the phenomenological parameter $`\lambda `$ given at the microscopic level by the expression $`(1ϵ)(1/2p)`$.
From Eq. (26), we can obtain the corresponding equation for the susceptibility. The solution in Fourier space is readily found to be
$$\chi (q,\stackrel{}{k})=\frac{1}{D_{}k^2+D_{}q^2+2i\lambda q+ϵ}.$$
(29)
Upon integration over $`\stackrel{}{k}`$ and $`q`$, one obtains the expression in real space
$$\chi (x_{},\stackrel{}{x}_{})=\frac{1}{(2\pi )^d}d^{d1}ke^{i\stackrel{}{k}\stackrel{}{x}_{}}_{\mathrm{}}^{\mathrm{}}𝑑q\frac{e^{iqx_{}}}{D_{}k^2+D_{}q^2+2i\lambda q+ϵ}.$$
(30)
This last integral can be performed analytically in $`d=1`$ and $`2`$ (see Appendix). For $`d>2`$, even though we do not have a closed expression, we can obtain the leading scaling behavior. To simplify the calculations, we set, without lack of generality, $`D_{}=D_{}=1`$. The integration in $`q`$ is done by the method of the residues (see Appendix). The integration of the $`\stackrel{}{k}`$ angular part yields
$$\chi (x_{},\stackrel{}{x}_{})=\frac{1}{2}\left(\frac{\gamma }{2\pi }\right)^{\nu +1}x_{}^\nu _0^{\mathrm{}}𝑑zz^{\nu +1}J_\nu (\gamma x_{}z)\frac{e^{x_{}(\gamma \sqrt{1+z^2}\lambda )}}{(1+z^2)^{1/2}}.$$
(31)
Here, $`J_\nu (z)`$ is the first kind Bessel function of order $`\nu `$, and we have defined the constants $`\nu =(d3)/2`$ and $`\gamma =(\lambda ^2+ϵ)^{1/2}`$. We are interested in the behavior of this integral for large distances, that is, in the limit $`x_{}x_{}1`$. In this limit, the weight of the integral is given by the region of small $`z`$, since the exponential suppresses large values. We can then approximate the integral in the interval $`0<z<1`$ and perform a Taylor expansion of the square root in the exponential and the denominator. In the denominator, we readily have $`(1+z^2)^{1/2}1`$. The term in the exponential, however, contains a constant term, and must be therefore expanded up to second order:
$$x_{}(\gamma \sqrt{1+z^2}\lambda )x_{}(\gamma [1+z^2/2]\lambda )=x_{}(\gamma \lambda )x_{}\gamma z^2/2.$$
(32)
In the limit $`ϵ0`$, we have $`\gamma \lambda `$, and the constant $`\gamma \lambda `$ can be expanded to give:
$$\gamma \lambda =(\lambda ^2+ϵ)^{1/2}\lambda \lambda \left(1+\frac{ϵ}{2\lambda ^2}\right)\lambda =\frac{ϵ}{2\lambda }.$$
(33)
Substituting these approximations into Eq. (31), we are led to the expression
$$\chi (x_{},\stackrel{}{x}_{})\frac{1}{2\lambda }\frac{1}{(2\pi )^{\nu +1}}x_{}^{(1d)}e^{x_{}ϵ/2\lambda }_0^{\mathrm{}}𝑑yy^{\nu +1}J_\nu (y)e^{\frac{x_{}}{2\lambda x_{}^2}y^2},$$
(34)
where we have performed the change of variables $`y=\gamma x_{}z`$ and extended again the upper limit of the integral to infinity (which is allowed given its exponential convergence). The integral in Eq. (34) yields :
$$\chi (x_{},\stackrel{}{x}_{})\frac{1}{2\lambda }\left(\frac{\lambda }{2\pi }\right)^{\nu +1}x_{}^{(1d)/2}e^{x_{}ϵ/2\lambda }e^{\lambda x_{}^2/2x_{}},$$
(35)
which as usual we can write in the scaling form,
$$\chi (x_{},\stackrel{}{x}_{})=x_{}^{(1d)/2}\mathrm{\Gamma }(\frac{x_{}}{\xi _{}},\frac{x_{}}{\xi _{}}).$$
(36)
From here, the correlation exponents read $`\nu _{}^{}=1`$ and $`\nu _{}^{}=1/2`$, as in the directed case. This implies again an affinity exponent $`\zeta =1/2`$.
The conclusion of the lengthy calculations developed in this section is that the presence of any amount of diffusion along the preferred direction of a directed sandpile model is completely irrelevant. As soon as there is anisotropy in a model (in our mathematical formulation, when $`\lambda 0`$, however small), it takes over and places the model in the universality class of completely directed sandpiles. In particular, we recover the result $`sL`$ for any anisotropic sandpile, in agreement with the numerical results in Ref. and the analytic results of Ref. . We remark, however, that our results do not rely in a particular model like those of , but only on symmetry arguments, and are therefore of a broader generality.
## VII Conclusions
In this paper we have presented a detailed numerical analysis of deterministic and stochastic directed sandpile models. We find definitive evidence for the existence of a new universality class, embracing directed sandpile models with stochastic rules. The origin of the different critical behavior can be traced back to the presence of multiple topplings in the latter case. An example of this feature is provided in Fig. 12, where we plot the local density of topplings in two avalanches corresponding to the DDS and ESDS models. From this figure it becomes evident that the stochastic dynamics induces multiple toppling events, which are forbidden in the deterministic models. This feature has been fruitfully exploited in Ref. to obtain an analytical solution of the stochastic model.
We have also studied the case of directed sandpiles with bulk dissipation. In this case, our results prove that the critical behavior is unchanged. This points out that the boundary dissipation does not play any particular role in the development of the critical behavior in directed sandpiles.
Finally, numerical results indicate that some critical exponents show a super-universal nature, assuming the same values independently of the universality class. We provide an analytical explanation of this feature by means of a continuous phenomenological equation that takes into account the energy balance condition imposed by the dynamical rules in sandpile models.
## Acknowledgements
This work has been supported by the European Network under Contract No. ERBFMRXCT980183. We thank D. Dhar, R. Dickman, M. A. Muñoz, A. Stella, and S. Zapperi for helpful comments and discussions.
##
In this Appendix, we work out the exact value of the integral (30), giving the static susceptibility for anisotropic sandpiles, in $`d=1`$ and $`2`$. We take again, without lack of generality, $`D_{}=D_{}=1`$.
### Static susceptibility in $`d=1`$
In $`d=1`$, Eq. (30) reduces to the simple form
$$\chi (x_{})=\frac{1}{2\pi }_{\mathrm{}}^{\mathrm{}}𝑑q\frac{e^{iqx_{}}}{q^2+2i\lambda q+ϵ}.$$
(37)
This integral can be done by the method of the residues. The zeroes of the denominator are
$$q_\pm =i[\pm (\lambda ^2+ϵ)^{1/2}\lambda ].$$
(38)
Closing the contour of the integral arround the upper complex plane, we obtain
$$\chi (x_{})=\frac{1}{2}\frac{e^{x_{}(\sqrt{\lambda ^2+ϵ}\lambda )}}{\sqrt{\lambda ^2+ϵ}}.$$
(39)
We are interested in the limit $`ϵ0`$. The denominator gives $`\sqrt{\lambda ^2+ϵ}\lambda `$, while the argument in the exponential yields
$$\sqrt{\lambda ^2+ϵ}\lambda =\lambda \left[\left(1+\frac{ϵ}{\lambda ^2}\right)^{1/2}1\right]\frac{ϵ}{2\lambda }.$$
(40)
From here, we write
$$\chi (x_{})\frac{1}{2\lambda }e^{x_{}ϵ/2\lambda },$$
(41)
which allows to determine the correlation length $`\xi _{}=ϵ^1`$.
### Static susceptibility in $`d=2`$
In the case $`d=2`$, we have
$$\chi (x_{},x_{})=\frac{1}{(2\pi )^2}_{\mathrm{}}^{\mathrm{}}𝑑ke^{ikx_{}}_{\mathrm{}}^{\mathrm{}}𝑑q\frac{e^{iqx_{}}}{k^2+q^2+2i\lambda q+ϵ}.$$
(42)
Integrating by residues the $`q`$ variable, similarly as in the previous subsection, we obtain
$$\chi (x_{},x_{})=\frac{1}{2\pi }e^{\lambda x_{}}_0^{\mathrm{}}𝑑k\mathrm{cos}(kx_{})\frac{e^{x_{}(\sqrt{k^2+\lambda ^2+ϵ})}}{\sqrt{k^2+\lambda ^2+ϵ}}.$$
(43)
The last integral is duable , giving a modified Bessel function of order $`0`$:
$$\chi (x_{},x_{})=\frac{1}{2\pi }e^{\lambda x_{}}K_0\left[(\lambda ^2+ϵ)^{1/2}(x_{}^2+x_{}^2)^{1/2}\right].$$
(44)
In the limit $`x_{}x_{}1`$, the modified Bessel function $`K_0`$ can be replaced by its asymptotic form $`K_0(z)e^z/(2\pi z)^{1/2}`$ . Then, we have
$$\chi (x_{},x_{})\frac{1}{(2\pi )^{3/2}}\frac{1}{(\lambda ^2+ϵ)^{1/4}}\frac{1}{(x_{}^2+x_{}^2)^{1/4}}\mathrm{exp}\left(\lambda x_{}[\lambda ^2+ϵ]^{1/2}[x_{}^2+x_{}^2]^{1/2}\right).$$
(45)
In the double limit $`x_{}x_{}1`$ and $`ϵ0`$, the argument in the last exponential can be expanded
$$\lambda x_{}[\lambda ^2+ϵ]^{1/2}[x_{}^2+x_{}^2]^{1/2}\lambda x_{}\left(\lambda +\frac{ϵ}{2\lambda }\right)\left(x_{}+\frac{x_{}^2}{2x_{}}\right)\frac{\lambda x_{}^2}{2x_{}}\frac{x_{}ϵ}{2\lambda },$$
(46)
where we have only kept terms linear in $`ϵ`$ and $`x_{}^2/x_{}`$.
The static susceptibility can be finally written
$$\chi (x_{},x_{})\frac{1}{(2\pi )^{3/2}}\frac{1}{\lambda ^{1/2}}x_{}^{1/2}e^{x_{}ϵ/2\lambda }e^{\lambda x_{}^2/2x_{}},$$
(47)
from which we immediately read the correlation lengths $`\xi _{}=ϵ^1`$ and $`\xi _{}=ϵ^{1/2}`$.
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# Dephasing of Electrons on Helium by Collisions with Gas Atoms
## I Introduction
The damping of quantum effects in a system coupled to external degrees of freedom is a fundamental problem of atomic physics, condensed matter physics and quantum optics. There is great interest in understanding and controlling such damping in well characterised systems. Here we study the damping of quantum effects in the transport properties of a two-dimensional electron gas deposited on the surface of a pool of liquid helium .
Electrons on the surface of helium are vertically confined by their image charges and an (optional) applied holding field. They constitute a two-dimensional electron gas similar to those in semiconductor devices but with different scattering and damping mechanisms. Electrons may scatter off ripples on the surface of the helium pool (“ripplons”) or off helium vapour atoms above the liquid surface. Above 1 K, gas atom scattering dominates , and we concentrate on this regime. On the electronic time-scale, the helium vapour atoms are almost stationary and hence similar to impurities in a metal film or a semiconductor device. Thus there are quantum interference corrections to the resistivity at low temperature familiar from studies of transport in metals and semiconductors . These corrections result from constructive interference between closed electron paths and their time-reversed counterparts, leading to a small enhancement of the resistance (“weak-localization correction”).
The slow movement of helium atoms leads to damping of weak-localization. There is an important distinction between the effect of vertical and horizontal motion of helium atoms. Roughly, horizontal movement produces damping by scrambling the phase of the interfering paths; vertical movement, by reducing the weight of contributing paths of long duration. The effect of horizontal movement has been analysed previously ; it is the purpose of this paper to study the effect of vertical motion.
The central result is that due to vertical motion of the helium atoms the interference contribution of paths of duration $`t`$ is reduced by a factor $`\mathrm{exp}(t/\tau _v)^3`$. Thus paths of duration greater than the damping time, $`\tau _v`$, are effectively cutoff. An interesting feature is that damping due to both vertical and horizontal movement of helium atoms is not a simple exponential; it cuts off more sharply as the exponential of $`t^3`$. In contrast, electron-electron and electron-phonon interactions in metals and semiconductors are supposed to produce simple exponential damping. Damping in atomic physics and nuclear magnetic resonance is also commonly a simple exponential; this is indicated by the Lorentzian shape of spectral and magnetic resonance lines <sup>*</sup><sup>*</sup>*Recall that the Fourier transform of $`e^{|t|}`$ is a Lorentzian.. As emphasized by Afonin et al. in context of quantum transport, the form of damping can be probed by measuring the magnetic field dependence of the weak-localization correction (“weak-localization lineshape”). In section III we exhibit some lineshapes corresponding to different forms of damping.
Weak-localization has been observed in a related system, electrons on a surface of solid hydrogen . In this system helium vapour was deliberately introduced above the solid hydrogen to scatter electrons; thus gas atom damping is relevant to this type of experiment. More recently, Karakurt et al. have systematically studied the dependence of the damping rate on various experimental parameters (electron density; gas vapour pressure, controlled via temperature; and holding field) for electrons on helium . In this way they have obtained quantitative information on the contributions of different mechanisms to the damping rate. It is the experiment of Karakurt et al. that prompted us to carry out the present investigation.
For orientation it is useful to recall some typical parameters for the experiment of Karakurt et al. In the absence of a holding field, the electron is bound to the surface by its image. The charge of the image is reduced from the bare charge of the electron by a factor $`(ϵ1)/(ϵ+1)=7\times 10^3`$ ; thus the vertical scale of the electronic wavefunction is 76 $`\mathrm{\AA }`$. The lowest vertical subband wavefunction is of the Fang-Howard form, $`\varphi (z)z\mathrm{exp}(z/b)`$, at zero holding field; this form remains an excellent variational ansatz with $`b`$ an adjustable parameter when a holding field is applied. Here $`z`$ denotes the distance of the electron above the helium pool. The subband spacing is 6 K; hence for sufficiently low temperaturesAt zero field for an ideal interface the subband spectrum is like that of atomic hydrogen. To calculate the ground state occupancy for such a spectrum it is neccessary to regulate the partition function as discussed by Fermi. Hence estimation of the temperature threshold below which the surface electrons are effectively two-dimensional involves some subtlety . and electron densities below $`2\times 10^{15}`$ /m<sup>2</sup> the surface electrons behave like a two-dimensional electron gas. Much of the data of Karakurt et al. is at temperatures around 2 K and at a typical density of $`2\times 10^{11}`$ /m<sup>2</sup> corresponding to a Fermi temperature of 0.6 mK. Note that their two-dimensional electron gas is therefore non-degenerate in contrast to the situation in metal films and typical semiconductor devices. Thus transport properties are not determined entirely by mono-energetic electrons on the Fermi surface; instead we must sum the Boltzman-weighted contribution of electrons of all energies. The electron-atom collision time inferred from mobility measurements was typically a few ps. The longest relevant electronic time scale is $`\tau _z`$, the time taken by a thermal electron to move a distance $`b`$ (see eq 5 below). At 2 K and zero holding field $`\tau _z=80`$ ps. In comparison, the atom-atom collision time is enormous, of the order of 10 ns.
## II Analysis of Damping
In this section we analyse the damping produced by the vertical motion of helium vapour atoms. First we analyse a simple model (model I) that captures some of the essential physics, but leads to the incorrect conclusion that the damping factor goes as the exponential of $`t^2`$ rather than $`t^3`$, a result obtained earlier by Stephen . We then identify a shortcoming of model I and in the next subsection introduce and analyse an improved version (model II) that leads to the correct answer.
### A Model I
In this model we assume that the Helium atoms are able to scatter electrons only if they are within a certain distance (denoted $`b`$) from the liquid helium surface. It is also assumed that the scattering is independent of the precise height of the atom so long as it lies within the prescribed distance.
Consider $`p(t)=`$ probability that an atom will remain within the scattering distance for a time $`t`$. At first let us assume the probability decays exponentially,
$$p(t)=\mathrm{exp}\left(\frac{t}{\sqrt{\pi }\tau _z}\right);$$
(1)
the numerical coefficient in the exponential is for later convenience. Since the motion of vapour atoms is essentially independent the probability that $`n`$ atoms will remain within the scattering distance for time $`t`$ is
$$[p(t)]^n=\mathrm{exp}\left(\frac{nt}{\sqrt{\pi }\tau _z}\right).$$
(2)
Weak-localization results from constructive interference between the history in which an electron traverses a particular closed path and the history in which it traverses the same path backwards. A path of duration $`t`$ involves $`t/\tau _e`$ collisions. For this path to contribute to weak-localization it is neccessary for all atoms to remain within the electron-scattering region for a duration of order $`t`$. Hence the fraction of paths of duration $`t`$ that contribute to weak-localization is
$$\gamma (t)=\mathrm{exp}\left(\frac{t^2}{\sqrt{\pi }\tau _e\tau _z}\right).$$
(3)
Here $`\gamma `$ is the damping factor; essentially this result is given in ref .
This argument must be improved in two ways. First, it is not neccessary for all the atoms to remain in place for the entire duration of a closed path. In particular, atoms encountered by the electron in the middle of a closed path are encountered by both the forward and backward path at essentially the same time. At the other extreme, atoms encountered early on the forward path are encountered towards the end on the backward path, a time $`t`$ later. Thus on average atoms need to remain in place for a time $`t/2`$. Hence the damping factor is really
$`\gamma (t)`$ $`=`$ $`\mathrm{exp}\left({\displaystyle \frac{t^2}{2\sqrt{\pi }\tau _e\tau _z}}\right)`$ (4)
$``$ $`\mathrm{exp}\left({\displaystyle \frac{t^2}{\tau _v^2}}\right)`$ (5)
with $`\tau _v=\pi ^{1/4}\sqrt{2\tau _e\tau _z/c}`$.
A second improvement is needed because eq (1) is incorrect. $`p(t)`$ is easily calculated and seen to not be exponential. Here we mention only the relevant features of $`p(t)`$; the details are relegated to appendix A. (i) As expected on dimensional grounds, $`p(t)`$ is a function of $`t/\tau _z`$ alone, where
$$\tau _z=\sqrt{\frac{Mb^2}{2kT}}.$$
(6)
Here $`M=`$ mass of a helium atom. Physically, $`\tau _z`$ is the time taken by a thermal atom to move a distance $`b`$. (ii) For short times, $`t\tau _z`$, we find
$$p(t)1\frac{t}{\sqrt{\pi }\tau _z}$$
(7)
(iii) For long times, $`t\tau _z`$, $`p(t)`$ vanishes in a manner not relevant to our purpose.
Now the probability that $`n`$ atoms remain near the surface is
$`[p(t)]^n`$ $`=`$ $`\mathrm{exp}[n\mathrm{ln}p(t)]`$ (8)
$`=`$ $`\mathrm{exp}\left[n\mathrm{ln}\left(1{\displaystyle \frac{t}{\sqrt{\pi }\tau _z}}+\mathrm{}\right)\right]`$ (9)
$``$ $`\mathrm{exp}\left({\displaystyle \frac{nt}{\sqrt{\pi }\tau _z}}\right).`$ (10)
This shows that for large $`n`$, $`[p(t)]^n`$ can be approximated as an exponential only for $`t\tau _z/\sqrt{n}`$; but since it becomes negligible in any case once $`t\tau _z/n`$, there is no significant error in taking $`[p(t)]^n`$ to be an exponential.
The upshot of this discussion is that although $`p(t)`$ is far from exponential, $`[p(t)]^n`$ is a simple exponential under appropriate circumstances; eq (2) is valid, although eq (1) is not. Similarly we see that eq (4) is also valid provided $`\tau _z\tau _e`$, a condition needed for weak-localization.
In summary, for model I the damping decays as the exponential of $`t^2`$. Provided $`\tau _z\tau _e`$, it is given by eq (4). The atomic time constant $`\tau _z`$ is given by eq (5). Evidently, the three time scales are arranged in the hierarchy $`\tau _z>\tau _v>\tau _e`$.
### B Model II
The shortcoming of model I is the assumption stated in the first paragraph of the previous subsection. It is more realistic to assume that the ability of an atom to scatter electrons turns off smoothly as it moves away from the liquid helium surface.
If we treat the atoms as hard core potentials, the contribution of a closed path to the return amplitude is a product of the amplitude for the electron to go to atom 1, multiplied by the amplitude to scatter off atom 1, multiplied by the amplitude to go to atom 2, multiplied by the amplitude to scatter off atom 2, and so on around the loop.
Let $`A(z)`$ be the amplitude to scatter from an atom at height $`z`$ above the helium surface. Model I can be described as the case in which $`A(z)`$ is a step function. Here we choose
$`A(z)`$ $`=`$ $`{\displaystyle \frac{4\lambda z^2}{b^3}}\mathrm{exp}\left({\displaystyle \frac{2z}{b}}\right)\mathrm{for}z>0;`$ (11)
$`=`$ $`0\mathrm{for}z<0.`$ (12)
This is derived by taking the vertical subband wavefunction of the electrons to be of the Fang-Howard form and treating the helium atom as a short-ranged hard-core potential.
If the helium atoms are only allowed to move vertically the forward and backward paths remain in phase; however the interference contribution to the return probability is still modified because the forward and backward paths have different amplitudes to scatter from each atom. We must consider
$$Q(t)=A(z)A(z+vt).$$
(13)
Here $`t`$ is the difference in the times at which the atom is encountered on the forward and return path. The atom is assumed to move ballistically at vertical speed $`v`$ for this time. $`\mathrm{}`$ denotes an average over all possible configurations of the Helium atom (vertical position is assumed to be uniformly distributed and vertical speed is given by the Maxwell-Boltzmann formula).
Introduce the normalization factor $`R(t)`$ defined by
$$R(t)=A(z)A(z+vt).$$
(14)
Here the average over vertical position is performed as in eq (9) but the velocity distribution is assumed to be a delta function peaked about zero. $`R(t)`$ is the value of $`Q(t)`$ when the atoms don’t move. Let
$$q(t)=Q(t)/R(t).$$
(15)
The contribution of paths of duration $`t`$ is then reduced roughly by the factor $`q(t)`$ raised to the power $`t/\tau _e`$, the number of atoms encountered.
$`q(t)`$ is analogous to $`p(t)`$ for model I. Again on dimensional grounds, $`q(t)`$ depends only on the ratio $`t/\tau _z`$ and again we are interested only in the short time behaviour. This is evaluated in Appendix B. The difference from the previous case is that
$$q(t)=1\frac{t^2}{3\tau _z^2}+\mathrm{}$$
(16)
for short times, $`t\tau _z`$. The behaviour is quadratic rather than linear (compare eq 6). Quadratic behaviour is generic; the linear behaviour for model I is an artifact of the discontinuous step in $`A(z)`$. Hence $`q(t)`$ raised to the power of $`n`$ is approximately Gaussian rather than exponential
$$[q(t)]^n\mathrm{exp}\left(\frac{nt^2}{3\tau _z^2}\right).$$
(17)
Eq (13) should be contrasted with eq (7) above for model I.
To obtain the damping factor, roughly we must replace $`n`$ in eq (13) by $`t/\tau _e`$, the number of atoms encountered in a path of duration $`t`$. Before that we must replace $`t^2`$ in eq (13) by $`t^2/3`$, its value averaged over the interval from 0 to $`t`$ with uniform weight. This is to take into account the range in the difference of times at which an atom is encountered along the forward and reversed histories.
The result for the damping factor is
$$\gamma (t)=\mathrm{exp}\left(\frac{t^3}{\tau _v^3}\right)$$
(18)
where $`\tau _v=(9\tau _e\tau _z^2)^{1/3}`$. Eq (14) is the central result of this paper. It is valid provided $`\tau _z\tau _e`$.
## III Lineshape
Karakurt et al. observed damping by vapour atom motion at low electron density and by electron-electron interaction at high density . At intermediate densities, damping by both mechanisms was substantial. Vapour atom scattering produces cubic exponential damping; electron-electron interaction is presumably simple exponential. Afonin et al. have pointed out that the weak-localization lineshape depends on the form of damping and they have given an expression for the lineshape in the extreme cases that the damping is entirely simple exponential or entirely cubic exponential. The purpose of this section is to study the lineshape in the intermediate regime and examine how it crosses over from one extreme form to the other.
For simplicity, first let us consider a degenerate electronic system. Assuming that the different damping mechanisms are independent the lineshape is given by
$`\delta g(E,B)`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}{\displaystyle \frac{e^2}{h}}\left({\displaystyle \frac{W}{L}}\right)\varphi (E,B);`$ (19)
$`\varphi (E,B)`$ $`=`$ $`{\displaystyle _{\tau _e}^{\mathrm{}}}𝑑t{\displaystyle \frac{4\pi eDB}{h}}{\displaystyle \frac{e^{t/\tau _1}e^{t^3/\tau _3^3}}{\mathrm{sinh}(4\pi eDtB/h)}}.`$ (20)
Here $`E`$ is the Fermi energy; $`W`$, the sample width; $`L`$, the sample length; $`D`$ the electron diffusion constant; $`1/\tau _1`$, the simple exponential damping rate; and $`1/\tau _3`$, the cubic exponential damping rate. Energy dependence enters the integrand in eq (15) through the diffusion constant $`D=E\tau _e/m`$ and through the energy dependence (if any) of the time constants $`\tau _1`$ and $`\tau _3`$. The $`\mathrm{sinh}`$ factor in eq (15) may be recognised as the Fourier transform of the directed area distribution for closed random walks on a plane .
It is useful to manipulate eq (15) into a more revealing form. To this end introduce the dimensionless variable $`u=8\pi eDtB/h`$ to obtain
$$\delta g=\frac{1}{\pi }\frac{e^2}{h}\left(\frac{W}{L}\right)_{B/B_e}^{\mathrm{}}𝑑u\frac{\mathrm{exp}u\left(\frac{1}{2}+\frac{B_1}{B}\right)\mathrm{exp}\left(u^3\frac{B_3^3}{B^3}\right)}{1e^u};$$
(21)
here $`B_e=h/(8\pi eD\tau _e)`$. Making use of the asymptotic formula
$$_ϵ^{\mathrm{}}𝑑u\frac{e^u}{u}\mathrm{ln}\frac{1}{ϵ}+\gamma $$
(22)
we obtain
$`\delta g/\left({\displaystyle \frac{e^2}{h}}{\displaystyle \frac{W}{L}}\right)`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}\left[\mathrm{ln}{\displaystyle \frac{B_e}{B_1}}+\mathrm{ln}{\displaystyle \frac{B_e}{B_3}}\right]`$ (24)
$`+{\displaystyle \frac{1}{\pi }}({\displaystyle \frac{B_1}{B}},{\displaystyle \frac{B_3}{B}})`$
where $`B_1=h/(8\pi eD\tau _1)`$, $`B_2=h/(8\pi eD\tau _3)`$, $`\gamma =0.577216\mathrm{}`$ is Euler’s constant and the function
$`(x,y)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{ln}x+{\displaystyle \frac{1}{2}}\mathrm{ln}y+\gamma `$ (26)
$`{\displaystyle _0^{\mathrm{}}}𝑑u\left[{\displaystyle \frac{e^u}{u}}{\displaystyle \frac{\mathrm{exp}u\left(\frac{1}{2}+x\right)\mathrm{exp}\left(u^3y^3\right)}{1e^u}}\right].`$
Eqs (18) and (19) constitute the generalisation of the standard weak-localisation lineshape to the case that both $`\tau _1`$ and $`\tau _3`$ damping are present. For the special case that there is no $`\tau _3`$ damping (hence $`y0`$) eqs (18,19) reduce to the familiar expression involving digamma functions by use of the integral representation
$$_0^{\mathrm{}}𝑑u\left(\frac{e^u}{u}\frac{e^{\left(\frac{1}{2}+x\right)u}}{1e^u}\right)=\psi (\frac{1}{2}+x).$$
(28)
A significant feature revealed by eqs (18,19) is that the lineshape is universal: $``$ does not depend on microscopic length scales. Note that the magnetic field dependence is entirely in the second term of eq (18); the first term is an additive constant. A practical advantage of eq (19) over eq (15) is that the integrand is well behaved for both large and small $`u`$. In contrast, the integrand in eq (15) diverges at the lower end.
To study the crossover in lineshape we fix the damping rate $`1/\tau _1+1/\tau _3=1/\tau _\varphi `$. Equivalently, we fix $`B_1+B_3=B_\varphi `$. $`\delta g`$ is plotted as a function of $`B`$ for several values of the ratio $`B_1/B_\varphi `$. Fig 1 shows that for the same damping rate the lineshape changes noticeably as damping shifts from simple exponential to cubic exponential.
Fig 2 shows the behaviour of the conductance minimum at $`B=0`$ for a fixed damping rate. It is given by
$$\delta g(B=0)=\frac{1}{\pi }\frac{e^2}{h}\left[\mathrm{ln}\left(\frac{B_e}{B_\varphi }\right)+u\left(\frac{B_3}{B_1}\right)\right]$$
(29)
with the crossover function
$$u(x)=\mathrm{ln}(1+x)+\gamma +_0^{\mathrm{}}𝑑s(1+3s^2x^3)\mathrm{ln}se^se^{s^3x^3}.$$
(30)
As implied by eqs (18) and (19) the crossover depends only on the ratio $`B_3/B_1`$. $`u`$ has the limiting values $`u(0)=0`$ and $`u(\mathrm{})=2\gamma /3`$.
Under experimental conditions the electron gas is non-degenerate. At finite temperature
$`\delta g(T,B)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑E{\displaystyle \frac{f}{E}}\delta g(E,B)`$ (31)
$``$ $`{\displaystyle \frac{n\pi \mathrm{}^2}{m(kT)^2}}{\displaystyle _{E_c}^{\mathrm{}}}𝑑E\delta g(E,B)e^{E/kT}.`$ (32)
$`n`$ is the area density of electrons. In the second line of eq (23) we have approximated the Fermi function by a Boltzman factor and imposed a lower cutoff $`E_c`$. Below the cutoff energy the electrons are presumed to be strongly localized and to make an insignificant contribution to the conductance. These finite temperature considerations make it more difficult to extract the form of damping from the lineshape.
## IV Conclusion
In summary, we have given a physical argument that due to vertical motion of helium atoms the interference of electron paths of duration $`t`$ is damped by a factor $`\mathrm{exp}(t/\tau _v)^3`$. We have derived a formula for the universal magnetoconductance lineshape for the case that both $`\tau _1`$ and $`\tau _3`$ damping are present. It should be possible to rederive these results via impurity averaged diagrams; this is left open for future work.
It is a pleasure to acknowledge helpful correspondence with M. Stephen. This work was supported in part by NSF Grants DMR 98-04983 (DH and HM) and DMR 97-01428 (AJD) and by the Alfred P. Sloan Foundation (HM). HM acknowledges the hospitality of the Aspen Center for Physics where this work was completed.
## A Asymptotics of $`p(t)`$
We wish to calculate $`p(t)`$, the probability that a vapour atom will remain within a vertical elevation $`b`$ of the liquid surface for a time $`t`$. We assume (i) the initial elevation of the atom is uniformly distributed between zero and $`b`$; (ii) the vertical velocity is Maxwell-Boltzman distributed; (iii) the atom moves ballistically; and (iv) if the atom strikes the liquid surface it sticks and does not reflect . Due to assumption (iii) the expression for $`p(t)`$ that we derive is valid only for times short compared to the atom-atom collision time; however this is not a serious restriction since we are interested only in the short time behaviour of $`p(t)`$.
Based on these assumptions we may write
$`p(t)`$ $`=`$ $`{\displaystyle _0^b}𝑑z{\displaystyle \frac{1}{b}}{\displaystyle _0^{(bz)/t}}𝑑v\sqrt{{\displaystyle \frac{M}{2\pi kT}}}\mathrm{exp}\left({\displaystyle \frac{Mv^2}{2kT}}\right)`$ (A2)
$`+{\displaystyle _0^b}𝑑z{\displaystyle \frac{1}{b}}{\displaystyle _{z/t}^0}𝑑v\sqrt{{\displaystyle \frac{M}{2\pi kT}}}\mathrm{exp}\left({\displaystyle \frac{Mv^2}{2kT}}\right).`$
The two contributions correspond to the atom moving up and down respectively.
By exchanging the order of integration we can perform the $`z`$ integral first to obtain
$$p(t)=\frac{2}{\sqrt{\pi }}_0^{1/\overline{t}}𝑑u(1u\overline{t})\mathrm{exp}(u^2).$$
(A3)
We have rescaled variables so that $`u=v/\sqrt{2kT/M}`$ and $`\overline{t}=t/\tau _z`$. Note that $`p(t=0)=1`$ and as $`t\mathrm{}`$, $`p(t)0`$. Eq (A2) is an exact expression for $`p(t)`$. The small time, $`t\tau _z`$, asymptotic behaviour is
$$p(t)\left(1\frac{\overline{t}}{\sqrt{\pi }}+\mathrm{}\right).$$
(A4)
## B Asymptotics of $`q(t)`$
To calculate $`q(t)`$ we assume that the initial elevation of the vapour atom is uniformly distributed between the liquid surface and an upper cutoff $`L`$. Ultimately we shall take $`L\mathrm{}`$. Aside from this we share the assumptions (ii), (iii) and (iv) of Appendix A.
Hence we obtain
$`Q(t)`$ $`=`$ $`A(z)A(z+vt)`$ (B1)
$`=`$ $`{\displaystyle \frac{1}{L}}\sqrt{{\displaystyle \frac{M}{2\pi kT}}}{\displaystyle _0^L}𝑑z{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑v\mathrm{exp}\left({\displaystyle \frac{Mv^2}{2kT}}\right)A(z)A(z+vt).`$ (B2)
Using eq (8) for $`A(z)`$ and rescaling we obtain
$`Q(t)`$ $`=`$ $`{\displaystyle \frac{16}{\sqrt{\pi }}}{\displaystyle \frac{\lambda ^2}{bL}}{\displaystyle _0^{\mathrm{}}}𝑑ue^{u^2}e^{2u\overline{t}}{\displaystyle _0^{L/b}}𝑑\zeta e^{4\zeta }\zeta ^2(\zeta +u\overline{t})^2`$ (B5)
$`+{\displaystyle \frac{16}{\sqrt{\pi }}}{\displaystyle \frac{\lambda ^2}{bL}}{\displaystyle _{\mathrm{}}^0}𝑑ue^{u^2}e^{2u\overline{t}}{\displaystyle _{u\overline{t}}^{L/b}}𝑑\zeta e^{4\zeta }\zeta ^2(\zeta +u\overline{t})^2`$
Here $`\overline{t}=t/\tau _z`$, $`u=v/\sqrt{(2kT)/M}`$ and $`\zeta =z/b`$.
Performing the $`\zeta `$ integral yields
$$Q(t)=\frac{\lambda ^2}{4\sqrt{\pi }bL}_0^{\mathrm{}}𝑑ue^{u^2}e^{2u\overline{t}}(3+6u\overline{t}+4u^2\overline{t}^2)$$
(B7)
and hence the normalization
$$R(t)=Q(0)=\frac{3\lambda ^2}{8bL}.$$
(B8)
The exact reduction factor is then
$`q(t)`$ $`=`$ $`{\displaystyle \frac{Q(t)}{R(t)}}`$ (B9)
$`=`$ $`{\displaystyle \frac{2}{3\sqrt{\pi }}}{\displaystyle _0^{\mathrm{}}}𝑑ue^{u^2}e^{2u\overline{t}}(3+6u\overline{t}+4u^2\overline{t}^2)`$ (B10)
with the small time, $`t\tau _z`$, asymptotic behaviour
$$q(t)1\frac{\overline{t}^2}{3}+\frac{\overline{t}^4}{2}+\mathrm{}$$
(B11)
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# Particle with internal dynamical asymmetry: chaotic self-propulsion and turning
(14th July 2000)
We consider model of a complex particle that consists of a rigid shell and a nucleus with spatial asymmetric interaction. The particle’s dynamics with the nucleus driven by a periodic excitation is considered. It is shown that unidirectional self-propulsed particle motion arises in the absence of spatial and temporary asymmetry of external potentials and influences. Transport modes are the general case of complex particle dynamics in the presence of nonlinear friction or periodic external potential. The changes of average transport velocity and direction of transport are determined by qualitative changes of internal dynamics regimes: local attractor bifurcations in the internal phase space of the complex particle. Finally, microbiological relevance of the proposed model is briefly discussed.
PACS number: 05.40.+a, 05.60.+k
Over the past years, the appearance and intensive development of some areas of nonlinear dynamics and statistical physics have been inspired by several microbiological phenomena \[1-4\]. First, it is the physical principles and functional mechanisms of molecular motors, nano- and microscale engines, which effectively convert chemical energy to mechanical work . Second, it is physics of motility, namely the unidirectional transport of microbiological objects . Both these research lines are closely related to each other, and the biological motor efficiency can be defined in the terms of average transport velocity.
One of the most interesting mechanisms of the directed transport of microobjects is the so-called ratchet-effect, related to particle motion in periodic potential with spatial asymmetry . As a rule, the problem of the ratchet mechanism is considered within the frame of stochastic approach, when the unidirectional motion of overdamped Brownian particle in asymmetrical potential takes places due to unbiased nonequilibrium fluctuations in the absence of any macroscopic forces and gradients \[4-5\]. However, recently a significant interest has been paid to the deterministic approach to this problem, taking into account the finiteness of inertial terms \[6-7\]. A deterministic description allows to expand dynamic modes spectrum, including regular and chaotic ones , and to study the role of interaction mechanisms in greater detail. Within the framework of this approach, similar to the deterministic diffusion phenomenon , low-dimensional dynamic chaos can effectively play the role of thermodynamic fluctuations. In , Mateos studied the event of the deterministic motion of a particle in smooth ratchet potential under the action of a symmetrical periodic time-dependent force. He showed that a wide spectrum of dynamic regimes, with bifurcation transitions between regular and chaotic regimes, arises when periodic force amplitude is varied. Moreover, he found, that current reversals are related with bifurcations from chaotic to periodic regimes. In a recent paper , Flach, Yevtushenko and Zolotaryuk have studied a more general case of a directed current appearance in deterministic nonlinear systems. It has been shown that the key factor for a transport regime appearance is the existence of spatial (”spatial ratchet”) and/or temporary (”temporal ratchet”) asymmetry of influences upon a particle.
Thus, one of the necessary triggers for directed transport appearance is the existence of asymmetry in the system. The basic idea of the present paper is that this asymmetry can be ”built” into the internal degrees of freedom of the complex particle and, in this way, the source of motion can be located inside of the object itself. At that, external potentials and the influences can be completely symmetrical . Such a particle is an active self-propulsed walker, which can control the direction and velocity of its own motion. Further we will formulate a simple one-dimensional model of a complex particle and we will consider its dynamics in: (i) a medium with nonlinear friction (that is true for real biological media nearly always ) and (ii) in a linear medium under external periodic symmetrical potential. We will analyse spectrum of dynamic regimes that arises under changing parameters of the system, and also the relations between the peculiarities of a complex particle internal dynamics and its directed transport characteristics. Finally, we will shortly discuss the biological motivation of the proposed model.
Let us consider a one-dimensional system, consisting of a rigid spherical shell with a mass $`M`$, whose position is determined by the coordinate of its centre $`X`$, and an internal particle (”nucleus”) with a mass $`m`$ and the coordinate $`x`$, attached to the internal walls of the shell (see inset in Fig1.). The elastic interaction between the shell and the nucleus is governed by an asymmetric nonlinear potential $`U(xX)`$, so that the force arising due to the nucleus displacement from the centre of the shell depends not only on the absolute displacement value $`\mathrm{\Delta }x=xX`$ , but also on its sign (Fig.1). The minimum of the potential is located at the origin ($`U^{^{}}(x)=0`$), which corresponds to the equilibrium nucleus position in the centre of the shell. Let us also presume that the nucleus is driven by a periodic symmetrical force $`f(t)=acos(\omega t)`$ which is defined, for example, by mechanochemical coupling to a periodic biochemical reaction, such as ATP hydrolisis . Then, generally, with an external potential and nonlinear friction present, the complex particle dynamics, after trivial mass normalisation, is defined by the pair of equations:
$`\ddot{x}=K_1(\dot{x}\dot{X}){\displaystyle \frac{U(xX)}{x}}+acos(\omega t),`$ (1)
$`\ddot{X}=K_2(\dot{X})+\mu {\displaystyle \frac{U(xX)}{x}}+W^{^{}}(X)`$ (2)
where $`K_i(\nu )`$ are nonlinear friction forces, $`\mu =\frac{m}{M}`$ is the mass factor, $`a,\omega `$ are driving force amplitude and frequency respectively. The potential of interaction is given by (Fig.1):
$`U(x)=\alpha (xx_0)^2+\beta (xx_0)^4\gamma (xx_0)^3`$ (3)
where $`\alpha ,\beta ,\gamma `$ are the parameters of interaction, $`x_0`$ is an appropriate shift, such that a minimum of potential (3) is located at the origin. Nonlinear friction forces are given by $`K_i=s_i\nu +q_i\nu ^3`$. The parameters $`\mu =2,\omega =1,\alpha =\beta =1,\gamma =2.5,x_01.553`$ are fixed throughout this paper.
The complete phase space of the system (1-2) is five-dimensional $`(x,X,\dot{x},\dot{X},t)`$, the dynamic equations are nonlinear and contain dissipative terms. So, in the phase space of system (1-2) both periodic (limiting cycles) and chaotic (strange attractors) attractors can exist . On the other hand, the complex particle structure induce, in the natural way, the splitting of the complete phase space into two subspaces, the external one, accessible to an external observer - $`(X,\dot{X},t)`$, and the internal one, corresponding to an observer located within the shell - $`(xX,\dot{x}\dot{X},t)`$. The first subspace corresponds to nonlocal transport of the complex particle as a whole, the second one corresponds to local dynamics, which is determined by the internal interaction. The most interesting question is how the correlation between these two scales of the particle’s complete dynamic is realised. While zero-mean driving symmetrical periodic force $`f(t)`$ does not contain a constant component, the appearance of the directed current can only be a consequence of the asymmetry of the internal interaction.
Let us consider now the motion of a particle in a nonlinear medium in the absence of an external potential ($`W(x)0,s_1=0.2,s_2=0.5,q_1=10^2,q_2=510^3`$). For the analysis of the particle’s dynamics in the external phase space ($`X,\dot{X},t`$) we used stroboscopic Poincare section with the period equal to the driving force period, $`T=2\pi /\omega `$. The bifurcation diagram for $`V=\dot{X}`$ in a limited range of the parameter $`a`$ is shown in Fig.2a . At first, the standard scenario of period-doubling route to chaos is realised ($`a[5.6,6.0347]`$), and then the bifurcation connected with the internal crisis of a chaotic attractor ($`a_16.2117`$) and the opposite tangent bifurcation connected with the birth of a steady period-three cycle take place ($`a_26.96441`$) . In Fig.2b we show the current $`J=\frac{1}{NM}_{j=1}^M_{t_0}^{t_0+N}V_j(t)𝑑t`$, $`(t_0=50,N=10^4,M=50)`$, as the function of parameter $`a`$ in the same range of values. As one can see, current reversal takes place exactly at the tangent bifurcation point $`a_2`$. As viewed from the internal phase space $`(xX,\dot{x}\dot{X},t)`$ it corresponds to transition from the chaotic attractor (Fig.3a) to the periodic one (Fig.3b). Besides, a smaller jump of the current value takes place at the internal crisis bifurcation point $`a_1`$, that corresponds to a sudden expansion of chaotic attractor . It can be shown that by a certain variation of the system parameters it is possible to achieve the current reversal taking place exactly at this bifurcation point. The investigation of the model’s dynamics in other ranges of parameter $`a`$ showed, that any abrupt change of the current $`J`$ is connected with sudden changes of the internal chaotic attractor - crisis (as interior crisis at $`a_1`$), or subduction (as tangent bifurcation at $`a_2`$) . Here, current reversal with the a greater probability occurs at bifurcations of the second type, that is tangent bifurcations (as in this case the jump of a current value is much greater).
Let us consider now the motion of a complex particle under external periodic potential $`W(x)=Acos(X)`$. To separate nonlinear friction mechanisms, let us presume that the medium is linear $`(q_i0)`$. We found that in this case, as well as in the one considered above, there also exist transport modes. However, unlike the free motion case, in this case the transport modes take place within certain parameter ranges. It is connected with the presence of external potential, which adds to the system some additional temporal and spatial characteristic scales, while some coherence between the parameters of external potential and internal force and interaction is needed for the appearance of transport. We found a very interesting effect which is connected with these coherence mechanism, the effect of sharp change (current ”switching”) (Fig.4b). In the parameter space this effect looks like full overbarrier reflection, when the particle flying in the ballistic mode above potential $`(V_m=_0^TV(t)𝑑t=1,a[5.1,5.206])`$ changes its flight direction to the opposite $`(V_m=_0^TV(t)𝑑t=1,a[5.363,5.57])`$ at a small variation of the driving force amplitude. The corresponding bifurcation diagram for $`V`$ is shown in Fig.4a. As one can see, the current switching, as it is in the nonlinear friction case, is connected with tangent bifurcation at $`a_c5.3633`$ leading to the birth of the limiting cycle from the chaotic attractor. The period-one limiting cycle of the particle internal phase space $`C_1`$ corresponds to the ballistic transport in the positive direction in the region, while another period-one cycle, $`C_2`$, corresponds to the transport in the negative direction (see inset in Fig.4a).
The existence of these limiting cycles, $`C_1`$ and $`C_2`$, allows to distinguish between two characteristic spatial scales and to separate the solution in the external phase space into two parts $`X(t)=X_s(t)+\xi (t)`$, where $`X_s(t)`$ is a slow part and $`\xi (t)`$ is a small fast part. After that, using perturbation theory ideology , for the slow variable we can construct a ”quasi” zero-order approximation:
$`\ddot{X}_s=s_2X_s+W^{^{}}(X_s)+S(t)`$ (4)
where
$`S(t)=G(x(t)X(t))={\displaystyle \frac{U(x(t)X(t))}{x(t)}}`$ (5)
is a periodic force acting on the shell that corresponds to the exact solution of the system (1-2). This forces $`S(t)`$ is zero-mean, $`_0^TS(t)𝑑t=0`$, with pronounced temporal asymmetry (Fig.5a and Fig.6a). Thus, the equation (4) describes the motion of a simple particle under a periodic potential, driving by a spatially homogeneous asymmetric periodic force . For the numerical integration of the equation (4) we used as excitation force the first six harmonics of the complete solution, $`S^{}(t)=_{k=1}^6S_kcos(k\omega t+\varphi _k)`$. The obtained results demonstrate very close agreement with the solutions of the entire system (1-2). In $`C_1`$ \- case in the in phase space of system (4) there exists the only one global attractor that corresponds to the particle ballistic transport in the positive direction (Fig.5a), and in $`C_2`$ \- case there is the only one global attractor that corresponds to the particle ballistic transport in the negative direction (Fig.5b). In the case of chaotic transport (e. g., $`a[2.208,5.363]`$) approximation (4) is not valid, because due to spectrum continuity there are no distinguished scales in the system.
It is necessary to note, that in each of the considered cases (free motion in nonlinear medium and motion under external periodic potential), there is only one global attractor (limit cycle or chaotic attractor) in the phase space of system (1-2) for the given value of parameter $`a`$. In those cases averaging over ensemble is equivalent to averaging over attractor invariant measure, and so we can talk about the directed motion of an individual self-propulsed complex particle rather than of current. The obtained current $`J`$ dependencies can be represented as dependencies of the average drift velocity $`V_m`$ of the particle, and we also can talk about turning instead of current reversal.
In summary, we have shown, that the internal dynamic asymmetry, ”dynamical chirality ”, can be a source of the directed motion of a complex particle. Transport modes are not exceptional, but occupy certain parameter ranges (it is a general case for the motion in a nonlinear medium). The direction and velocity of the motion are determined by the particle internal dynamics properties, and their changes are connected with the internal attractor bifurcations. It is interesting to note, that a similar situation will take place, if we consider a particle with symmetrical nonlinear internal interaction $`U(x)`$ and a temporal asymmetrical driving periodic force $`f(t)`$ . Such a force may be stimulated by a limiting cycle of some biochemical reaction , which is conformally coupled with the internal particle structure. Here, the force temporal asymmetry will be a general case, while symmetry is possible as an exceptional case in the space of reaction parameters.
The assumption of absolute rigidity of the shell is not a strict requirement. The validity of such assumption corresponds to the validity of the inequality $`\theta =\omega _s/\omega 1`$, where $`\omega `$ is the driving force frequency and $`\omega _s`$ is the characteristic frequency of large-scale deformation vibrations of the shell. Moreover, the viscosity of the external and internal medium also impedes the excitation of such oscillations. The energy dissipation by the small shell vibrations can be taken into account within nonlinear friction terms.
The considered mechanism can shed light upon physics of some non-canonical forms of microbiological motility. For example, swimming Cyanobacteria, having the shape of almost exact spheroid and having no external appendages perform directed motion in a liquid without any observable shape changes . The mechanism of internal dynamic asymmetry can be additional to the hydrodynamic mechanism accounting for the transport of Cyanobacteria through to small-scale periodic modulations of the shell shape . In the case of more complex cells, the internal structure of the particle can be a concentrated image of the cytoskeletal structure . Cytoskeletal polymers all have complex nonlinear viscoelastic characteristics . Moreover, interactions between cytoskeletal polymer systems in living cell may result in composite material with unique nontrivial dynamic properties . Furthermore, the proposed model provides for dynamic modelling of such collective phenomena, as microorganisms interaction and self-organisation at the level of microscopic equations of the motion of a single biological particle. The microscopic approach to these phenomena that consider a microorganism as an active object in active environment , makes it possible to effectively use the ideas about interaction and synchronization of internal attractors of individual microobjects .
For valuable discussions, we are indebted to Prof. V. V. Yanovsky, Prof. O. V. Usatenko and Dr. Hab. S. Flach.
FIGURE CAPTIONS
1. The interaction potential $`U(x)`$ between shell and nucleus. Inset shows the complex particle structure.
2. (a) Bifurcation diagram for $`V`$ as a function of $`a`$ and (b) the current $`J`$ as a function $`a`$ for $`\mu =2,\omega =1,\alpha =\beta =1,\gamma =2.5,W(X)0,s_1=0.2,s_2=0.5,q_1=10^2,q_2=510^3`$.
3. Projections of complex particle internal phase space attractors: (a) chaotic attractor for $`a=6.96`$ and (b) the period-three limit cycle just below tangent bifurcation at $`a6.9644`$.
4. (a) Bifurcation diagram for $`V`$ as a function of $`a`$ and (b) the current $`J`$ as a function $`a`$ for $`A=3,s_1=0.5,s_2=0.8,q_1=q_20`$. Inset shows (a) internal phase space period-one limit cycles: $`C_1`$ ($`a=5.18`$, positive current) and $`C_2`$ ($`a=5.41`$, negative current) and (b) corresponding particle trajectories.
5. (a) The shape of periodic zero-mean asymmetric force $`S(t)`$, corresponding to limit cycle $`C_1`$ ($`a=5.18`$, positive current): solid curve - exact solution for system (1 - 2), dotted curve - first six harmonics sum; (b) corresponding limit cycle in the reduced external phase space ($`V_m=1`$): solid curve for full system (1), dotted curve - for approximation (4) with the first six harmonics as driving force.
6. Same as in Fig.5 for cycle $`C_2`$ ($`a=5.41`$, negative current).
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# LAPTH-806/00 DTP-MSU/00-10 hep-th/0007228 Solitons and black holes in Einstein-Born-Infeld-dilaton theory
## 1 Introduction
Born-Infeld (BI) type generalisations of both Abelian and non-Abelian gauge theories have recently attracted much attention in the context of superstring theory. An Abelian BI action was obtained as the disc open string effective action in an external constant vector field (for references and a recent review see ). The computation is valid both for the bosonic string and the superstring and is exact in terms of the sigma-model coupling $`\alpha ^{}`$. Alternatively, the BI action was derived as an effective action governing the dynamics of vector fields on $`D`$-branes . The BI lagrangian introduces a natural bound for the field strength – a Born-Infeld ‘critical field’ – which should now be regarded as originating from the non-locality of the underlying fundamental theory. A direct consequence of this is the smoothing of the pointlike singularities of the vector field. As was shown already in the original papers by Born and Infeld , the Coulomb field of a point charge has a finite energy in this theory.
It is expected that a similar regularization of gravitational singularities should follow from the non-locality of closed strings . However no explicit closed string effective action summing up all terms in $`\alpha ^{}`$ was obtained so far. A somewhat simpler question is how the singularity of the charged black hole gets modified when gravity is still treated classically, while the dynamics of the vector field is governed by the strings. The guess is that the singularity, although not eliminated, will be smoothed somehow. As was shown by Gibbons and Rasheed , the Einstein-Born-Infeld charged black holes are less singular indeed as compared with their Einstein-Maxwell counterpart, the Reissner-Nordström (RN) solution. Namely, there is no RN-type divergent term in the metric near the singularity $`g_{00}Q^2/r^2`$, but only a Schwarzschild type term $`g_{00}m_0/r`$. But, contrary to the original Schwarzschild solution, both signs of $`m_0`$ are possible now for different values of the electric charge, thus both timelike and spacelike singularities may appear. In the latter case an internal Cauchy horizon is present as in the RN metric.
In fact, within the context of string theory two types of charged black holes are known: besides the RN one there is also a dilatonic black hole with a rather different internal structure. The extremal dilatonic black holes are particularly different from the RN ones: the horizon is moved to the singularity which is therefore lightlike. Moreover, in the string frame the metric associated with the extremal magnetic dilaton black hole turns out to be perfectly regular. Thus the dilaton is also able to produce a regularizing effect on gravitational singularities. A natural question arises concerning the effect of both the BI non-linearity and the dilaton, this is the main subject of the present paper. Since already the flat space BI action gives rise to regular particle-like solutions, we also expect such solutions in the Born-Infeld-dilaton theory and its gravitating generalization. As we will see, there is a one-parameter family of such solutions which, for a limiting value of the parameter, tend (after suitable coordinate rescaling) to the near horizon limit of the extremal dilatonic black hole. The limiting solution saturates the BPS bound. Another class of solutions is black holes. We show that a generic Einstein-Born-Infeld-dilaton black hole has a timelike singularity with a power-low behavior of the local mass function.
The choice of the lagrangian is worth to be discussed in detail. We adopt the $`SL(2,R)`$ invariant version of the dilaton-axion coupled Born-Infeld action , which is a direct Born-Infeld type generalisation of the toroidally compactified heterotic string effective theory ($`N=4,D=4`$ supergravity). The main reason is that we would like to make contact with the dilatonic black holes . This version of the theory is also distinguished as the unique dilaton gereralization of the original Born-Infeld theory exhibiting a continuous electric-magnetic duality symmetry. In fact, the original BI theory as well as its gravitating generalization is symmetric under the $`SO(2)`$ electric-magnetic duality rotations. This duality can be extended to the $`SL(2,R)`$ symmetry when a dilaton and an axion are suitably added. However, the open string version of the BI action coupled to a dilaton does not enjoy electric-magnetic duality. This type of theory was discussed along similar lines in a recent paper with results differing somewhat from those presented below. Meanwhile the $`SL(2,Z)S`$-duality is now regarded as an exact symmetry of the superstring/M theory, so it is reasonable to concentrate on this version of the dilaton-coupled BI theory.
As was shown in , the unique BI-dilaton-axion action generating the $`SL(2,R)`$ invariant equations of motion reads (with $`G=\mathrm{}=c=1`$)
$`S`$ $`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle \left\{R+2(\varphi )^2+\frac{1}{2}\mathrm{e}^{4\varphi }(\kappa )^2\kappa F_{\mu \nu }\stackrel{~}{F}^{\mu \nu }+4\beta ^2(1)\right\}\sqrt{g}d^4x}`$ (1.1)
$`=`$ $`{\displaystyle L\sqrt{g}d^4x},`$
where $`\varphi `$ is a dilaton, $`\kappa `$ is an axion and
$$=\sqrt{1+\frac{F^2\mathrm{e}^{2\varphi }}{2\beta ^2}\frac{(F\stackrel{~}{F})^2\mathrm{e}^{4\varphi }}{16\beta ^4}},$$
(1.2)
with $`F^2=F_{\mu \nu }F^{\mu \nu },F\stackrel{~}{F}=\stackrel{~}{F}_{\mu \nu }F^{\mu \nu }`$. In the limit $`\beta \mathrm{}`$ this action reduces to the (truncated) heterotic string effective action in four dimensions.
In the BI theory, as in electrodynamics in media, one must distinguish between the field tensor $`F_{\mu \nu }`$ incorporating the electric strength and the magnetic induction, and the induction tensor $`G_{\mu \nu }`$ combining the electric induction and the magnetic field strength
$$G^{\mu \nu }=\frac{1}{2}\frac{L}{F_{\mu \nu }}=\frac{\mathrm{e}^{2\varphi }}{}\left(F^{\mu \nu }\frac{(F\stackrel{~}{F})}{4\beta ^2}\stackrel{~}{F}^{\mu \nu }\mathrm{e}^{2\varphi }\right)+\kappa \stackrel{~}{F}^{\mu \nu }.$$
(1.3)
The Maxwell equations in terms of $`G_{\mu \nu }`$ are sourceless
$$dG=0,G=G_{\mu \nu }dx^\mu dx^\nu ,$$
(1.4)
thus coinciding with the Bianchi identity for the field tensor
$$dF=0,F=F_{\mu \nu }dx^\mu dx^\nu .$$
(1.5)
Therefore the linear transformations
$$FaF+b\stackrel{~}{G},GcGd\stackrel{~}{F},$$
(1.6)
do not change this system of equations. Moreover, one can check that the equations for $`\varphi ,\kappa `$ following from the action (1.1) as well as the corresponding Einstein equations
$$R_{\mu \nu }\frac{1}{2}g_{\mu \nu }R=8\pi T_{\mu \nu },$$
(1.7)
$$T_{\mu \nu }=G_{\mu \lambda }F_{}^{\lambda }{}_{\nu }{}^{}+\varphi _{,\mu }\varphi _{,\nu }+\frac{1}{4}\kappa _{,\mu }\kappa _{,\nu }\mathrm{e}^{4\varphi }g_{\mu \nu }L$$
(1.8)
remain also unchanged , provided the complex axidilaton $`z=\kappa +i\mathrm{e}^{2\varphi }`$ undergoes the transformation
$$z\frac{czd}{bz+a}$$
(1.9)
with real $`a,b,c,d`$ subject to $`acbd=1`$. Note that these transformations change the action (1.1) by a total divergence.
In what follows we will consider either purely electric or purely magnetic configurations. In this case the axion decouples and can be consistently set to zero, with the term $`F\stackrel{~}{F}`$ in the lagrangian omitted. This truncated version of the theory still has a discrete electric-magnetic duality $`(b=d=1,a=c=0)`$:
$$F\stackrel{~}{G},G\stackrel{~}{F},\varphi \varphi .$$
(1.10)
The plan of the paper is as follows. In Sec. 2 we perform the dimensional reduction of the action for static spherically symmetric purely electric (magnetic) configurations, and present the system in two alternative gauges. Sec. 3 is devoted to everywhere regular solutions which are studied both analytically and numerically. Black hole solutions are discussed in Sec. 4. Our results are summarized in Sec. 5. In the Appendix, the behavior of the solutions near the singularities is discussed in more detail. The numerical calculations presented in this paper were performed in collaboration with V.V. Dyadichev.
## 2 Basic equations
We start with the action, which is a truncated (for purely electric or magnetic configurations) version of (1.1)
$$S=\frac{1}{4}\left\{\frac{R}{G}+2(\varphi )^2+4\beta ^2(1)\right\}\sqrt{g}d^4x,$$
(2.1)
where we have restored the Newton constant $`G`$, and assume for $``$ a more general form with an arbitrary dilaton coupling constant $`\gamma `$:
$$=\sqrt{1+\frac{F^2\mathrm{e}^{2\gamma \varphi }}{2\beta ^2}}.$$
(2.2)
The discrete electric-magnetic duality (1.10) continues to hold for an arbitrary $`\gamma `$. This theory has altogether three dimensional parameters: $`G`$, $`\beta `$ and $`\gamma `$ with dimensionalities (in units $`\mathrm{}=c=1`$)
$$[G]=\mathrm{L}^2,[\beta ]=\mathrm{L}^2,[\gamma ]=\mathrm{L},$$
(2.3)
We will be dealing with static spherically symmetric configurations, for which the equations of motion reduce to one-dimensional equations. Assuming the general static spherically symmetric parametrization of the metric
$$ds^2=N\sigma ^2dt^2\frac{dr^2}{N}R^2\left(d\theta ^2+\mathrm{sin}^2\theta d\phi ^2\right),$$
(2.4)
where $`N,R,\sigma `$ are functions of $`r`$ only, we perform a dimensional reduction leading to the Lagrangian
$$L=L_g+L_m$$
(2.5)
where the gravitational part is
$$L_g=\frac{1}{G}\left(N\sigma ^{}RR^{}+\frac{\sigma }{2}((NR)^{}R^{}+1)\right),$$
(2.6)
while the matter part reads
$$L_m=\frac{1}{2}\sigma NR^2\varphi ^2+\sigma \beta ^2R^2(1).$$
(2.7)
The dual induction two-form is reduced to
$$\stackrel{~}{G}=\stackrel{~}{G}_{\mu \nu }dx^\mu dx^\nu =\frac{a^{}R^2\mathrm{e}^{2\gamma \varphi }}{\sigma }\mathrm{sin}\theta d\theta d\phi $$
(2.8)
with
$$=\sqrt{1\left(\frac{a^{}\mathrm{e}^{\gamma \varphi }}{\beta \sigma }\right)^2}.$$
(2.9)
Integrating the Maxwell equation $`d\stackrel{~}{G}=0`$ one obtains
$$a^{}\mathrm{e}^{2\gamma \varphi }=\frac{QV\sigma }{R^2},=V\left(1+\frac{Q^2\mathrm{e}^{2\gamma \varphi }}{\beta ^2R^4}\right)^{1/2}.$$
(2.10)
The dilaton equation reads
$$\left(\sigma NR^2\varphi ^{}\right)^{}=\sigma R^2\gamma \beta ^2\left(\frac{1}{V}V\right).$$
(2.11)
Variation of the action over $`R,N,\sigma `$ gives the full set of Einstein equations
$$R(N\sigma ^{})^{}+\frac{1}{2}[\sigma (NR)^{}]^{}+\frac{1}{2}N(\sigma R^{})^{}=G\left(2\beta ^2\sigma R(1V)\sigma NR\varphi ^2\right)$$
(2.12)
$$R(\sigma R^{})^{}2\sigma ^{}RR^{}=G\sigma R^2\varphi ^2$$
(2.13)
$$(NRR^{})^{}\frac{1}{2}\left(R^{}(NR)^{}+1\right)=G\left(\beta ^2R^2\frac{(1V)}{V}+\frac{1}{2}NR^2\varphi ^2\right).$$
(2.14)
Actually, the metric (2.4) is invariant under reparametrizations $`r=r(\rho )`$ of the radial coordinate, so that a gauge condition may be imposed. We will use two particular gauges in what follows: one is $`R=r`$, and another is $`\sigma =1`$.
In the first gauge ($`\underset{¯}{R=r}`$),
$$ds^2=N\sigma ^2dt^2\frac{dr^2}{N}r^2\left(d\theta ^2+\mathrm{sin}^2\theta d\phi ^2\right).$$
(2.15)
Introducing a new variable $`\psi =\beta ^2\gamma ^4Q^2\mathrm{e}^{2\gamma \varphi }`$, and performing rescalings
$$\beta \gamma rr,\gamma \varphi \varphi ,$$
(2.16)
(leading to a dimensionless independent variable $`r`$), the system of equations (2.11), (2.13), (2.14) may be rewritten as
$`Nr(r\varphi ^{})^{}+r\varphi ^{}`$ $`=`$ $`{\displaystyle \frac{\psi }{\sqrt{\psi +r^4}}}+2gr\varphi ^{}\left(\sqrt{\psi +r^4}r^2\right)`$ (2.17)
$`(rN)^{}1`$ $`=`$ $`g\left(2r^22\sqrt{\psi +r^4}Nr^2\varphi ^2\right)`$ (2.18)
$`\sigma ^{}`$ $`=`$ $`g\sigma r\varphi ^2,`$ (2.19)
with the only dimensionless coupling constant
$$g=\frac{G}{\gamma ^2}.$$
(2.20)
Another useful relation is the combination $`2R`$(2.12) $`2N`$(2.13) which gives, after rescaling,
$$\left(\frac{(N\sigma ^2)^{}r^2}{\sigma }\right)^{}=4g\sigma r^2(1V).$$
(2.21)
In the second gauge ($`\underset{¯}{\sigma =1}`$),
$$ds^2=\lambda ^2dt^2\frac{d\rho ^2}{\lambda ^2}R^2\left(d\theta ^2+\mathrm{sin}^2\theta d\phi ^2\right),$$
(2.22)
the radial coordinate $`\rho `$ being related to the radial coordinate $`r`$ in the first gauge by $`|d\rho /dr|=\sigma (r)`$. We note that a suitable combination of Eqs. (2.11), (2.12) and (2.14) eliminates the source terms, leading to the equation
$$2(\lambda ^2RR^{})^{}\frac{1}{2}(\lambda ^2R^2)^{\prime \prime }+\frac{2G}{\gamma }(\lambda ^2R^2\varphi ^{})^{}1=0.$$
(2.23)
The first integral of this equation, together with equations (2.12) and (2.11), leads (after rescaling $`\rho `$ and $`R`$, as well as $`\varphi `$, as before), to the system
$`{\displaystyle \frac{R^{}}{R}}{\displaystyle \frac{\lambda ^{}}{\lambda }}+2g\varphi ^{}`$ $`=`$ $`{\displaystyle \frac{\rho \rho _0}{\lambda ^2R^2}},`$ (2.24)
$`R^{\prime \prime }`$ $`=`$ $`gR\varphi ^2,`$ (2.25)
$`(\lambda ^2R^2\varphi ^{})^{}`$ $`=`$ $`uV{\displaystyle \frac{\psi }{R^2}},`$ (2.26)
with $`\rho _0`$ an integration constant.
## 3 Regular solutions
In the gauge $`R=r`$ the conditions of asymptotic flatness read, in terms of rescaled quantities,
$`N`$ $`=`$ $`1{\displaystyle \frac{2M}{r}}+O({\displaystyle \frac{1}{r^2}}),`$ (3.1)
$`\varphi `$ $`=`$ $`\varphi _{\mathrm{}}+{\displaystyle \frac{D}{r}}+O({\displaystyle \frac{1}{r^2}}),`$ (3.2)
$`\sigma `$ $`=`$ $`1{\displaystyle \frac{gD^2}{2r^2}}+O({\displaystyle \frac{1}{r^3}}).`$ (3.3)
Here $`M`$ and $`D`$ are the mass and the dilaton charge in the rescaled variables. Alternatively, one can introduce the unscaled quantities $`\overline{M}`$ (via $`N=12G\overline{M}/r+\mathrm{}`$) and $`\overline{D}`$, then $`M=G\beta \gamma \overline{M}`$, $`D=\beta \overline{D}`$.
For these quantities one can derive useful sum rules. Define a local mass function $`m(r)`$ through
$$N=1\frac{2m(r)}{r},$$
(3.4)
with $`m(r)M+O(1/r)`$ as $`r\mathrm{}`$. Integrating Eq. (2.21) from some finite $`r_1`$ to infinity and taking into account the conditions of asymptotic flatness one obtains
$$M=2g_{r_1}^{\mathrm{}}\sigma r^2(1V)𝑑r+\frac{(N\sigma ^2)^{}r^2}{2\sigma }|_{r=r_1}.$$
(3.5)
For solutions regular at the origin the last term vanishes as $`r_10`$, so one obtains
$$M=2g_0^{\mathrm{}}\sigma r^2(1V)𝑑r.$$
(3.6)
In a similar way one can derive from the dilaton equation
$$D=_{r1}^{\mathrm{}}\sigma r^2\left(VV^1\right)𝑑r\sigma Nr\xi |_{r=r_1},$$
(3.7)
with $`\xi =r\varphi ^{}`$. Thus for the regular solutions we have
$$D=_0^{\mathrm{}}\sigma r^2\left(VV^1\right)𝑑r.$$
(3.8)
Combining these two formulas we find the relation
$$M+gD=g_0^{\mathrm{}}\sigma r^2\frac{(1V)^2}{V}𝑑r,$$
(3.9)
showing that the dilaton charge for the electric solution is negative ($`V<1`$) and that its absolute value is greater than the mass.
To analyse the behavior of regular solutions near the origin it is convenient to introduce the logarithmic variable $`\tau =\mathrm{ln}r`$. The system (2.17)–(2.19) may be rewritten as the system of first order equations
$`\dot{\psi }`$ $`=`$ $`2\xi \psi ,`$ (3.10)
$`N\dot{\xi }`$ $`=`$ $`u\xi v,`$ (3.11)
$`\dot{N}`$ $`=`$ $`vN(1+g\xi ^2),`$ (3.12)
$`\dot{\sigma }`$ $`=`$ $`g\sigma \xi ^2,`$ (3.13)
where
$$u=\frac{\psi }{\sqrt{\psi +r^4}},v=1+2g\left(r^2\sqrt{\psi +r^4}\right).$$
(3.14)
By definition, the regular solution should be flat in the vicinity of the origin $`r=0`$ ($`\tau \mathrm{}`$), i.e. $`N=1`$ at the origin. Under this condition one can find the series solution in powers of $`1/\tau `$ with logarithmic coefficients
$`\psi `$ $`=`$ $`{\displaystyle \frac{1}{\tau ^2}}\left(1+{\displaystyle \frac{2L}{\tau }}+{\displaystyle \frac{3L^22(2g1)L+4(1g)}{\tau ^2}}+\mathrm{}\right),`$ (3.15)
$`\xi `$ $`=`$ $`{\displaystyle \frac{1}{\tau }}\left(1+{\displaystyle \frac{L+2g1}{\tau }}{\displaystyle \frac{L^23(2g1)L+4(1g)^2+1}{\tau ^2}}+\mathrm{}\right),`$ (3.16)
$`N`$ $`=`$ $`1+{\displaystyle \frac{2g}{\tau }}\left(1+{\displaystyle \frac{L+1/2}{\tau }}+{\displaystyle \frac{L^2+2(1g)(L+3/2)}{\tau ^2}}+\mathrm{}\right).`$ (3.17)
In these formulas
$$L=(2g1)\mathrm{ln}(\tau )+c,$$
(3.18)
where $`c`$ is a free parameter. One can also show from Eq. (2.21) that $`(N\sigma ^2)\dot{}(4g/3)\sigma e^{2\tau }`$ is vanishingly small, so that $`\lambda ^2=N\sigma ^2`$ is constant to all orders in this expansion. Two particular values of $`g`$ should be mentioned. The first one is $`g=0`$ which corresponds to a vanishing Newton constant. In this case the expansion (3.17) generates $`N1`$, so we get the flat space dilaton-BI solution. The second is $`g=1/2`$, in this case there are no logarithmic terms in the expansions.
These local solutions being continued (numerically) to large $`r`$ meet the conditions of asymptotic flatness for continuously varying $`c`$ subject to the condition
$$c>c_{cr}(g).$$
(3.19)
The dependence of $`c_{cr}`$ on the effective coupling constant $`g`$ is shown in Fig.1. The limit $`g=0`$ corresponds to the flat space solutions, in this case $`N1`$. The corresponding value is approximately $`c_{cr}=1.52`$. When $`g`$ is large enough the critical value moves to the negative half-plane. The variation with $`c_{cr}`$ of the mass, dilaton charge and the asymptotic value of $`\psi `$ are shown in Fig. 2. It is worth emphasizing that we deal with a one-parameter family, i.e. for a given mass the values of the dilaton charge and the electric charge (proportional to the square root of $`\psi _{\mathrm{}}`$) are fixed. Physically these particle-like solutions can be regarded as the original Born-Infeld field distributions of the point charge deformed by gravity and by the dilaton. Thus the mere existence of regular particle-like solutions in the present theory is by no means surprising.
When the parameter $`c`$ decreases and approaches the critical value, both the mass and the absolute value of the dilaton charge increase, with an approximate power law dependence $`M(cc_{cr})^p`$, where $`p=1/(1+g)`$ (see Fig. 2), and at the same time
$$g|D|M\mathrm{as}cc_{cr}.$$
(3.20)
This can be understood qualitatively using the sum rules above. When $`c`$ is close to its boundary value the main contribution to the integral comes from the region of large $`r`$ where $`V1`$. Then from Eqs. (3.6), (3.8), (3.9) it is clear that if the integral for $`M`$ diverges for $`r\mathrm{}`$, the integral for the sum $`M+gD`$ is less divergent than both $`M`$ and $`|D|`$. The numerical results also show that the ratio $`\psi (\mathrm{})/MD`$ goes to a finite limit as $`cc_{cr}`$. To understand this, note that a generic asymptotically flat solution of the non linear theory is asymptotic for $`r\mathrm{}`$ ($`V1`$) to an electric solution of the linear Einstein-Maxwell-dilaton theory (obtained from the magnetic solution of by transforming $`\varphi \varphi `$ and rescaling $`\rho `$, $`R`$ and $`\varphi `$)
$`\psi `$ $`=`$ $`{\displaystyle \frac{\rho _+\rho _{}}{g+1}}\left(1{\displaystyle \frac{\rho _{}}{\rho }}\right)^{2/(g+1)},R=\rho \left(1{\displaystyle \frac{\rho _{}}{\rho }}\right)^{1/(g+1)},`$
$`\lambda ^2`$ $`=`$ $`\left(1{\displaystyle \frac{\rho _+}{\rho }}\right)\left(1{\displaystyle \frac{\rho _{}}{\rho }}\right)^{(g1)/(g+1)}`$ (3.21)
($`g=\gamma ^2`$ for $`G=1`$). Comparing with (3.1), we see that $`M+gD=(\rho _+\rho _{})/2`$, so that the limit $`M+gDM`$ corresponds to the BPS limit $`\rho _{}\rho _+`$. In this limit, we obtain from (3)
$$\psi (\mathrm{})(g+1)\left(\frac{M}{g}\right)^2.$$
(3.22)
This relation is well confirmed by the numerical observations (Fig. 2). Relations (3.22) and (3.20) together imply
$$M^2+g^2D^2=Q_{eff}^2,Q_{eff}=g\left(\frac{2}{g+1}\psi (\mathrm{})\right)^{1/2},$$
(3.23)
which corresponds to the BPS saturation of the Einstein-Maxwell-dilaton theory .
The transition point in the parameter space $`c=c_{cr}`$ corresponds to a limiting solution which is not asymptotically flat. To find its asymptotic behaviour, consider the asymptotic solution (3) in the BPS limit $`\rho _{}\rho _+`$. The limit in which both $`\rho `$ and $`M=g\rho _+/(g+1)`$ go to infinity corresponds to the near horizon limit $`|\rho \rho _+|\rho _+`$. Putting $`(\rho \rho _+)/\rho _+=(r/\rho _+)^{g+1}`$, rescaling $`t`$, and taking the limit $`\rho _+\mathrm{}`$, we arrive at the asymptotic critical solution
$$ds_{cr}^2=a^2r^{2g}dt^2(g+1)^2dr^2r^2d\mathrm{\Omega }^2,\psi _{cr}=r^2/(g+1).$$
(3.24)
This asymptotic behaviour has not been directly tested, as in all our numerical computations the parameter $`c`$ was at best only approximately equal to its critical value. Nevertheless the critical solution can be approached numerically as the envelope of asymptotically flat regular solutions with $`cc_{cr}`$. In this respect, we note that the curves for $`\xi (\tau )`$ show increasingly large and flat maxima near $`\xi =1`$ (before dropping again to 0 as $`\tau \mathrm{}`$), in agreement with the asymptotic critical behaviour $`\xi _{cr}(\mathrm{})=1`$ from (3.24).
For $`c<c_{cr}`$ the solution starting at the origin as (3.15,3.16,3.17) develops a coordinate singularity at some radius $`r=r^{}`$ . In the gauge $`R=r`$ one observes that for $`rr^{}`$ the metric function $`N`$ approaches zero while the dilaton has a square root singularity
$$N(r^{}r),\varphi \sqrt{r^{}r},\xi \frac{1}{\sqrt{r^{}r}}.$$
(3.25)
Such a behaviour can be derived by defining locally a new radial coordinate $`x`$:
$$r=r^{}\mathrm{exp}(gx^2/2),$$
(3.26)
and assuming $`\psi \psi ^{}(12bx)`$, where $`b`$ is some constant. Then $`\xi b/gx`$, and Eq. (3.11) yields $`Ngv^{}x^2`$, with $`v^{}<0`$ ($`g\psi ^{}>r^2+1/4g`$) by continuity, while Eq. (3.12) yields $`b^2=1`$, so that $`b=\pm 1`$. The complete local solution
$`\xi `$ $`=`$ $`\pm (1/gx)+d+\mathrm{},`$
$`N`$ $`=`$ $`gx^2(v^{}g(u^{}+dv^{})x+\mathrm{}),`$ (3.27)
$`\sigma `$ $`=`$ $`k(1/x2gd+\mathrm{})`$
depends on the four integration constants $`r^{}`$, $`\psi ^{}`$ $`d`$ and $`k`$, and so is generic.
Obviously the coordinate patch $`r<r^{}`$ fails to describe the full solution. The extension may be achieved by going to the second gauge (2.22), as from (3.26) and the last equation (3) $`d\rho dx`$ near $`Rr=r^{}`$. From Eq. (2.26), $`R^2\lambda ^2\varphi ^{}`$ is an increasing function of $`\rho `$. So if one starts from a solution regular near $`R=0`$ ($`\varphi _\rho ^{}\varphi _r^{}>0`$) and increases $`R`$ and $`\rho `$, one must have $`\varphi _\rho ^{}>0`$ at $`R=r^{}`$, with $`R_\rho ^{}=0`$. This is a maximum of $`R`$ from Eq. (2.25). As $`\rho `$ is further increased, $`R`$ must decrease to zero, while $`\varphi _\rho ^{}`$ stays positive so that $`\varphi `$ continues to increase, and a regular solution at $`R=0`$ ($`\varphi \mathrm{}`$) cannot be achieved. The resulting solution is compact and singular (‘bag of gold’), with the power law behaviours near the singularity
$`\psi r^{2\delta },`$ $`Nr^{1g\delta ^2},`$ $`\sigma r^{g\delta ^2}(g>1/4),`$ (3.28)
$`\psi r^{1/g},`$ $`Nr^{1/2g},`$ $`\sigma r^{1/4g}(g<1/4),`$ (3.29)
where $`\delta <0`$ is bounded by
$$\sqrt{4+3g^1}2<\delta <\sqrt{g^1}.$$
(3.30)
(details on the derivation of these behaviours and bounds are given in the Appendix).
Let us consider now magnetically charged solitons. Reverting to unscaled variables and introducing a magnetic charge via
$$F=P\mathrm{sin}\theta d\theta d\phi $$
(3.31)
we obtain
$$=\sqrt{1+\frac{P^2\mathrm{e}^{2\gamma \varphi }}{\beta ^2R^2}}$$
(3.32)
Rescaling the variables in the gauge $`R=r`$ we obtain the same system (2.17)–(2.19) with $`\xi \xi ,QP`$. This is in a perfect agreement with the electric-magnetic duality discussed in the Introduction. Accordingly, the electric solitons described in this section can be reinterpreted as magnetic solitons. In this case the parameter $`\delta `$ in (3.28) is positive, and it is easy to see that the bounds above still hold if $`\delta `$ is replaced by $`|\delta |`$. It is worth noting that the above symmetry holds in the Einstein frame, in the conformally related ‘string’ frame the metrics of electric and magnetic solutions are essentially different.
## 4 Black holes
To study black holes, let us assume the existence of a non-degenerate horizon at some $`r=r_h`$, i.e. $`N(r_h)=0`$, with $`N^{}(r_h)>0`$. It is again convenient to work with the rescaled first order differential system (3.10-3.13) with the logarithmic variable $`\tau =\mathrm{ln}(r/r_h)`$. From this system we find the following series solution
$`N`$ $`=`$ $`v_h\tau +{\displaystyle \frac{1}{2v_h}}(v_1v_hv_h^2gu_h^2)\tau ^2+O(\tau ^3),`$ (4.1)
$`\xi `$ $`=`$ $`{\displaystyle \frac{u_h}{v_h}}+{\displaystyle \frac{1}{2v_h^2}}(v_hu_1v_1u_h)\tau +O(\tau ^2),`$ (4.2)
$`\psi `$ $`=`$ $`b+{\displaystyle \frac{2bu_h}{v_h}}\tau +O(\tau ^2),`$ (4.3)
Here $`u_h=u(r_h)`$,
$$v_1=2g\left(2r_h^2\frac{bu_h+2v_hr_h^4}{v_h\sqrt{b+r_h^4}}\right),u_1=\frac{b}{(b+r_h^4)^{3/2}}\left(2r_h^4\left(\frac{u_h}{v_h}1\right)+b\frac{u_h}{v_h}\right),$$
(4.4)
and $`b=\psi (r_h)`$ is a free parameter varying in the finite interval
$$r_h^4<b<\frac{1}{g}\left(\frac{1}{4g}+r_h^2\right).$$
(4.5)
The left bound corresponds to the condition of positivity under the square root in $`V`$, while the right bound comes from the assumption $`v_h>0`$.
Actually the right bound on the parameter $`b`$ for asymptotically flat solutions is lower, namely
$$b<b_{cr}(g)<\frac{1}{g}\left(\frac{1}{4g}+r_h^2\right),$$
(4.6)
otherwise the bag of gold type singularity is met. The critical $`b`$ depends on the horizon radius as well. On Fig. 4 the curves $`b_{cr}(g)`$ are shown for various values of $`r_h`$. For small $`r_h`$ the exterior black hole solutions approach regular solutions. As in the case of regular solutions with $`c`$ close to $`c_{cr}`$, both the mass and the dilaton charge grow with $`b`$. From the Eqs. (3.5) and (3.7) one can see that the boundary term in the black hole case vanishes for the dilaton but remains finite for the mass:
$$M=M_0+2g_{r_h}^{\mathrm{}}\sigma r^2(1V)𝑑r,M_0=\frac{v_h\sigma _hr_h^2}{2}$$
(4.7)
Here the first term may be regarded as the ’bare’ mass of the black hole, and the second as a contribution from the black hole hair. The sum rule for the dilaton charge preserves its form (3.8), where the integration now is performed over the exterior space:
$$D=_{r_h}^{\mathrm{}}\sigma r^2\left(VV^1\right)𝑑r$$
(4.8)
The combined sum rule now reads
$$MM_0+gD=g_{r_h}^{\mathrm{}}\sigma r^2\frac{(1V)^2}{V}𝑑r.$$
(4.9)
When $`b_{cr}`$ is approached this quantity remains finite while $`M`$ and $`D`$ diverge, so that asymptotically $`g|D|`$ tends to the field mass $`MM_0`$.
From Eq. (2.25), $`R(\rho )`$ must vanish for some $`\rho _s<\rho _h`$, leading to a singularity. There can be no inner Cauchy horizon, so that this singularity must be timelike. To show this, note that from Eq. (2.26), $`\lambda ^2R^2\varphi ^{}`$ must decrease when $`\rho `$ decreases inside the event horizon $`\rho =\rho _h`$, and thus must stay negative ($`\varphi _h^{}`$ is positive, while $`\lambda ^2<0`$ for $`r<r_h`$), so that $`\lambda ^2`$ cannot vanish again. Reversing the sign of $`\tau `$ in the Eqs. (3.10)–(3.13) one can integrate inside the black hole up to the singularity. The leading terms near the singularity are the same as in Eq. (3.28), with now $`\delta `$ positive for electric black holes and negative for magnetic black holes,
$$m\frac{\mu }{r^{g\delta ^2}},\xi \delta ,\psi \nu ^2r^{2\delta },\sigma ar^{g\delta ^2}$$
(4.10)
where $`\delta `$, $`\mu >0,\nu >0`$ and $`a>0`$ are free parameters. The solution (4.10) is generic, so starting from the expansions (4.1)–(4.3) inside the horizon one always meets a member of the local family (4.10) with certain $`\mu ,\nu ,a,\delta `$.
The parameter $`\delta `$ is constrained by various bounds (see the Appendix). One of these is obtained from the combined sum rule similar to (4.9), but with the integration covering the internal space,
$$\mu a(1g\delta ^2+2g|\delta |)M_0=g_0^{\rho _h}R^2\frac{(1V)^2}{V}𝑑\rho .$$
(4.11)
On account of the positivity of $`M_0`$ and $`\mu a`$, this implies $`1g\delta ^2+2g|\delta |>0`$, i.e.
$$|\delta |<1+\sqrt{1+g^1}.$$
(4.12)
Another bound, obtained by combining the sum rule (4.11) with the first integral (2.24) evaluated at the singularity, is
$$g|\delta |<\rho _h/2\mu a1.$$
(4.13)
In the limit $`g\mathrm{}`$, corresponding to $`\gamma 0`$, the dilaton decouples, and the bound (4.13) leads, for fixed ($`\rho _h`$, $`\mu `$, $`a`$) (note that the ratio $`\rho _h/\mu a`$ is invariant under the rescaling (2.16)), to $`\delta =O(g^1)`$, so that the exponents in (4.10) go to zero, and the mass function goes to a constant at the singularity $`r=0`$, consistently with the results of for Einstein-Born-Infeld black holes <sup>1</sup><sup>1</sup>1Note that (4.11) implies $`\mu >0`$ in the limit $`g\mathrm{}`$. Here there is only one horizon for finite $`g`$, so the two-horizon case ($`\mu <0`$) cannot be obtained as a limit.. Thus the dilaton reinforces a divergence of the mass function at the singularity.
## 5 Conclusion
Let us summarize our results. We have shown that the Einstein-Born-Infeld-dilaton theory admits a one-parameter family of particle-like globally regular solutions characterized by their mass, electric (magnetic) charge, and dilaton charge. The unique parameter determining these three quantities is bounded from below, and when the boundary is approached, the charges exhibit a BPS saturation similar to that of the BPS black holes of Einstein-Maxwell-dilaton theory (recall that the latter does not possess particle-like solutions). In this limit the absolute values of charges rise indefinitely, as well as the effective radius of the particle, so that the main contribution to the total charges comes from large radii. Since at large radii the Born-Infeld non-linearity is small, the correspondence with the Maxwellian counterpart of the theory is by no means surprising. However the BPS saturation property of regular solutions at large masses is not obvious a priori.
Another type of solutions regular at the origin are compact, with a point singularity at the ‘other end’ (such configurations are commonly called ‘bags of gold c’). For these solutions, the areas of two-sphere sections increase up to some finite limiting value and then decrease again up to the final curvature singularity. These solutions also form a one-parameter sequence, corresponding to values of the parameter below the threshold for particle-like solutions.
We have also found a two-parameter family of black-hole solutions for arbitrary values of the horizon radius and the second free parameter varying on a finite interval. Alternatively, one can consider as independent parameters the mass and the electric (magnetic) charge of the black hole. The dilaton charge is a dependent quantity like in the linear Einstein-Maxwell-dilaton theory. Altogether these three quantities satify a BPS inequality which is saturated in the limit of the infinite mass and the dilaton charge. The black holes do not possess internal Cauchy horizons and the singularity inside them is always spacelike. This has to be contrasted with the Born-Infeld black holes without dilaton in which case the solutions possessing an internal Cauchy horizon still exist. Like in the linear Maxwell-dilaton case, the scalar field prevents the formation of an internal horizon. The local mass function has a power-low behaviour near the singularity with a power index depending on the global black hole parameters. Solutions with a regular event horizon exist also in the bag of gold (compact) form, these are doubly singular.
The above picture holds in the Einstein frame. One should inquire how the regular solutions look in the ’string’ frame which is related to the present one by a conformal transformation with the conformal factor $`\psi `$. According to (3.15), this conformal factor vanishes logarithmically as $`r0`$ for electric solutions (and diverges for magnetic ones), so the ’string’ metric in the vicinity of the origin is singular. Nevertheless, the ADM mass remains finite in this frame both for electric and magnetic solutions.
It is worth comparing our results with those of the paper where another type of coupling of the Born-Infeld theory to dilaton was assumed. In the latter case there is no symmetry between electric and magnetic solutions, in particular, the magnetic non-singular particle-like solutions were found in the string frame (contrary to electric ones). Our version of the theory is (classically) $`S`$-dual and, moreover, it inherits the structure of the $`𝒩=4,D=4`$ supergravity. Therefore in the Einstein frame magnetic solutions are related to electric ones simply by reversing the sign of the dilaton. The string metric is singular at the origin in both cases.
Acknowledgements
We would like to thank V.V. Dyadichev for assistance in numerical calculations. D.G. is grateful to Laboratoire de Gravitation et Cosmologie Relativiste, Universitè Paris-6 for hospitality and to CNRS for support while the work was initiated. His work was also partially supported by the Russian Foundation for Basic Research under the grant 00-02-16306.
## Appendix: Power law behaviours near the singularity
Let us first consider the electric bag of gold of Sect. 3 with a regular origin and a point singularity. Assume the power law behaviours
$$\psi \nu ^2r^{2\delta },\sigma ar^{g\delta ^2}$$
(A.1)
(with the integration constants $`\delta <0`$, $`\nu >0`$, $`a>0`$). Then, if
$$g\delta ^2+\delta +1>0,$$
(A.2)
Eqs. (2.17)-(2.18) are solved near the singularity $`r=0`$ by
$$N2\mu r^{1g\delta ^2},$$
(A.3)
with $`\mu <0`$. Several bounds on the exponent $`\delta `$ can be derived. First, Eq. (2.24) where $`\rho _0=(\lambda ^2R/\sigma )(1+2g\xi )|_{r=0}=0`$ for regular solutions gives, near the singularity $`\rho =\rho _1`$,
$$g\delta ^24g\delta 3=\rho _1/\mu a>0,$$
(A.4)
leading to the bound
$$\delta >\sqrt{4+3g^1}2.$$
(A.5)
Other bounds can be obtained from sum rules obtained in a manner similar to (3.6), (3.8), (3.9). These sum rules are
$`\mu a(g\delta ^21)`$ $`=`$ $`2g{\displaystyle _0^{\rho _1}}R^2(1V)𝑑\rho ,`$ (A.6)
$`2\mu a\delta `$ $`=`$ $`{\displaystyle _0^{\rho _1}}R^2(V^1V)𝑑\rho ,`$ (A.7)
$`\mu a(1g\delta ^2+2g\delta )`$ $`=`$ $`g{\displaystyle _0^{\rho _1}}R^2{\displaystyle \frac{(1V)^2}{V}}𝑑\rho ,`$ (A.8)
where all the right-hand sides are finite and positive. From Eqs. (A.6) and (A.8) we obtain the bounds
$$\sqrt{1+g^1}1<\delta <\sqrt{g^1},$$
(A.9)
while Eq. (A.7) gives nothing new. It is easy to show that the lower bound (A.9) is weaker than the bound (A.5), so that our final bounds are
$$\sqrt{4+3g^1}2<\delta <\sqrt{g^1}.$$
(A.10)
A corollary of Eq. (A.10) is that the singular behaviour (3.28) can exist only for
$$g>1/4,$$
(A.11)
which ensures the inequality (A.2). Also, the combination of (A.4) and (A.8) leads to the inequality
$$g\delta <1\rho _1/2\mu a.$$
(A.12)
Assume now the power law behaviours (A.1) with
$$g\delta ^2+\delta +1<0,$$
(A.13)
which is possible only for
$$g<1/4.$$
(A.14)
Then the behaviour
$$N\alpha r^\delta $$
(A.15)
solves Eq. (2.18) if $`\alpha =2g\nu /(1+\delta +g\delta ^2)>0`$, and Eq. (2.17) if $`1+2g\delta =0`$, leading to the behaviours (3.29). Eq. (2.24) does not give new information, neither do the sum rules analogous to (A.7) and (A.8) whose both sides diverge with the same power law and sign, while the sum rule analogous to (A.6) shows that the behaviour of $`\lambda ^2`$ near the singularity must be $`\lambda ^2\alpha a^2\beta (\rho _1\rho )^\gamma `$, with $`\gamma =(14g)/(1+4g)>0`$, $`\beta >0`$. For $`g=1/4`$, (A.15) is replaced by $`N(\nu /2)r^2\mathrm{ln}r`$.
Consider now the black hole case. Let us assume again the behaviours (A.1), with $`\delta >0`$ for electric black holes. In this case $`g\delta ^2+|\delta |+1>0`$ for all $`g0`$, so that the behaviour of $`N`$ is always given by (A.3) with $`\mu >0`$. As in the case of the bag of gold singularity, we can derive various bounds on $`|\delta |`$. On the horizon Eq. (2.24) gives
$$\rho _0\rho _h=r_h^2(\lambda ^2)^{}(\rho _h)/2=M_0>0.$$
(A.16)
Applying now the same equation (2.24) near the singularity $`\rho r^{g\delta ^2+1}0`$ and inserting the preceding result, we obtain to leading order
$$g\delta ^24g\delta 3=\frac{(\rho _h+M_0)}{\mu a}.$$
(A.17)
The negativity of the right-hand side leads to the bound
$$\delta <2+\sqrt{4+3g^1}.$$
(A.18)
On the other hand, the sum rules obtained as in the bag of gold case give
$`\mu a(g\delta ^21)+M_0`$ $`=`$ $`2g{\displaystyle _0^{\rho _h}}R^2(1V)𝑑\rho ,`$ (A.19)
$`2\mu a\delta `$ $`=`$ $`{\displaystyle _0^{\rho _h}}R^2(V^1V)𝑑\rho ,`$ (A.20)
$`\mu a(1g\delta ^2+2g\delta )M_0`$ $`=`$ $`g{\displaystyle _0^{\rho _h}}R^2{\displaystyle \frac{(1V)^2}{V}}𝑑\rho .`$ (A.21)
This last sum rule (with $`M_0`$ and $`\mu a`$ positive) leads to the bound
$$\delta <1+\sqrt{1+g^1},$$
(A.22)
which is stronger than the bound (A.18). Again, the combination of (A.17) and (A.21) leads to the inequality
$$g\delta <\rho _h/2\mu a1.$$
(A.23)
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# Simple predictions from ALCORc for rehadronisation of charmed quark matter
(July 22. 2000)
## Abstract
We study the production of charmed hadrons with the help of ALCOR<sub>c</sub>, the algebraic coalescence model for rehadronisation of charmed quark matter. Mesonic ratios are introduced as factors connecting various antibaryon to baryon ratios. The resulting simple relations could serve as tests of quark matter formation and coalescence type rehadronization in heavy ion collisions.
Charm hadron production has gained an enhanced attention in relativistic heavy ion collisions at CERN SPS (Super Proton Synchrotron). The measured anomalous suppression of the $`J/\psi `$ in Pb+Pb collision is considered as one of the strongest candidates for an evidence of quark-gluon plasma (QGP) formation in Pb+Pb collision at 158 GeV/nucleon bombarding energy . So far only the $`J/\psi `$ and $`\psi ^{}`$ production was measured in heavy ion collisions through dilepton decay channels. However, recent efforts to measure D-meson support the theoretical investigation of charm production from a different point of view. Namely, it is interesting to search for predictions on the total numbers of charmed hadrons and their ratios. The answer to this question may become very important at the RHIC accelerator, where a large number of charmed quark-antiquark pairs will be produced and a number of different charmed hadrons could be detected.
In this paper we assume that quark matter is formed in heavy ion collisions and the charm hadrons will be produced directly from this state via quark coalescence. Quark coalescence was successfully applied to describe direct hadron production from deconfined quark matter phase (see. the ALCOR , the Transchemistry and the MICOR models). In these models the hadronic rescatterings are assumed to be weak and they are neglected. Thus the results of quark coalescence processes were compared directly to the experimental data - and the agreement was remarkably good.
Multicharm hadron production was already investigated in a simplified quark coalescence model and first results were obtained at RHIC and LHC energies, where an appreciable number of charm quark may appear . Here we summarize simple predictions from a non-linear algebraic coalescence model ALCOR<sub>c</sub>, the extension of the ALCOR model of algebraic coalescence of strange quark matter for the inclusion of charmed quarks, mesons and baryons.
The description of the charmed baryons has to deal with the fact that two possible $`(1/2)^+`$ baryon multiplets exist containing $`c`$, $`s`$ and $`u`$ (or $`d`$) quarks, one being flavor symmetric under $`s`$ and $`d`$ (or $`u`$) exchange and the other being antisymmetric . The heavier (symmetric) states decay into the lighter (antisymmetric) one by emission of a $`\gamma `$ or a $`\pi `$ meson. However, if quark clusterization is the basic hadronization process, then the effect of these decay processes will be cancelled from charmed antibaryon to baryon ratios. Neglecting the difference between the light $`u`$ and $`d`$ quarks and using the notation $`q`$ for them, the 10 different types of produced quark clusters can be connected to the 40 lowest lying SU(4)-flavor baryon species in the following way (see e.g. Ref. for precise quark content, hadron names and masses):
$`N(qqq)`$ $`:=`$ $`p,n,\mathrm{\Delta }^{++},\mathrm{\Delta }^+,\mathrm{\Delta }^0,\mathrm{\Delta }^{};`$
$`Y(qqs)`$ $`:=`$ $`\mathrm{\Lambda }^0,\mathrm{\Sigma }^+,\mathrm{\Sigma }^0,\mathrm{\Sigma }^{},\mathrm{\Sigma }^+,\mathrm{\Sigma }^0,\mathrm{\Sigma }^{};`$
$`\mathrm{\Xi }(qss)`$ $`:=`$ $`\mathrm{\Xi }^0,\mathrm{\Xi }^{},\mathrm{\Xi }^0,\mathrm{\Xi }^{};`$
$`\mathrm{\Omega }(sss)`$ $`:=`$ $`\mathrm{\Omega }^{};`$
$`Y_c(qqc)`$ $`:=`$ $`\mathrm{\Lambda }_c^+,\mathrm{\Sigma }_c^{++},\mathrm{\Sigma }_c^+,\mathrm{\Sigma }_c^0,\mathrm{\Sigma }_c^{++},\mathrm{\Sigma }_c^+,\mathrm{\Sigma }_c^0;`$
$`\mathrm{\Xi }_c(qsc)`$ $`:=`$ $`\mathrm{\Xi }_c^+,\mathrm{\Xi }_c^0,\mathrm{\Xi }_{c}^{}{}_{}{}^{}{}_{}{}^{+},\mathrm{\Xi }_{c}^{}{}_{}{}^{}{}_{}{}^{0},\mathrm{\Xi }_c^+,\mathrm{\Xi }_c^0;`$
$`\mathrm{\Omega }_c(ssc)`$ $`:=`$ $`\mathrm{\Omega }_c^0,\mathrm{\Omega }_c^0;`$
$`\mathrm{\Xi }_{cc}(qcc)`$ $`:=`$ $`\mathrm{\Xi }_{cc}^+,\mathrm{\Xi }_{cc}^{++},\mathrm{\Xi }_{cc}^+,\mathrm{\Xi }_{cc}^{++};`$
$`\mathrm{\Omega }_{cc}(scc)`$ $`:=`$ $`\mathrm{\Omega }_{cc}^+,\mathrm{\Omega }_{cc}^+;`$
$`\mathrm{\Omega }_{ccc}(ccc)`$ $`:=`$ $`\mathrm{\Omega }_{ccc}^{++}`$ (1)
In ALCOR, the algebraic coalescence model of rehadronization it is assumed that the number of directly produced hadrons is given by the product of the the number of quarks (or anti-quarks) from which those hadrons are produced, multiplied by coalescence coefficients $`C_h`$ and by non-linear normalization coefficients $`b_q`$, that take into account conservation of quark numbers during quark coalescence, as will be explained subsequently. The number of various hadrons and quarks is denoted by the symbol usual for that type of particles, e.q. $`q`$, $`s`$ and $`c`$ denote the number of light, strange and charmed quarks, respectively, $`N`$ denotes the number of protons, neutrons and deltas etc.
In this way the baryons and antibaryons can be described through the following clustering relations:
$`N(qqq)`$ $`=`$ $`C_N(b_qq)^3\overline{N}(\overline{q}\overline{q}\overline{q})=C_{\overline{N}}(b_{\overline{q}}\overline{q})^3`$
$`Y(qqs)`$ $`=`$ $`C_Y(b_qq)^2(b_ss)\overline{Y}(\overline{q}\overline{q}\overline{s})=C_{\overline{Y}}(b_{\overline{q}}\overline{q})^2(b_{\overline{s}}\overline{s})`$
$`\mathrm{\Xi }(qss)`$ $`=`$ $`C_\mathrm{\Xi }(b_qq)(b_ss)^2\overline{\mathrm{\Xi }}(\overline{q}\overline{s}\overline{s})=C_{\overline{\mathrm{\Xi }}}(b_{\overline{q}}\overline{q})(b_{\overline{s}}\overline{s})^2`$
$`\mathrm{\Omega }(sss)`$ $`=`$ $`C_\mathrm{\Omega }(b_ss)^3\overline{\mathrm{\Omega }}(\overline{s}\overline{s}\overline{s})=C_{\overline{\mathrm{\Omega }}}(b_{\overline{s}}\overline{s})^3`$
$`Y_c(qqc)`$ $`=`$ $`C_Y^c(b_qq)^2(b_cc)\overline{Y}_c(\overline{q}\overline{q}\overline{c})=C_{\overline{Y}}^c(b_{\overline{q}}\overline{q})^2(b_{\overline{c}}\overline{c})`$
$`\mathrm{\Xi }_c(qsc)`$ $`=`$ $`C_\mathrm{\Xi }^c(b_qq)(b_ss)(b_cc)\overline{\mathrm{\Xi }}_c(\overline{q}\overline{s}\overline{c})=C_{\overline{\mathrm{\Xi }}}^c(b_{\overline{q}}\overline{q})(b_{\overline{s}}\overline{s})(b_{\overline{c}}\overline{c})`$
$`\mathrm{\Omega }_c(ssc)`$ $`=`$ $`C_\mathrm{\Omega }^c(b_ss)^2(b_cc)\overline{\mathrm{\Omega }}_c(\overline{s}\overline{s}\overline{c})=C_{\overline{\mathrm{\Omega }}}^c(b_{\overline{s}}\overline{s})^2(b_{\overline{c}}\overline{c})`$
$`\mathrm{\Xi }_{cc}(qcc)`$ $`=`$ $`C_\mathrm{\Xi }^{cc}(b_qq)(b_cc)^2\overline{\mathrm{\Xi }}_{cc}(\overline{q}\overline{c}\overline{c})=C_{\overline{\mathrm{\Xi }}}^{cc}(b_{\overline{q}}\overline{q})(b_{\overline{c}}\overline{c})^2`$
$`\mathrm{\Omega }_{cc}(scc)`$ $`=`$ $`C_\mathrm{\Omega }^{cc}(b_ss)(b_cc)^2\overline{\mathrm{\Omega }}_{cc}(\overline{s}\overline{c}\overline{c})=C_{\overline{\mathrm{\Omega }}}^{cc}(b_{\overline{s}}\overline{s})(b_{\overline{c}}\overline{c})^2`$
$`\mathrm{\Omega }_{ccc}(ccc)`$ $`=`$ $`C_\mathrm{\Omega }^{ccc}(b_cc)^3\overline{\mathrm{\Omega }}_{ccc}(\overline{c}\overline{c}\overline{c})=C_{\overline{\mathrm{\Omega }}}^{ccc}(b_{\overline{c}}\overline{c})^3`$ (2)
Mesons in the pseudoscalar and vector SU(4)-flavor multiplets are grouped in the following way:
$`\pi (q\overline{q})`$ $`:=`$ $`\pi ^+,\pi ^0,\pi ^{},\eta ,\rho ^+,\rho ^0,\rho ^{},\omega ;`$
$`K(q\overline{s})`$ $`:=`$ $`K^+,K^0,K^+,K^0;`$
$`\overline{K}(\overline{q}s)`$ $`:=`$ $`K^{},\overline{K}^0,K^{},\overline{K}^0;`$
$`\varphi (s\overline{s})`$ $`:=`$ $`\eta ^{},\varphi ;`$
$`D(\overline{q}c)`$ $`:=`$ $`D^+,D^0,D^+,D^0;`$
$`\overline{D}(q\overline{c})`$ $`:=`$ $`D^{},\overline{D}^0,D^{},\overline{D}^0;`$
$`D_s(\overline{s}c)`$ $`:=`$ $`D_s^+,D_s^+;`$
$`\overline{D}_s(s\overline{c})`$ $`:=`$ $`D_s^{},D_s^{};`$
$`J/\psi (\overline{c}c)`$ $`:=`$ $`\eta _c,J/\psi ;`$ (3)
Thus the number of directly produced mesons reads as
$`\pi (q\overline{q})`$ $`=`$ $`C_\pi (b_qq)(b_{\overline{q}}\overline{q})J/\psi (c\overline{c})=C_{J/\psi }(b_cc)(b_{\overline{c}}\overline{c})`$
$`K(q\overline{s})`$ $`=`$ $`C_K(b_qq)(b_{\overline{s}}\overline{s})D(\overline{q}c)=C_D(b_{\overline{q}}\overline{q})(b_cc)`$
$`\overline{K}(\overline{q}s)`$ $`=`$ $`C_{\overline{K}}(b_{\overline{q}}\overline{q})(b_ss)\overline{D}(q\overline{c})=C_{\overline{D}}(b_qq)(b_{\overline{c}}\overline{c})`$
$`\varphi (s\overline{s})`$ $`=`$ $`C_\varphi (b_ss)(b_{\overline{s}}\overline{s})D_s(\overline{s}c)=C_D^s(b_{\overline{s}}\overline{s})(b_cc)`$ (4)
$`\overline{D}_s(s\overline{c})=C_{\overline{D}}^s(b_ss)(b_{\overline{c}}\overline{c})`$
As a straightforward extension to the ALCOR model, the non-linear coalescence factors $`b_q`$, $`b_s`$, $`b_c`$ and the $`b_{\overline{q}}`$, $`b_{\overline{s}}`$, $`b_{\overline{c}}`$ are determined unambiguously from the requirement that the number of the constituent quarks and anti-quarks do not change during the hadronization, and that all initially available quarks and anti-quarks have to end up in the directly produced hadrons. This constraint is a basic assumption in all models of quark coalescence. The correct quark counting yields to the following equations, expressing the conservation of the number of quarks:
$`q`$ $`=`$ $`3N(qqq)+2Y(qqs)+\mathrm{\Xi }(qss)+K(q\overline{s})+\pi (q\overline{q})+`$ (5)
$`+2Y_c(qqc)+\mathrm{\Xi }_c(qsc)+\mathrm{\Xi }_{cc}(qcc)+\overline{D}(q\overline{c})`$
$`\overline{q}`$ $`=`$ $`3\overline{N}(\overline{q}\overline{q}\overline{q})+2\overline{Y}(\overline{q}\overline{q}\overline{s})+\overline{\mathrm{\Xi }}(\overline{q}\overline{s}\overline{s})+\overline{K}(\overline{q}s)+\pi (q\overline{q})+`$ (6)
$`+2\overline{Y}_c(\overline{q}\overline{q}\overline{c})+\overline{\mathrm{\Xi }}_c(\overline{q}\overline{s}\overline{c})+\overline{\mathrm{\Xi }}_{cc}(\overline{q}\overline{c}\overline{c})+D(\overline{q}c)`$
$`s`$ $`=`$ $`3\mathrm{\Omega }(sss)+2\mathrm{\Xi }(qss)+Y(qqs)+\overline{K}(\overline{q}s)+\varphi (s\overline{s})+`$ (7)
$`+2\mathrm{\Omega }_c(ssc)+\mathrm{\Xi }_c(qsc)+\mathrm{\Omega }_{cc}(scc)+\overline{D}_s(s\overline{c})`$
$`\overline{s}`$ $`=`$ $`3\overline{\mathrm{\Omega }}(\overline{s}\overline{s}\overline{s})+2\overline{\mathrm{\Xi }}(\overline{q}\overline{s}\overline{s})+\overline{Y}(\overline{q}\overline{q}\overline{s})+K(q\overline{s})+\varphi (s\overline{s})+`$ (8)
$`+2\overline{\mathrm{\Omega }}_c(\overline{s}\overline{s}\overline{c})+\overline{\mathrm{\Xi }}_c(\overline{q}\overline{s}\overline{c})+\overline{\mathrm{\Omega }}_{cc}(\overline{s}\overline{c}\overline{c})+D_s(\overline{s}c)`$
$`c`$ $`=`$ $`3\mathrm{\Omega }_{ccc}(ccc)+2\mathrm{\Xi }_{cc}(qcc)+\mathrm{\Lambda }_c(qqc)+D(\overline{q}c)+J/\psi (c\overline{c})+`$ (9)
$`+2\mathrm{\Xi }_{cc}(scc)+\mathrm{\Lambda }_c(qsc)+\mathrm{\Lambda }_c(ssc)+D_s(\overline{s}c)`$
$`\overline{c}`$ $`=`$ $`3\overline{\mathrm{\Omega }}_{ccc}(\overline{c}\overline{c}\overline{c})+2\overline{\mathrm{\Xi }}_{cc}(\overline{q}\overline{c}\overline{c})+\overline{\mathrm{\Lambda }}_c(\overline{q}\overline{q}\overline{c})+\overline{D}(q\overline{c})+J/\psi (c\overline{c})+`$ (10)
$`+2\overline{\mathrm{\Xi }}_{cc}(\overline{s}\overline{c}\overline{c})+\overline{\mathrm{\Lambda }}_c(\overline{q}\overline{s}\overline{c})+\overline{\mathrm{\Lambda }}_c(\overline{s}\overline{s}\overline{c})+\overline{D}_s(s\overline{c})`$
These equations for $`q`$, $`s`$, $`c`$ and $`(\overline{q}`$, $`\overline{s}`$, $`\overline{c})`$ determine the six $`b_i`$ normalization factors — which are not free parameters. These constraints, together with the prescription of the coalescence factors $`C_i`$, complete the description of hadron production from charmed quark matter by quark coalescence, and define the ALCOR<sub>c</sub> model.
In this paper, we will evaluate only the simplest predictions from ALCOR<sub>c</sub>, by considering ratios of the number of particles to the number of anti-particles and by assuming the symmetry of the coalescence process for charge conjugation, extending the results of ref. to the case of charmed quarks, mesons and baryons.
Assuming that the coalescence coefficients $`C`$ for hadrons are equal to that for the corresponding anti-particles, e.g. $`C_\mathrm{\Lambda }=C_{\overline{\mathrm{\Lambda }}}`$, the following relations were obtained for the ratio of light and strange antibaryons and baryons :
$`{\displaystyle \frac{\overline{N}(\overline{q}\overline{q}\overline{q})}{N(qqq)}}`$ $`=`$ $`\left[{\displaystyle \frac{b_{\overline{q}}\overline{q}}{b_qq}}\right]^3`$ (11)
$`{\displaystyle \frac{\overline{Y}(\overline{q}\overline{q}\overline{s})}{Y(qqs)}}`$ $`=`$ $`\left[{\displaystyle \frac{b_{\overline{q}}\overline{q}}{b_qq}}\right]^2\left[{\displaystyle \frac{b_{\overline{s}}\overline{s}}{b_ss}}\right]`$ (12)
$`{\displaystyle \frac{\overline{\mathrm{\Xi }}(\overline{q}\overline{s}\overline{s})}{\mathrm{\Xi }(qss)}}`$ $`=`$ $`\left[{\displaystyle \frac{b_{\overline{q}}\overline{q}}{b_qq}}\right]\left[{\displaystyle \frac{b_{\overline{s}}\overline{s}}{b_ss}}\right]^2`$ (13)
$`{\displaystyle \frac{\overline{\mathrm{\Omega }}(\overline{s}\overline{s}\overline{s})}{\mathrm{\Omega }(sss)}}`$ $`=`$ $`\left[{\displaystyle \frac{b_{\overline{s}}\overline{s}}{b_ss}}\right]^3`$ (14)
Inspecting eqs. (11)-(14) one can recognize, that the kaon to anti-kaon ratio $`𝒮^{qs}`$ has a special role as it acts as a stepping factor that connects various antibaryon to baryon rations,
$$𝒮^{qs}\frac{K(q\overline{s})}{\overline{K}(\overline{q}s)}=\left[\frac{b_qq}{b_{\overline{q}}\overline{q}}\right]\left[\frac{b_{\overline{s}}\overline{s}}{b_ss}\right].$$
(15)
This factor $`𝒮^{qs}`$ substitutes a light quark with a strange quark in the antibaryon to baryon ratios. Thus it shifts the antibaryon to baryon ratios and changes their strangeness content by one unit, as the following relations display:
$`𝒮^{qs}\left[{\displaystyle \frac{\overline{N}}{N}}\right]`$ $`=`$ $`{\displaystyle \frac{\overline{Y}}{Y}}`$ (16)
$`𝒮^{qs}𝒮^{qs}\left[{\displaystyle \frac{\overline{N}}{N}}\right]`$ $`=`$ $`{\displaystyle \frac{\overline{\mathrm{\Xi }}}{\mathrm{\Xi }}}`$ (17)
$`𝒮^{qs}𝒮^{qs}𝒮^{qs}\left[{\displaystyle \frac{\overline{N}}{N}}\right]`$ $`=`$ $`{\displaystyle \frac{\overline{\mathrm{\Omega }}}{\mathrm{\Omega }}}`$ (18)
The inverse factor, $`𝒮^{sq}=(𝒮^{qs})^1`$ decreases the strangeness content and increases the number of light quarks in the antibaryon to baryon ratios. Note that these relations hold between the ratios of the directly produced anti-baryons to baryons and that the number of observed anti-baryons and baryons have to be corrected to the various chains of resonance decays .
Extending the above ALCOR model to the case of charmed baryons and antibaryons, further relations are obtained:
$`{\displaystyle \frac{\overline{Y}_c(\overline{q}\overline{q}\overline{c})}{Y_c(qqc)}}`$ $`=`$ $`\left[{\displaystyle \frac{b_{\overline{q}}\overline{q}}{b_qq}}\right]^2\left[{\displaystyle \frac{b_{\overline{c}}\overline{c}}{b_cc}}\right]{\displaystyle \frac{\overline{\mathrm{\Xi }}_c(\overline{q}\overline{s}\overline{c})}{\mathrm{\Xi }_c(qsc)}}=\left[{\displaystyle \frac{b_{\overline{q}}\overline{q}}{b_qq}}\right]\left[{\displaystyle \frac{b_{\overline{s}}\overline{s}}{b_ss}}\right]\left[{\displaystyle \frac{b_{\overline{c}}\overline{c}}{b_cc}}\right]`$
$`{\displaystyle \frac{\overline{\mathrm{\Xi }}_{cc}(\overline{q}\overline{c}\overline{c})}{\mathrm{\Xi }_{cc}(qcc)}}`$ $`=`$ $`\left[{\displaystyle \frac{b_{\overline{q}}\overline{q}}{b_qq}}\right]\left[{\displaystyle \frac{b_{\overline{c}}\overline{c}}{b_cc}}\right]^2{\displaystyle \frac{\overline{\mathrm{\Omega }}_c(\overline{s}\overline{s}\overline{c})}{\mathrm{\Omega }_c(ssc)}}=\left[{\displaystyle \frac{b_{\overline{s}}\overline{s}}{b_ss}}\right]^2\left[{\displaystyle \frac{b_{\overline{c}}\overline{c}}{b_cc}}\right]`$
$`{\displaystyle \frac{\overline{\mathrm{\Omega }}_{ccc}(\overline{c}\overline{c}\overline{c})}{\mathrm{\Omega }_{ccc}(ccc)}}`$ $`=`$ $`\left[{\displaystyle \frac{b_{\overline{c}}\overline{c}}{b_cc}}\right]^3{\displaystyle \frac{\overline{\mathrm{\Omega }}_{cc}(\overline{s}\overline{c}\overline{c})}{\mathrm{\Omega }_{cc}(scc)}}=\left[{\displaystyle \frac{b_{\overline{s}}\overline{s}}{b_ss}}\right]\left[{\displaystyle \frac{b_{\overline{c}}\overline{c}}{b_cc}}\right]^2`$ (19)
These ratios and the ratios from eqs. (11)-(14) can be organized into a special structure displayed in Fig.1. We can introduce two more factors $`𝒮^{sc}`$ and $`𝒮^{cq}`$ constructed as in eq.(15) but from the ratios of charmed mesons:
$`𝒮^{sc}{\displaystyle \frac{\overline{D}_s(s\overline{c})}{D_s(\overline{s}c)}}`$ $`=`$ $`\left[{\displaystyle \frac{b_ss}{b_{\overline{s}}\overline{s}}}\right]\left[{\displaystyle \frac{b_{\overline{c}}\overline{c}}{b_cc}}\right]`$ (20)
$`𝒮^{cq}{\displaystyle \frac{D(\overline{q}c)}{\overline{D}(q\overline{c})}}`$ $`=`$ $`\left[{\displaystyle \frac{b_cc}{b_{\overline{c}}\overline{c}}}\right]\left[{\displaystyle \frac{b_{\overline{q}}\overline{q}}{b_qq}}\right]`$ (21)
The factor $`𝒮^{sc}`$ substitutes a strange quark with a charm one and the factor $`𝒮^{cq}`$ changes the charm quark into a light one. These properties lead to the following identity:
$$𝒮^{qs}𝒮^{sc}𝒮^{cq}1$$
(22)
This identity can be rewritten as an identity between the mesonic ratios:
$$\frac{\overline{D}_s/D_s}{\overline{D}/D}=\overline{K}/K$$
(23)
A comparison of this simple relation with experimental data could serve as test of quark matter formation and coalescence type rehadronization in heavy ion collisions.
The inverse of the step factors is defined as $`𝒮^{ji}=(𝒮^{ij})^1`$. The structure of the antibaryon to baryon ratios in ALCOR<sub>c</sub> is visualized in a geometric manner in Fig. 1. This way, more complicated but definitely interesting relations can be obtained. Since the baryons with one charm quark (or antiquark) can be measured most easily, one may consider the following relations as candidates for an experimental test:
$$\frac{\overline{\mathrm{\Xi }}_c(\overline{q}\overline{s}\overline{c})}{\mathrm{\Xi }_c(qsc)}=𝒮^{qs}\left[\frac{\overline{Y}_c}{Y_c}\right]=𝒮^{qc}\left[\frac{\overline{Y}}{Y}\right]=𝒮^{sc}\left[\frac{\overline{\mathrm{\Xi }}}{\mathrm{\Xi }}\right]=𝒮^{sq}\left[\frac{\overline{\mathrm{\Omega }}_c}{\mathrm{\Omega }_c}\right].$$
(24)
These yield the following simple relation between baryonic and mesonic ratios:
$`{\displaystyle \frac{\overline{Y}/Y}{\overline{Y}_c/Y_c}}`$ $`=`$ $`D_s/\overline{D}_s,`$ (25)
$`{\displaystyle \frac{\overline{N}/N}{\overline{Y}_c/Y_c}}`$ $`=`$ $`D/\overline{D},`$ (26)
$`{\displaystyle \frac{\overline{N}/N}{\overline{Y}/Y}}`$ $`=`$ $`\overline{K}/K.`$ (27)
Fig. 1. The application of mesonic step factors $`𝒮^{qs}`$, $`𝒮^{sc}`$ and $`𝒮^{cq}`$ on the antibaryon to baryon ratios. The arrows are indicating the three corresponding directions of shifting.
A number of similar expressions can be derived from Figure 1, picking up a given ratio and following all the paths to reach that from its neighbors.
In summary , we have made simple predictions from the ALCOR<sub>c</sub> model, extending the ALCOR model of algebraic coalescence and rehadronization of quark matter to the case when charmed quarks and final state hadrons are present in a significant number. We found that the various $`\overline{M}/M`$ mesonic ratios connect different $`\overline{B}/B`$ ratios. The agreement between the obtained theoretical relations and those in the measured data could serve as proof or disproof of the formation of quark matter in heavy ion collisions followed by a fast hadronization via quark coalescence. The predictions made in this paper are independent from the detailed values of coalescence coefficients, we have assumed only their symmetry for charge conjugation. The calculations of the absolute numbers of produced particles from ALCOR<sub>c</sub> requires the specification of these coalescence coefficients from calculations of cross-sections.
Acknowledgment: This work was supported by the OTKA Grants No. T029158, T025579 and T024094.
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# ELECTRON-PHONON COUPLING ORIGIN OF THE RESISTIVITY IN YNI2B2C SINGLE CRYSTALS
## 1 Introduction
In the past seven years, a great interest has been aroused by the discovery of the family of superconducting quaternary rare-earth borocarbide intermetallic compounds RNi<sub>2</sub>B<sub>2</sub>C (R = rare earth), of which YNi<sub>2</sub>B<sub>2</sub>C is one of the most studied non-magnetic members. In some way similarly to high-$`T_\mathrm{c}`$ cuprates, the borocarbides have a layered crystal structure, even if band-structure calculations have shown the presence of a three-dimensional electronic structure (see Ref. for a short overview of selected properties). Initially, a conventional BCS description of the superconducting properties in these materials has been supported by many of the experimental and theoretical results. More recently some particular features of the electronic specific heat (as a function of temperature and magnetic field) and of the upper critical field in the non-magnetic borocarbides LuNi<sub>2</sub>B<sub>2</sub>C and YNi<sub>2</sub>B<sub>2</sub>C have been interpreted as signs of an unconventional pairing, possibly of $`d`$-wave type $`^\mathrm{?}`$. On the other hand, Shulga *et al.* have recently explained the positive curvature of the upper critical field of YNi<sub>2</sub>B<sub>2</sub>C as a function of the temperature near $`T_\mathrm{c}`$, its magnitude and shape in the framework of $`s`$-wave Migdal-Eliashberg theory by considering the presence of two bands, one of them being more deeply involved in the transport properties of the compound $`^\mathrm{?}`$.
In the present work we demonstrate the complete agreement of the experimental resistivity data obtained in YNi<sub>2</sub>B<sub>2</sub>C with the predictions of the classic single-band theory for the electron-phonon (e-p) interaction (Bloch-Grüneisen theory). We obtain a value of the transport e-p coupling constant that well agrees with previous experimental and theoretical results. The transport e-p spectral function $`\alpha _{\mathrm{tr}}^2F(\mathrm{\Omega })`$ is also obtained by the fit of the experimental data in the whole temperature range and used for the calculation of the normalized tunneling conductance by directly solving the real-axis Eliashberg equations.
## 2 Experiment
High-quality YNi<sub>2</sub>B<sub>2</sub>C single crystals were grown by using the rf - zone melting process $`^\mathrm{?}`$. The critical temperature of the crystals, measured by AC susceptibility and perfectly confirmed by resistivity measurements, is $`T_\mathrm{c}=15.5`$K with $`\mathrm{\Delta }T_\mathrm{c}(1090\%)=0.5`$K. The imaginary part of the susceptibility shows a single, very high and narrow peak at $`TT_\mathrm{c}`$ which confirms the purity and crystallographic quality of the samples.
We measured the resistivity of these crystals as a function of the temperature, on cooling from 300 K down to 4.2 K, by using the AC version of the standard four-probe technique. The current leads were directly soldered to the opposite sides of the samples, which have parallelepipedal shape. The voltage leads, made from very thin gold wires, were glued with a conducting paste to the surface of the crystals. We improved the sensitivity of the measurement by injecting in the crystals an AC current of typically 10 mA at 133 and detecting the voltage by the standard lock-in technique. Due to a very slow cooling-down procedure, we were able to collect nearly three thousand resistivity values for every curve between 4.2 and 300 K.
In Figure 1 the resistivity $`\rho (T)`$ of one of the YNi<sub>2</sub>B<sub>2</sub>C crystals is shown (open circles). For clarity, only a reduced number of points is reported. In the inset of the same figure we show an enlargement of the low-temperature part of the resistivity curve that contains all the measured points at $`T<30`$ K (open circles).
The resistivity of Fig. 1 shows a perfect Bloch-Grüneisen (BG) behavior that has already been observed in previous experiments $`^\mathrm{?}`$. The small residual value of the resistivity $`\rho (0)`$= 3 $`\mu \mathrm{\Omega }`$cm and its perfectly linear high-temperature part (with a slope $`\mathrm{d}\rho /\mathrm{d}T`$ = 0.12 $`\mu \mathrm{\Omega }`$cm/K) indicate the high quality and low impurity content of the samples. Quite similar results were obtained in various YNi<sub>2</sub>B<sub>2</sub>C samples.
## 3 Discussion
The resistivity of a normal Fermi-liquid metal follows the well known Matthiessen’s rule: $`\rho (T)`$=$`\rho _0`$+$`\rho _{\mathrm{ph}}(T)`$, where $`\rho _0`$ and $`\rho _{\mathrm{ph}}(T)`$ are the residual and the phonon resistivities, respectively.
According to the BG theory, the high-temperature part of the $`\rho _{\mathrm{ph}}(T)`$ shows a linear behavior given by the following expression:
$$\rho _{\mathrm{ph}}(T)=(2\pi \epsilon _0k_\mathrm{B}/\mathrm{}\omega _\mathrm{p}^2)\lambda _{\mathrm{tr}}T$$
where $`\omega _\mathrm{p}`$ is the plasma frequency, $`\lambda _{\mathrm{tr}}`$=$`2_0^{\mathrm{}}[\alpha _{\mathrm{tr}}^2F(\mathrm{\Omega })/\mathrm{\Omega }]𝑑\mathrm{\Omega }`$ is the transport e-p coupling constant, and $`\alpha _{\mathrm{tr}}^2F(\mathrm{\Omega })`$ is the transport e-p spectral function. Once known the experimental value of the plasma energy $`\mathrm{}\omega _\mathrm{p}`$ and the temperature coefficient of the linear part of the resistivity at $`T>`$ 100 K (see Fig. 1), we can determine the transport coupling constant $`\lambda _{\mathrm{tr}}=(\mathrm{}\omega _\mathrm{p}^2/2\pi \epsilon _0k_\mathrm{B})\mathrm{d}\rho _{\mathrm{ph}}/\mathrm{d}T=(\mathrm{}\omega _\mathrm{p}^2/2\pi \epsilon _0k_\mathrm{B})\mathrm{d}\rho /\mathrm{d}T`$. The values $`\omega _\mathrm{p}=`$ 4.25 eV determined in Ref. by reflectance and EELS measurements and $`\mathrm{d}\rho /\mathrm{d}T=`$ 0.12 $`\mu \mathrm{\Omega }`$cm/K from Fig. 1 lead to the result $`\lambda _{\mathrm{tr}}=`$ 0.53.
Full information on the e-p coupling in YNi<sub>2</sub>B<sub>2</sub>C can be obtained by fitting the resistivity of Fig. 1 in the whole temperature range. Again following the standard BG theory, we use the most general expression for $`\rho _{\mathrm{ph}}(T)`$, given by:
$$\rho _{\mathrm{ph}}(T)=(4\pi \epsilon _0k_\mathrm{B}T/\mathrm{}\omega _\mathrm{p}^2)_0^{\mathrm{\Omega }_{\mathrm{max}}}[\alpha _{\mathrm{tr}}^2F(\mathrm{\Omega })/\mathrm{\Omega }][\mathrm{}\mathrm{\Omega }/2k_\mathrm{B}T\mathrm{sinh}(\mathrm{}\mathrm{\Omega }/2k_\mathrm{B}T)]^2𝑑\mathrm{\Omega }.$$
Our goal is to determine the function $`\alpha _{\mathrm{tr}}^2F(\mathrm{\Omega })`$ from the fit. Actually, as a first approximation, for the $`\alpha _{\mathrm{tr}}^2F(\mathrm{\Omega })`$ we use the phonon spectral density $`G(\mathrm{\Omega })`$ obtained in inelastic neutron scattering experiments $`^\mathrm{?}`$ multiplied by a two-step weighting function, whose constant values for $`\mathrm{\Omega }<37.5`$ meV and $`37.5<\mathrm{\Omega }<70`$ meV are to be determined by the fit. These energy ranges are chosen because they correspond to the two most-distinguishable structures of the $`G(\mathrm{\Omega })`$ which shows peaks at about 20 and 50 meV and a value close to zero just at 37.5 meV. We neglected in the fit the presence of high energy phonons at about 100 meV.
The results of the fit are shown as a solid line in Fig. 1. The experimental $`\rho (T)`$ is *perfectly* fitted by the theoretical BG curve in the whole temperature range. Fig. 2 shows the relative deviations of the experimental curve from the fit $`(\rho \rho _{\mathrm{FIT}})/\rho `$ as function of temperature. They never exceed $`\pm 5\%`$.
Fig. 3 shows the spectral function $`\alpha _{\mathrm{tr}}^2F(\mathrm{\Omega })`$ obtained as the product of the $`G(\mathrm{\Omega })`$ by the two-step weighting function $`\alpha _{\mathrm{tr}}^2`$ determined from the fit. The latter function is shown in the inset of Fig. 3. Of course, the $`\lambda _{\mathrm{tr}}`$ calculated from the $`\alpha _{\mathrm{tr}}^2F(\mathrm{\Omega })`$ perfectly coincides with the value previously obtained from the linear part of $`\rho (T)`$.
It is well known that, to a first approximation, $`\lambda _{\mathrm{tr}}\lambda `$, where $`\lambda `$ is the e-p coupling factor involved in the BCS coupling of the Cooper’s pairs. In the hypothesis of an e-p coupling origin of the superconductivity in YNi<sub>2</sub>B<sub>2</sub>C we calculate the quasiparticle density of states in this compound in both $`s`$\- and $`d`$-wave symmetry by solving in direct way the real-axis Eliashberg equations for the strong e-p coupling $`^\mathrm{?}`$ and using as $`\alpha _{\mathrm{tr}}^2F(\mathrm{\Omega })`$ the function shown in Fig. 3.
Fig. 4 shows the SIN tunneling conductances in $`s`$\- and $`d`$-wave symmetry calculated at 4.2 K from the density of states (solid and dash lines, respectively). From the solution of the Eliashberg equations we obtained both the correct $`T_\mathrm{c}`$ and the superconducting gap $`\mathrm{\Delta }`$ 2 meV observed in tunneling experiments $`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$, by using a value $`\lambda `$= 0.57 slightly greater than $`\lambda _{\mathrm{tr}}`$. This fact is consistent with the conventional relation between $`\lambda `$ and $`\lambda _{\mathrm{tr}}`$. From the curves of Fig. 4 and the tunneling experimental data in STM configuration $`^{\mathrm{?},\mathrm{?}}`$ it is very difficult to determine the possible symmetry of the order parameter in YNi<sub>2</sub>B<sub>2</sub>C due to the appreciable broadening of the SIN data.
A more clear indication can be obtained by the comparison of the calculated SIS tunneling conductance with the break-junction experimental data present in literature. This comparison is shown in Fig. 5 where the theoretical SIS curves at 4 K determined from the solution of Eliashberg equation for pure $`d`$-wave (dash line), pure $`s`$-wave (dot line) and $`s`$-wave plus magnetic impurities in quasi non-unitary limit $`^\mathrm{?}`$ (solid line) are presented together with the break-junction data of Ekino et al $`^\mathrm{?}`$ (inset). The $`s`$-wave curve in presence of a small amount of magnetic impurities (for details see the caption of the figure and Ref. ) reproduces all the main features of the experimental data including the well pronounced dip at about twice the energy of the gap. On the other hand, the $`d`$-wave tunneling conductance is quite different from the experimental curve. These results suggest a possible $`s`$-wave symmetry (or, at least, a dominant $`s`$-wave component) in YNi<sub>2</sub>B<sub>2</sub>C and give evidence for the essential role played by the electron-phonon coupling in the pairing mechanism in this compound.
## 4 Conclusions
We have demonstrated that the resistivity of YNi<sub>2</sub>B<sub>2</sub>C has a temperature dependence *perfectly* described by the standard BG theory for the e-p coupling in conventional metals. The value of $`\lambda _{\mathrm{tr}}`$ here determined from the resistivity is representative of an intermediate e-p coupling strength and is consistent with the value used in Ref. for the discussion of the superconducting properties of YNi<sub>2</sub>B<sub>2</sub>C in the framework of the isotropic single-band model. The transport e-p spectral function $`\alpha _{\mathrm{tr}}^2F(\mathrm{\Omega })`$ determined from the fit of the resistivity was used, to a first approximation, in the direct solution of the Eliashberg equations both in $`s`$\- and $`d`$-wave symmetry. The resulting $`T_\mathrm{c}`$ and $`\mathrm{\Delta }`$ are quite in agreement with the experimental data, while the comparison of the calculated SIS tunneling conductance with the break-junction data supports a conventional $`s`$-wave symmetry for the order parameter in YNi<sub>2</sub>B<sub>2</sub>C.
## References
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# Rapid solution of problems by nuclear-magnetic resonance quantum computation
## INTRODUCTION
Recently there has been interest in trying to use the thermal state as a starting point for NMR computation. We note two efforts to pursue this, one by Zhou, Leung, and Chuang , the other by Woodward and Brüschweiler . We come at the problem from a different point of view to obtain results slightly stronger than those of , as well as showing some different ways to proceed, and much more explicit than those claimed in .
Computational complexity brings the idea of cost vs. problem size into problems solvable by use of computers. For certain problems, cost grows with problem size more slowly for quantum computers than it does for a Turing machine , showing that the complexity of a problem depends on the computer used to solve it. With the Turing machine no longer the only game in town, the question is opened: what problems are natural to one or another computer design ?
Are all quantum computers alike with respect to the problems that they solve efficiently? Three types of quantum computer will be discussed in connection with problems of function classification, the prototype of which is the Deutsch-Jozsa (DJ) problem , which concerns determining a property of an $`n`$-bit function $`f:𝐙_N𝐙_2`$, given an oracle that evaluates $`f`$, where $`N`$ is written as shorthand for $`2^n`$.
In theory, which is all this paper deals with, a quantum computer yields the solution to a problem as the outcome of a quantum measurement , and can be called an outcome quantum computer (OQC) to distinguish it from an expectation-value quantum computer (EVQC), which in place of an outcome yields, to some finite precision, the expectation value for a measurement operator and a (possibly mixed) state . A nuclear-magnetic-resonance quantum computer (NMRQC) is a restricted EVQC, the restriction stemming from facts of NMR spectrometers. The restriction on an NMRQC relative to a general EVQC has consequences which seem to have gone unnoticed. Attention to them shows better how an NMR spectrometer can act as a quantum computer, stimulates a generalization of the DJ problem, and shows the way to solving the original DJ problem without exponential loss of signal as the number of bits $`n`$ increases .
## DJ PROBLEM FOR THE OQC AND THE EVQC
As stated originally, the DJ problem is this: given any function $`f:𝐙_N𝐙_2`$, show at least one of following: (A) $`f`$ is not constant, or (B) $`f`$ is not balanced, where a balanced function has the value 0 for just half of its $`N`$ arguments and 1 for the other half.
We review briefly the history of methods for use of an OQC to solve this problem. For later generalization, it is convenient to organize the method of solution in three steps, the middle one of which is a compound step that may be repeated: (1) Prepare a quantum register in a starting state; (2) apply operators including one for an Oracle for the function $`f`$; and (3) make a quantum measurement defined by a projection.
The method as first presented required a work bit and hence a quantum register of $`n+1`$ bits; it also required two invocations of the oracle (repetition of step (2)). Later Cleve et al. showed how to solve the problem invoking the Oracle only once ; building on this, Collins et al. showed how to skip the work bit so the register is only $`n`$ bits ; this calls for a Hilbert space spanned by $`N`$ orthonormal vectors $`|j`$, $`j=0,1,\mathrm{},N1`$. In this version, the method consists of the following steps:
1. Prepare the starting state
$$|w\stackrel{\mathrm{def}}{=}N^{1/2}\underset{j=0}{\overset{N1}{}}|j.$$
(1)
2. Apply the operator $`U_f`$ for the Oracle for the function $`f`$ defined by its effect on basis vectors $`|j`$:
$$U_f|j=(1)^{f(j)}|j$$
(2)
(no repetition and no other operators).
3. Make the measurement defined by the projection $`|ww|`$, which has eigenvalues 0 and 1, and hence two possible outcomes.
If the outcome is 1 the function is not balanced, while if the outcome is 0 the function is not constant, as follows from the probability of the outcome being 1:
$`\mathrm{Pr}\{\mathrm{outcome}=1\}`$ $`=`$ $`\mathrm{Tr}(|ww|U_f|ww|U_f^{})`$ (3)
$`=`$ $`\{\begin{array}{cc}1,\hfill & \text{if }f\text{ is constant,}\hfill \\ 0,\hfill & \text{if }f\text{ is balanced.}\hfill \end{array}`$
(In case $`f`$ is neither balanced nor constant, the OQC outcome can be either 0 or 1 with probabilities determined by the usual rules of quantum mechanics, but the outcome varies from one trial to another.)
Another version of the DJ problem restricts the class of functions to be the union of constant and balanced functions; we shall have occasion to introduce analogs to this version.
Turn now to the use of an EVQC, which in place of an outcome yields the expectation value $`\mathrm{Tr}(M\rho )`$ for a measurement $`M`$ of a density matrix $`\rho `$ . An EVQC is characterized by a parameter of resolution $`ϵ`$: Two density matrices $`\rho _1`$ and $`\rho _2`$ are taken to be distinguishable by a measurement described by an operator $`M`$ if and only if the difference in the expectation values exceeds the minimum resolution:
$$|\mathrm{Tr}(M\rho _1)\mathrm{Tr}(M\rho _2)|>ϵ\mathrm{\Lambda }(M),$$
(4)
where $`\mathrm{\Lambda }(M)`$ is the difference between the minimum and the maximum eigenvalue of the measurement operator $`M`$. (The factor $`\mathrm{\Lambda }(M)`$ makes limitations of resolution immune to the mere analytic trick of multiplying the measurement operator by a constant.) For an Oracle exercising $`U_f`$ on a density matrix $`\rho `$, a measurement operator $`M`$ yields an expectation value
$$E(f)=\mathrm{Tr}(MU_f\rho U_f^{}).$$
(5)
Using the measurement operator $`|ww|`$, an EVQC measuring the state $`U_f|w`$ obtains the expectation value
$`E(f)`$ $`=`$ $`\mathrm{Tr}(|ww|U_f|ww|U_f^{})`$ (6)
$`=`$ $`N^2{\displaystyle \underset{j,k=0}{\overset{N1}{}}}(1)^{f(j)+f(k)}.`$
For this case, it follows that
$$E(f)=\left(N^1\underset{j=0}{\overset{N1}{}}(1)^{f(j)}\right)^2.$$
(7)
This expectation $`E(f)`$ has the nice property of invariance under permutations of the arguments of $`f`$, and hence depends only on what might be called the “imbalance” of $`f`$, defined by
$`I(f)`$ $`\stackrel{\mathrm{def}}{=}`$ $`{\displaystyle \frac{1}{2}}[\text{(Number of values of}j\text{ for which}f(j)=1)`$ (8)
$`(\text{Number of values of}j\text{ for which}f(j)=0)].`$
Depending on $`f`$, $`I(f)`$ takes on integral values $`N/2I(f)N/2`$. (Recall $`N=2^n`$, and $`n>0`$, so $`N`$ is even.) That is, one has for this case
$$E(f)=4N^2I^2(f).$$
(9)
For example, if $`f`$ is balanced, one sees $`I(f)=0`$, so it follows that $`E(f)=0`$, while if $`f`$ is constant, $`I(f)=\pm N/2`$ so $`E(f)=1`$; the two cases are resolvable by an EVQC for any $`ϵ<1`$.
Drastic sensitivity to $`ϵ`$ is seen in the satisfiability problem of distinguishing the unsatisfiable function $`f_0`$ having the zero value for all arguments from any function $`f_1`$ that takes the value 1 for just one argument. One can check to see that $`I(f_0)=N/2`$ and $`I(f_1)=1N/2`$, so that
$`E(f_0)E(f_1)`$ $`=`$ $`4N^2[(N/2)^2(1N/2)^2]`$ (10)
$`=`$ $`2^{2n}(12^n).`$
This becomes exponentially small as the number $`n`$ of bits increases, so that $`|E(f_0)E(f_1)|>ϵ\mathrm{\Lambda }(|ww|)`$ only for
$$n<\mathrm{log}_2(4/ϵ).$$
(11)
## GENERALIZATION
Equation (6) suggests the following generalization. Given any $`N\times N`$ matrix $`B`$, define a mapping $`S_B`$ from the set of functions to numbers by
$`S_B(f)`$ $`\stackrel{\mathrm{def}}{=}`$ $`{\displaystyle \underset{j,k=0}{\overset{N1}{}}}(1)^{f(j)+f(k)}B_{jk}`$ (12)
$`=`$ $`\mathrm{Tr}(B)+{\displaystyle \underset{j=0}{\overset{N2}{}}}{\displaystyle \underset{k=j+1}{\overset{N1}{}}}(1)^{f(j)+f(k)}(B_{jk}+B_{kj}).`$
Then Eq. (6) is equivalent to $`E(f)=S_B^{}(f)`$, where $`B^{}`$ is the matrix defined by
$$(j,k)B_{jk}^{}=N^2.$$
(13)
In the general case defined by (5), one finds
$`E(f)`$ $`=`$ $`\mathrm{Tr}(MU_f\rho U_f^{})`$ (14)
$`=`$ $`{\displaystyle \underset{j,k=0}{\overset{N1}{}}}M_{jk}(1)^{f(j)}\rho _{kj}(1)^{f(k)}`$
$`=`$ $`S_B(f)`$
for a matrix $`B(\rho ,M)`$ having elements
$$B_{jk}=M_{jk}\rho _{kj}(\mathrm{no}\mathrm{sum}).$$
(15)
One is thus led to explore generalizations of EVQC computations that implement $`S_B`$ for matrices $`B(\rho ,M)`$ of the general form of Eq. (15) rather than the special form of Eq. (13). In particular, if $`S_B(f)=0`$, we shall say that $`f`$ is balanced with respect to $`B`$. By inspection, one arrives at the following:
###### Proposition 1
For $`\{c_j\}`$ any set of constants and $`\{B_j\}`$ any set of $`N\times N`$ matrices, if $`f`$ is balanced with respect to $`B_1,B_2,\mathrm{},`$ then $`f`$ is balanced with respect to $`_jc_jB_j`$.
It follows from Eq. (12) that
###### Proposition 2
If the matrix $`B`$ is written as the sum of symmetric and antisymmetric parts, only the symmetric part contributes to $`S_B`$.
For any $`f:𝐙_N𝐙_2`$, let $`\overline{f}`$ be the logical complement of $`f`$, so $`(j)\overline{f}(j)=1f(j)`$. Then it follows immediately from Eq. (2) that
###### Proposition 3
If $`\overline{f}`$ is the logical complement of $`f`$, then
$$(B)S_B(\overline{f})=S_B(f).$$
(16)
The three-step procedure for solving the DJ problem readily generalizes to execute $`S_B(f)`$ for a variety of matrices $`B`$, as will be illustrated in connection with the NMRQC.
## NMR SPECTROMETER USED AS A COMPUTER
We review the use of a nuclear-magnetic-resonance spectrometer as an NMRQC for solving the Deutsch-Jozsa problem, in order to point out obstacles that impede it (relative to a general EVQC. For step (1) on an NMR spectrometer, a liquid sample begins in a mixed state of thermal equilibrium and is manipulated one way or another into a starting state. The thermal-equilibrium density matrix is proportional to $`\mathrm{exp}(/k_BT)`$, where $``$ is the hamiltonian for the $`n`$-spin molecule (in the liquid sample) used as a quantum register, $`k_B`$ is Boltzmann’s constant, and $`T`$ is the temperature. In the high-temperature approximation the thermal density matrix is given adequately well by the first two terms in the Taylor expansion:
$$\rho _{\mathrm{eq}}=2^n(\mathrm{𝟏}/k_BT)N^1\mathrm{𝟏}\frac{\mathrm{}}{Nk_BT}\underset{i}{}\omega _iI_z^i,$$
(17)
where $`\omega _i`$ is the resonant angular frequency of the $`i`$-th nucleus, and $`I_z^i`$ is defined by a tensor product over all $`n`$ spins in which all the factors are unit operators except for $`\frac{1}{2}\mathrm{Diag}(1,1)`$ as the $`i`$-th factor of the tensor product. This state, being diagonal, is invariant under the action of the Oracle and so must be manipulated into some other density matrix to serve as a starting state.
How to produce a starting density matrix has been much discussed. One way to prepare a starting density matrix is to produce a pseudopure state using gradient pulses , resulting in a starting density matrix of the form
$$\rho =(1\alpha /N)N^1\mathrm{𝟏}+\frac{\alpha }{N}|ww|,$$
(18)
for some (usually small) coefficient $`\alpha `$; a cost is a reduction exponential in $`n`$ of $`\alpha `$ and hence of the available spectrometer signal. (The small size of $`\alpha `$ compounds the exponential loss of polarization expressed by the explicit appearance of $`N`$ in the formula for the pseudopure state.) Another way to deal with a starting state is temporal averaging, which avoids the signal loss of a pseudopure state, but requires repetitions of the whole procedure and addition of the resulting spectra, costing much time . A third way uses extra qubits as ancilla , and a fourth advocates another use of extra bits . All these methods are elaborate and expensive of signal or time or number of bits required. A ray of hope is the simplified use of the equilibrium density matrix, which has been shown to work for the DJ problem for functions of one bit and two bits , but has not been developed into an algorithm applicable to the general case of $`n`$ bits.
Whatever method prepares a starting state, in step (2) a unitary transformation on the density matrix is implemented by use of r.f. pulses combined with waiting periods during which spin-spin couplings inherent in the molecule of the liquid sample exercise their effect. This results in some density matrix $`\rho ^{}`$ at some time $`t^{}`$.
Step (3), which we particularly want to notice, is modified in NMR to result in a spectrum conventionally expressed as the time evolution of the measurement of $`F^+`$, which is equivalent to the simultaneous measurement of $`F_x`$ and $`F_y`$, defined by
$$F_{x,y}=\underset{j=1}{\overset{n}{}}I_{x,y}^j,$$
(19)
where $`I_x^j`$ is a tensor product over all $`n`$ spins in which all the factors are unit operators except the $`j`$-th factor, which is
$$I_x=\frac{1}{2}\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right);$$
(20)
$`I_y^j`$ has instead of $`I_x`$ the $`j`$-th factor
$$I_y=\frac{1}{2}\left(\begin{array}{cc}0& i\\ i& 0\end{array}\right).$$
(21)
(If the resonances of individual spins $`j`$ are well resolved (e.g. if spin $`j`$ has a unique gyromagnetic ratio),the corresponding $`I_{x,y}^j`$ can be measured using analogue or digital filters, and not just the sum over all $`j`$.) The spectrometer signal for $`F_x`$ starts at $`t^{}`$ and is a sequence of expectation values obtained at measurement times $`t_k=t^{}+k\mathrm{\Delta }t`$, $`k=0,1,2`$, …, where $`\mathrm{\Delta }t`$ is the sampling interval. In the Heisenberg picture, the density matrix $`\rho ^{}`$ is fixed and the $`k`$-th expectation value is $`\mathrm{Tr}(\rho ^{}M_k)`$, where, for example with $`M=F_x`$,
$$M_k=\mathrm{exp}\left(\frac{i}{\mathrm{}}k\mathrm{\Delta }t\right)F_x\mathrm{exp}\left(\frac{i}{\mathrm{}}k\mathrm{\Delta }t\right).$$
(22)
Analogous time sequences can be defined for $`F_y`$ and, in the well resolved case, for $`I_{x,y}^j`$.
Typically, the signal (which is damped by relaxation in a way not shown in Eq. (22)) is Fourier transformed into an NMR spectrum. Either one deals with complications from a less than general coupling, e.g., by use of swap operations , or one must use a molecule and a spectrometer which exhibit distinct frequencies for all single-spin transitions.
Remark: For a molecule in which all transitions between basis states have distinct frequencies, to see them one must resolve all $`n2^{n1}`$ peaks of the Fourier spectrum, which requires a time-bandwidth product exponential in the number of spins.
This makes it desirable to avoid Fourier transforms of the time-domain signal, leading us to focus on single-time measurements which involve no Fourier transform. The requirement that single-time measurement operators in NMR be unitarily equivalent to $`F_{x,y}`$ or to $`I_{x,y}^j`$ now becomes an obstacle, because the operators $`F_{x,y}`$ and $`I_{x,y}^j`$ all have spectra with multiple eigenvalues, so that no single-time operator is nondegenerate. Thus no single-time operator has the power of a nondegenerate operator to resolve states; this constraint limits the NMRQC.
The original method for solving the DJ problem used for its measurement the projection operator $`|w,0w,0|`$ while the streamlined method used $`|ww|`$, also a projection operator. Less important than it seems at first glance but still provoking of thought is the following:
###### Proposition 4
No single-time measurement of any NMR operator can implement any nontrivial projection.
Proof: Any single-time operator $`M_k`$ is some unitary transform of some weighted sum of operators $`I_{x,y}^j`$, all of which are traceless. Trace is preserved under unitary transform, so all the candidates for $`M_k`$ are traceless (which indeed they must be if the large term proportional to the unit matrix in Eq. 17 is not to drown out all the effects of interest). A nontrivial projection has nonzero trace. Q.E.D.
One can get around Proposition 4 by invoking a nonprojective operator to distinguish balanced from constant functions, but questions remain that are less easily disposed of. One requires in place of $`|ww|`$ an operator $`M`$ that (a) works with the starting density operator of the form of Eq. (18), and (b) via Eq. (6) produces an expectation value that is invariant under permutation of the arguments of $`f`$. It is proved in Appendix A that:
###### Proposition 5
(i) Given a density operator $`\rho `$ of the form of Eq. (18), the expectation value $`\mathrm{Tr}(MU_f\rho U_f^{})`$ is invariant under permutation of the arguments of $`f`$ if and only if $`M=c|ww|+D+A`$, where $`D`$ is any diagonal matrix, $`A`$ is any antisymmetric matrix, and $`c`$ is any scalar; (ii) the resulting expectation value $`E(f)`$ is independent of the antisymmetric matrix A; (iii) if $`M`$ is hermitian, $`c`$ and $`D`$ are real and $`A`$ is pure imaginary.
Using a measurement operator $`c|ww|+D+A`$ unitarily equivalent to $`F_x`$ and the starting density matrix of Eq. (18), the expectation value for constant functions differs from that for balanced functions by $`c\alpha /N`$, with the result that the two classes of functions are distinguishable if and only if the resolution satisfies
$$ϵ<\frac{\alpha }{N}|c|/\mathrm{\Lambda }(F_x)=\frac{\alpha |c|}{nN},$$
(23)
where the second equality follows from $`\mathrm{\Lambda }(F_x)=n`$ as the difference between the minimum and maximum eigenvalues of $`F_x`$. Thus a small value of $`|c|`$ demands fine resolution.
The straightforward way to produce a measurement operator in NMR spectrometry is by unitary transform of $`F_x`$. (It adds nothing to allow unitary transforms of $`F_y`$, which is unitarily equivalent to $`F_x`$.) This and Eq. (23) raise the question of how large a value of $`|c|`$ is possible for an operator of the form $`M=c|ww|+D+A`$ that is constrained to be unitarily equivalent to $`F_x`$. As follows from the invariance of eigenvalues under unitary transform, the constraint is that $`F_x`$ and $`M`$ have the same eigenvalues with the same multiplicities. This implies
###### Proposition 6
For the matrix $`M`$ of Proposition 5 to be unitarily equivalent to $`F_x`$, it is necessary that $`D`$, $`A`$, and $`c`$ be such that for all eigenvalues $`\lambda _k`$ of $`F_x`$, $`k=1`$, …, $`N`$, $`det(M\lambda _k)=0`$ and $`\mathrm{Tr}(M)=\mathrm{Tr}(F_x)=0`$.
It is instructive to look at the first two cases, $`n=1`$ and $`n=2`$. For $`n=1`$, one has $`F_x=I_x`$ and $`M=I_x`$ produces the largest possible value of $`|c|/\mathrm{\Lambda }(M)`$ consistent with the eigenvalues of $`\pm 1/2`$, namely, $`|c|/\mathrm{\Lambda }(M)=1`$. For the two-spin case ($`n=2`$), an analysis of the restriction that the eigenvalues be those of $`F_x`$ shows
$$\mathrm{For}n=2,|c|/\mathrm{\Lambda }(M)<(2/3)^{1/2};$$
(24)
with more work, somewhat lower bounds can be demonstrated. (We found an $`M`$ for two spins unitarily equivalent to $`F_x`$ for which $`|c|/\mathrm{\Lambda }(M)=3^{1/2}`$, but we do not know if this is the best that can be done.) Thus there is a drop-off in $`|c|/\mathrm{\Lambda }`$ between the case $`n=1`$ and the case $`n=2`$, and hence an increase in the fineness of required resolution relative to $`\alpha /N`$ (see Eq. (23)). This drop-off suggests the following question for future analysis:
Question: For a measurement operator of the form $`M=c|ww|+D+A`$, unitarily equivalent to $`F_x`$, how does the largest possible value of $`|c|/\mathrm{\Lambda }(M)`$ vary with the number of bits $`n`$?
Via Eq. (23), the answer to this question will determine as a function of $`n`$ the resolution necessary for an NMRQC to solve the DJ problem using this method. The cases examined suggest a decreasing function; if confirmed this poses a serious obstacle of signal loss beyond that already known to the use of this method to solve the DJ problem.
## EXAMPLES OF FUNCTIONS NATURAL TO NMR
With this background, we ask: are there function classes for which a single-time measurement suffices to distinguish a function of that class from the constant function? Here are some such classes, the definitions of which depend on the concept of a Hamming distance. Let the argument $`j`$ of a function $`f:𝐙_N𝐙_2`$ be written as an $`n`$-bit string, padded with 0’s to the left. Given two integers $`j,k`$ (with $`0j,k2^n1`$), let $`d`$ be the number of bits of $`j`$ that are different from the corresponding bits of $`k`$. This is the Hamming distance between $`j`$ and $`k`$, denoted $`d(j,k)`$. Consider functions $`f:𝐙_N𝐙_2`$ such that: a) $`f(j)=1`$ for $`N/4`$ values of $`j`$, and b) if $`f(j)=f(k)=1`$, then $`d(j,k)1`$. Let $`𝒞_N`$ be the set of all such functions together with all their binary complements.
###### Proposition 7
For all $`j`$, every function of $`𝒞_N`$ is balanced with respect to $`I_x^j`$.
Proof: Suppose $`f𝒞_N`$; let $`g`$ be $`f`$ or $`\overline{f}`$, whichever function takes the value 1 for $`N/4`$ of its arguments. By Proposition 16 it suffices to show that, for all $`j`$, $`U_gI_x^jU_g^{}=0`$. The ($`l,m`$)-element of $`I_x^j`$ is nonzero if and only if the ($`n`$-bit representations of) $`l`$ and $`m`$ differ at just bit $`j`$. Hence this element is nonzero only if the Hamming distance $`d(l,m)=1`$. Because $`I_x^j`$ is proportional to a permutation matrix, it has one nonzero element in each row, so $`U_g`$ acting on its left changes the sign of half the nonzero elements of $`I_x^j`$. $`U_g`$ multiplied on the right also changes half the nonzero elements of $`U_gI_x^j`$. If elements negated by multiplication on the right are distinct from those negated by multiplication on the operation on the left, then half the elements change sign and we are done. For this to fail, at the element $`(l,m)`$, it must be that $`g(l)=g(m)=1`$ and $`(I_x^j)_{l,m}0`$. But that can happen only for functions not in $`𝒞_N`$. Q.E.D.
From this proposition, it will be shown that functions of $`𝒞_N`$ can be efficiently distinguished from constant functions by use of an NMRQC; moreover, solving this problem by use of a classical computer requires a number of function evaluations that grows exponentially with $`n`$.
Consider a starting state $`\rho `$ prepared from the equilibrium density operator by a hard $`90^{}`$ $`y`$-pulse:
$$\rho \stackrel{\mathrm{def}}{=}U_{90y}\rho _{\mathrm{eq}}U_{90y}^{}N^1\mathrm{𝟏}\frac{\mathrm{}}{Nk_BT}\underset{i=1}{\overset{n}{}}\omega _iI_x^i;$$
(25)
suppose the Oracle executes $`U_f`$ and the measurement operator is $`F_x=_iI_x^i`$. Applying Eq. (15) to this case and using $`(I_x^j)_{l,m}(I_x^k)_{m,l}=\delta _{jk}/4`$, one finds
$$B(\rho ,F_x)=\frac{\mathrm{}}{2Nk_BT}\underset{i}{}\omega _iI_x^i.$$
(26)
By Propositions 1 and 7, every function $`f𝒞_N`$ is balanced with respect to $`B`$, so $`S_B(f)E(f)=0`$, regardless of $`\omega _i`$. In contrast, for the constant functions $`f_0(j)=0`$ and $`f_1(j)=1`$ (for all $`j`$), one finds
$$E(f_{0,1})=\underset{l,m=0}{\overset{N1}{}}B_{lm}=\frac{\mathrm{}}{k_BT}\underset{i=1}{\overset{n}{}}\omega _i.$$
(27)
Notice the absence of a factor of $`N`$ in the denominator, removed by summing over the $`N`$ elements of $`I_x^i`$, each 1/2. Hence, neglecting effects beyond reach of this theory, an NMRQC operating with the starting density matrix defined in Eq. (25) can distinguish, for any $`n`$, functions of class $`𝒞_N`$ from constant functions for any resolution
$$ϵ<\frac{\mathrm{}}{nk_BT}\underset{i=1}{\overset{n}{}}\omega _i.$$
(28)
A striking feature of this result is the appearance in the denominator of $`n`$, the number of nuclear spins, rather than $`N2^n`$. Hence we have a method that avoids the much lamented exponential loss of signal.
## THE THERMAL STATE AND $`|ww|`$-BALANCED FUNCTIONS
At the expense of an extra bit and a more complex Oracle, the balanced functions (i.e., balanced with respect to $`|ww|`$) can be distinguished from constant functions using the starting state obtained merely by a hard $`90^{}`$ $`y`$-pulse applied to the thermal state (see Eq. (25)), requiring neither the pseudopure state of Eq. (18) nor temporal averaging. One requires an Oracle for a function $`f:𝐙_{N/2}𝐙_2`$ that implements $`U_f^{}`$ not for $`f`$ but for $`f^{}:𝐙_N𝐙_2`$, related to $`f`$ by
$$f^{}(j)=\{\begin{array}{cc}f(j),\hfill & \text{if }0jN/21\text{,}\hfill \\ 0,\hfill & \text{if }N/2jN1\text{.}\hfill \end{array}$$
(29)
Thus while the function $`f`$ is a function on $`n1`$ bits balanced with respect to $`|ww|`$, $`f^{}`$ is a function on $`n`$ bits balanced with respect to $`I_x^1`$. The three steps of execution to decide if $`f`$ is balanced or constant are then: (1) apply a hard 90 $`y`$-pulse on the thermal state, (2) apply $`U_f^{}`$ for $`f^{}`$ related to $`f`$ as above, (3) measure $`I_x^1`$ in the time domain, in the limit of small times, to obtain a signal that is substantial if $`f`$ is constant but vanishes if $`f`$ is balanced with respect to $`|ww|`$.
It is easy to check that given this more complex Oracle, an NMRQC decides between balanced functions and constant functions for any resolution
$$ϵ<\frac{\mathrm{}}{k_BT}\omega _1.$$
(30)
There is no exponential growth in the demand for resolution; indeed there is no growth at all, an advantage over the procedure described in . (The factor of $`n`$ in Ex. (23) vanishes when $`F_x`$ is replaced by $`I_x^1`$.) This shows a way to solve the original DJ problem on NMR with no loss of signal as $`n`$ increases.
It should be remarked that the operations discussed do not display the couplings needed to make any kind of quantum computer serve to distinguish the function classes discussed. These couplings are required, however, in the NMR implementation of the Oracle.
## APPENDIX A. PROOF OF PROPOSITION 5
Assume for some $`\alpha `$ that $`\rho =N^1(1\frac{\alpha }{N})\mathrm{𝟏}+\frac{\alpha }{N}|ww|`$. Then we have for the expectation value, $`E(f)=\mathrm{Tr}(M\rho )=N^1[(1\frac{\alpha }{N})\mathrm{Tr}(M)+\frac{\alpha }{N}S_M(f)]`$, whence it follows that $`(M,f,g)[E(f)=E(g)S_M(f)=S_M(g)]`$. To see the condition imposed on $`M`$ by the required invariance of $`S_M`$ under permutations, let $`P_{lm}`$ be the matrix obtained by permuting rows $`l`$ and $`m`$ of the $`N\times N`$ identity matrix. Define an operation of $`P_{lm}`$ on $`f`$ by
$$(P_{lm}f)(j)=\{\begin{array}{cc}f(m),\hfill & \text{if }j=l\text{,}\hfill \\ f(l),\hfill & \text{if }j=m\text{,}\hfill \\ f(j),\hfill & \text{otherwise.}\hfill \end{array}$$
(31)
Because general permutations are compositions of elementary permutations, the necessary and sufficient condition for $`\mathrm{Tr}(M\rho )`$ to be invariant under all permutations is that
$$(l,m,f)S_M(f)=S_M(P_{l,m}f).$$
(32)
It follows from Eqs. (12) and (31) that for any $`f`$ such that $`f(l)f(m)`$,
$`S_M(P_{lm}f)`$ (33)
$`=`$ $`\mathrm{Tr}(M)+{\displaystyle \underset{j=0}{\overset{N2}{}}}{\displaystyle \underset{k=j+1}{\overset{N1}{}}}(1)^{P_{lm}f(j)+P_{lm}f(k)}(M_{jk}+M_{kj})`$
$`=`$ $`\mathrm{Tr}(M)+{\displaystyle \underset{j=0}{\overset{N2}{}}}{\displaystyle \underset{k=j+1}{\overset{N1}{}}}(1)^{f(j)+f(k)+\delta jl+\delta _{jm}+\delta _{kl}+\delta _{km}}`$
$`\times (M_{jk}+M_{kj}),`$
where $`\delta _{jl}`$ is the Kronecker $`\delta `$, equal to 1 if $`j=l`$ and otherwise equal to 0.
For any $`l<m`$, assume any $`f`$ such that $`f(l)f(m)`$; set $`\widehat{M}_{jk}=M_{jk}+M_{kj}`$, require $`S_M(f)=S_M(P_{lm}f)`$, use Eq. (33), and eliminate terms that are the same on the two sides of the equation to show:
$`0`$ $`=`$ $`{\displaystyle \underset{k=m+1}{\overset{N1}{}}}(1)^{f(k)}(\widehat{M}_{lk}\widehat{M}_{mk})`$ (34)
$`+{\displaystyle \underset{j=0}{\overset{l1}{}}}(1)^{f(j)}(\widehat{M}_{jl}\widehat{M}_{jm})`$
$`+{\displaystyle \underset{k=l+1}{\overset{m1}{}}}(1)^{f(k)}\widehat{M}_{lk}{\displaystyle \underset{j=l+1}{\overset{m1}{}}}(1)^{f(j)}\widehat{M}_{jm}.`$
(The convention is used that if the upper limit of a sum is less than the lower limit, the sum is zero.) On relabeling some indices in the sums, this becomes
$`(l<m)(f\mathrm{s}.\mathrm{t}.f(l)f(m))`$
$`0`$ $`=`$ $`{\displaystyle \underset{k=0,kl,km}{\overset{N1}{}}}(1)^{f(k)}(\widehat{M}_{lk}\widehat{M}_{mk}).`$ (35)
This can hold for all admissible $`f`$ only if
$$(l<m)(kl,m)\widehat{M}_{lk}=\widehat{M}_{mk}.$$
(36)
This, together with the symmetry from its definition that $`\widehat{M}_{jk}=\widehat{M}_{kj}`$, implies that for $`(jk)\widehat{M}_{jk}`$ is independent of $`j`$ and $`k`$. From this, part (i) of the proposition follows immediately. Part (ii) is an immediate consequence of Proposition 2, and part (iii) depends only on the definition of a hermitian matrix. Q.E.D.
## APPENDIX B. DECIDING BETWEEN $`𝒞_N`$ AND CONSTANT <br>FUNCTIONS CLASSICALLY
###### Proposition 8
The number of invocations of a classical Oracle required to decide with certainty that a function $`f𝒞_N`$ is not constant is at least $`2^{n1}+1.`$
Remark: The issue in proving this is to rule out the possibility that the constraint on the Hamming distance associated with the function class $`𝒞_N`$ can greatly reduce the number of invocations required of the Oracle.
Proof: Given an Oracle that, on demand, takes an argument $`j`$ and computes the function value $`f(j)`$, how many invocations of the Oracle are sufficient to assure a decision between “$`f`$ is constant” and “$`f𝒞_N`$”? Suppose one has obtained from the Oracle the values $`f(j)`$ for any $`K`$ values of $`j`$, with $`KN/2`$, and suppose for all these arguments, $`f(j)=0`$. Then the possibility that $`f`$ is constant is not excluded. What about the possibility that $`f𝒞_N`$? We show that under these conditions there exists an $`f^{}𝒞_N`$ that satisfies $`f^{}(j)=0`$ for all the $`K`$ arguments tested, so the possibility that $`f𝒞_N`$ is also not excluded. This follows as soon as we show that in any set of $`N/2`$ arguments there exist a subset of $`N/4`$ arguments each separated by a Hamming distance greater than 1 from all the others of the subset. To see this is so, observe that out of the at least $`N/2`$ arguments unchecked by the Oracle, any one can be chosen and called $`j_0`$. Partition all the arguments of the unchecked subset into classes $`W_m`$ where $`kW_m`$ if and only if $`m=d(j_0,k)`$. Observe that $`(j,k)d(j,k)=1[(m=m^{}\pm 1)`$ such that $`jW_m`$ and $`kW_m^{}]`$. From this it follows that for any pair $`j`$ and $`k`$ both in $`W_0W_2`$, …, $`d(j,k)1`$; similarly for any pair $`j,k`$ both in $`W_1W_3`$, …, $`d(j,k)1`$. At least one of these unions of $`W`$-classes has $`N/4`$ elements, and hence holds arguments for some $`f^{}𝒞_N`$. Q.E.D.
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# Convergence of zeta functions of graphs
## Introduction
Associated to any finite graph $`X`$ there is a zeta function $`Z(X,u)`$, $`u`$. It is defined as an infinite product but shown (in various different cases) by Ihara, Hashimoto, and Bass to be a polynomial. Indeed the rationality formula for a $`q+1`$ regular $`X`$ states that:
$$Z(X,u)=(1u^2)^{\chi (X)}\mathrm{Det}(I\delta u+qu^2).$$
(0.1)
Here $`\delta `$ is the adjacency operator of $`X`$.
In , an $`L^2`$-zeta function is defined for noncompact graphs with symmetries, using the machinery of von Neumann algebras. A rationality formula similar to (0.1) expresses the relationship between the zeta function and the von Neumann determinant of a Laplace operator. The results of this paper focus on an especial case. Let $`Y`$ be an infinite graph which covers a finite graph $`B=\pi \backslash Y`$. The $`L^2`$-zeta function $`Z_\pi (Y,u)`$ is defined in only in a small neighborhood of zero. The first result of this paper is to extend the $`L^2`$-zeta function to the interior of $`C=\{u:|u|=q^{1/2}\}[1,\frac{1}{q}][\frac{1}{q},1]`$.
In the second part of the paper, we consider a tower of finite graphs $`B_i`$ covered by $`Y`$. Put $`N_i=|B_i|/|B|`$. In Theorem 2.1 we show that the zeta functions for the $`B_i`$, renormalized by taking $`N_i^{\text{th}}`$ roots, converge to the $`L^2`$-zeta function for $`Y`$. The argument is inspired by, and uses, work of Lück .
In the first section we recall the definitions of the zeta functions of finite and infinite graphs. One of the main results of this paper is Theorem 1.5 in this section. In the second section of this paper we prove the convergence theorem and exhibit interesting examples. Theorem 2.4 generalizes work of Deitmar .
## 1. Zeta functions
In this section we recall the definition of the zeta function and related material. We first recall the definition of the zeta function for finite graphs.
### 1.1. Finite Graphs
For a graph $`X`$, let $`VX`$ and $`EX`$ denote the sets of vertices and edges, respectively, of $`X`$. If each vertex has the same degree then $`X`$ is *regular*.
###### Definition 1.1.
Let $`X`$ be a finite graph. A closed path in $`X`$ is *primitive* if it is not a nontrivial power of another path inside the fundamental group of $`X`$. Let $`P`$ be the set of free homoptopy classes of primitive closed paths of $`X`$. Then the zeta function of $`X`$ is
$$Z(X,u)=\underset{\gamma P}{}(1u^{\mathrm{}(\gamma )}),$$
where $`\mathrm{}(\gamma )`$ is the minimum length of paths in the class of $`\gamma `$.
Let $`\delta `$ be the adjacency operator of the graph $`X`$ acting on $`l^2(VX)`$. For $`fl^2(VX)`$ let $`Qf(x)=q(x)f(x)`$ where $`q(x)+1`$ is the degree of the vertex $`x`$. Put $`\mathrm{\Delta }(X,u)=I\delta u+Qu^2`$. The Ihara rationality formula says that $`Z(X,u)`$ is a polynomial:
$$Z(X,u)=(1u^2)^{\chi (X)}\mathrm{Det}(\mathrm{\Delta }(X,u)).$$
(1.1)
The zeta function satisfies the following functional equation: \[1, Corollary 3.10\]
$$Z(X,(qu)^1)=\left(\frac{1u^2}{q^2u^21}\right)^{\chi (X)}q^{v2e}u^{2e}Z(X,u)$$
(1.2)
where $`e=|EX|`$ and $`v=|VX|`$.
For more details and examples see .
The following proposition is well known. A proof can be found in \[6, Page 59\].
###### Proposition 1.2.
The zeros of the zeta function for any finite $`q+1`$ regular graph lie on the set $`C`$ where
$$C=\{u:|u|=q^{1/2}\}[1,\frac{1}{q}][\frac{1}{q},1].$$
Let $`\mathrm{\Omega }`$ be the interior of $`C`$. Then:
###### Proposition 1.3.
If $`u\mathrm{\Omega },\lambda [(q+1),q+1]`$ then $`1\lambda u+qu^2(\mathrm{},0]`$.
###### Proof.
Let $`u=a+bi`$. For $`1\lambda u+qu^2`$ to be real we must have $`1\lambda u+qu^2=1\lambda \overline{u}+q\overline{u}^2.`$ This implies that either $`b=0`$ or $`a=\frac{\lambda }{2q}`$. If $`a=\frac{\lambda }{2q}`$ then $`1\lambda u+qu^2=1q(a^2+b^2)`$ which is negative or zero if $`u`$ is on or outside of the circle $`|u|=q^{1/2}`$. In the case that $`b=0`$ put $`f(a):=1\lambda a+qa^2=1\lambda u+qu^2`$. If $`f`$ has no real root then $`f`$ is always positive. Otherwise $`|\lambda |2\sqrt{q}.`$ Then $`f(\frac{1}{q})`$ and $`f(\frac{1}{q})`$ are non-negative, and $`f^{}`$ is not zero in $`(\frac{1}{q},\frac{1}{q})`$. Therefore $`f`$ will be positive on $`(\frac{1}{q},\frac{1}{q})`$. ∎
###### Corollary 1.4.
For a $`q+1`$ regular finite graph X, the polynomial $`\mathrm{Det}\mathrm{\Delta }(X,u)`$ has an analytic $`N^{\text{th}}`$ root on $`\mathrm{\Omega }`$ for all $`N`$.
###### Proof.
We know $`\mathrm{Det}\mathrm{\Delta }(X,u)=\mathrm{\Pi }_\lambda (1\lambda u+qu^2)`$ where $`\lambda `$ varies over eigenvalues of the adjacency operator of the graph. Then $`\mathrm{\Pi }_\lambda (1\lambda u+qu^2)^{\frac{1}{N}}`$ is an analytic $`N^{\text{th}}`$ root for $`\mathrm{Det}\mathrm{\Delta }(X,u)`$. ∎
### 1.2. Infinite Graphs
For $`\pi `$ a countable discrete group, the *von Neumann algebra* of $`\pi `$ is the algebra $`𝒩(\pi )`$ of bounded $`\pi `$-equivariant operators from $`l^2(\pi )`$ to $`l^2(\pi )`$. The *von Neumann trace* of an element $`f𝒩(\pi )`$ is defined by
$$\mathrm{Tr}_\pi f=f(e),e$$
for $`e\pi `$ the unit element. For $`H=_{i=1}^nl^2(\pi )`$ and a bounded $`\pi `$-equivariant operator $`f:HH`$, define
$$\mathrm{Tr}_\pi f=\underset{i=1}{\overset{n}{}}\mathrm{Tr}_\pi f_{ii}.$$
The trace as defined is independent of the decomposition of $`H`$. The von Neumann trace extends to bounded $`\pi `$-equivariant operators on Hilbert $`𝒩(\pi )`$-modules, but they will not be needed.
Now let $`Y=(VY,EY)`$ be an infinite graph. Suppose the group $`\pi `$ acts freely on $`Y`$ with finite quotient $`B`$.
Let $`P`$ denote the set of free homotopy classes of primitive closed paths in $`Y`$. For $`\gamma P`$, $`\mathrm{}(\gamma )`$ is the length of the shortest representative of $`\gamma `$. The group $`\pi _\gamma `$ is the stabilizer of $`\gamma `$ under the action of $`\pi `$. The *$`L^2`$-zeta function* of $`Y`$ is defined in as the infinite product
$$Z_\pi (Y,u)=\underset{\gamma \pi \backslash P}{}\left(1u^{\mathrm{}(\gamma )}\right)^{\frac{1}{|\pi _\gamma |}}.$$
The adjacency operator $`\delta `$ and Laplace operator $`\mathrm{\Delta }(Y,u)=I\delta u+Qu^2`$ are $`\pi `$-equivariant operators on $`l^2(VY)`$. Choosing lifts of vertices of $`B`$ yields a decomposition $`l^2(VY)=l^2(\pi )`$. Then from \[2, Theorem 0.3\],
$$Z_\pi (Y,u)=(1u^2)^{\chi ^{(2)}(Y)}\mathrm{Det}_\pi \mathrm{\Delta }(Y,u).$$
(1.3)
In this formula, $`\chi ^{(2)}(Y)`$ is the $`L^2`$-Euler characteristic of $`Y`$, which in our setting is simply equal to $`\chi (B)`$. The determinant $`\mathrm{Det}_\pi \mathrm{\Delta }(Y,u)`$ is defined via formal power series as $`(\mathrm{Exp}Tr_\pi \mathrm{Log})\mathrm{\Delta }(Y,u)`$ and converges for small $`u`$. More precisely, if $`Y`$ is $`q+1`$ regular then the radius of convergence of $`Z_\pi (Y,u)`$ is greater than or equal to $`\frac{1}{q}.`$
###### Theorem 1.5.
Let $`Y`$ be a $`q+1`$ regular graph. Then $`Z_\pi (Y,u)`$ has a holomorphic extension to $`\mathrm{\Omega }.`$
###### Proof.
By (1.3) it is enough to show that $`\mathrm{Det}_\pi \mathrm{\Delta }(Y,u)`$ has a holomorphic extension on $`\mathrm{\Omega }`$. Let $`g_u(\lambda )=\mathrm{log}(1\lambda u+qu^2)`$. Here and in the rest of the paper $`\mathrm{log}`$ is the principal branch of the logarithm, defined and analytic on $`(\mathrm{},0]`$. Fix $`u\mathrm{\Omega }`$. Then using Proposition 1.3, there exists an open set $`V_u[(q+1),q+1]`$ on which $`g_u`$ is analytic. Since $`\delta `$ is self-adjoint and $`\delta q+1`$, the spectrum $`\sigma (\delta )[(q+1),q+1]`$. By the spectral theorem for self-adjoint operators we can write:
$$\delta =_k^{}^k\lambda 𝑑E(\lambda ).$$
(1.4)
Now
$$g_u(\delta )=_k^{}^k\mathrm{log}(1\lambda u+qu^2)𝑑E(\lambda )$$
(1.5)
is well defined, and
$$\mathrm{Tr}_\pi g_u(\delta )=_k^{}^k\mathrm{log}(1\lambda u+qu^2)d(\mathrm{Tr}_\pi E(\lambda ))$$
(1.6)
is a holomorphic function of $`u`$ on $`\mathrm{\Omega }`$. Now for small $`u`$,
$`\mathrm{Det}_\pi (\mathrm{\Delta }(Y,u))`$ $`=\mathrm{Exp}\mathrm{Tr}_\pi \mathrm{Log}(\mathrm{\Delta }(Y,u))`$ (1.7)
$`=\mathrm{Exp}\mathrm{Tr}_\pi \mathrm{Log}(I\delta u+qu^2)`$
$`=\mathrm{Exp}\mathrm{Tr}_\pi g_u(\delta ).`$
###### Remark.
Using the functional equation (1.2) it is possible to extend the $`L^2`$-zeta function to the exterior of $`C`$. Defining the zeta function on $`C`$ itself presents some problems. In examples, the function on the interior of $`C`$ and the exterior of $`C`$ do not match continuously on $`C`$. However the absolute value of the resulting function is the Fuglede-Kadison determinant, and may extend continuously to $`C`$.
## 2. Convergence of Zeta Functions for Towers of Graphs
In this section we prove that the zeta functions for a tower of finite graphs, suitably renormalized, converge to the $`L^2`$-zeta function for an infinite covering graph. The argument uses an idea from .
### 2.1. The Convergence Theorem
###### Theorem 2.1.
Let $`Y`$ be a $`q+1`$ regular graph. Suppose the group $`\pi `$ acts freely on $`Y`$, and that $`B=\pi \backslash Y`$ is a finite graph. Suppose $`\pi =\pi _1\pi _2\mathrm{}`$ is a tower of finite index normal subgroups and $`\pi _i=\{e\}`$. Let $`[\pi :\pi _i]=N_i`$ and $`B_i=\pi _i\backslash Y`$. Then for $`u\mathrm{\Omega }`$, we have
$$\underset{i\mathrm{}}{lim}Z(B_i,u)^{\frac{1}{N_i}}=Z_\pi (Y,u).$$
(2.1)
The convergence is uniform on compact subsets of $`\mathrm{\Omega }`$.
###### Remark.
Notice
$`Z(B_i,u)`$ $`=(1u^2)^{\chi (B_i)}\mathrm{Det}\mathrm{\Delta }(B_i,u)`$ (2.2)
$`=(1u^2)^{N_i\chi (B)}\mathrm{Det}\mathrm{\Delta }(B_i,u).`$
So we let
$$Z(B_i,u)^{1/N_i}=(1u^2)^{\chi (B)}(\mathrm{Det}\mathrm{\Delta }(B_i,u))^{1/N_i}$$
(2.3)
The $`N^{\text{th}}`$ root in (2.3) is taken in the sense of Corollary 1.4.
###### Proof of Theorem 2.1.
From the remark we need to show that
$$\underset{i\mathrm{}}{lim}(\mathrm{Det}\mathrm{\Delta }(B_i,u))^{1/N_i}=\mathrm{Det}\mathrm{\Delta }(B,u).$$
(2.4)
Let $`F_i(\lambda )=\frac{1}{N_i}|\{\mu \text{ eigenvalue of }\delta _i\text{}\mu \lambda \}|`$. Let $`F(\lambda )=\mathrm{Tr}_\pi E(\lambda )`$, where $`\{E(\lambda )\}_\lambda `$ is the spectral decomposition of $`\delta `$ acting on $`l^2(Y).`$ We now set
$$\begin{array}{cccccc}\hfill \overline{F}(\lambda )& =& lim\; sup_i\mathrm{}F_i(\lambda );\hfill & \hfill \underset{¯}{F}(\lambda )& =& lim\; inf_i\mathrm{}F_i(\lambda );\hfill \\ \hfill \overline{F}^+(\lambda )& =& lim_{\epsilon 0^+}\overline{F}(\lambda +\epsilon );\hfill & \hfill \underset{¯}{F}^+(\lambda )& =& lim_{\epsilon 0^+}\underset{¯}{F}(\lambda +\epsilon ).\hfill \end{array}$$
Then from \[7, Theorem 2.3.1\], for all $`\lambda [k,k]`$,
$$F(\lambda )=\overline{F}^+(\lambda )=\underset{¯}{F}^+(\lambda )$$
We know
$$\mathrm{log}\mathrm{Det}\mathrm{\Delta }(Y_i,u)=_k^{}^k\mathrm{log}(1u\lambda +qu^2)𝑑F_i(\lambda ).$$
(2.5)
By (1.6)
$$\mathrm{log}\mathrm{Det}_\pi \mathrm{\Delta }(Y,u)=_k^{}^k\mathrm{log}(1u\lambda +qu^2)𝑑F(\lambda ).$$
(2.6)
If $`K\mathrm{\Omega }`$ is compact then $`\mathrm{log}(1u\lambda +qu^2)`$ is bounded uniformly for $`uK`$ and $`\lambda `$ in an open interval containing $`[k,k].`$ Now integration by parts shows that indeed
$$_k^{}^k\mathrm{log}(1u\lambda +qu^2)𝑑F_i(\lambda )_k^{}^k\mathrm{log}(1u\lambda +qu^2)𝑑F(\lambda ).$$
(2.7)
as $`i\mathrm{}`$. ∎
### 2.2. Examples
###### Example 2.2.
Let $`Y`$ be the line, with vertices $``$ and edges connecting $`n`$ to $`n+1`$. The group $`\pi =`$ acts on $`Y`$ with quotient having one vertex and one edge. Let $`\pi _n=n`$, so $`B_n=\pi _n\backslash Y`$ is an $`n`$-cycle. As $`B_n`$ has only two primitive loops and $`Y`$ has none, we have $`Z(B_n,u)=(1u^n)^2`$ and $`Z_\pi (Y,u)=1`$. The graphs are 2-regular, so $`\mathrm{\Omega }`$ is the (open) unit disk. For $`u\mathrm{\Omega }`$,
$$\underset{n\mathrm{}}{lim}(1u^n)^{2/n}=1.$$
Notice that the functional equation (1.2) gives $`_\pi (Y,u)=u^2`$ outside of the unit disk.
###### Example 2.3.
Let $`Y`$ be the $`q+1`$ regular tree and $`B=\pi \backslash Y`$ finite. Since $`Y`$ has no closed loops, the $`L^2`$-zeta function of $`Y`$ is the constant function 1. So if $`B_n=\pi _n\backslash Y`$ is a tower of finite graphs covering $`B`$,
$$\underset{n\mathrm{}}{lim}Z(B_n,u)^{1/[\pi :\pi _n]}=1$$
for $`u\mathrm{\Omega }`$. This result is contained in for $`|u|`$ small.
###### Theorem 2.4.
Let $`Y`$ and $`B`$ be as in the previous example. We have:
$$Z(B,u)=\frac{\mathrm{Det}\mathrm{\Delta }(B,u)}{\mathrm{Det}_\pi \mathrm{\Delta }(Y,u)}$$
for $`u\mathrm{\Omega }`$.
For small $`u`$, the above theorem is the main result of .
###### Proof.
From \[2, Theorem 0.3\] we know that
$$Z_\pi (Y,u)=(1u^2)^{\chi ^{(2)}(Y)}\mathrm{Det}_\pi \mathrm{\Delta }(Y,u).$$
As in the previous example, $`Z_\pi (Y,u)=1`$. Now we have
$$Z(B,u)=\frac{Z(B,u)}{Z_\pi (Y,u)}=\frac{(1u^2)^{\chi (B)}\mathrm{Det}\mathrm{\Delta }(B,u)}{(1u^2)^{\chi ^{(2)}(Y)}\mathrm{Det}_\pi \mathrm{\Delta }(Y,u)}.$$
As $`\chi (B)=\chi ^{(2)}(Y)`$, the theorem follows. ∎
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# Gravity with a dynamical preferred frame
## I Introduction
The Lorentz group is non-compact, since the boost parameter is unbounded. This makes exact Lorentz invariance impossible to test uniformly. Lorentz invariance has thus been tested only up to some maximum boost and beyond that lies an infinite volume of uncharted territory in the Lorentz group. Contrast this with the rotation group. Rotation invariance can be tested by filling in the compact $`SO(3)`$ group manifold more and more densely with data points, or by checking a few randomly selected rotations. The rotation group can be and has been uniformly explored.
There is also reason to doubt exact Lorentz invariance: it leads to divergences in quantum field theory associated with states of arbitrarily high energy and momentum. This problem can be cured with a short distance cutoff which, however, breaks Lorentz invariance.
For these reasons we entertain the possibility that there is a preferred rest frame at each spacetime point. In particular, we seek a viable effective field theory incorporating a breaking of local Lorentz invariance.
If the preferred frame were to be a fixed external structure, then it would violate general covariance, which would require us to abandon general relativity (or any generally covariant modification thereof). General covariance ordinarily implies that the divergence of the matter energy-momentum tensor $`T_{ab}`$ vanishes when the matter field satisfies its equation of motion. This is required for consistency of the Einstein field equation<sup>*</sup><sup>*</sup>*We use units with $`c=1`$ and the metric signature $`(+)`$. $`G_{ab}=8\pi GT_{ab}`$, since the divergence of the Einstein tensor $`G_{ab}`$ is identically zero by virtue of the contracted Bianchi identity. If a fixed preferred frame is introduced into the matter action, for example, general covariance is lost, $`T_{ab}`$ is not divergenceless, and the Einstein equation is inconsistent. We therefore seek to incorporate the preferred frame while preserving general covariance, which requires that the preferred frame be dynamical.
It would be most appealing if the preferred frame were somehow determined by the spacetime metric itself. As discussed below, in spacetimes with an initial singularity the metric can be used to define a cosmological time function, the gradient of which then determines a preferred frame (which is by construction timelike). However, the nonlocal relationship between this frame and the metric results in infinitely nonlocal field equations if this frame is incorporated into the action principle. Theories with even a finite amount of nonlocality are generally horribly unstable, so we do not consider this a viable approach.
To avoid such unacceptable behavior the preferred frame should arise from local conditions, which of course reflect conditions at earlier times but only through dynamics. For example, this dynamical frame could be defined by a vector field or by the gradient of a scalar field. In these approaches the presence or absence of a preferred frame depends on the field configuration, since the preferred vector may vanish or may not be timelike.
Since our motivation arises from doubts about the fundamental validity of exact Lorentz invariance, we are not interested in a theory possessing a Lorentz invariant phase. We wish to study instead an effective theory in which there is always a preferred frame. This frame is defined by a timelike direction or, equivalently, by a unit timelike (contravariant or covariant) vector field.A unit timelike vector contains a discrete piece of information that a frame by itself does not have, namely, a time orientation. The theory may or may not depend on this orientation. Such a field carries a nonlinear representation of the local Lorentz group since the field takes values not in a vector space but on the unit hyperboloid in the tangent space. This could therefore be called a theory of nonlinearly realized Lorentz invariance. It is analogous to a gauge theory with a non-linear sigma model Higgs field of fixed norm. There seems to be no generally accepted terminology for this sort of symmetry structure. Since the symmetry breaking unit vector field is not a state-dependent expectation value but rather breaks this symmetry in all states, it may be misleading to say the symmetry is “spontaneously broken”. For lack of a better idea we shall take refuge in ambiguity and just call it ”broken”. What this really means is that in order to implement the local Lorentz symmetry one must transform not only the matter fields but also the background unit vector, so for practical purposes it will appear as what would normally be called broken Lorentz symmetry.
The theory described here was devised about a decade ago by J.C. Dell together with one of us, and we have since learned that similar ideas have independently been studied previously. The non-gravitational part of the theory (and generalizations thereof) was considered by Dirac in the early 1950’s as a new theory of electrons (in which the unit timelike vector played the dual role of gauge-fixed vector potential and flow vector of a stream of charged dust). A class of generally covariant theories breaking Lorentz invariance was studied by Gasperini in many papers. In this work the tetrad formalism was used, and the local Lorentz symmetry was broken by including in the action terms referring to a fixed “internal” unit timelike vector. This is equivalent to our formulation in terms of the metric and a unit timelike vector. To see the equivalence, note that the tetrad defines a metric and associates to the fixed internal vector a unit vector field on spacetime. The only other information in the tetrad is the gauge freedom parametrized by the local rotations leaving invariant the preferred timelike vector. Eliminating this gauge freedom leads to the formalism used in this paper. Gasperini has studied both cosmological and central field solutions, with various choices for the specific form of the second derivative terms in the Lagrangian. In the present paper we focus on a different Lagrangian.
The particular Lagrangian studied here was also considered by Kostelecký and Samuel, as a simplified model of the spontaneous Lorentz symmetry breaking that might occur in string theory, although in Ref. the preferred vector was not necessarily timelike. More generally, those authors argued that spontaneous Lorentz symmetry breaking in string theory may produce vacuum expectation values of more than one tensor field. In this case, rather than having just a single “preferred frame” there might be several background Lorentz tensors which collectively break part or all of the Lorentz symmetry.
The remainder of this paper is organized as follows. In section II we explain the nonlocality problem encountered if a cosmological time is used to define the preferred frame. In section III our proposed field theory of a preferred frame is formulated and its general properties are investigated. It is seen that the solutions for our theory comprise a subset of the solutions to the coupled Einstein-Maxwell-charged dust equations. Several types of exact solutions to the field equations are characterized in section IV, and the linearized theory is studied in section V. Coupling of the preferred frame to matter fields is discussed in section VI. Both the dimension $`4`$ couplings and some higher dimension ones are examined. The paper concludes with a brief discussion in section VII.
Among the dimension $`>4`$ couplings are included theories involving Lorentz non-invariant dispersion at high wavevectors. This is motivated by recent work in which high frequency dispersion is invoked to avoid the role of trans-Planckian modes in the Hawking effect (for a review see ). In this framework such theories can be formulated in a generally covariant manner so that gravitational effects can be consistently incorporated. Higher dimension couplings also also provide an alternative generally covariant formulation of variable speed of light models, in which different fields propagate at different speeds possibly at different cosmological epochs. Such models have recently been of interest as potential alternatives to standard cosmology, and have been given generally covariant formulations using additional fields to define the preferred frame.
## II Cosmological times
In this section we briefly describe the construction of a cosmological time function determined purely by the metric, and the reason for rejecting it for the purposes of an effective theory of local Lorentz symmetry breaking.
The cosmological metric of our universe, by virtue of its (approximate) homogeneity, (approximately) defines a preferred spacelike foliation of spacetime. However, this particular definition of the time function relies on the symmetry of the spacetime. For a workable theory with general covariance what is needed is a definition of cosmological time that can be used independently of symmetry.
It is difficult to think of a notion of cosmological time that would make sense for all possible cosmologies. However, if we restrict attention to spacetimes with a “beginning”, then two notions of cosmological time at $`x`$ present themselves: ($`i`$) volume time, the spacetime volume (or perhaps the fourth root thereof) of the past set $`I^{}[x]`$, and ($`ii`$) maximal time, the maximal proper time along a causal curve going back to the initial singularity.The maximal time function has been discussed in Refs. . In particular, a powerful theorem proved in Ref. establishes a number of properties of this function under the further assumption that the initial singularity is the only place past directed causal curves can end. Other possibilities are combinations or smoothed averages of these times. Both of these time functions are determined non-locally but causally by the spacetime to the past of $`x`$. They may or may not be sufficiently smooth functions to enter meaningfully into a local action principle.<sup>§</sup><sup>§</sup>§One of the results of the theorem of Ref. referred to in the previous footnote is that the maximal time function is locally Lipschitz and its first and second derivatives exist almost everywhere. The volume time function may well be even better behaved. If we assume that they are indeed sufficiently smooth we find that there is in any case a fatal problem with using them in this manner, as will now be explained.
Suppose that to the usual action for gravity and matter fields is added a term involving one of the above cosmological times,
$$S=S_{\mathrm{local}}+S_\tau .$$
(1)
We assume that the equations of motion are obtained as usual by requiring that the action is stationary with respect to variations of the fields. The variational derivative of the action with respect to $`g_{ab}(x)`$ has the form
$$\frac{\delta S}{\delta g_{ab}(x)}=\frac{\delta S_{\mathrm{local}}}{\delta g_{ab}(x)}+\frac{\delta ^{}S_\tau }{\delta g_{ab}(x)}+d^4x^{}\frac{\delta S_\tau }{\delta \tau (x^{})}\frac{\delta \tau (x^{})}{\delta g_{ab}(x)},$$
(2)
where the prime on $`\delta ^{}`$ indicates that the metric dependence of $`\tau (x)`$ is not included in the variation. Since $`\delta \tau (x^{})/\delta g_{ab}(x)`$ has support when $`x^{}`$ lies to the future of $`x`$, the field equation $`\delta S/\delta g_{ab}(x)=0`$ involves the values of the fields to the future of $`x`$. Indeed the metric field equation is infinitely non-local in time, since the time function at any point to the future can be affected by a metric variation at $`x`$. Even finite nonlocality in time leads to unphysical instability , so this approach to incorporating a preferred frame must be rejected. If the action depends on $`\tau `$ only through its derivative $`_a\tau `$, then the equation of motion would be causal if $`\delta _c\tau (x^{})`$ depended only on $`g_{ab}(x)`$ at $`x=x^{}`$. However, this is not the case for the either the volume time or the maximal time.
## III Aether dynamics
We now turn to a class of theories in which there is a preferred frame which is determined by a local field. It is convenient to give a name to this field, and “aether” seems as good a name as any. Let us take the aether field to be a unit timelike vector field $`u^a`$, which is dimensionless, like the metric. To handle the condition that $`u^a`$ is a unit vector, we include in the action a Lagrange multiplier term. Note that we are implicitly assuming that the spacetime admits a globally defined unit timelike vector field which is the case if and only if the spacetime is time orientable.
### A Action
The most general Lagrangian involving the metric and the aether with two or fewer derivatives is, up to a total divergence,
$`_{g,u}`$ $`=`$ $`a_0a_1Ra_2R_{ab}u^au^b`$ (4)
$`b_1F^{ab}F_{ab}b_2(_au_b)(^au^b)b_3\dot{u}^a\dot{u}_a,`$
where $`\dot{u}^a:=u^m_mu^a`$, and $`F_{ab}`$ is defined in analogy to the electromagnetic field strength,
$$F_{ab}:=2_{[a}u_{b]}.$$
(5)
The term $`(_au^a)^2`$ is equivalent, via integration by parts, to the combination $`(_au_b)(^au^b)(1/2)F^{ab}F_{ab}+R_{ab}u^au^b`$, so has not been included in (4).
The Lagrangian (4) is similar to the one discussed in Ref. as the most general Lagrangian for a vector-tensor theory of gravity including terms up to second order in derivatives and quadratic in the vector field. The differences are that (i) the terms $`u^2`$ and $`Ru^2`$ are missing from our action since the vector field is constrained to be a unit vector, and (ii) we have included the quartic term $`b_3\dot{u}^2`$, which was omitted in Ref. because it is not quadratic in $`u^a`$. Note that, even without the last term, our theory is not a special case of the vector-tensor theories discussed in Ref. , since the constraint $`u^2=1`$ affects the field equations.
The coefficient $`a_0`$ in the action (4) has mass dimension 4 while $`a_{1,2}`$ and $`b_{1,2,3}`$ have mass dimension 2. Lacking the underlying fundamental theory we do not try to assign a priori the values of these coefficients. A partial analysis of the observational consequences and limits on them has been done for the vector-tensor theories, however that analysis does not apply directly to our case due to the presence of the constraint term. It is fairly clear nevertheless that whatever values $`b_{1,2,3}`$ take, agreement with observation will require that $`a_2a_1`$, and that $`a_0/a_1`$ (which is basically the cosmological constant) must not be much larger than the squared Hubble constant.
In this initial foray we shall restrict attention to the simple case in which the only terms with non-zero coefficients are $`R`$ and $`F^2`$. That is, we set $`a_0=a_2=b_2=b_3=0`$. The minimal theory we consider is thus defined by the action
$`S_{min}[`$ $`g_{ab},u^a,\lambda ]={\displaystyle }d^4x\sqrt{g}`$ (7)
$`\left(a_1Rb_1F^{ab}F_{ab}+\lambda (g_{ab}u^au^b1)\right).`$
This minimal theory is one of the models considered by Kostelecký and Samuel in the paper mentioned in section I. Those authors studied a broader class of models in which $`u^a`$ is not necessarily constrained to be a unit vector but rather possesses a Lorentz-invariant potential energy with a minimum at some fixed value of $`u_au^a`$. They also allowed for extra, compact spatial dimensions of spacetime, and examined cases where the symmetry breaking vector lies in the extra dimensions as well as cases where it lies in the four ordinary spacetime dimensions. Our paper by contrast is restricted to four dimensions and to a timelike vector of fixed norm. Later in this paper we shall also add matter terms to the action, including terms which couple the aether field to the matter.
Note that $`F_{ab}`$ is invariant under the “gauge transformation”
$$u_au_a+_af,$$
(8)
however the constraint $`u^2=1`$ does not share this symmetry (nor do the additional couplings in general), so the theory is certainly not “gauge invariant”. The constraint does have a limited version of this symmetry however, namely for those functions $`f`$ satisfying $`(u_a+_af)(u^a+^af)=u_au^a=1`$. The general solution to this equation is
$$u^a_af=1\pm \sqrt{1+q^{ab}_af_bf},$$
(9)
where
$$q_{ab}:=g_{ab}+u_au_b$$
(10)
is the (positive definite) spatial metric orthogonal to $`u^a`$. Thus the action (7) is invariant under the gauge transformation (8) if $`f`$ is chosen arbitrarily on a spacelike surface and then determined uniquely elsewhere (up to a discrete choice of sign) by integration of (9) along the flow of $`u^a`$.
### B Field equations
The equations of motion arising from the action (7) are
$$G_{ab}=\frac{2b_1}{a_1}(F_{am}F_b{}_{}{}^{m}\frac{1}{4}F^2g_{ab})+\frac{\lambda }{a_1}u_au_b,$$
(11)
$$_aF^{ab}=\frac{\lambda }{2b_1}u^b,$$
(12)
$$g_{ab}u^au^b=1.$$
(13)
The metric equation (11) has the form of the Einstein equation $`G_{ab}=8\pi GT_{ab}`$, where $`G=1/16\pi a_1`$, and the stress tensor receives contributions from both the $`F^2`$ term and the constraint term in the action. (The constraint equation (13) has been used to drop the contribution to (11) that would have come from the variation of $`\sqrt{g}`$ in the constraint term.) The contribution from the constraint term looks like that of a (pressureless) dust with rest energy density $`2\lambda `$, and that from the $`F^2`$ term is the usual Maxwell tensor familiar from electromagnetism, if we identify the vector potential as
$$\text{}A_m\text{}2\sqrt{b_1}u_m.$$
(14)
The stress tensor thus satisfies the usual energy conditions provided $`b_1/a_1`$ and $`\lambda /a_1`$ are positive.
In terms of the vector potential $`A_m`$ (14) the constraint equation (13) becomes
$$A_mA^m=4b_1$$
(15)
which can be interpreted as a gauge condition. The aether field equation (12) becomes the Maxwell equation with source equal to the current of a charged dust fluid with 4-velocity $`u_b`$ and charge density $`(\lambda /\sqrt{b_1})`$. The evolution of $`\lambda `$ is determined by the current conservation equation which follows from divergence of the aether field equation (12) upon using the identity $`_a_bF^{ab}0`$. Thus $`\lambda `$ satisfies a first order ordinary differential equation along the flow lines of $`u^a`$:
$$u^a_a\lambda =\lambda _au^a.$$
(16)
In particular, if $`\lambda `$ vanishes on a Cauchy surface, it must vanish everywhere. Also, the sign of $`\lambda `$ on a given flow line cannot change, since if $`\lambda =0`$ at any point on a flow line it must vanish everywhere on that line.
#### Relation to Einstein-Maxwell-charged dust system
We have just seen that the field equations of the minimal theory take the form of the coupled Einstein-Maxwell equations, with a charged dust matter source possessing charge to mass ratio $`1/2\sqrt{b_1}`$. There is no explicit equation of motion for the dust, however the normalization condition (13) provides such an equation. Taking the gradient of $`u^2=1`$ we have
$`0`$ $`=`$ $`_a(u^bu_b)`$ (17)
$`=`$ $`2u^b_au_b`$ (18)
$`=`$ $`2(u^b_bu_a+u^bF_{ab}).`$ (19)
Let us define $`\stackrel{~}{F}_{ab}=2_{[a}A_{b]}=2\sqrt{b_1}F_{ab}`$. Then (19) becomes
$$u^b_bu_a=\frac{1}{2\sqrt{b_1}}\stackrel{~}{F}_{ab}u^b,$$
(20)
which is the equation of motion for a particle in the electromagnetic field $`\stackrel{~}{F}_{ab}`$, with charge to mass ratio $`1/2\sqrt{b_1}`$, the same ratio we inferred from the Einstein and Maxwell equations! Thus any solution of our minimal theory is a solution of the Einstein-Maxwell-charged dust equations (although the converse is not true). The equivalence to a subset of the charged dust solutions demonstrates that the equations of our theory admit an initial value formulation, and it provides some useful intuition about the nature of the solutions.
Our theory is not equivalent to the Einstein-Maxwell-charged dust system because in the general solution of that system the dust 4-velocity is not proportional to the vector potential in some gauge. That is, although there is always a gauge transformation that will make $`A_m/2\sqrt{b_1}`$ a unit vector, it cannot in general be made to coincide with the dust 4-velocity.The general form of the discrepancy between these two 4-vectors was found by Dirac (see the second paper of Ref.), who showed that (in four spacetime dimensions) there is always a gauge in which $`A_m/2\sqrt{b_1}=u_m+\xi _m\eta `$, where $`\xi `$ and $`\eta `$ are functions that are constant along the flow lines of $`u^a`$. Dirac included the functions $`\xi `$ and $`\eta `$ as dynamical variables in order to obtain a theory in which arbitrary electron streams were admitted. In the third paper of Ref. he allowed for multiple streams.
## IV Solutions
In this section we characterize a few types of solutions to the field equations.
### A Solutions with $`\lambda =0`$
If $`\lambda =0`$, then the two field equations (11,12) are just the Einstein-Maxwell equations. Any solution to these equations is a solution in our theory provided a gauge can be chosen so that the constraint equation (15) is satisfied. Such a gauge always exists, at least locally.
### B Solutions with $`F_{ab}=0`$
A special class of solutions to the field equations with $`\lambda =0`$ are those with $`F_{ab}=0`$. For such fields, (12) implies that $`\lambda =0`$, and the field equations (11-13) reduce to the ordinary vacuum Einstein equation together with the constraint $`u^2=1`$. When $`F_{ab}=0`$ it follows, at least locally, that $`u_a=_a\tau `$ for some function $`\tau `$, and the constraint then implies that $`_a\tau ^a\tau =1`$. The general solution for such a function $`\tau `$ can specified by assigning the value $`\tau =0`$ to an arbitrary spacelike surface, and determining $`\tau `$ elsewhere by “uniform normal extension”, i.e. by the differential equation $`n^a_a\tau =1`$, where $`n^a`$ is the unit normal to the surface.
Another way to think of this construction is in terms of the congruence of integral curves of $`u^a`$. When $`F_{ab}=0`$, eqn. (19) implies that these curves are geodesics. Moreover, if $`u^a`$ is the unit tangent field to a congruence of geodesics, then $`F_{ab}=2_{[a}u_{b]}=0`$ if and only if the congruence is hypersurface-orthogonal. Hence the general solution of this type is just an arbitrary solution to the Einstein equation, together with $`u^a`$ given by the unit tangent field of any hypersurface-orthogonal congruence of timelike geodesics in this metric. A special case is flat spacetime, where the $`u^a`$ congruence consists of straight lines normal to an arbitrary initial spacelike hypersurface.
#### Singular aether evolution
This characterization of the $`F_{ab}=0`$ solutions shows that, at least for such solutions, the evolution of $`u_a`$ is generally singular. The geodesics launched normally from a spacelike surface will typically cross. Where they do, the quantity $`_au_b`$ will diverge.
The existence of such singular evolutions for $`u_a`$ signals a breakdown of the effective theory we are using. Perhaps it would be cured by including the term $`(_{(a}u_{b)})(^{(a}u^{b)})`$ in the action. (Without this term the action is insensitive to gradients for which the antisymmetrized derivative $`_{[a}u_{b]}`$ vanishes.) For the purposes of the present paper we shall not pursue this question, but it should be addressed.
#### Cosmological solutions
If $`u^a`$ shares the symmetry of a homogeneous isotropic cosmological metric, then $`F_{ab}=0`$. The presence of the aether field therefore has no influence on the cosmological evolution unless there are additional terms in the action beyond the minimal model. In Ref. we examine some cosmological effects of coupling to a scalar field though a fourth spatial derivative term as discussed in Sect. VI B.
#### Black hole solutions
For a spherically symmetric black hole, a suitable congruence of geodesics is given by the radial free-fall trajectories that all have the same Killing energy, i.e. the same asymptotic velocity at spatial infinity. The same construction can even be applied in the case of a Kerr black hole, at least for the geodesics that are at rest at spatial infinity. This follows from the work of Ref. , in which this congruence is employed to construct a coordinate system for the Kerr metric using the time function $`\tau `$ mentioned above.
### C Spherically symmetric, static solutions
Here we seek to characterize the general spherically symmetric, static solution. We shall find that, besides the mass, the metric in these solutions has an additional free parameter, the “aether charge”.
Some of the linearized static, spherically symmetric solutions were previously studied in Ref. . Those authors examined the case where $`u^a`$ is spacelike, and while in four spacetime dimensions they restricted attention to vanishing Lagrange multiplier $`\lambda `$ and vanishing field strength $`F_{ab}`$. In the present work we treat the nonlinear case, considering only timelike $`u^a`$ and imposing no further restrictions on the fields.
Coordinates can be chosen so the line element takes the form
$$ds^2=g_{tt}dt^2+g_{rr}dr^2r^2(d\theta ^2+\mathrm{sin}^2\theta d\phi ^2),$$
(21)
and the aether field takes the form
$$u=u_t(r)dt+u_r(r)dr.$$
(22)
The only potentially nonzero component of $`F_{ab}`$ is then $`F_{rt}=_ru_t`$, and the constraint (13) implies
$$g^{tt}u_t^2+g^{rr}u_r^2=1.$$
(23)
The aether field equation (12) in coordinate form reads
$$\frac{1}{\sqrt{g}}_\alpha \left(\sqrt{g}g^{\alpha \mu }g^{\beta \nu }F_{\mu \nu }\right)=\frac{\lambda }{2b_1}u^\beta ,$$
(24)
or, taking into account the form of the metric (21),
$$\frac{1}{\sqrt{g}}_r\left(\sqrt{g}g^{rr}g^{\beta t}F_{rt}\right)=\frac{\lambda }{b_1}u^\beta .$$
(25)
The left hand side vanishes when $`\beta =r`$, hence the field equation implies that $`\lambda u_r=0`$, which in turn implies that either $`u_r=0`$ or $`\lambda =0`$. In the former case, $`u^a`$ is proportional to the timelike Killing field itself. There are thus two cases to consider.
If $`\lambda 0`$ there are in fact no static solutions, unless the coefficients in the action are such that the charge to mass ratio of the dust is extremal. Recall that any solution to our theory is a solution to the charged dust theory. However under the influence of gravitational and electric forces, the non-extremal charged dust cannot remain static, since there is no pressure.
If $`\lambda =0`$ then (cf. Sect. IV A) these are just the spherically symmetric static Einstein-Maxwell solutions, i.e. the Reissner-Nordstrom solutions, in a spherically symmetric, static gauge with fixed norm (15). Such a gauge always exists, at least locally. Consider the gauge transform $`A_\mu =A_\mu ^{}+\alpha _{,\mu }`$ of an arbitrary vector potential $`A_\mu ^{}`$. To maintain spherical symmetry and time independence we must have $`\alpha =\beta t+\gamma (r)`$ (using the coordinates in (21)), so that $`A_t=A_t^{}+\beta `$ and $`A_r=A_r^{}+\gamma _{,r}`$. The normalization (15) then implies
$$\gamma _{,r}=A_r^{}\pm \sqrt{g^{tt}\left[g^{tt}(A_t^{}+\beta )^24b_1\right]},$$
(26)
where we have used the fact that for the Reissner-Nordstrom metrics $`g_{rr}=1/g_{tt}`$. In any region where $`g^{tt}`$ and $`A_t`$ are bounded one can always choose $`\beta `$ large enough so that the radical is real, and then $`\gamma (r)`$ can be found by integration.
If a horizon is present $`g^{tt}`$ diverges and it is not clear from the preceding discussion whether the unit timelike gauge can be accessed in a smooth way across the horizon. Indeed it can, however the maximal extension of the region over which such a gauge can be accessed depends on the parameter $`b_1`$. This can also be understood from the equivalence with a charged dust solution. The radial congruence $`u^a`$ must satisfy the Lorentz force equation (20), and this congruence can be nonsingular and time-independent only if the trajectories are monotonic in the $`r`$ coordinate. In general, however, the trajectories bounce inside the black hole.
#### Comparison with observation
To compare with observation it would be necessary to determine which of the above solutions to use in the presence of a spherically symmetric static source such as a planet, star, or black hole. The metric associated with one of these objects depends on its “aether charge” —the charge of the Reissner-Nordstrom solution—which is determined by the “charge” of the “aether dust” that fell in when the object condensed. The choice is determined by the initial conditions on $`\lambda `$, which are presumably cosmological in origin. We have no theory of these initial conditions at this stage, but agreement with observations can put a bound on the amount of aether charge. If this charge is zero, then we have the usual Schwarzschild solution of general relativity (and, as discussed in section IV B, the aether field is the tangent field to a hypersurface orthogonal congruence of timelike geodesics), which of course agrees with observations.
## V Linearized theory
In this section we study the linearized equations defined by expanding about a background solution,
$`g_{ab}`$ $`=`$ $`g_{ab}^{(0)}+h_{ab}`$ (27)
$`u_a`$ $`=`$ $`u_a^{(0)}+v_a,`$ (28)
$`\lambda `$ $`=`$ $`\lambda ^{(0)}+\lambda ^{(1)}.`$ (29)
For the background we take the flat metric $`g_{ab}^{(0)}=\eta _{ab}`$ and a constant $`u_a^{(0)}`$. In this background solution the equations of motion imply that the Lagrange multiplier $`\lambda ^{(0)}`$ must vanish, hence we shall use the letter $`\lambda `$ for the perturbation $`\lambda ^{(1)}`$. In this section we use the flat background metric to raise and lower indices. Note that we use the perturbation of the covariant vector $`u_a`$ to define the perturbation $`v_a`$.
We choose cartesian coordinates $`(x^0,x^i)`$, $`i=1,2,3`$, in which the components of $`\eta _{ab}`$ are $`\text{diag}(1,1,1,1)`$ and those of $`u_a^{(0)}`$ are $`(1,0,0,0)`$. The linearized field equations for this theory were first written down in a general gauge in Ref., who also pointed out that the Lorentz gauge can be accessed using the linearized diffeomorphism invariance of the action. The Lorentz gauge condition for the metric perturbation is $`\overline{h}_{ab}{}_{}{}^{,b}=0`$, where $`\overline{h}_{ab}h_{ab}\frac{1}{2}h\eta _{ab}`$ is the trace-reversed metric perturbation. In this gauge, the linearized equations of motion are
$$\mathrm{}\overline{h}_{ab}=2\frac{\lambda }{a_1}u_a^{(0)}u_b^{(0)}$$
(30)
$$\mathrm{}v_b_b(^av_a)=\frac{\lambda }{2b_1}u_b^{(0)}$$
(31)
$$h_{ab}u^{(0)a}u^{(0)b}+2u^{(0)b}v_b=0$$
(32)
In a source-free region, the residual gauge freedom is usually employed to set $`\overline{h}_{0i}=h_{0i}=0`$ and $`\overline{h}=h=0`$. The possibility of doing so depends on the fact that these quantities satisfy the wave equation in general relativity. In our case, the components $`h_{0i}`$ satisfy the wave equation, however $`h`$ does not, due to the source term on the right hand side of (30) that is always present (even outside matter) unless $`\lambda =0`$. This source corresponds to the energy density of the “charged dust”, and we wish to allow for the presence of this term. Therefore, rather than setting $`h`$ to zero, we choose to set to zero the trace of the spatial part, $`\overline{h}_i^i`$, which does satisfy the wave equation. The proof that this can be done follows the same logic as in the usual case. This gauge condition implies $`\overline{h}=\overline{h}_{00}`$, hence $`h_{00}=\overline{h}_{00}\frac{1}{2}\overline{h}\eta _{00}=\frac{1}{2}\overline{h}_{00}`$, and the Lorentz gauge condition implies $`_0\overline{h}_{00}=0`$ and $`_i\overline{h}_{ij}=0`$. Thus $`h_{00}`$ is time-independent and the spatial part $`\overline{h}_{ij}`$ is a transverse traceless solution to the wave equation. This is not quite the same as the usual transverse traceless gauge in general relativity however, since $`h_{ij}=\overline{h}_{ij}h_{00}\delta _{ij}`$, so $`h_{ij}`$ is not transverse unless $`_ih_{00}=0`$, and $`h=2h_{00}0`$.
It remains to consider the linearized equations for $`h_{00}`$, $`v_a`$, and $`\lambda `$. The $`00`$-component of the metric equation (30) and the constraint equation (32) determine $`\lambda `$ and $`v_0`$ in terms of $`h_{00}`$:
$`\lambda `$ $`=`$ $`a_1^2h_{00}`$ (33)
$`v_0`$ $`=`$ $`\frac{1}{2}h_{00}.`$ (34)
Using the time-independence of $`v_0`$, the time and space components of the aether equation (31) read
$`^2v_0_0(^iv_i)`$ $`=`$ $`{\displaystyle \frac{\lambda }{2b_1}}`$ (35)
$`\mathrm{}v_i_i(^jv_j)`$ $`=`$ $`0.`$ (36)
Let us decompose $`v_i`$ into transverse and longitudinal parts, $`v_i=v_i^T+v_i^L`$, where $`^iv_i^T=0`$ and $`v_i^L=_if`$ for some scalar field $`f`$. Then (36) implies that the transverse part satisfies the wave equation, $`\mathrm{}v_i^T=0`$, so the the aether field has two transverse massless modes.
As for the longitudinal part $`v_i^L`$, equation (35) implies
$$_0(^iv_i^L)=\frac{b_1a_1}{2b_1}^2h_{00},$$
(37)
so in particular $`v_i^L`$ has at most linear time dependence. (The same conclusion follows from the divergence of Eq. (36).) Thus $`v_i^L=_ia(x)t+_ib(x)`$, where $`a(x)`$ is determined by $`h_{00}`$ and $`b(x)`$ is arbitrary.
In summary, the perturbation spectrum consists of two massless transverse traceless modes of $`\overline{h}_{ij}`$, and two massless transverse modes of $`v_i`$. In addition there is a mode in which $`h_{00}`$ is an arbitrarily specified time-independent function, which determines $`v_0`$, $`\lambda `$, and the time derivative of the longitudinal part $`v_i^L`$. The time-independent part of $`v_i^L`$ is also arbitrary. This last freedom corresponds to the linearization of the restricted gauge symmetry (9).
The longitudinal mode looks very strange at first sight. In the charged dust interpretation, the dust energy density is adjusted to produce an arbitrary gravitational potential $`h_{00}`$, and the perturbed metric, electromagnetic field, and charge density are all time independent, while the perturbed dust world lines are time dependent. This is a peculiarity of the first order perturbative solution however. No exact solution shares this property, as can be easily seen from the aether field equation (12). If the left hand side is invariant with respect to a timelike Killing field, and if $`\lambda `$ is also invariant, then so must be $`u^b`$. Evidently the higher order terms in the equations of motion induce time dependence into the solution. A similar phenomenon can be seen upon expanding the simpler Einstein-neutral dust system about the flat space solution with constant dust 4-velocity and vanishing density. As in our case, the dust density perturbation can set up any static metric perturbation, and the linearized geodesic equation for the dust yields a time-dependent dust velocity perturbation.
## VI Matter couplings
We have so far considered only the terms in the action involving the metric and the aether field and up to two derivatives. Suppose a matter term $`S_{mat}[g_{ab},u^a,\psi ]`$ is added to the action $`S_{min}[g_{ab},u^a,\lambda ]`$ of the minimal theory (7), where $`\psi `$ stands for a generic matter field. The variation of $`S_{mat}`$ with respect to the metric produces an additional contribution to the stress-energy tensor, and the variation with respect to $`u^a`$ produces an additional term in the current on the right hand side of (12). The resulting field equation takes the form
$$^aF_{ab}=\frac{1}{2b_1}\left(\lambda u_b+\frac{1}{2}\frac{\delta S_{mat}}{\delta u^b}\right).$$
(38)
The identity $`^a^bF_{ab}=0`$ then implies
$$u^a_a\lambda =\lambda _au^a\frac{1}{2}^a\frac{\delta S_{mat}}{\delta u^a},$$
(39)
which shows that now even if $`\lambda `$ is initially zero it need not remain zero. In the presence of such matter couplings the equivalence to charged dust is lost.
We now consider specific types of matter couplings, first of dimension less than or equal to four, and next of dimension greater than four.
### A Couplings of dimension $`4`$
A complete classification of Lorentz violating, gauge-invariant extensions of the SU(3)$`\times `$SU(2)$`\times `$U(1) minimal standard model has been given by Colladay and Kostelecký, restricting attention to operators whose mass dimension is less than or equal to four, so as to preserve power-counting renormalizability. This class of low energy effective actions includes both $`CPT`$-even and $`CPT`$-odd terms, and involves various coupling tensors with “generation” indices allowing for mixing of fermions from different generations. These coupling tensors are supposed to be Lorentz violating vacuum expectation values arising in a theory with a fundamental underlying Lorentz symmetry.
Here we consider the above class of Lorentz violating terms, keeping only those couplings that can be constructed with the aether field $`u^a`$. With this restriction the antisymmetric tensor couplings are excluded, which rules out Lorentz-violating Yukawa couplings and couplings of gauge field strengths to Higgs bilinears, and limits the form of modifications of the gauge field kinetic terms. Invariance under time reversal $`u^au^a`$ would be required if the physical significance of the aether is only to define a preferred frame and not a preferred local time orientation. If we accordingly further assume this symmetry, all the $`CPT`$-odd terms are excluded, which rules out terms with a vector coupled to fermion or Higgs currents, gauge field Chern-Simons currents, and the U(1) potential.
The only possibilities remaining after all these restrictions have been imposed are the modifications of the fermion, gauge field, and Higgs kinetic terms:
$$\frac{1}{2}i(c_L)_{IJ}u^au^b\overline{L}_I\gamma _aD_bL_J+\text{h.c.}+\mathrm{},$$
(40)
$$\frac{1}{4}c_Bu^au^bg^{mn}B_{am}B_{bn}+\mathrm{},$$
(41)
$$\frac{1}{2}c_\mathrm{\Phi }u^au^b(D_a\mathrm{\Phi })^{}D_b\mathrm{\Phi }.$$
(42)
The indices $`I,J`$ in (40) are generational indices, and the coupling constants $`(c_L)_{IJ}`$, $`c_B`$, and $`c_\mathrm{\Phi }`$ are all dimensionless. The ellipses in (40) stand for similar terms for the other fermions, while those in (41) stand for similar terms for the other gauge fields.
Such additional kinetic terms modify the propagation speed of the various fields. For example, the propagation speed for the Higgs field, with respect to the preferred frame, is $`(1+c_\mathrm{\Phi })^{1/2}`$, which is less or more than the speed of light if $`c_\mathrm{\Phi }`$ is positive or negative respectively. The coupling constants must therefore be small numbers for fields whose propagator has been measured accurately. It would be interesting to determine what limits can be placed on these coefficients, particularly for fields such as the Higgs or gluons whose propagators are presumably not yet so well measured.
### B Couplings of dimension $`>4`$
Once the restriction to terms of dimension 4 or less is dropped, the possibilities for Lorentz violating terms—like those for Lorentz invariant ones—are endless. Here we would like to consider just two types, which illustrate different possibilities that arise in the presence of Lorentz symmetry breaking.
#### Modified kinetic terms
If the coupling coefficient for a Lorentz violating kinetic term like (40-42) is field dependent and polynomial, rather than a constant, then the term is a dimension $`>4`$ operator. In this case it is possible that the coefficient was larger in the early universe than it is today, due to the cosmological evolution of the field(s) on which the coupling function depends. This provides an alternate approach to constructing generally covariant, variable speed of light cosmologies. Approaches using a vector or a scalar to define the preferred frame have been the subject of some recent papers.
#### Modified dispersion
Next we consider a deviation from Lorentz invariance that becomes strong only at high wavevectors. In the early universe, when the fields were highly excited at large wavevectors, the gravitational effects of such a deviation could have been of paramount importance. The study of a model incorporating such effects is left to another paper. Here we indicate only an example of a term in the Lagrangian that produces high frequency dispersion in the propagation of a matter field, and we display the the form of the resulting contribution to the energy-momentum tensor.
Consequences of non-Lorentz invariant high frequency dispersion for the Hawking effect have previously been studied using 1+1 dimensional model field theories in which higher spatial derivative terms are added to the action (for a review see ), and recently such models have been generalized to field theory in the background of a 3+1 dimensional Robertson-Walker spacetime in order to study the consequences for the spectrum of primordial density fluctuations in inflationary cosmology. These models can be extended to an arbitrary 3+1 dimensional setting, preserving general covariance as well as spatial rotation symmetry in the local preferred frame. As an example consider the Lagrangian
$$_\phi =\frac{1}{2}\left(^a\phi _a\phi +k_0^2(D^2\phi )^2\right).$$
(43)
Here $`k_0`$ is a constant with the dimensions of inverse length which sets the scale for deviations from Lorentz invariance, and $`D^2`$ is the covariant spatial Laplacian, i.e.,
$$D^2\phi =D^aD_a\phi =q^{ac}_a(q_c{}_{}{}^{b}_{b}^{}\phi ),$$
(44)
where $`D_a`$ is the spatial covariant derivative operator and $`q_{ab}`$ is the spatial metric (10).
The $`u^a`$-dependence of the Lagrangian (43) produces a “matter” term in the aether field equation (38). The energy-momentum tensor for this Lagrangian is
$`T_{ab}=_a\phi _b\phi _\phi g_{ab}`$ (45)
$`k_0^2[2D^2\phi u^mu_{(a}_{|m|}D_{b)}\phi +2_m(D^2\phi q_{(a}{}_{}{}^{m})_{b)}\phi `$ (46)
$`^m(q_{ab}D^2\phi D_m\phi )].`$ (47)
In Ref. we evaluate the expectation value of this energy-momentum tensor in a thermal state in flat spacetime, which allows us to determine the modification of the equation of state produced by the fourth derivative term. This equation of state is then be used to study how the cosmological evolution is affected by the high frequency dispersion.
## VII Discussion
We have made an initial attempt to study the possible consequences of incorporating a preferred frame—the aether—into a generally covariant theory. With the action adopted in this paper the aether vector generically develops gradient singularities even when the metric is perfectly regular. We take this as a sign that the theory is unphysical as an effective theory (although if the aether sector is ignored the theory can be made to agree with observations with an appropriate choice of initial conditions, i.e. by setting $`F_{ab}`$ to zero.) The primary open questions are ($`i`$) what determines the initial values of the aether field and the Lagrange multiplier field, and ($`ii`$) are the gradient singularities, which appear to be generic in the evolution of the aether, eliminated by including a symmetrized derivative term $`(_{(a}u_{b)})(^{(a}u^{b)})`$ in the action along with the antisymmetrized derivative term used in this paper? It is plausible that adding the symmetrized derivative term will have a significant effect, since with it the action is sensitive to the existence of any large gradients.
## Acknowledgements
We are grateful to V.A. Kostelecký, M. Luty, and R.P. Woodard for helpful discussions. This work was supported in part by the National Science Foundation under grant No. PHY98-00967.
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# Second order families of special Lagrangian 3-folds
## 1. Introduction
The study of special Lagrangian submanifolds was introduced by Harvey and Lawson in §III of their fundamental paper on calibrated geometries. They analyzed the local and global geometry of these submanifolds in flat complex $`m`$-space, constructing many interesting examples and proving local existence theorems.
Several important classes of examples of these submanifolds have been constructed since then, mainly with an eye to applications in the theory of calibrations or minimizing submanifolds. Particularly important was the construction by Lawlor of a special Lagrangian manifold asymptotic to a pair of planes that violate the angle criterion, thus proving that such a pair of planes is not area-minimizing. See also Harvey for a thorough account of this example and its applications.
The deformation theory of compact examples in Calabi-Yau manifolds was studied in the 1990 thesis of R. McLean , who showed that the moduli space of compact special Lagrangian submanifolds of a given Calabi-Yau manifold is always a disjoint union of smooth manifolds.
Special Lagrangian geometry received renewed attention in 1996 when its role in mirror symmetry was discovered by Strominger, Yau, and Zaslow . Since then, interest in special Lagrangian geometry has grown quite rapidly. The reader might consult (for what is currently known about the moduli of compact special Lagrangian submanifolds), (for some examples arising from algebraic geometry), (for further information about mirror symmetry), (for examples with large symmetry groups), (for embedding a given real-analytic Riemannian $`3`$-manifold as a special Lagrangian submanifold of a Calabi-Yau $`3`$-fold), (for some interesting speculations about how one might count the isolated special Lagrangian submanifolds), and (for information about special Lagrangian cones in $`^{\mathrm{\hspace{0.17em}3}}`$).
Still, the systematic exploration of special Lagrangian geometry seems to have hardly begun. The known explicit examples have largely been found by the well-known Ansatz of symmetry reduction or other special tricks.
The research that lead to this article was an attempt to classify families of special Lagrangian submanifolds that are characterized by invariant, differential geometric conditions, in particular, conditions on the second fundamental form of the special Lagrangian submanifold.
At least when the ambient space is flat, the lowest order invariant of a special Lagrangian submanifold is its second fundamental form. Now, for a Lagrangian submanifold of a linear symplectic vector space, the second fundamental form, usually defined as a quadratic form with values in the normal bundle, has a natural interpretation as a symmetric cubic form $`C`$ on the submanifold, called the *fundamental cubic*. When the submanifold is special Lagrangian, it turns out that the trace of this cubic form with respect to the first fundamental form vanishes, but there are no further pointwise conditions on this cubic that are satisfied for all special Lagrangian submanifolds.
It is natural to ask whether one can obtain nontrivial families of special Lagrangian submanifolds by imposing pointwise conditions on the fundamental cubic. In the language of overdetermined systems of PDE, one would like to be able to say whether there are any second order systems of PDE that are ‘compatible’ with the (first order) system that represents the special Lagrangian condition.<sup>1</sup><sup>1</sup>1 Here, ‘compatibility’ is not strictly defined, but, roughly speaking, means that there exist at least as many (local) solutions to the overdetermined system as one would expect from a naïve ‘equation counting’ argument. A more precise description would involve concepts from exterior differential systems, such as involutivity, that will not be needed in this article.
The first task is to understand the space of pointwise invariants of a traceless cubic form under the special orthogonal group. For example, in dimension $`3`$ (which is the case this article mainly considers), the space of traceless cubics is an irreducible $`\mathrm{SO}(3)`$-module of dimension $`7`$, so one would expect there to be four independent polynomial invariants.
However, the relations that one gets by imposing conditions on these invariants are generally singular at the cubics that have a nontrivial stabilizer under the action of $`\mathrm{SO}(3)`$. For comparison, consider the classical case of hypersurfaces in Euclidean space. The fundamental invariants are the principal curvatures, i.e., the eigenvalues of the second fundamental form with respect to the first fundamental form. These are smooth away from the (generalized) umbilic locus, i.e., the places where two or more of the principal curvatures come together. It is exactly at these places that the stabilizer of the second fundamental form in the orthogonal group is larger than the minimum possible stabilizer. Of course, the umbilic locus is also the place where moving frame adaptations generally run into trouble, unless one assumes that the multiplicities of the principal curvatures are constant.
There is a similar phenomenon in special Lagrangian geometry. In place of the umbilic locus, one looks that the places where the fundamental cubic has a nontrivial stabilizer,<sup>2</sup><sup>2</sup>2 In contrast to the familiar case of hypersurfaces in Euclidean space, where the stabilizer, though generically finite, is always nontrivial, it turns out that the stabilizer at a generic point of the fundamental cubic of a ‘generic’ special Lagrangian $`3`$-fold is trivial. and at the special Lagrangian submanifolds where the stabilizer of the cubic is nontrivial at the generic point. These are the special special Lagrangian submanifolds.
In this article, after making some general remarks to introduce the structure equations of special Lagrangian geometry, I classify the possible nontrivial $`\mathrm{SO}(3)`$-stabilizers of traceless cubics in three variables. It turns out that the $`\mathrm{SO}(3)`$-stabilizer of a nontrivial traceless cubic is isomorphic to either a copy of $`\mathrm{SO}(2)`$, the group $`A_4`$ of order $`12`$, the group $`S_3`$ of order $`6`$, the group $`_3`$, the group $`_2`$, or is trivial. I then consider, for each of the nontrivial subgroups $`G`$ on this list, the problem of classifying the special Lagrangian $`3`$-folds whose cubic form at each point has its stabilizer contain a copy of $`G`$.
For example, it turns out that the only special Lagrangian $`3`$-folds in $`^{\mathrm{\hspace{0.17em}3}}`$ whose cubic form has a continuous stabilizer at each point are the $`3`$-planes and the $`\mathrm{SO}(3)`$-invariant examples discovered by Harvey and Lawson.
There are no special Lagrangian $`3`$-folds whose cubic stabilizer at at generic point is of type $`A_4`$, but the ones whose stabilizer at a generic point is of type $`S_3`$ turn out to be the austere special Lagrangian $`3`$-folds and these are known to be the orthogonal products, the special Lagrangian cones, and the ‘twisted’ special Lagrangian cones.
The special Lagrangian $`3`$-folds with cubic stabilizer at a generic point isomorphic to $`_2`$ turn out to be the examples discovered by Lawlor, extended by the work of Harvey and then Joyce.<sup>3</sup><sup>3</sup>3 I am indebted to Joyce for suggesting (by private communication) that the family that I had shown to exist in this case might be the Lawlor-Harvey-Joyce family. He was correct, and this saved me quite a bit of work in integrating the corresponding structure equations.
The special Lagrangian $`3`$-folds whose cubic stabilizer at a generic point is isomorphic to $`_3`$ turn out to be asymptotically conical and, indeed, turn out to be deformations, in a certain sense, of the special Lagrangian cones, as explained in the thesis of Haskins .
The above results are explained more fully in §3.
At the conclusion of §3, I consider a different type of invariant condition on the fundamental cubic, namely that, at every point $`xL`$, the degree $`3`$ curve in $`(T_xL)`$ defined by the fundamental cubic have a real singular point. This is one semi-algebraic condition on the fundamental cubic. I show that the special Lagrangian $`3`$-folds with this property are exactly the ruled special Lagrangian $`3`$-folds.
Moreover, while the most general known family of ruled special Lagrangian $`3`$-folds up until now was one discovered by Borisenko and that depends on four functions of one variable (in the sense of exterior differential systems), I show that the full family depends on six functions of one variable. Moreover, I show that, when one interprets the ruled special Lagrangian $`3`$-folds as surfaces in the space of lines in $`^{\mathrm{\hspace{0.17em}3}}`$, the surfaces that one obtains are simply the ones that are holomorphic with respect to a canonical Levi-flat, almost CR-structure on the space of lines. This interpretation has several implications for the structure of ruled special Lagrangian $`3`$-folds, among them being that any real-analytic ruled surface in $`^{\mathrm{\hspace{0.17em}3}}`$ on which the Kähler form vanishes lies in a (essentially unique) ruled special Lagrangian $`3`$-fold. Moreover, a special Lagrangian $`3`$-fold is ruled if and only if it contains a ruled surface.
It has to be said that the results of this article are only the first step in understanding the compatibility of the special Lagrangian condition with higher order conditions. Now that the ‘umbilic’ cases are understood, the serious work on the ‘generic’ case can be undertaken. This will be reported on in a subsequent work.
Also, while, for the sake of brevity, this work has concerned itself (essentially exclusively) with the $`3`$-dimensional case, there are obvious higher dimensional generalizations that need to be investigated and that should yield to the same or similar techniques.<sup>4</sup><sup>4</sup>4 My student, Marianty Ionel, has recently completed a study of the special Lagrangian $`4`$-folds in $`^{\mathrm{\hspace{0.17em}4}}`$ whose fundamental cubic has nontrivial symmetries.
### 1.1. Special Lagrangian geometry
In this article, a slightly more general notion of special Lagrangian geometry is adopted than is customary. The reader might compare this discussion with Harvey and Lawson’s original article or Harvey’s more recent book .
#### 1.1.1. Special Kähler structures
Let $`M`$ be a complex $`m`$-manifold endowed with a Kähler form $`\omega `$ and a holomorphic volume form $`\mathrm{{\rm Y}}`$. It is *not* assumed that $`\mathrm{{\rm Y}}`$ be parallel, or even of constant norm, with respect to the Levi-Civita connection associated to $`\omega `$. The pair $`(\omega ,\mathrm{{\rm Y}})`$ is said to define a *special Kähler structure* on $`M`$.
#### 1.1.2. Special Lagrangian submanifolds
A submanifold $`LM`$ of real dimension $`m`$ is said to be *Lagrangian*<sup>5</sup><sup>5</sup>5 or *$`\omega `$-Lagrangian* if there is any danger of confusion if the pullback of $`\omega `$ to $`L`$ vanishes.
Harvey and Lawson show \[12, §III, Theorem 1.7\] that for any Lagrangian submanifold $`LM`$, the pullback of $`\mathrm{{\rm Y}}`$ to $`L`$ can never vanish. A Lagrangian submanifold $`L`$ is said to be *special Lagrangian* if the pullback of $`\mathrm{Im}(\mathrm{{\rm Y}})`$ to $`L`$ vanishes. When $`L`$ is special Lagrangian, it has a canonical orientation for which $`\mathrm{Re}(\mathrm{{\rm Y}})`$ pulls back to $`L`$ to be a positive volume form, and this is the orientation that will be assumed throughout this article.
More generally, if $`\lambda `$ is a complex number of unit modulus, one says that an oriented Lagrangian submanifold $`LM`$ has *constant phase* $`\lambda `$ if $`\overline{\lambda }\mathrm{{\rm Y}}`$ pulls back to $`L`$ to be a (real-valued) positive volume form. Obviously, for any fixed $`\lambda `$, this notion is not significantly more general than the notion of special Lagrangian, so I will usually consider only special Lagrangian submanifolds in this article.
#### 1.1.3. The Calabi-Yau case
When $`\mathrm{{\rm Y}}`$ is parallel with respect to the Levi-Civita connection associated to $`\omega `$, the Kähler metric has vanishing Ricci tensor and the pair $`(\omega ,\mathrm{{\rm Y}})`$ is said to define a *Calabi-Yau* structure on $`M`$. In this case, Harvey and Lawson show that any special Lagrangian submanifold $`LM`$ is minimal. Moreover, if $`L`$ is compact, it is absolutely minimizing in its homology class since it is then calibrated by $`\mathrm{Re}(\mathrm{{\rm Y}})`$.
#### 1.1.4. Local existence
Assume that $`\omega `$ is real-analytic with respect to the standard real-analytic structure on $`M`$ that underlies its complex analytic structure. Harvey and Lawson show\[12, §III, Theorem 5.5\] that any real-analytic submanifold $`NM`$ of dimension $`m1`$ on which $`\omega `$ pulls back to be zero lies in a unique special Lagrangian submanifold $`LM`$. (Although their result is stated only for the case of the standard flat special Kähler structure on $`^m`$, their proof is valid in the general case, provided one makes the necessary trivial notational changes.)
Thus, there are many special Lagrangian submanifolds locally, at least in the real-analytic category. By adapting arguments from \[12, §III.2\], one can also prove local existence of special Lagrangian submanifolds even without the assumption of real-analyticity. Instead, one uses local existence for an elliptic second order scalar equation.
#### 1.1.5. Deformations
R. McLean proved that, in the Calabi-Yau case, a compact special Lagrangian submanifold $`LM`$ is a point in a smooth, finite dimensional moduli space $``$ consisting of the special Lagrangian deformations of $`L`$ and that the tangent space to $``$ at $`L`$ is isomorphic to the space of harmonic 1-forms on $`L`$.
McLean’s argument makes no essential use of the assumption that $`\mathrm{{\rm Y}}`$ be $`\omega `$-parallel. Instead, it is sufficient for the conclusion of McLean’s theorem that $`\mathrm{{\rm Y}}`$ be closed (in fact, one only really needs that the imaginary part of $`\mathrm{{\rm Y}}`$ be closed.) For a related result, see .
### 1.2. Special Kähler reduction
One reason for considering the slightly wider notion of special Lagrangian geometry adopted here is that it is stable under the process of *reduction*, as explained in , , and .
Let $`(\omega ,\mathrm{{\rm Y}})`$ be a special Kähler structure on $`M`$. A vector field $`X`$ on $`M`$ will be said to be an *infinitesimal symmetry* of the structure if the (locally defined) flow of $`X`$ preserves both $`\omega `$ and $`\mathrm{{\rm Y}}`$.
Suppose that $`X`$ is an infinitesimal symmetry of $`(\omega ,\mathrm{{\rm Y}})`$ and that $`X`$ is, moreover, $`\omega `$-Hamiltonian, i.e., that there exists a function $`H`$ on $`M`$ satisfying $`X\text{ }\text{ }\omega =\mathrm{d}H`$. The flow lines of $`X`$ are tangent to the level sets of $`H`$.
Say that a value $`h`$ is a *good* value for $`H`$ if it is a regular value of $`H`$ and if the flow of $`X`$ on the level set $`H^1(h)M`$ is simple, i.e., there is a smooth manifold structure on the set $`M_h`$ of flow lines of $`X`$ in the level set $`H^1(h)`$ so that the natural projection $`\pi _h:H^1(h)M_h`$ is a smooth submersion. The (real) dimension of $`M_h`$ is necessarily $`2m2`$.
When $`h`$ is good, there exists a unique $`2`$-form $`\omega _h`$ on $`M_h`$ for which $`\pi _h^{}(\omega _h)`$ is the pullback of $`\omega `$ to $`H^1(h)`$ and there exists a unique complex-valued $`(m1)`$ form $`\mathrm{{\rm Y}}_h`$ on $`M_h`$ for which $`\pi _h^{}(\mathrm{{\rm Y}}_h)`$ is the pullback to $`H^1(h)`$ of $`X\text{ }\text{ }\mathrm{{\rm Y}}`$.
It is trivial to verify that $`(\omega _h,\mathrm{{\rm Y}}_h)`$ defines a special Kähler structure on $`M_h`$. Note, however, that, even if $`(\omega ,\mathrm{{\rm Y}})`$ is Calabi-Yau, its reductions will generally *not* be Calabi-Yau. In fact, this happens only when the length of $`X`$ is constant along the level set $`H^1(h)`$.
If $`LH^1(h)`$ is a special Lagrangian submanifold that is tangent to the flow of $`X`$, then $`L=\pi _h^1(L_h)`$ where $`L_hM_h`$ is also special Lagrangian. Conversely, if $`L_hM_h`$ is special Lagrangian, then $`L=\pi _h^1(L_h)`$ is special Lagrangian in $`M`$.
This method of special Kähler reduction allows one to construct many examples of special Lagrangian submanifolds by starting with a Hamiltonian $`(m1)`$-torus action and doing a series of reductions, leading to a 1-dimensional special Kähler manifold, where the integration problem is reduced to integrating a holomorphic 1-form on a Riemann surface.
## 2. The Structure Equations
The structure equations of a Kähler manifold adapted for special Lagrangian geometry can be found in , and will only be reviewed briefly here.
### 2.1. The special coframe bundle
The *standard special Kähler structure* on $`^m`$ is the one defined by
(2.1)
$$\omega _0=\frac{\mathrm{i}}{2}\left(\mathrm{d}z_1\mathrm{d}\overline{z_1}+\mathrm{}+\mathrm{d}z_m\mathrm{d}\overline{z_m}\right)\text{and}\mathrm{{\rm Y}}_0=\mathrm{d}z_1\mathrm{}\mathrm{d}z_m,$$
where $`z_1,\mathrm{},z_m`$ are the usual complex linear coordinates on $`^m`$. The corresponding Kähler metric is, of course
(2.2)
$$g_0=\mathrm{d}z_1\mathrm{d}\overline{z_1}+\mathrm{}+\mathrm{d}z_m\mathrm{d}\overline{z_m}.$$
Note that $`^m^m`$ is a special Lagrangian subspace.
Let $`(\omega ,\mathrm{{\rm Y}})`$ be a special Kähler structure on the complex $`m`$-manifold $`M`$. There is a unique positive function $`B`$ on $`M`$ that satisfies
(2.3)
$$\mathrm{{\rm Y}}\overline{\mathrm{{\rm Y}}}=\frac{2^m(\mathrm{i})^{m^2}}{m!}B^2\omega ^m.$$
A linear isomorphism $`u:T_xM^m`$ will be said to be a *special Kähler coframe* at $`x`$ if it satisfies $`\omega _x=u^{}(\omega _0)`$ and $`\mathrm{{\rm Y}}_x=B(x)u^{}(\mathrm{{\rm Y}}_0)`$. Such a coframe is necessarily complex linear, i.e., satisfies $`u(J_xv)=\mathrm{i}u(v)`$ for all $`vT_xM`$, where $`J_x:T_xMT_xM`$ is the complex structure map.
The set of special Kähler coframes at $`x`$ will be denoted $`P_x`$ and is the fiber of a principal right $`\mathrm{SU}(m)`$-bundle $`\pi :PM`$, with right action given by $`R_a(u)=a^1u`$ for $`a\mathrm{SU}(m)`$.
As usual, the $`^m`$-valued, tautological 1-form $`\zeta `$ on $`P`$ is defined by requiring that $`\zeta _u=u\pi ^{}(u):T_uP^m`$ for $`uP`$. It satisfies $`R_a^{}(\zeta )=a^1\zeta `$ for $`a\mathrm{SU}(m)`$.
The components of $`\zeta `$ will be written as $`\zeta _i`$ for $`1im`$. The equations
(2.4)
$$\omega =\frac{\mathrm{i}}{2}\left(\zeta _1\overline{\zeta _1}+\mathrm{}+\zeta _m\overline{\zeta _m}\right)\text{and}\mathrm{{\rm Y}}=B\zeta _1\mathrm{}\zeta _m$$
hold on $`P`$, where, as is customary, I have omitted the $`\pi ^{}`$, thus implicitly embedding the differential forms on $`M`$ into the differential forms on $`P`$ via pullback.
Finally, there are functions $`𝐞=(𝐞_i)`$, where $`𝐞_i:PTM`$ is a bundle mapping satisfying $`\zeta _i(𝐞_j)=\delta _{ij}`$. In other words $`\pi ^{}(v)=\zeta _i(v)𝐞_i(u)`$ for all $`vT_uP`$.
###### Remark 1 (The flat case).
When $`M=^m`$ and $`(\omega ,\mathrm{{\rm Y}})=(\omega _0,\mathrm{{\rm Y}}_0)`$, it is customary to use the vector space (parallel) trivialization of the tangent bundle of $`^m`$ to identify all of the tangent spaces to the vector space $`^m`$ itself. In this case, the functions $`𝐞_i`$ will be regarded as vector-valued functions on $`P^m\times \mathrm{SU}(m)`$ and the basepoint projection will be denoted as $`𝐱:P^m`$. Then the above relations take on the more familiar ‘moving frame’ form
$$\mathrm{d}𝐱=𝐞_i\zeta _i,$$
and so on. The reader should have no trouble figuring out what is meant in context.
### 2.2. The structure equations
The Levi-Civita connection associated to the underlying Kähler structure on $`M`$ is represented on $`P`$ by a $`𝔲(m)`$-valued 1-form $`\psi =\psi ^{}=(\psi _{i\overline{ȷ}})`$ that satisfies the *first structure equation*
(2.5)
$$\mathrm{d}\zeta _i=\psi _{i\overline{ȷ}}\zeta _j.$$
The equation $`\mathrm{d}\mathrm{{\rm Y}}=0`$ implies
(2.6)
$$(\overline{})\left(\mathrm{log}B\right)=\mathrm{i}\mathrm{d}^c(\mathrm{log}B)=\mathrm{tr}(\psi )=\psi _{iı}.$$
Note that $`(\omega ,\mathrm{{\rm Y}})`$ is Calabi-Yau if and only if $`B`$ is constant, i.e., if and only if $`\psi `$ takes values in $`𝔰𝔲(m)`$.
In the Calabi-Yau case, where $`\mathrm{{\rm Y}}`$ is parallel with respect to the Levi-Civita connection of $`\omega `$, the relation $`\psi _{ii}=0`$ holds. Moreover, the Calabi-Yau structure is locally equivalent to the standard structure if and only if the Levi-Civita connection of $`\omega `$ vanishes, i.e., if and only if $`\mathrm{d}\psi =\psi \psi `$, which is known as the *second structure equation* of a flat Calabi-Yau space.
### 2.3. Special Lagrangian submanifolds
For the study of special Lagrangian submanifolds, it is convenient to separate the structure equations into real and imaginary parts. Thus, set $`\zeta _i=\omega _i+\mathrm{i}\eta _i`$ and $`\psi _{i\overline{ȷ}}=\alpha _{ij}+\mathrm{i}\beta _{ij}`$. The first structure equations can then be written in the form
(2.7)
$$\begin{array}{cc}\hfill \mathrm{d}\omega _i& =\alpha _{ij}\omega _j+\beta _{ij}\eta _j,\hfill \\ \hfill \mathrm{d}\eta _i& =\beta _{ij}\omega _j\alpha _{ij}\eta _j.\hfill \end{array}$$
where $`\alpha _{ij}=\alpha _{ji}`$ and $`\beta _{ij}=\beta _{ji}`$. Note that (2.6) becomes $`\beta _{ii}=\mathrm{d}^c\left(\mathrm{log}B\right)`$.
Let $`LM`$ be a special Lagrangian submanifold. For $`xL`$, a special Kähler coframe $`u:T_xM^m`$ is said to be *$`L`$-adapted* if $`u(T_xL)=^m^m`$ and $`u:T_xL^m`$ is orientation preserving. The space of $`L`$-adapted coframes forms a principal right $`\mathrm{SO}(m)`$-subbundle $`P_L\pi ^1(L)P`$ over $`L`$.
The equations $`\eta _i=0`$ hold on $`P_L`$. Thus, by the structure equations (2.7), the relations $`\beta _{ij}\omega _j=0`$ hold on $`P_L`$ while $`\omega _1\mathrm{}\omega _m`$ is nowhere vanishing. It follows from Cartan’s Lemma that there are functions $`h_{ijk}=h_{jik}=h_{ikj}`$ on $`P_L`$ so that
(2.8)
$$\beta _{ij}=h_{ijk}\omega _k.$$
The second fundamental form of $`L`$ can then be written as
(2.9)
$$\mathrm{I}\mathrm{I}=h_{ijk}J𝐞_i\omega _j\omega _k=J𝐞_iQ_i$$
where $`Q_i=h_{ijk}\omega _j\omega _k`$. The information in the second fundamental form is thus contained in the symmetric cubic form
(2.10)
$$C=h_{ijk}\omega _i\omega _j\omega _k=\omega _iQ_i,$$
which is well-defined on $`L`$. This symmetric cubic form will be referred to as the *fundamental cubic* of the special Lagrangian submanifold $`L`$.
Note that the trace of $`C`$ with respect to the induced metric $`g=\omega _{1}^{}{}_{}{}^{2}+\mathrm{}+\omega _{n}^{}{}_{}{}^{2}`$ on $`L`$ satisfies
(2.11)
$$\mathrm{tr}_gC=h_{iik}\omega _k=\beta _{ii}=\mathrm{d}^c(\mathrm{log}B)\text{ }{}_{L}{}^{},$$
which is the restriction to $`L`$ of an ambient $`1`$-form. In the Calabi-Yau case, $`0=\psi _{ii}=\mathrm{i}\beta _{ii}`$, so the fundamental cubic $`C`$ is traceless.
Finally, in the flat case, the curvature vanishing condition $`\mathrm{d}\psi =\psi \psi `$ can be separated into real and imaginary parts. The result will be referred to as the *second structure equations*:
(2.12a) $`\mathrm{d}\alpha _{ij}`$ $`=\alpha _{ik}\alpha _{kj}+\beta _{ik}\beta _{kj},`$
(2.12b) $`\mathrm{d}\beta _{ij}`$ $`=\beta _{ik}\alpha _{kj}\alpha _{ik}\beta _{kj}.`$
#### 2.3.1. A Bonnet-type result
Given an $`m`$-manifold $`L`$ endowed with a Riemannian metric $`g`$ and a symmetric cubic form $`C`$ that is traceless with respect to $`g`$, one can ask whether there is an isometric imbedding of $`(L,g)`$ into $`^m`$ as a special Lagrangian submanifold that induces $`C`$ as the fundamental cubic.
It is easy to see that there is a Bonnet-style theorem derivable from the above structure equations. Namely, there are analogs of the Gauss and Codazzi equations that give necessary and sufficient conditions for the solution of this problem.
To see this, first choose a $`g`$-orthonormal coframing $`\omega =(\omega _i)`$ on an open subset $`UL`$. Then define $`\eta _i=0`$ and let $`\alpha _{ij}=\alpha _{ji}`$ be the unique $`1`$-forms on $`U`$ that satisfy the equations $`\mathrm{d}\omega _i=\alpha _{ij}\omega _j`$. (The existence and uniqueness of such $`\alpha `$ is just the Fundamental Lemma of Riemannian geometry.) Write $`C=h_{ijk}\omega _i\omega _j\omega _k`$ and set $`\beta _{ij}=h_{ijk}\omega _k`$.
The equation (2.12b) then expresses the fact that $`C`$ must satisfy a Codazzi-type equation with respect to $`g`$, namely, that the covariant derivative of $`C`$ with respect to the Levi-Civita connection of $`g`$ is fully symmetric. The equation (2.12a) then expresses the fact that that $`C`$ must satisfy a Gauss-type equation with respect to $`g`$, namely, that $`C`$ satisfies an algebraic equation of the form $`Q_g(C)=\text{Riem}(g)`$, where $`Q_g`$ is a certain quadratic mapping (depending on $`g`$) from symmetric cubic forms into tensors of the same algebraic type as the Riemann curvature.
Thus, when $`L`$ is simply connected, these Codazzi and Gauss equations are the necessary and sufficient conditions on $`g`$ and $`C`$ for there to be an isometric immersion of $`(L,g)`$ into $`^m`$ as a special Lagrangian submanifold inducing $`C`$ as its fundamental cubic. Moreover, such an isometric immersion will be unique up to rigid motion.
## 3. Second Order Families
### 3.1. The second fundamental form as a cubic
It was already explained in §2 how the second fundamental form of a special Lagrangian submanifold $`L^m`$ can be regarded as a symmetric cubic form that is traceless with respect to the first fundamental form.
Thus, the polynomial second order invariants of such a submanifold correspond to the $`\mathrm{SO}(m)`$-invariant polynomials on the space $`_3(^m)`$ of harmonic polynomials on $`^m`$ that are homogeneous of degree $`3`$. Moreover, the $`\mathrm{SO}(m)`$-stabilizer of a given cubic in this space corresponds to the ambiguity in the choice of an adapted coframe for the corresponding special Lagrangian submanifold. In particular, points on special Lagrangian submanifolds at which the $`\mathrm{SO}(m)`$-stabilizer of the second fundamental form is nontrivial can be regarded as analogs of umbilic points in the classical theory of surfaces in 3-space.
The space $`_3(^m)`$ is an irreducible $`\mathrm{SO}(m)`$-module when $`\mathrm{SO}(m)`$ acts in the natural way by pullback. Thus, there are no invariant linear functions and, up to multiples, exactly one invariant quadratic polynomial, which is, essentially, the squared norm of the polynomial. It is not difficult to show that there are no invariant cubic polynomials on this space, that the space of invariant quartics is of dimension $`2`$ (one of which is the square of the invariant quadratic), and so on. The exact structure of the ring of invariants for general $`m`$ is complicated, however, and I will not discuss this further.
Since I will only be using the results of the case $`m=3`$, I am going to be assuming this from now on. The space $`_3(^{\mathrm{\hspace{0.17em}3}})`$ has dimension $`7`$, and one would expect that the ‘generic’ $`\mathrm{SO}(3)`$-orbit in this vector space to have dimension $`3`$. In fact, I am now going to determine the orbits that have non-trivial stabilizers.
#### 3.1.1. Special orbits
The main goal of this section is to prove the following result, which is undoubtedly classical even though I have been unable to locate a proof in the literature.
###### Proposition 1.
The $`\mathrm{SO}(3)`$-stabilizer of $`h_3(^{\mathrm{\hspace{0.17em}3}})`$ is nontrivial if and only if $`h`$ lies on the $`\mathrm{SO}(3)`$-orbit of exactly one of the following polynomials
1. $`0_3(^{\mathrm{\hspace{0.17em}3}})`$, whose stabilizer is $`\mathrm{SO}(3)`$.
2. $`r(2z^33zx^23zy^2)`$ for some $`r>0`$, whose stabilizer is $`\mathrm{SO}(2)`$.
3. $`6sxyz`$ for some $`s>0`$, whose stabilizer is the subgroup $`A_4\mathrm{SO}(3)`$ of order $`12`$ generated by the rotations by an angle of $`\pi `$ about the $`x`$-, $`y`$-, and $`z`$-axes and by rotation by an angle of $`\frac{2}{3}\pi `$ about the line $`x=y=z`$.
4. $`s(x^33xy^2)`$ for some $`s>0`$, whose stabilizer is the subgroup $`S_3\mathrm{SO}(3)`$ of order $`6`$ generated by the rotation by an angle of $`\pi `$ about the $`x`$-axis and the rotation by an angle of $`\frac{2}{3}\pi `$ about the $`z`$-axis.
5. $`r(2z^33zx^23zy^2)+6sxyz`$ for some $`r,s>0`$ satisfying $`sr`$, whose stabilizer is the $`_2`$-subgroup of $`\mathrm{SO}(3)`$ generated by rotation by an angle of $`\pi `$ about the $`z`$-axis.
6. $`r(2z^33zx^23zy^2)+s(x^33xy^2)`$ for some $`r,s>0`$ satisfying $`sr\sqrt{2}`$, whose stabilizer is the $`_3`$-subgroup of $`\mathrm{SO}(3)`$ generated by rotation by an angle of $`\frac{2}{3}\pi `$ about the $`z`$-axis.
###### Remark 2 (Special Values).
The reader may wonder about the conditions $`sr`$ and $`sr\sqrt{2}`$ in the last two cases. It is not difficult to verify that the polynomial $`(2z^33zx^23zy^2)+6xyz`$ lies on the $`\mathrm{SO}(3)`$-orbit of $`2(x^33xy^2)`$ and that the polynomial $`(2z^33zx^23zy^2)+\sqrt{2}(x^33xy^2)`$ lies on the $`\mathrm{SO}(3)`$-orbit of $`6\sqrt{3}xyz`$.
###### Proof.
Suppose that $`h_3(^{\mathrm{\hspace{0.17em}3}})`$ has a nontrivial stabilizer $`G\mathrm{SO}(3)`$. Obviously $`G=\mathrm{SO}(3)`$ if and only if $`h=0`$, so suppose that $`h0`$ from now on. Since $`G`$ is closed in $`\mathrm{SO}(3)`$, it is compact and has a finite number of components.
Suppose first that $`G`$ is not discrete. Then the identity component of $`G`$ must be a closed $`1`$-dimensional subgroup and hence conjugate to the subgroup $`\mathrm{SO}(2)\mathrm{SO}(3)`$ consisting of the rotations about the $`z`$-axis. Thus, $`h`$ lies on the orbit of a cubic polynomial that is invariant under this rotation group. By replacing $`h`$ by such an element, it can be supposed that the identity component of $`G`$ is $`\mathrm{SO}(2)`$.
Consider the following four subspaces of $`_3(^{\mathrm{\hspace{0.17em}3}})`$: Let $`V_0`$ be the $`1`$-dimensional space spanned by $`z(2z^23x^23y^2)`$; let $`V_1`$ be the $`2`$-dimensional space spanned by $`x(4z^2x^2y^2)`$ and $`y(4z^2x^2y^2)`$; let $`V_2`$ be the $`2`$-dimensional space spanned by $`(x^2y^2)z`$ and $`xyz`$; and let $`V_3`$ be the $`2`$-dimensional space spanned by $`(x^33xy^2)`$ and $`(3x^2yy^3)`$. Each of these subspaces is preserved by the elements of $`\mathrm{SO}(2)`$. Moreover, $`\mathrm{SO}(2)`$ acts trivially on $`V_0`$, while the element $`R_\alpha \mathrm{SO}(2)`$ that represents rotation by an angle $`\alpha `$ about the $`z`$-axis, acts as rotation by the angle $`k\alpha `$ on the $`2`$-dimensional space $`V_k`$ for $`k=1,2,3`$.
Obviously, the only nonzero elements of $`_3(^{\mathrm{\hspace{0.17em}3}})`$ that are fixed by $`\mathrm{SO}(2)`$ are those of the form $`rz(2z^23x^23y^2)`$ for some nonzero $`r`$. Moreover, since $`z`$ is the unique linear factor of this polynomial, it follows that $`G`$, the stabilizer of this polynomial, must preserve the $`z`$-axis. Also, since this polynomial is positive on exactly one of the two rays in the $`z`$-axis emanating from the origin, it follows that $`G`$ must also fix the orientation of the $`z`$-axis. Thus, $`G=\mathrm{SO}(2)`$. Moreover, note that by a rotation that reverses the $`z`$-axis, the element $`rz(2z^23x^23y^2)`$ is carried into the element $`rz(2z^23x^23y^2)`$. Thus, one can assume that $`r>0`$.
Now suppose that $`G`$ is discrete (and hence finite). Let $`AG`$ be an element of finite order $`p>1`$. Then $`A`$ is rotation about a line by an angle of the form $`(2q/p)\pi `$ for some integer $`q`$ relatively prime to $`p`$ and satisfying $`0<q<p`$. Replacing $`h`$ by an element in its $`\mathrm{SO}(3)`$-orbit, I can assume that the fixed line of $`A`$ is the $`z`$-axis. Since the action of $`A`$ on $`V_k`$ is a rotation by the angle $`(2kq/p)\pi `$ for $`k=1,2,3`$, it follows that, unless either $`2q/p`$ or $`3q/p`$ are integers, then the only elements of $`_3(^{\mathrm{\hspace{0.17em}3}})`$ that are fixed by $`A`$ are the elements of $`V_0`$. Since these elements have a continuous symmetry group, and so, by hypothesis, cannot be $`h`$, it follows that either $`2q/p`$ or $`3q/p`$ are integers, i.e., that $`p=2`$ or $`p=3`$.
If $`p=2`$, then $`h`$ must lie in $`V_0+V_2`$, i.e., there must be constants $`r`$, $`s`$, and $`t`$, so that
$$h=rz(2z^23x^23y^2)+3\left(s(2xy)+t(x^2y^2)\right)z.$$
By a rotation that reverses the $`z`$-axis, if necessary, I can assume that $`r0`$ and then, by applying a rotation in $`\mathrm{SO}(2)`$, I can assume that $`t=0`$ and $`s0`$. Since $`G`$ is discrete, $`s`$ cannot be zero, so $`s>0`$. Note that $`A`$ is a rotation by an angle of $`\pi `$ about the $`z`$ axis, and that this certainly preserves any $`h`$ in the above form. Note also that every such $`h`$ has a linear factor. In particular, to each element $`A`$ of order $`2`$ in $`G`$, there corresponds a linear factor of $`h`$ that is fixed (up to a sign) by $`A`$.
If $`r=0`$, then $`h=6sxyz`$, and it is clear that the elements of $`G`$ must permute the planes $`x=0`$, $`y=0`$, and $`z=0`$. It follows that $`G`$ must be the group $`A_4`$ of order $`12`$ described in the proposition.
If $`r>0`$, then it is still true that $`h`$ has a linear factor, i.e.,
$$h=\left(2rz^23rx^23ry^2+6sxy\right)z.$$
When $`rs`$, the quadratic factor in the above expression is irreducible (since $`r`$ and $`s`$ are positive), so $`G`$ must stabilize the $`z`$-axis. In fact, since $`h`$ is positive on the positive ray of the $`z`$-axis, $`G`$ must actually be a subgroup of $`\mathrm{SO}(2)`$. Since $`s>0`$, $`G`$ must therefore be isomorphic to $`_2`$, generated by the rotation by $`\pi `$ about the $`z`$-axis. On the other hand, when $`r=s`$, the polynomial $`h`$ factors as
$$h=r(\sqrt{2}z\sqrt{3}x+\sqrt{3}y)(\sqrt{2}z+\sqrt{3}x\sqrt{3}y)z.$$
These three linear factors of $`h`$ are linearly dependent, so that $`h`$ vanishes on the union of three coaxial planes that meet pairwise at an angle of $`\pi /3`$. Consequently, $`h`$ lies on the $`\mathrm{SO}(3)`$-orbit of an element of the form
$$p(x^33xy^2)=px\left(x\sqrt{3}y\right)\left(x+\sqrt{3}y\right)$$
where $`p>0`$. Since $`G`$ must preserve these factors up to a sign, $`G`$ is isomorphic to $`S_3`$ and is generated as claimed in the proposition.
Finally, assume that $`G`$ has no element of order $`2`$. Then, by the above argument, all of the nontrivial elements of $`G`$ have order $`3`$. By the well-known classification of the finite subgroups of $`\mathrm{SO}(3)`$, <sup>6</sup><sup>6</sup>6 Up to conjugation, these subgroups consist of the cyclic subgroups, the dihedral subgroups, and the symmetry groups of the Platonic solids. it follows that $`G`$ must be isomorphic to $`_3`$.
Let $`A`$ be a generator of $`G`$ and assume (as one may, by replacing $`h`$ by an element in its $`\mathrm{SO}(3)`$-orbit) that $`A`$ is rotation by an angle of $`2\pi /3`$ about the $`z`$-axis. Then the elements of $`_3(^{\mathrm{\hspace{0.17em}3}})`$ that are fixed by $`A`$ are the elements in $`V_0+V_3`$, i.e., those of the form
$$h=rz(2z^23x^23y^2)+s(x^33xy^2)+t(3x^2yy^3).$$
By a rotation about the $`z`$-axis, $`h`$ can be replaced by an element in its orbit that is of the above form but that satisfies $`t=0`$ and $`s0`$. Now, $`s>0`$, since, otherwise the stabilizer of $`h`$ would contain $`\mathrm{SO}(2)`$. After rotation by an angle of $`\pi `$ about the $`x`$-axis if necessary, I can further assume that $`r0`$. In fact, $`r>0`$, since, otherwise, $`G`$ would be isomorphic to $`S_3`$, contrary to hypothesis.
It remains to determine those positive values of $`r`$ and $`s`$ (if any) for which
$$h=rz(2z^23x^23y^2)+s(x^33xy^2)$$
has a symmetry group larger than $`_3`$.
If the symmetry group $`G`$ is to be larger than $`_3`$, then, by the aforementioned classification, either $`G`$ contains an element of order $`2`$ or $`G`$ is infinite. In either case, by the above arguments, $`h`$ must have a linear factor. Now, it is straightforward to verify that $`h`$ has no linear factor unless $`s=r\sqrt{2}`$. Thus, the stabilizer is $`_3`$ except in this case. On the other hand
$$z(2z^23x^23y^2)+\sqrt{2}(x^33xy^2)=(z+\sqrt{2}x)\left(x+\sqrt{3}y\sqrt{2}z\right)\left(x\sqrt{3}y\sqrt{2}z\right),$$
and the three linear factors vanish on three mutually orthogonal $`3`$-planes. It follows immediately that this $`h`$ lies on the orbit of $`6pxyz`$ for some $`p>0`$. ∎
The argument for Proposition 1 can be used to prove two more easy results:
###### Proposition 2.
A cubic $`h_3(^{\mathrm{\hspace{0.17em}3}})`$ is reducible if and only if it has a symmetry of order $`2`$. It factors into three linear factors if and only if it is either the zero cubic, has symmetry $`A_4`$, or has symmetry $`S_3`$.
###### Proof.
By Proposition 1, any cubic that has a symmetry of order $`2`$ has a linear factor. Conversely, suppose that $`h_3(^{\mathrm{\hspace{0.17em}3}})`$ has a linear factor and is nonzero. By applying an $`\mathrm{SO}(3)`$ symmetry, it can be assumed that $`z`$ divides $`h`$, implying that $`h`$ has the form
$$h=z\left(r(2z^23x^23y^2)+3p(x^2y^2)+3q(2xy)\right),$$
which clearly has a symmetry of order $`2`$ that fixes $`z`$. The quadratic factor is reducible if and only if either $`r=0`$, in which case a rotation in the $`xy`$-plane reduces $`p`$ to zero, so that the symmetry group is $`A_4`$, or else $`p^2+q^2=r^2`$, in which case $`h`$ factors into three linearly dependent factors, so that the symmetry group is $`S_3`$. ∎
Before stating the next proposition, it will be useful to establish some notation. For any given linear function $`w:^{\mathrm{\hspace{0.17em}3}}^{}`$, the subgroup $`G_w\mathrm{SO}(3)`$ of rotations that preserve $`w`$ is isomorphic to $`\mathrm{SO}(2)=S^1`$. The induced representation of $`G_w`$ on $`_3(^{\mathrm{\hspace{0.17em}3}})^{\mathrm{\hspace{0.17em}7}}`$ is the sum of four $`\mathrm{G}_w`$-irreducible subspaces, $`V_0^w`$, $`V_1^w`$, $`V_2^w`$, and $`V_3^w`$, where $`V_0^w`$ has dimension $`1`$ and is the trivial representation and, for $`k>0`$, $`V_k^w`$ has dimension $`2`$ and is the representation on which a rotation by an angle of $`\alpha `$ in $`G_w`$ acts as a rotation by an angle of $`k\alpha `$.
For example, the proof of Proposition 1 lists an explicit basis for $`V_k^z`$ for $`0k3`$. Note that a cubic $`h_3(^{\mathrm{\hspace{0.17em}3}})`$ is linear in $`z`$ if and only if it lies in $`V_2^z+V_3^z`$. By symmetry, it follows that a cubic in $`_3(^{\mathrm{\hspace{0.17em}3}})`$ is linear in a variable $`w`$ if and only if it lies in $`V_2^w+V_3^w`$.
###### Proposition 3.
The set of cubics in $`_3(^{\mathrm{\hspace{0.17em}3}})`$ that are linear in some variable is a closed semi-analytic variety of codimension $`1`$ in $`_3(^{\mathrm{\hspace{0.17em}3}})`$ and consists of the cubics $`h_3(^{\mathrm{\hspace{0.17em}3}})`$ for which the projective plane curve $`h=0`$ has a real singular point.
Any cubic that is linear in two distinct variables is reducible and is on the $`\mathrm{SO}(3)`$-orbit of
$$h=3sxyz+sp(x^33xy^2),$$
which, in addition to being linear in $`z`$, is linear in $`w=y+pz`$ as well. When $`s`$ and $`p`$ are nonzero, this cubic is not linear in any other variables.
Any cubic that is linear in three distinct variables is on the $`\mathrm{SO}(3)`$-orbit of $`3sxyz`$ for some $`s0`$.
###### Proof.
A cubic $`h`$ is linear in a direction $`w`$ if and only if the direction generated by $`w`$ is a singular point of the projectivized curve $`h=0`$ in $`^2`$. Thus, the first statement follows, since the set of real cubic curves with a real singular point is a semi-analytic set of codimension $`1`$. If there are two distinct singular points, then the curve $`h=0`$ must be a union of a line with a conic. If there are three distinct singular points, then the curve $`h=0`$ must be the union of three nonconcurrent lines. Further details are left to the reader. ∎
### 3.2. Continuous symmetry
I now want to consider those special Lagrangian submanifolds $`L^{\mathrm{\hspace{0.17em}3}}`$ whose cubic second fundamental form has an $`\mathrm{SO}(2)`$ symmetry at each point.
###### Example 1 ($`\mathrm{SO}(3)`$-invariant special Lagrangian submanifolds).
By looking for special Lagrangian submanifolds of $`^{\mathrm{\hspace{0.17em}3}}=^{\mathrm{\hspace{0.17em}3}}+\mathrm{i}^{\mathrm{\hspace{0.17em}3}}`$ that are invariant under the ‘diagonal’ action of $`\mathrm{SO}(3)`$ on the two $`^{\mathrm{\hspace{0.17em}3}}`$-summands, Harvey and Lawson found the following examples:
$$L_c=\{(s+\mathrm{i}t)𝐮𝐮S^2^{\mathrm{\hspace{0.17em}3}},t^33s^2t=c^3\}.$$
Here, $`c`$ is a (real) constant. Note that $`L_0`$ is the union of three special Lagrangian $`3`$-planes. When $`c0`$, the submanifold $`L_c`$ has three components and each one is smooth and complete. In fact, these three components are isometric, as scalar multiplication in $`^{\mathrm{\hspace{0.17em}3}}`$ by a nontrivial cube root of unity permutes them cyclically. Each of these components is asymptotic to one pair of $`3`$-planes drawn from $`L_0`$. The $`\mathrm{SO}(3)`$-stabilizer of a point of $`L_c`$ is isomorphic to $`\mathrm{SO}(2)`$, so it follows that the fundamental cubic at each point has at least an $`\mathrm{SO}(2)`$-symmetry. It is not difficult to verify that this cubic is nowhere vanishing on $`L_c`$. Note also that, for $`\lambda `$ real and nonzero, $`\lambda L_c=L_{\lambda c}`$, so that, up to scaling, all of the $`L_c`$ with $`c0`$ are isometric.
###### Theorem 1.
If $`L^{\mathrm{\hspace{0.17em}3}}`$ is a connected special Lagrangian submanifold whose cubic fundamental form has an $`\mathrm{SO}(2)`$ symmetry at each point, then either $`L`$ is a $`3`$-plane or else $`L`$ is, up to rigid motion, an open subset of one of the Harvey-Lawson examples.
###### Proof.
Let $`L^{\mathrm{\hspace{0.17em}3}}`$ satisfy the hypotheses of the theorem. If the fundamental cubic $`C`$ vanishes identically, then $`L`$ is a $`3`$-plane, so assume that it does not. The locus where $`C`$ vanishes is a proper real-analytic subset of $`L`$, so its complement $`L^{}`$ is open and dense in $`L`$. Replace $`L`$ by a component of $`L^{}`$, so that it can be assumed that $`C`$ is nowhere vanishing on $`L`$.<sup>7</sup><sup>7</sup>7 In principle, this strategy could cause problems, but, as it is eventually going to be shown that $`L=L^{}`$ in the general case anyway, no problems ensue.
By Proposition 1, since the stabilizer of $`C_x`$ is $`\mathrm{SO}(2)`$ for all $`xL`$, there is a positive (real-analytic) function $`r:L^+`$ with the property that the equation
(3.1)
$$C=r\omega _1\left(2\omega _{1}^{}{}_{}{}^{2}3\omega _{2}^{}{}_{}{}^{2}3\omega _{3}^{}{}_{}{}^{2}\right)$$
defines an $`\mathrm{SO}(2)`$-subbundle $`FP_L`$ of the adapted coframe bundle $`P_LL`$. On the subbundle $`F`$, the following identities hold:
(3.2)
$$\left(\begin{array}{ccc}\beta _{11}& \beta _{12}& \beta _{13}\\ \beta _{21}& \beta _{22}& \beta _{23}\\ \beta _{31}& \beta _{32}& \beta _{33}\end{array}\right)=\left(\begin{array}{ccc}2r\omega _1& r\omega _2& r\omega _3\\ r\omega _2& r\omega _1& 0\\ r\omega _3& 0& r\omega _1\end{array}\right).$$
Moreover, because $`F`$ is an $`\mathrm{SO}(2)`$-bundle, relations of the form
(3.3)
$$\begin{array}{cc}\hfill \alpha _{21}& =t_{21}\omega _1+t_{22}\omega _2+t_{23}\omega _3\hfill \\ \hfill \alpha _{31}& =t_{31}\omega _1+t_{32}\omega _2+t_{33}\omega _3\hfill \end{array}$$
hold on $`F`$ for some functions $`t_{ij}`$. Moreover, for $`i=1`$, $`2`$, $`3`$ there exist functions $`r_i`$ on $`F`$ so that
(3.4)
$$\mathrm{d}r=r_i\omega _i.$$
Substituting the relations (3.2), (3.3), and (3.4) into the identities
(3.5)
$$\mathrm{d}\beta _{ij}=\beta _{ik}\alpha _{kj}\alpha _{ik}\beta _{kj}$$
and using the identities $`\mathrm{d}\omega _i=\alpha _{ij}\omega _j`$ then yields polynomial relations among these quantities that can be solved, leading to relations of the form
(3.6)
$$\begin{array}{cc}\hfill \alpha _{21}& =t\omega _2,\hfill \\ \hfill \alpha _{31}& =t\omega _3,\hfill \\ \hfill \mathrm{d}r& =4rt\omega _1,\hfill \end{array}$$
where, for brevity, I have written $`t`$ for $`t_{22}`$.
Note that (3.6) implies that $`d\omega _1=0`$. Differentiating the last equation in (3.6) implies that there exists a function $`u`$ on $`F`$ so that
(3.7)
$$\mathrm{d}t=u\omega _1.$$
Substituting (3.6) and (3.7) into the identities
(3.8)
$$\mathrm{d}\alpha _{ij}=\alpha _{ik}\alpha _{kj}+\beta _{ik}\beta _{kj}$$
and expanding, again using the identities $`\mathrm{d}\omega _i=\alpha _{ij}\omega _j`$, yields the relations
(3.9)
$$u=(3r^2t^2),d\alpha _{23}=(t^2+r^2)\omega _2\omega _3.$$
Differentiating these last equations yields only identities.
The structure equations found so far can be summarized as follows: $`FL`$ is an $`\mathrm{SO}(2)`$ bundle on which the 1-forms $`\omega _1,\omega _2,\omega _3,\alpha _{23}(=\alpha _{32})`$ are a basis. They satisfy the structure equations
(3.10)
$$\begin{array}{cc}\hfill \mathrm{d}\omega _1& =0,\hfill \\ \hfill \mathrm{d}\omega _2& =t\omega _1\omega _2\alpha _{23}\omega _3,\hfill \\ \hfill \mathrm{d}\omega _3& =t\omega _1\omega _3+\alpha _{23}\omega _3,\hfill \\ \hfill \mathrm{d}\alpha _{23}& =(t^2+r^2)\omega _2\omega _3,\hfill \\ \\ \hfill \mathrm{d}r& =4rt\omega _1,\hfill \\ \hfill \mathrm{d}t& =(3r^2t^2)\omega _1,\hfill \end{array}$$
and the exterior derivatives of these equations are identities.
These equations imply that $`\mathrm{d}\left(r^{3/2}+r^{1/2}t^2\right)=0`$. Since $`L`$ and $`F`$ are connected, it follows that there is a constant $`c>0`$ so that $`r^{3/2}+r^{1/2}t^2=c^{3/2}`$. Consequently, there is a function $`\theta `$ that is well-defined on $`L`$ that satisfies
$$r^{3/4}=c^{3/4}\mathrm{cos}3\theta ,r^{1/4}t=c^{3/4}\mathrm{sin}3\theta .$$
and the bound $`|\theta |<\pi /6`$. It then follows from the last two equations of (3.10) that
$$\omega _1=c\frac{\mathrm{d}\theta }{(\mathrm{cos}3\theta )^{4/3}}.$$
Moreover, setting $`\eta _i=c^1(\mathrm{cos}3\theta )^{1/3}\omega _i`$ for $`i=2`$ and $`3`$ yields
$$\mathrm{d}\eta _2=\alpha _{23}\eta _3,\mathrm{d}\eta _3=\alpha _{23}\eta _2,\mathrm{d}\alpha _{23}=\eta _2\eta _3,$$
which are the structure equations of the metric of constant curvature $`1`$ on $`S^2`$.
Conversely, if $`d\sigma ^2`$ is the metric of constant curvature $`1`$ on $`S^2`$, then, on the product $`L=(\pi /6,\pi /6)\times S^2`$, consider the quadratic and cubic forms defined by
$$g=c^2\frac{\mathrm{d}\theta ^2+\mathrm{cos}^23\theta d\sigma ^2}{(\mathrm{cos}3\theta )^{8/3}}\text{and}C=c^2\frac{2\mathrm{d}\theta ^33\mathrm{cos}^23\theta \mathrm{d}\theta d\sigma ^2}{(\mathrm{cos}3\theta )^{8/3}}.$$
The metric $`g`$ is complete and the pair $`(g,C)`$ satisfy the Gauss and Codazzi equations that ensure that $`(L,g)`$ can be isometrically embedded as a special Lagrangian $`3`$-fold in $`^{\mathrm{\hspace{0.17em}3}}`$ inducing $`C`$ as the fundamental cubic. Thus, for each value of $`c`$, there exists a corresponding special Lagrangian $`3`$-fold that is complete and unique up to special Lagrangian isometries of $`^{\mathrm{\hspace{0.17em}3}}`$.
Since the parameter $`c`$ is accounted for by dilation in $`^{\mathrm{\hspace{0.17em}3}}`$, it now follows that these special Lagrangian $`3`$-folds are the Harvey-Lawson examples, as desired. Note that since these are complete and since $`r`$ is nowhere vanishing, it follows that $`L^{}=L`$ for the Harvey-Lawson examples, and hence for all examples. ∎
### 3.3. $`A_4`$ symmetry
Now consider those special Lagrangian submanifolds $`L^{\mathrm{\hspace{0.17em}3}}`$ whose fundamental cubic has an $`A_4`$-symmetry at each point. Unfortunately, I cannot begin the discussion by providing a nontrivial example.
###### Theorem 2.
The only special Lagrangian submanifold of $`^{\mathrm{\hspace{0.17em}3}}`$ whose fundamental cubic has an $`A_4`$-symmetry at each point is a special Lagrangian $`3`$-plane.
###### Remark 3.
It is interesting to compare the results of Theorems 1 and 2. The $`\mathrm{SO}(3)`$-orbits consisting of the $`h_3(^{\mathrm{\hspace{0.17em}3}})`$ that have an $`\mathrm{SO}(2)`$ symmetry form a cone of dimension $`3`$ in $`_3(^{\mathrm{\hspace{0.17em}3}})`$, while the ones with an $`A_4`$-symmetry form a cone of dimension $`4`$ in $`_3(^{\mathrm{\hspace{0.17em}3}})`$. Thus, one might expect, based on ‘equation counting’, that the the condition of having all cubics have a $`\mathrm{SO}(2)`$-symmetry would have fewer solutions than the condition of having all cubics have an $`A_4`$-symmetry. However, just the opposite is true.
###### Proof.
Let $`L^{\mathrm{\hspace{0.17em}3}}`$ be a connected special Lagrangian submanifold with the property that its fundamental cubic $`C`$ has an $`A_4`$-symmetry at each point. If $`C`$ vanishes identically, then $`L`$ is an open subset of a special Lagrangian $`3`$-plane, so assume that it does not. Let $`L^{}L`$ be the dense open subset where $`C`$ is nonzero.
By Proposition 1, since the stabilizer of $`C_x`$ is $`A_4`$ for all $`xL^{}`$, there is a positive (real-analytic) function $`r:L^+`$ for which the equation
(3.11)
$$C=6r\omega _1\omega _2\omega _3$$
defines an $`A_4`$-subbundle $`FP_L`$ over $`L^{}`$ of the adapted coframe bundle $`P_LL`$. On $`F`$, the following identities hold:
(3.12)
$$\left(\begin{array}{ccc}\beta _{11}& \beta _{12}& \beta _{13}\\ \beta _{21}& \beta _{22}& \beta _{23}\\ \beta _{31}& \beta _{32}& \beta _{33}\end{array}\right)=\left(\begin{array}{ccc}0& r\omega _3& r\omega _2\\ r\omega _3& 0& r\omega _1\\ r\omega _2& r\omega _1& 0\end{array}\right).$$
Since $`F`$ is an $`A_4`$-bundle over $`L^{}`$, there are relations
(3.13)
$$\begin{array}{cc}\hfill \alpha _{23}& =t_{11}\omega _1+t_{12}\omega _2+t_{13}\omega _3\hfill \\ \hfill \alpha _{31}& =t_{21}\omega _1+t_{22}\omega _2+t_{23}\omega _3\hfill \\ \hfill \alpha _{12}& =t_{31}\omega _1+t_{32}\omega _2+t_{33}\omega _3\hfill \end{array}$$
holding on $`F`$ for some functions $`t_{ij}`$. Moreover, there exist functions $`r_i`$ for $`i=1,2,3`$ on $`F`$ so that
(3.14)
$$\mathrm{d}r=r_i\omega _i.$$
Substituting the relations (3.12), (3.13), and (3.14) into the identities
(3.15)
$$\mathrm{d}\beta _{ij}=\beta _{ik}\alpha _{kj}\alpha _{ik}\beta _{kj}$$
and using the identities $`\mathrm{d}\omega _i=\alpha _{ij}\omega _j`$ then yields polynomial relations among these quantities that can be solved, leading to relations
(3.16)
$$\alpha _{ij}=0,\mathrm{d}r=0.$$
Substituting (3.12) and (3.16) into the identities
(3.17)
$$\mathrm{d}\alpha _{ij}=\alpha _{ik}\alpha _{kj}+\beta _{ik}\beta _{kj}$$
yields $`r=0`$, contrary to hypothesis. ∎
### 3.4. $`S_3`$ symmetry
Now consider the special Lagrangian submanifolds of $`^{\mathrm{\hspace{0.17em}3}}`$ whose fundamental cubic has an $`S_3`$-symmetry at every point. In contrast to the case of $`A_4`$-symmetry, there clearly are nontrivial examples of this type.
###### Example 2 (Products).
This example is fairly trivial: Write $`^{\mathrm{\hspace{0.17em}3}}=^{\mathrm{\hspace{0.17em}1}}\times ^{\mathrm{\hspace{0.17em}2}}`$ and look for special Lagrangian submanifolds of the form $`L=^{}\times \mathrm{\Sigma }`$, where $`\mathrm{\Sigma }^{\mathrm{\hspace{0.17em}2}}`$ is a surface. It is not difficult to see that there is a unique complex structure on $`^{\mathrm{\hspace{0.17em}2}}`$ (not the given one!) with the property that $`L`$ is special Lagrangian if and only if $`\mathrm{\Sigma }`$ is a complex curve with respect to this structure.
Explicitly, writing $`z_k=x_k+\mathrm{i}y_k`$, one sees $`L=^{}\times \mathrm{\Sigma }`$ is special Lagrangian for $`\mathrm{\Sigma }^{\mathrm{\hspace{0.17em}2}}`$ if and only if the $`2`$-forms $`\mathrm{d}x_2\mathrm{d}y_2+\mathrm{d}x_3\mathrm{d}y_3`$ and $`\mathrm{d}x_2\mathrm{d}y_3+\mathrm{d}y_2\mathrm{d}x_3`$ each vanish when pulled back to $`\mathrm{\Sigma }`$. Since
$$(\mathrm{d}x_2\mathrm{d}y_2+\mathrm{d}x_3\mathrm{d}y_3)+\mathrm{i}(\mathrm{d}x_2\mathrm{d}y_3+\mathrm{d}y_2\mathrm{d}x_3)=(\mathrm{d}x_2\mathrm{i}\mathrm{d}x_3)(\mathrm{d}y_2+\mathrm{i}\mathrm{d}y_3),$$
these $`2`$-forms vanish on $`\mathrm{\Sigma }`$ if and only if $`\mathrm{\Sigma }`$ is a complex curve in $`^{\mathrm{\hspace{0.17em}2}}`$ endowed with the complex structure for which $`u=x_2\mathrm{i}x_3`$ and $`v=y_2+\mathrm{i}y_3`$ are holomorphic.
Now, each of these special Lagrangian $`3`$-folds is easily seen to have its fundamental cubic be expressible as a cubic polynomial in a pair of $`1`$-forms, from which it follows from Proposition 1 that the $`\mathrm{SO}(3)`$-symmetry group of the cubic at each point is either everything (if the cubic vanishes at the given point) or else isomorphic to $`S_3`$.
###### Example 3 (Special Lagrangian cones).
A more interesting example is to consider the special Lagrangian cones. Suppose that $`\mathrm{\Sigma }S^5`$ is a (possibly immersed) surface with the property that the cone $`C(\mathrm{\Sigma })^{\mathrm{\hspace{0.17em}3}}`$ is special Lagrangian. Then it is not difficult to show that the fundamental cubic of $`C(\mathrm{\Sigma })`$ has an $`S_3`$-stabilizer at those points where it is not zero. (This is because the cubic form uses only two of the directions.)
The necessary and sufficient conditions on $`\mathrm{\Sigma }`$ that $`C(\mathrm{\Sigma })`$ be special Lagrangian are easily stated: Let $`𝐮:S^5^{\mathrm{\hspace{0.17em}3}}`$ be the inclusion mapping. Define a $`1`$-form $`\theta `$ on $`S^5`$ by $`\theta =J𝐮\mathrm{d}𝐮`$ and define a $`2`$-form $`\mathrm{\Psi }`$ on $`S^5`$ by $`\mathrm{\Psi }=𝐮\text{ }\text{ }\mathrm{Im}(\mathrm{{\rm Y}})`$. Then $`\mathrm{\Sigma }S^5`$ has the property that $`C(\mathrm{\Sigma })`$ is special Lagrangian if and only if $`\theta `$ and $`\mathrm{\Psi }`$ vanish when pulled back to $`\mathrm{\Sigma }`$. An elementary application of the Cartan-Kähler theorem shows that any real-analytic curve $`\gamma S^5`$ to which $`\theta `$ pulls back to be zero lies in an irreducible real-analytic surface $`\mathrm{\Sigma }`$ that satisfies these conditions. Thus, there are many such surfaces. (In the terminology of exterior differential systems, these surfaces depend on two functions of one variable.)
In addition, many explicit examples of such surfaces are now known. For example, in , a thorough study is done of the special Lagrangian cones that are invariant under a circle action. In fact, the differential equation for these surfaces admits a Bäcklund transformation and can be formulated as an integrable system. In principle, the compact torus solutions can be described explicitly in terms of $`\vartheta `$-functions via loop group constructions.
###### Example 4 (Twisted special Lagrangian cones).
The special Lagrangian cones can be generalized somewhat, using a construction found in \[6, §4\].
Again, let $`𝐱:\mathrm{\Sigma }S^5`$ be an immersion of a simply connected surface for which the cone on $`𝐱(\mathrm{\Sigma })`$ is special Lagrangian. Endow $`\mathrm{\Sigma }`$ with the metric and orientation that it inherits from this immersion and let $`:\mathrm{\Omega }^p(\mathrm{\Sigma })\mathrm{\Omega }^{2p}(\mathrm{\Sigma })`$ be the associated Hodge star operator. Since $`\mathrm{\Sigma }`$ is minimal, it follows that
(3.18)
$$\mathrm{d}(\mathrm{d}𝐱)+2𝐱=0.$$
Now, let $`b:\mathrm{\Sigma }`$ be any solution to the second order, linear elliptic equation
(3.19)
$$\mathrm{d}(\mathrm{d}b)+2b=0.$$
(For example, $`b`$ could be one of the components of $`𝐱`$.) Equations (3.18) and (3.19) imply that the vector-valued $`1`$-form
(3.20)
$$\beta =𝐱\mathrm{d}bb\mathrm{d}𝐱$$
is closed. Thus, there exists a $`^{\mathrm{\hspace{0.17em}3}}`$-valued function $`𝐛:\mathrm{\Sigma }^{\mathrm{\hspace{0.17em}3}}`$ so that $`\mathrm{d}𝐛=\beta `$.
Now, consider the immersion $`X:^{}\times \mathrm{\Sigma }^{\mathrm{\hspace{0.17em}3}}`$ defined by
(3.21)
$$X=𝐛+t𝐱.$$
Since $`dX=𝐱(dt+\mathrm{d}b)+t\mathrm{d}𝐱b\mathrm{d}𝐱`$, it follows that $`X`$ immerses $`^{}\times \mathrm{\Sigma }`$ as a special Lagrangian $`3`$-fold in $`^{\mathrm{\hspace{0.17em}3}}`$, at least away from the locus $`t=b=0`$ in $`^{}\times \mathrm{\Sigma }`$, where $`X`$ fails to be an immersion. Moreover, at those places where the fundamental cubic of this immersed submanifold is nonzero, it has $`S_3`$-symmetry.
It turns out that the image $`X(^{}\times \mathrm{\Sigma })`$ determines the data $`𝐱:\mathrm{\Sigma }S^5`$ and $`b:\mathrm{\Sigma }^{}`$ up to a replacement of the form $`(𝐱,b)(𝐱,b)`$, except in the case that $`𝐱(\mathrm{\Sigma })`$ lies in a special Lagrangian $`3`$-plane, in which case, $`X(^{}\times \mathrm{\Sigma })`$ lies in a parallel $`3`$-plane.
Note that when $`b=0`$, the function $`𝐛`$ is constant, so that $`X(^{}\times \mathrm{\Sigma })`$ is just a translation of the cone on $`\mathrm{\Sigma }`$. Thus, these examples properly generalize the special Lagrangian cones. I will refer to these examples as *twisted special Lagrangian cones*.
As explained in , this example can be generalized somewhat by allowing $`𝐱:\mathrm{\Sigma }S^5`$ to be a branched immersion that is an integral manifold of $`\theta `$ and $`\mathrm{\Psi }`$, but then one must allow $`b`$ to have ‘pole-type’ singularities at the branch points of the immersion $`𝐱`$.
###### Theorem 3.
Suppose that $`L^{\mathrm{\hspace{0.17em}3}}`$ is a connected special Lagrangian $`3`$-fold with the property that its fundamental cubic at each point has an $`S_3`$-symmetry. Then either $`L`$ is congruent to a product $`^{}\times \mathrm{\Sigma }`$ as in Example 2, or else $`L`$ contains a dense open set $`L^{}L`$ such that every point of $`L^{}`$ has a neighborhood that lies in a twisted special Lagrangian cone $`X(^{}\times \mathrm{\Sigma })`$, as in Example 4.
###### Proof.
Suppose that $`L^{\mathrm{\hspace{0.17em}3}}`$ satisfies the hypotheses of the theorem. If the fundamental cubic $`C`$ vanishes identically on $`L`$, then $`L`$ is a $`3`$-plane and there is nothing to show, so suppose that $`C0`$. Let $`L^{}L`$ be the open dense subset where $`C0`$.
The hypothesis that $`C_x`$ has $`S_3`$-symmetry at every $`xL^{}`$ implies that there is a positive function $`s:L^{}^{}`$ and an $`S_3`$-subbundle $`FP_L`$ over $`L^{}`$ with projection $`𝐱:FL^{}^{\mathrm{\hspace{0.17em}3}}`$ on which the identity
(3.22)
$$C=s\left(\omega _{2}^{}{}_{}{}^{3}3\omega _2\omega _{3}^{}{}_{}{}^{2}\right)$$
holds. In particular, the second fundamental form of $`L^{}`$ has the form
(3.23)
$$\text{II}=J𝐞_2s(\omega _{2}^{}{}_{}{}^{2}\omega _{3}^{}{}_{}{}^{2})+J𝐞_3s(2\omega _2\omega _3),$$
where $`𝐞_1,𝐞_2,𝐞_3`$ are the vector-valued functions defined by the moving frame relation $`\mathrm{d}𝐱=𝐞_1\omega _1+𝐞_2\omega _2+𝐞_3\omega _3`$.
It follows from (3.23) that $`L^{\mathrm{\hspace{0.17em}3}}`$ is an austere submanifold of dimension $`3`$. By Theorem 4.1 of , it follows that either $`L^{}`$ is locally the product of a line in $`^{\mathrm{\hspace{0.17em}3}}`$ with a minimal surface $`\mathrm{\Sigma }`$ in the orthogonal $`5`$-plane, or else there exists a dense open subset $`L^{}L^{}`$ so that every point of $`L^{}`$ has an open neighborhood in $`L^{}`$ that lies in a twisted cone constructed as in Example 4 from a minimal immersion $`𝐱:\mathrm{\Sigma }S^5`$ and an auxiliary function $`b:\mathrm{\Sigma }`$ satisfying (3.19).
Since the group of translations and $`\mathrm{SU}(3)`$-rotations in $`^{\mathrm{\hspace{0.17em}3}}`$ acts transitively on the space of lines, it follows that if $`L`$ is locally an orthogonal product $`^{}\times \mathrm{\Sigma }`$ and is special Lagrangian, then, up to translation by a constant, $`\mathrm{\Sigma }`$ must be a complex curve in the complex $`2`$-plane $`P`$ orthogonal to the linear factor, where the complex structure on $`P`$ is taken to be as defined in Example 2.
On the other hand, if $`L`$ is not locally an orthogonal product and so is a twisted cone as described above, then one sees from the formula for $`dX`$ derived in Example 4 that the immersion $`𝐱:\mathrm{\Sigma }S^5`$ must not only be minimal, but must have the property that $`𝐱^{}\theta =𝐱^{}\mathrm{\Psi }=0`$ as well, as desired. ∎
###### Remark 4 (Singular behavior).
The reader may be annoyed by the apparent need to restrict to the open dense subset $`L^{}L`$. However, there are subtle singularity issues that seem to require this. For more discussion, see the final pages of .
###### Remark 5 (Austerity).
Theorem 3 implies that the austere special Lagrangian 3-folds in $`^{\mathrm{\hspace{0.17em}3}}`$ are completely described by Examples 2 and 4.
###### Remark 6 (Generality).
The reader knowledgeable about exterior differential systems may wonder about the generality of the austere special Lagrangian $`3`$-folds in the sense of Cartan-Kähler theory. While I have avoided this approach to the analysis of these examples in this treatment, I should confess that I first understood the local geometry of these examples by doing a Cartan-Kähler analysis. The obvious exterior differential system that one writes down for these examples is involutive, with Cartan characters $`s_1=4`$ and $`s_2=s_3=0`$. The characteristic variety of the involutive prolongation consists of two complex conjugate points, each of multiplicity $`2`$.
#### 3.4.1. Structure equations
For use in the next section, I will record here the structure equations that one derives for systems of this kind. I will maintain the notation established in the proof of Theorem 3 for $`L^{}L`$, the function $`s`$, and the $`S_3`$-bundle $`\pi :FL^{}`$. Thus, the fundamental cubic factors as
(3.24)
$$C=s\left(\omega _{2}^{}{}_{}{}^{3}3\omega _2\omega _{3}^{}{}_{}{}^{2}\right)=s\omega _2(\omega _2+\sqrt{3}\omega _3)(\omega _2\sqrt{3}\omega _3).$$
In particular, the equations
(3.25)
$$\left(\begin{array}{ccc}\beta _{11}& \beta _{12}& \beta _{13}\\ \beta _{21}& \beta _{22}& \beta _{23}\\ \beta _{31}& \beta _{32}& \beta _{33}\end{array}\right)=\left(\begin{array}{ccc}0& 0& 0\\ 0& s\omega _2& s\omega _3\\ 0& s\omega _3& s\omega _2\end{array}\right)$$
hold on $`F`$. Moreover, because $`F`$ is an $`S_3`$-bundle, relations of the form
(3.26)
$$\begin{array}{cc}\hfill \alpha _{23}& =t_{11}\omega _1+t_{12}\omega _2+t_{13}\omega _3\hfill \\ \hfill \alpha _{31}& =t_{21}\omega _1+t_{22}\omega _2+t_{23}\omega _3\hfill \\ \hfill \alpha _{12}& =t_{31}\omega _1+t_{32}\omega _2+t_{33}\omega _3\hfill \end{array}$$
hold on $`F`$ for some functions $`t_{ij}`$. Also, for $`i=1,2,3`$ there exist functions $`s_i`$ on $`F`$ so that
(3.27)
$$\mathrm{d}s=s_i\omega _i.$$
Substituting the relations (3.25), (3.26), and (3.27) into the identities
(3.28)
$$\mathrm{d}\beta _{ij}=\beta _{ik}\alpha _{kj}\alpha _{ik}\beta _{kj}$$
and using the identities $`\mathrm{d}\omega _i=\alpha _{ij}\omega _j`$ then yields polynomial relations among these quantities that can be solved, leading to relations of the form
(3.29)
$$\begin{array}{cc}\hfill \alpha _{23}& =r_2\omega _1t_2\omega _2t_3\omega _3,\hfill \\ \hfill \alpha _{31}& =3r_2\omega _23r_3\omega _3,\hfill \\ \hfill \alpha _{12}& =3r_3\omega _23r_2\omega _3,\hfill \\ \\ \hfill \mathrm{d}s& =3s\left(r_3\omega _1t_3\omega _2+t_2\omega _3\right),\hfill \end{array}$$
where I have renamed the covariant derivative variables in the solution for simplicity and symmetry of notation.
It is worth mentioning that not all of these functions on $`F`$ are invariant under the action of the group $`S_3`$ on the fibers. The functions $`s`$ and $`r_2`$ are invariant, the function $`r_3`$ and the $`1`$-form $`\omega _1`$ are invariant under the odd order elements of $`S_3`$ but switch sign under an element of order $`2`$, and the complex function $`t=(t_2,t_3):F^{\mathrm{\hspace{0.17em}2}}`$ is $`S_3`$-equivariant when $`^{\mathrm{\hspace{0.17em}2}}`$ is appropriately identified with the nontrivial irreducible representation of dimension $`2`$ of $`S_3`$. Thus, $`s`$, $`r_2`$, $`r_{3}^{}{}_{}{}^{2}`$, and $`t_{2}^{}{}_{}{}^{2}+t_{3}^{}{}_{}{}^{2}`$ are all well-defined on $`L`$, but $`r_3`$, for example, is only well-defined up to a sign.
Because the exterior differential system mentioned above is involutive, it can be shown that one can prescribe the functions $`t_2`$, $`t_3`$, $`r_2`$, and $`r_3`$ essentially arbitrarily along any curve on which $`\omega _{2}^{}{}_{}{}^{2}+\omega _{3}^{}{}_{}{}^{2}`$ is nonzero (the curves defined by the differential equations $`\omega _2=\omega _3=0`$ are characteristic) and generate a solution.
The functions $`r_2`$ and $`r_3`$ vanish identically if and only if $`L`$ is an orthogonal product. Otherwise, $`L`$ is (locally) a twisted cone.
The structure equations derived so far imply that
$$\omega _2\mathrm{d}\omega _2=\frac{1}{4}(\omega _2\pm \sqrt{3}\omega _3)\mathrm{d}(\omega _2\pm \sqrt{3}\omega _3)=2r_2\omega _1\omega _2\omega _3,$$
so it follows that the three linear factors of $`C`$ define integrable $`2`$-plane fields on $`L^{}`$ if and only if $`r_20`$ (and that if any one of the three is integrable, then so are the other two).
If one considers the differential system with the additional condition $`r_20`$, one sees that it implies the structure equation $`dr_3=3r_{3}^{}{}_{}{}^{2}\omega _1`$ and that the reduced system, with this condition added, is still involutive, but now with Cartan characters $`s_1=2`$ and $`s_2=s_3=0`$. In fact, the condition $`r_2=0`$ characterizes the special Lagrangian cones and (under the additional condition $`r_3=0`$) the orthogonal products.
Finally, note that, when $`r_2=0`$, the $`2`$-dimensional leaves of the $`2`$-plane field defined by $`\omega _2=0`$ will not lie in $`3`$-planes unless $`t_30`$. Since the condition $`t_3=0`$ is not $`S_3`$-invariant unless $`t_2=0`$ as well, it follows that, except in the very special case $`r_2=t_3=t_2=0`$, at most one of the three foliations has its leaves lying in $`3`$-planes.
It is not difficult to show that, up to congruence, there is only one example that satisfies $`r_2t_3t_20`$, namely, the Harvey-Lawson example $`L^{\mathrm{\hspace{0.17em}3}}`$ defined in the standard coordinates by the equations $`|z_1|^2=|z_2|^2=|z_3|^2`$ and $`\mathrm{Im}(z_1z_2z_3)=0`$. This cone is cut into surfaces by three distinct families of Lagrangian planes. For example, each element of the circle of Lagrangian planes defined by the relations
$$z_1\mathrm{e}^{\mathrm{i}\theta }\overline{z_2}=z_3\mathrm{e}^{2\mathrm{i}\theta }\overline{z_3}=0$$
meets $`L`$ in a $`2`$-dimensional cone. One gets the other two families by permuting the coordinates $`z_i`$.
In the case that only one of the three linear divisors of $`C`$ defines a foliation by surfaces that lie in $`3`$-planes, one can reduce to a $`_2`$-subbundle of $`F`$ by imposing the condition that $`t_30`$. Then, by pursuing the calculation of the integrability conditions, one finds that the remaining quantities $`r_3`$ and $`t_2`$ must satisfy the equations
(3.30)
$$\begin{array}{cc}\hfill \mathrm{d}r_3& =3r_{3}^{}{}_{}{}^{2}\omega _1\hfill \\ \hfill \mathrm{d}t_2& =3t_2r_3\omega _1+(t_{2}^{}{}_{}{}^{2}+9r_{3}^{}{}_{}{}^{2}2s^2)\omega _3\hfill \end{array}$$
Note that if $`r_3`$ vanishes anywhere, it vanishes identically. As already mentioned, this is the case of a product. It is not difficult to show that any connected example of this kind is congruent to an open subset of the special Lagrangian $`3`$-fold $`L_c`$ defined by the equations
$$y_1=(x_2\mathrm{i}x_3)^2(y_2+\mathrm{i}y_3)^2c^2=0,$$
where $`c>0`$ is a real parameter. This meets the circle of Lagrangian planes defined by
$$y_1=\mathrm{cos}\theta x_2\mathrm{sin}\theta y_2=\mathrm{cos}\theta x_3+\mathrm{sin}\theta y_3=0$$
in congruent surfaces that are hyperbolic cylinders.
On the other hand, if $`r_3`$ is nonzero, one can reduce the structure bundle to a parallelization of $`L`$ by imposing the conditions $`t_3=0`$ and $`r_3>0`$, so assume this. By the structure equations, the expression $`G=(s^2+t_{2}^{}{}_{}{}^{2}+9r_{3}^{}{}_{}{}^{2})s^{2/3}r_{3}^{}{}_{}{}^{4/3}`$ is constant on $`L`$. Moreover, one easily sees from the structure equations that the vector field $`X`$ that satisfies $`\omega _1(X)=\omega _3(X)=0`$ and $`\omega _2(X)=s^{1/3}r_{3}^{}{}_{}{}^{2/3}`$ is a symmetry vector field of the system and hence must correspond to an ambient symmetry of the corresponding solution. Since this symmetry must fix the vertex of the cone, it follows that it is a rotation. Pursuing this observation, it is not difficult to show that all of these solutions can be described as follows: Let $`\lambda _1\lambda _2>0>\lambda _3`$ be real numbers satisfying $`\lambda _1+\lambda _2+\lambda _3=0`$, and consider the $`3`$-fold $`L_\lambda ^{\mathrm{\hspace{0.17em}3}}`$ consisting of the points of the form
$$(r_1\mathrm{e}^{\mathrm{i}(\pi /6+\lambda _1t)},r_2\mathrm{e}^{\mathrm{i}(\pi /6+\lambda _2t)},r_3\mathrm{e}^{\mathrm{i}(\pi /6+\lambda _3t)})$$
where $`t`$$`r_1`$, $`r_2`$, and $`r_3`$ are real numbers satisfying $`\lambda _1r_{1}^{}{}_{}{}^{2}+\lambda _2r_{2}^{}{}_{}{}^{2}+\lambda _3r_{3}^{}{}_{}{}^{2}=0`$. Then $`L_\lambda `$ is a special Lagrangian cone with a foliation by $`2`$-dimensional $`3`$-plane slices given by $`\mathrm{d}t=0`$. (These $`3`$-plane slices are all congruent and are Euclidean cones.) Moreover, every $`L`$ of the type under discussion is congruent to $`L_\lambda `$ for some $`\lambda `$.
Note, by the way, that $`L_\lambda `$ is not closed unless the ratios of the $`\lambda _i`$ are rational. Thus, for ‘generic’ $`\lambda `$, the cone $`L_\lambda `$ is dense in the $`4`$-dimensional cone in $`^{\mathrm{\hspace{0.17em}3}}`$ defined by the equations
$$\lambda _1|z_1|^2+\lambda _2|z_2|^2+\lambda _3|z_3|^2=\mathrm{Re}(z_1z_2z_3)=0.$$
Part of the significance of these examples will be explained in the next section.
### 3.5. $`_2`$ symmetry
Now consider a special Lagrangian submanifold $`L^{\mathrm{\hspace{0.17em}3}}`$ whose fundamental cubic $`C`$ has a $`_2`$-symmetry at each point. Equivalently, by Proposition 2, this is the same as assuming that the fundamental cubic $`C`$ is reducible at each point.
Several nontrivial examples have already been seen: In fact, if $`C_x`$ has a continuous stabilizer at each $`x`$ or if $`C_x`$ has an $`S_3`$-stabilizer at each $`x`$, then Proposition 1 shows that $`C_x`$ must be reducible at each point. In the first case, the examples are classified by Theorem 1 and in the second case, the examples are classified by Theorem 3. However, these examples have stabilizer groups strictly larger than $`_2`$, so the interesting question is whether there exist any other examples. By Proposition 1 and Theorem 2, any such example $`L`$ will have to have the property that the $`\mathrm{SO}(3)`$-stabilizer of $`C_x`$ is exactly $`_2`$ for generic $`xL`$.
Before discussing explicit examples, I will describe a geometrically interesting condition that forces there to be a $`_2`$-symmetry of $`C_x`$ for all $`xL`$.
###### Proposition 4.
Let $`L^{\mathrm{\hspace{0.17em}3}}`$ be a special Lagrangian submanifold that supports a smooth codimension $`1`$ foliation $`𝒮`$ with the property that each $`𝒮`$-leaf $`SL`$ lies in a $`3`$-plane. Then $`C_x`$ is reducible for all $`xL`$. In particular, the $`\mathrm{SO}(3)`$-stabilizer of $`C_x`$ contains an element of order $`2`$.
###### Proof.
It suffices to assume that $`L`$ is connected, so do this.
If any $`𝒮`$-leaf $`S`$ is planar, even locally, then this plane must be $`\omega `$-isotropic and Harvey and Lawson’s Theorem 5.5 of §III in implies that $`L`$ itself must contain an open subset of a special Lagrangian $`3`$-plane. By real-analyticity, it follows that $`L`$ itself is planar and hence that $`C_x`$ vanishes identically for all $`xL`$. Thus, from now on, I can assume that none of the $`𝒮`$-leaves are planar and that $`L`$ itself is nonplanar.
Choose $`xL`$ and restrict $`L`$ to a neighborhood $`U`$ on which the foliation can be expressed a product, i.e., $`U=X\left((0,1)\times D\right)`$ for some open domain $`D^2`$, and the $`𝒮`$-leaves in $`U`$ are of the form $`X(t,D)`$ for $`t(0,1)`$. Then, by hypothesis, for each $`t(0,1)`$, there exists a unique real $`3`$-plane $`P(t)^{\mathrm{\hspace{0.17em}3}}`$ so that $`UP(t)=X(t,D)`$, and the surface $`UP(t)`$ is $`\omega `$-isotropic. Since the surface $`UP(t)`$ is non-planar, the plane $`P(t)`$ itself must Lagrangian, although it cannot be special Lagrangian, since, otherwise, the uniqueness aspect of Harvey and Lawson’s Theorem $`5.5`$ would imply that $`UP(t)`$, contradicting the assumption that $`L`$ is not planar. It is not difficult to see that the curve $`tP(t)`$ must be smooth, since the foliation $`𝒮`$ is assumed to be smooth.
Now, consider the $`\mathrm{SO}(2)`$-subbundle $`FP_L`$ over $`U`$ with the property that the vector-valued functions $`𝐞_2`$ and $`𝐞_3`$ are an oriented basis of the tangent space to the $`𝒮`$-leaves. Then $`\omega _1`$ is well-defined on $`U`$ and vanishes when pulled back to any $`𝒮`$-leaf.
Now, the set of Lagrangian planes that contain $`𝐞_2`$ and $`𝐞_3`$ is the circle of $`3`$-planes that contain $`𝐞_2`$, $`𝐞_3`$ and that are contained in the span of $`𝐞_1,J𝐞_1,𝐞_2,𝐞_3`$. In particular, $`P(t)`$ lies in this plane for each leaf $`X(t,D)U`$. Since each leaf $`\omega _1=0`$ lies in $`P(t)`$, it follows that the second fundamental form
$$\text{II}=J𝐞_1Q_1+J𝐞_2Q_2+J𝐞_3Q_3$$
has the property that $`Q_2`$ and $`Q_3`$ must vanish when restricted to the 2-planes defined by $`\omega _1=0`$, i.e., it must be true that $`Q_2`$ and $`Q_3`$ are multiples of $`\omega _1`$. However, by Euler’s homogeneity relation
$$C=\omega _1Q_1+\omega _2Q_2+\omega _3Q_3,$$
it now follows that $`C`$ itself must be a multiple of $`\omega _1`$, i.e., $`C`$ is reducible at every point of $`U`$, as desired.
Finally, by Proposition 2, the $`\mathrm{SO}(3)`$-stabilizer of $`C_x`$ must contain an element of order $`2`$ for all $`xL`$. ∎
###### Remark 7 (Non-integrable factors and non-planar foliations).
It is worth pointing out that there are examples of special Lagrangian $`3`$-folds $`L^{\mathrm{\hspace{0.17em}3}}`$ for which the fundamental cubic $`C`$ is reducible, but for which the factors of $`C`$ do not define codimension $`1`$ foliations of $`L`$. In fact, by the discussion in 3.4.1, it follows that, for the generic special Lagrangian $`3`$-fold $`L`$ for which the fundamental cubic $`C`$ has an $`\mathrm{SO}(3)`$-stabilizer isomorphic to $`S_3`$, the cubic $`C`$ factors into three linear factors, no one of which defines an integrable $`2`$-plane field.
Moreover, even in the case where $`r_20`$ (in which case, $`L`$ is a cone), so that the three factors are each integrable, the leaves of the three foliations will not lie in $`3`$-planes unless $`t_30`$, which does not hold for the general special Lagrangian cone.
###### Example 5 (Lawlor-Harvey).
This example was first found by Lawlor , and was subsequently generalized and extended by Harvey \[11, 7.78–9\]. While their results are valid in all dimensions, I will only discuss the dimension $`3`$ case.
They show that, for any compact $`2`$-dimensional ellipsoid $`EP`$ where $`P^{\mathrm{\hspace{0.17em}3}}`$ is a Lagrangian (but not special Lagrangian) $`3`$-plane, the special Lagrangian extension $`L`$ of $`E`$ is foliated in codimension $`1`$ by a $`1`$-parameter family of $`2`$-dimensional ellipsoids, each of which lies in a $`3`$-plane. By Proposition 4, it follows that the fundamental cubic of the Lawlor-Harvey examples must be reducible at each point, and thus have a symmetry of order $`2`$.
It is not difficult to see that, except when the ellipsoid is a round $`2`$-sphere, the Lawlor-Harvey examples are not special cases of either the $`\mathrm{SO}(2)`$-symmetry examples or of the $`S_3`$-symmetry examples. Thus, it follows that, at least at a generic point $`xL`$, the $`\mathrm{SO}(3)`$-stabilizer of $`C_x`$ must be isomorphic to $`_2`$.
###### Remark 8 (Joyce’s extension).
Dominic Joyce has informed<sup>8</sup><sup>8</sup>8 private communication, 3 July 2000 me that, in fact, the Lawlor-Harvey foliation result continues to hold when $`E`$ is any quadric surface in $`P`$, not necessarily an ellipsoid, or even a non-singular quadric.
###### Theorem 4.
Suppose that $`L^{\mathrm{\hspace{0.17em}3}}`$ is a connected special Lagrangian $`3`$-fold whose fundamental cubic $`C`$ is of $`_2`$-stabilizer type on an open dense subset $`L^{}L`$. Then $`L^{}`$ has a codimension $`1`$ foliation $`𝒮`$ such that each $`𝒮`$-leaf lies in a $`3`$-plane and is, moreover, a quadric surface in that $`3`$-plane. The space of maximally extended special Lagrangian $`3`$-folds of this type is finite dimensional and, in fact, coincides with the space of Lawlor-Harvey examples, as extended by Joyce.
###### Proof.
By assumption, at a generic point $`xL`$, the $`\mathrm{SO}(3)`$-stabilizer subgroup of $`C_x`$ is isomorphic to $`_2`$. Let $`L^{}L`$ be the open, dense subset where this holds. Then by Proposition 1, there exist positive functions $`r,s:L^{}`$ with $`rs`$ and a $`_2`$-subbundle $`FP_L`$ over $`L^{}`$ on which the following identity holds:
(3.31)
$$C=r\omega _1(2\omega _{1}^{}{}_{}{}^{2}3\omega _{2}^{}{}_{}{}^{2}3\omega _{3}^{}{}_{}{}^{2})+6s\omega _1\omega _2\omega _3.$$
(Of course, $`\pi :FL^{}`$ is a double cover and the reader can just think of the coframing $`\omega `$ as being well-defined on $`L^{}`$ up to the ambiguity of replacing $`\omega _2`$ and $`\omega _3`$ by $`\omega _2`$ and $`\omega _3`$.)
Consequently, on the subbundle $`F`$, the following identities hold:
(3.32)
$$\left(\begin{array}{ccc}\beta _{11}& \beta _{12}& \beta _{13}\\ \beta _{21}& \beta _{22}& \beta _{23}\\ \beta _{31}& \beta _{32}& \beta _{33}\end{array}\right)=\left(\begin{array}{ccc}2r\omega _1& s\omega _3r\omega _2& s\omega _2r\omega _3\\ s\omega _3r\omega _2& r\omega _1& s\omega _1\\ s\omega _2r\omega _3& s\omega _1& r\omega _1\end{array}\right).$$
Moreover, because $`F`$ is a $`_2`$-bundle, relations of the form
(3.33)
$$\begin{array}{cc}\hfill \alpha _{23}& =t_{11}\omega _1+t_{12}\omega _2+t_{13}\omega _3\hfill \\ \hfill \alpha _{31}& =t_{21}\omega _1+t_{22}\omega _2+t_{23}\omega _3\hfill \\ \hfill \alpha _{12}& =t_{31}\omega _1+t_{32}\omega _2+t_{33}\omega _3\hfill \end{array}$$
hold on $`F`$ for some functions $`t_{ij}`$. Moreover, for $`i=1,2,3`$ there exist functions $`r_i`$ and $`s_i`$ on $`F`$ so that
(3.34)
$$\mathrm{d}r=r_i\omega _i,\mathrm{d}s=s_i\omega _i.$$
Substituting the relations (3.32), (3.33), and (3.34) into the identities
(3.35)
$$\mathrm{d}\beta _{ij}=\beta _{ik}\alpha _{kj}\alpha _{ik}\beta _{kj}$$
and using the identities $`\mathrm{d}\omega _i=\alpha _{ij}\omega _j`$ then yields polynomial relations among these quantities that can be solved,<sup>9</sup><sup>9</sup>9 During the derivation of (3.36), one uses the assumptions that $`r`$, $`s`$ and $`r^2s^2`$ are all nonzero. leading to relations of the form
(3.36)
$$\begin{array}{cc}\hfill \mathrm{d}r& =2(s^2+2r^2)t_1\omega _1+(2rst_3+s^2t_2)\omega _2(2rst_2+s^2t_3)\omega _3,\hfill \\ \hfill \mathrm{d}s& =s\left(6rt_1\omega _1+(2st_3+rt_2)\omega _2(2st_2+rt_3)\omega _3\right),\hfill \\ \\ \hfill \alpha _{23}& =\frac{1}{2}(st_2rt_3)\omega _2+\frac{1}{2}(st_3rt_2)\omega _3\hfill \\ \hfill \alpha _{31}& =st_2\omega _1+st_1\omega _2rt_1\omega _3,\hfill \\ \hfill \alpha _{12}& =st_3\omega _1+rt_1\omega _2st_1\omega _3,\hfill \end{array}$$
where, for brevity, I have introduced the notation
$$t_1=t_{23}/r,t_2=t_{21}/s,t_3=t_{31}/s.$$
Using (3.36) to expand out the identities
$$0=\mathrm{d}(\mathrm{d}\omega _1)=\mathrm{d}(\mathrm{d}\omega _2)=\mathrm{d}(\mathrm{d}\omega _3)=\mathrm{d}(\mathrm{d}r)=\mathrm{d}(\mathrm{d}s)$$
yields relations on the exterior derivatives of $`t_1`$, $`t_2`$, and $`t_3`$. These can be expressed by the condition that there exist functions $`u_1`$, $`u_2`$, and $`u_3`$ so that the equations
(3.37)
$$\begin{array}{cc}\hfill \mathrm{d}t_1& =(su_13r3r^2t_{1}^{}{}_{}{}^{2})\omega _1\hfill \\ \hfill \mathrm{d}t_2& =3t_1(rt_2st_3)\omega _1+(u_2\frac{3}{2}rt_{2}^{}{}_{}{}^{2})\omega _2+(u_3+\frac{3}{2}st_{2}^{}{}_{}{}^{2})\omega _3,\hfill \\ \hfill \mathrm{d}t_3& =3t_1(rt_3st_2)\omega _1(u_3+\frac{3}{2}st_{3}^{}{}_{}{}^{2})\omega _2(u_2\frac{3}{2}rt_{3}^{}{}_{}{}^{2})\omega _3\hfill \end{array}$$
hold. Substituting (3.36) and (3.37) into the identities
(3.38)
$$\mathrm{d}\alpha _{ij}=\alpha _{ik}\alpha _{kj}+\beta _{ik}\beta _{kj}$$
and expanding, again using the identities $`\mathrm{d}\omega _i=\alpha _{ij}\omega _j`$, yields
(3.39)
$$\begin{array}{cc}\hfill u_2& =\frac{1}{2}\left(2rt_{1}^{}{}_{}{}^{2}+rt_{2}^{}{}_{}{}^{2}3st_2t_3+rt_{3}^{}{}_{}{}^{2}\right)rsu_1,\hfill \\ \hfill u_3& =\frac{1}{2}\left(2st_{1}^{}{}_{}{}^{2}st_{2}^{}{}_{}{}^{2}+3rt_2t_3st_{3}^{}{}_{}{}^{2}\right)+s+ru_1.\hfill \end{array}$$
Finally, expanding out the identities $`\mathrm{d}(\mathrm{d}t_1)=\mathrm{d}(\mathrm{d}t_2)=\mathrm{d}(\mathrm{d}t_3)=0`$ shows that they are equivalent to the formula
(3.40)
$$\begin{array}{cc}\hfill \mathrm{d}u_1=& 2t_1\left(3ru_1+s(t_{1}^{}{}_{}{}^{2}+2t_{2}^{}{}_{}{}^{2}+2t_{3}^{}{}_{}{}^{2})\right)\omega _1\hfill \\ & \left(u_1(rt_2+st_3)+3(rt_3+st_2)(1+t_{1}^{}{}_{}{}^{2})\right)\omega _2\hfill \\ & +\left(u_1(st_2+rt_3)3(rt_2+st_s)(1+t_{1}^{}{}_{}{}^{2})\right)\omega _3.\hfill \end{array}$$
The exterior derivative of (3.40) is an identity.
For future use, I record the formulae
(3.41)
$$\begin{array}{cc}\hfill \mathrm{d}\omega _1& =s(t_3\omega _2t_2\omega _3)\omega _1,\hfill \\ \hfill \mathrm{d}\omega _2& =t_1(r\omega _2s\omega _3)\omega _1+\frac{1}{2}(rt_3st_2)\omega _2\omega _3,\hfill \\ \hfill \mathrm{d}\omega _3& =t_1(r\omega _3s\omega _2)\omega _1+\frac{1}{2}(rt_2st_3)\omega _2\omega _3.\hfill \end{array}$$
which follow from the identities $`\mathrm{d}\omega _i=\alpha _{ij}\omega _j`$ coupled with (3.36).
At this point, it is worthwhile taking stock of what has been accomplished. Consider the system of quantities
$$\omega _1,\omega _2,\omega _3,r,s,t_1,t_2,t_3,u_1.$$
The formulae (3.41), (3.36), (3.37), and (3.40) express the exterior derivatives of these quantities as polynomials in these quantities. Moreover, the relation $`\mathrm{d}(\mathrm{d}q)=0`$ for $`q`$ any one of these quantities follows by formal expansion and use of the given exterior derivative formulae.
By a theorem<sup>10</sup><sup>10</sup>10 This was originally part of Cartan’s general theory of intransitive pseudo-groups. In more recent times, this theorem has been subsumed into the theory of *Lie algebroids*. For an introduction, the reader could try the Appendix of . of Élie Cartan, for any six constants $`\overline{r},\overline{s},\overline{t}_1,\overline{t}_2,\overline{t}_3,\overline{u}_1`$, there exists an open neighborhood $`U`$ of $`0^3`$ that is endowed with three linearly independent $`1`$-forms $`\omega _i`$ and six functions $`r,s,t_1,t_2,t_3,u_1`$ that satisfy the equations (3.41), (3.36), (3.37), and (3.40) and also satisfy
$$r(0)=\overline{r},s(0)=\overline{s},t_1(0)=\overline{t}_1,t_2(0)=\overline{t}_2,t_3(0)=\overline{t}_3,u_1(0)=\overline{u}_1.$$
Moreover these functions and forms are real-analytic and unique in a neighborhood of $`0`$, up to a real-analytic local diffeomorphism fixing $`0`$.
Now, given such a system $`(\omega ,r,s,t,u)`$ on a simply connected $`3`$-manifold $`L`$, one can set $`\eta _i=0`$, define $`\alpha _{ij}=\alpha _{ji}`$ by the last three equations of (3.36), define $`\beta _{ij}=\beta _{ji}`$ by the equations (3.32), and see that the affine structure equations
(3.42)
$$\begin{array}{cc}\hfill \mathrm{d}\omega _i& =\alpha _{ij}\omega _j+\beta _{ij}\eta _j,\hfill \\ \hfill \mathrm{d}\eta _i& =\beta _{ij}\omega _j\alpha _{ij}\eta _j,\hfill \\ \hfill \mathrm{d}\beta _{ij}& =\beta _{ik}\alpha _{kj}\alpha _{ik}\beta _{kj},\hfill \\ \hfill \mathrm{d}\alpha _{ij}& =\alpha _{ik}\alpha _{kj}+\beta _{ik}\beta _{kj}\hfill \end{array}$$
are identities. Thus, there is an immersion of $`L`$, unique up to translation and $`\mathrm{SU}(3)`$-rotation, as a special Lagrangian $`3`$-manifold in $`^{\mathrm{\hspace{0.17em}3}}`$ that induces these structure equations.
In particular, it follows that the space of germs of special Lagrangian $`3`$-manifolds in $`^{\mathrm{\hspace{0.17em}3}}`$ whose fundamental cubics are of the form (3.31) is of dimension $`6`$. Moreover, any two that agree to order $`4`$ at a single point must be equal in a neighborhood. It is not difficult to argue from this that the space one gets by reducing modulo the equivalence relation defined by analytic continuation is a $`3`$-dimensional singular space.
Now, the first of the equations (3.41) shows that the $`2`$-plane field $`\omega _1=0`$ is integrable, moreover, the structure equations found so far imply
(3.43)
$$\mathrm{d}\left(𝐞_2𝐞_3(J𝐞_1t_1𝐞_1)\right)0mod\omega _1.$$
In particular, the $`3`$-plane $`𝐞_2𝐞_3(J𝐞_1t_1𝐞_1)`$ is constant along each leaf of $`\omega _1`$ and, moreover each such leaf lies in an affine $`3`$-plane parallel to this $`3`$-plane. Thus, all of these examples are foliated in codimension $`1`$ by $`3`$-plane sections.
Moreover, an examination of the structure equations shows that the space of congruence classes of such $`3`$-plane sections is of dimension $`3`$, the same as the dimension of quadric surfaces in $`3`$-space. In fact, using the structure equations, it is not difficult to show that these $`3`$-plane sections are, in fact, quadric surfaces. For the sake of brevity, I will not include the details of this routine calculation here.
It follows that these special Lagrangian $`3`$-folds all belong to the class of Lawlor-Harvey examples, as extended by Joyce. ∎
###### Corollary 1.
Any connected special Lagrangian $`3`$-fold $`L^{\mathrm{\hspace{0.17em}3}}`$ that is foliated in codimension $`1`$ by $`3`$-plane sections is an open subset of a Lawlor-Harvey-Joyce example.
###### Proof.
By Proposition 4, any such $`L`$ must have a reducible fundamental cubic $`C`$. Thus, the $`\mathrm{SO}(3)`$-stabilizer of $`C_x`$ at each point contains a $`_2`$ and so is either isomorphic to $`\mathrm{SO}(2)`$, $`S_3`$, or $`_2`$.
If this stabilizer is isomorphic to $`\mathrm{SO}(2)`$ at a generic point, then Theorem 1 applies, showing that $`L`$ is a Lawlor-Harvey-Joyce example.
If this stabilizer is isomorphic to $`S_3`$ at a generic point, the discussion at the end of §3.4.1 shows that the only such examples that are foliated in codimension $`1`$ by $`3`$-plane sections have the property that these sections are necessarily (possibly singular) quadric surfaces, so that such an $`L`$ is, again, a Lawlor-Harvey-Joyce example.
Finally, if the stabilizer is isomorphic to $`_2`$ at a generic point, then Theorem 4 applies. ∎
###### Remark 9 (Harvey’s characterization).
In his proof of Theorem 7.78 in , Harvey gives a characterization of the Lawlor-Harvey examples that is closely related to Proposition 4. What he shows is that any special Lagrangian $`m`$-fold $`L^m`$ that meets a certain concurrent family of Lagrangian $`m`$-planes in a codimension $`1`$ foliation whose leaves are compact must belong to the family that they construct. When $`m=3`$, Corollary 1 is more general than this, since it makes no assumption about the family of Lagrangian planes that cut $`L`$ to produce the foliation and makes no assumption about compactness (or even completeness) of the leaves.
Of course, one expects that the higher dimensional analog of Corollary 1 holds, i.e., that any connected special Lagrangian $`m`$-fold $`L^m`$ that is foliated in codimension $`1`$ by $`m`$-plane sections is necessarily an open subset of a Lawlor-Harvey-Joyce example. I have not tried to prove this, but it should be straightforward.
### 3.6. $`_3`$ symmetry
Now consider those special Lagrangian submanifolds $`L^{\mathrm{\hspace{0.17em}3}}`$ whose cubic second fundamental form has an $`_3`$-symmetry at each point.
###### Example 6.
Let $`\mathrm{\Sigma }S^5`$ be a surface such that the cone on $`\mathrm{\Sigma }`$ is special Lagrangian, and consider the $`3`$-fold
$$L_\mathrm{\Sigma }=\{(s+\mathrm{i}t)𝐮𝐮\mathrm{\Sigma },t^33s^2t=c\},$$
where $`c`$ is a (real) constant. This $`L_\mathrm{\Sigma }`$ is special Lagrangian. For example, see , where a more general result for special Lagrangian cones in $`^n`$ is proved.
Note that $`L_\mathrm{\Sigma }`$ is diffeomorphic to the disjoint union of three copies of $`\times \mathrm{\Sigma }`$. In fact, each ‘end’ of each component of $`L_\mathrm{\Sigma }`$ is asymptotic to the cone on $`\lambda \mathrm{\Sigma }S^5`$ for some $`\lambda `$ satisfying $`\lambda ^6=1`$.
When $`\mathrm{\Sigma }`$ is not totally geodesic in $`S^5`$ the $`\mathrm{SO}(3)`$-stabilizer of the fundamental cubic at a generic point of point of $`L_\mathrm{\Sigma }`$ is isomorphic to $`_3`$.
###### Theorem 5.
If $`L^{\mathrm{\hspace{0.17em}3}}`$ is a connected special Lagrangian submanifold whose fundamental cubic has $`_3`$-symmetry at each point of a dense open subset of $`L`$, then $`L`$ contains a dense open set $`L^{}`$ such that every point of $`L^{}`$ has an open neighborhood in $`L`$ that is an open subset of one of the special Lagrangian $`3`$-folds of Example 6.
###### Proof.
Let $`L^{\mathrm{\hspace{0.17em}3}}`$ satisfy the hypotheses of the theorem. The locus of points $`xL`$ for which the $`\mathrm{SO}(3)`$-stabilizer of $`C`$ is larger than $`_3`$ is a proper real-analytic subset of $`L`$, so its complement $`L^{}`$ is open and dense in $`L`$. Thus, I can, without loss of generality, replace $`L`$ by a component of $`L^{}`$. In other words, I can assume that the $`\mathrm{SO}(3)`$-stabilizer of $`C_x`$ is isomorphic to $`_3`$ for all $`xL`$.
By Proposition 1, since the stabilizer of $`C_x`$ is $`_3`$ for all $`xL`$, there are positive (real-analytic) functions $`r`$ and $`s`$ on $`L`$ with the property that the equation
(3.44)
$$C=r\omega _1\left(2\omega _{1}^{}{}_{}{}^{2}3\omega _{2}^{}{}_{}{}^{2}3\omega _{3}^{}{}_{}{}^{2}\right)+s\left(\omega _{2}^{}{}_{}{}^{3}3\omega _2\omega _{3}^{}{}_{}{}^{2}\right)$$
defines a $`_3`$-subbundle $`FP_L`$ of the adapted coframe bundle $`P_LL`$. Moreover, the expression $`sr\sqrt{2}`$ is nowhere vanishing on $`L`$.
Now, on the subbundle $`F`$, the following identities hold:
(3.45)
$$\left(\begin{array}{ccc}\beta _{11}& \beta _{12}& \beta _{13}\\ \beta _{21}& \beta _{22}& \beta _{23}\\ \beta _{31}& \beta _{32}& \beta _{33}\end{array}\right)=\left(\begin{array}{ccc}2r\omega _1& r\omega _2& r\omega _3\\ r\omega _2& r\omega _1+s\omega _2& s\omega _3\\ r\omega _3& s\omega _3& r\omega _1s\omega _2\end{array}\right).$$
Moreover, because $`F`$ is a $`_3`$-bundle, relations of the form
(3.46)
$$\begin{array}{cc}\hfill \alpha _{23}& =t_{11}\omega _1+t_{12}\omega _2+t_{13}\omega _3\hfill \\ \hfill \alpha _{31}& =t_{21}\omega _1+t_{22}\omega _2+t_{23}\omega _3\hfill \\ \hfill \alpha _{12}& =t_{31}\omega _1+t_{32}\omega _2+t_{33}\omega _3\hfill \end{array}$$
hold on $`F`$ for some functions $`t_{ij}`$. Moreover, for $`i=1,2,3`$ there exist functions $`r_i`$ and $`s_i`$ on $`F`$ so that
(3.47)
$$\mathrm{d}r=r_i\omega _i,\mathrm{d}s=s_i\omega _i.$$
Substituting the relations (3.45), (3.46), and (3.47) into the identities
(3.48)
$$\mathrm{d}\beta _{ij}=\beta _{ik}\alpha _{kj}\alpha _{ik}\beta _{kj}$$
and using the identities $`\mathrm{d}\omega _i=\alpha _{ij}\omega _j`$ then yields polynomial relations among these quantities that can be solved,<sup>11</sup><sup>11</sup>11 During the derivation of (3.49), one uses the assumptions that $`r`$ and $`s`$ are nonzero but not the assumption that $`sr\sqrt{2}`$ is nonzero. leading to relations of the form
(3.49)
$$\begin{array}{cc}\hfill \mathrm{d}r& =4rt_1\omega _1,\hfill \\ \hfill \mathrm{d}s& =s\left(t_1\omega _1+3t_3\omega _23t_2\omega _3\right),\hfill \\ \\ \hfill \alpha _{23}& =t_2\omega _2t_3\omega _3\hfill \\ \hfill \alpha _{31}& =t_1\omega _3,\hfill \\ \hfill \alpha _{12}& =t_1\omega _2,\hfill \end{array}$$
where, for brevity, I have introduced the notation
$$t_1=t_{23},t_2=t_{12},t_3=t_{13}.$$
Using (3.49) to expand out the identities
$$0=\mathrm{d}(\mathrm{d}\omega _1)=\mathrm{d}(\mathrm{d}\omega _2)=\mathrm{d}(\mathrm{d}\omega _3)=\mathrm{d}(\mathrm{d}r)=\mathrm{d}(\mathrm{d}s)$$
and also the identities
$$\mathrm{d}\alpha _{ij}=\alpha _{ik}\alpha _{kj}+\beta _{ik}\beta _{kj}$$
yields relations on the exterior derivatives of $`t_1`$, $`t_2`$, and $`t_3`$. When these are solved, one finds that there are functions $`u_2`$ and $`u_3`$ so that the equations
(3.50)
$$\begin{array}{cc}\hfill \mathrm{d}t_1& =(3r^2t_{1}^{}{}_{}{}^{2})\omega _1\hfill \\ \hfill \mathrm{d}t_2& =t_1t_2\omega _1+u_2\omega _2+(u_3+v)\omega _3,\hfill \\ \hfill \mathrm{d}t_3& =t_1t_3\omega _1u_2\omega _3+(u_3v)\omega _2\hfill \end{array}$$
hold where
$$v=s^2\frac{1}{2}(r^2+t_{1}^{}{}_{}{}^{2}+t_{2}^{}{}_{}{}^{2}+t_{3}^{}{}_{}{}^{2}).$$
Observe that, if one sets $`r=0`$ in the current structure equations, then these become, up to a trivial change of notation, the same structure equations as those for the special Lagrangian cones discussed in §3.4. This is a first hint that these examples must be related to the special Lagrangian cones.<sup>12</sup><sup>12</sup>12 Also, if one now computes the Cartan characters of the naïve exterior differential system that models these structure equations, one finds that $`s_1=2`$ while $`s_2=s_3=0`$ and that this exterior differential system is involutive. The characteristic variety is a pair of complex conjugate points, each of multiplicity $`1`$.
The next observation is that the structure equations
(3.51)
$$\mathrm{d}r=4rt_1\omega _1,\text{and}\mathrm{d}t_1=(3r^2t_{1}^{}{}_{}{}^{2})\omega _1$$
are identical (after replacing $`t_1`$ by $`t`$) to the last two equations of (3.10). In particular, there must exist a constant $`c>0`$ and a function $`\theta `$ on $`L`$ satisfying the bound $`|\theta |<\pi /6`$ so that
(3.52)
$$r^{3/4}=c^{3/4}\mathrm{cos}3\theta ,r^{1/4}t_1=c^{3/4}\mathrm{sin}3\theta .$$
It then follows from (3.51) that
(3.53)
$$\omega _1=\frac{\mathrm{d}\theta }{c(\mathrm{cos}3\theta )^{4/3}}.$$
By dilation in $`^{\mathrm{\hspace{0.17em}3}}`$, one can reduce to the case $`c=1`$, so assume this from now on.
Consider the following expressions:
(3.54)
$$\begin{array}{cc}\hfill p& =r^{1/4}s,\hfill \\ \hfill q_2& =r^{1/4}t_2,q_3=r^{1/4}t_3,\hfill \\ \hfill v_2& =r^{1/2}u_2,v_3=r^{1/2}u_3,\hfill \\ \hfill \eta _2& =r^{1/4}\omega _2,\eta _3=r^{1/4}\omega _3.\hfill \end{array}$$
The structure equations derived above show that
(3.55)
$$\begin{array}{cc}\hfill \mathrm{d}\eta _2& =q_2\eta _2\eta _3\hfill \\ \hfill \mathrm{d}\eta _3& =q_3\eta _2\eta _3\hfill \\ \hfill \mathrm{d}p& =3p(q_3\eta _2q_2\eta _3)\hfill \\ \hfill \mathrm{d}q_2& =v_2\eta _2+(v_3+w)\eta _3\hfill \\ \hfill \mathrm{d}q_3& =v_2\eta _3+(v_3w)\eta _2\hfill \end{array}$$
where $`w=\frac{1}{2}(1+q_{2}^{}{}_{}{}^{2}+q_{3}^{}{}_{}{}^{2})p^2`$.
In particular, $`d(p^{1/3}\eta _2)=d(p^{1/3}\eta _3)=0`$. Let $`xL`$ be fixed and let $`UL`$ be an $`x`$-neighborhood on which there exist functions $`y_2`$ and $`y_3`$ vanishing at $`x`$ that satisfy $`p^{1/3}\eta _2=\mathrm{d}y_2`$ and $`p^{1/3}\eta _3=\mathrm{d}y_3`$. Then the functions $`(\theta ,y_2,y_3)`$ are independent on $`U`$ and, by shrinking $`U`$ if necessary, I can assume that $`(\theta ,y_2,y_3)(U)^{\mathrm{\hspace{0.17em}3}}`$ is a product open set of the form $`I\times D`$ where $`I(\frac{\pi }{6},\frac{\pi }{6})`$ is a connected interval and $`D^2`$ is a disc centered on the origin. Of course, the functions $`p`$, $`q_2`$, $`q_3`$, $`v_2`$ and $`v_3`$ can be regarded as functions on $`D`$, since their differentials are linear combinations of $`\mathrm{d}y_2`$ and $`\mathrm{d}y_3`$. In fact, these functions and forms can now be regarded as defined on the open set $`(\frac{\pi }{6},\frac{\pi }{6})\times D`$ by simply reading the formulae above backwards. Thus, for example
$$s=r^{1/4}p=(\mathrm{cos}3\theta )^{1/4}p$$
and so forth. This gives quantities $`\omega _i`$, $`r`$, $`s`$, $`t_i`$, and $`u_i`$ that are well-defined on all of $`(\frac{\pi }{6},\frac{\pi }{6})\times D`$ and that satisfy the originally derived structure equations. It follows that there is an immersion of $`(\frac{\pi }{6},\frac{\pi }{6})\times D`$ into $`^{\mathrm{\hspace{0.17em}3}}`$ as a special Lagrangian $`3`$-fold that extends $`U`$ and pulls back the constructed forms and quantities to agree with the given ones on $`U`$. The chief difference is that each of the $`\theta `$-curves in $`(\frac{\pi }{6},\frac{\pi }{6})\times D`$ is mapped to a complete curve in $`^{\mathrm{\hspace{0.17em}3}}`$.
Next, observe that the equations
(3.56)
$$\begin{array}{cc}\hfill \mathrm{d}𝐱& 𝐞_1\omega _1\hfill \\ \hfill \mathrm{d}𝐞_1& J𝐞_1(2r\omega _1)\hfill \\ \hfill \mathrm{d}(J𝐞_1)& 𝐞_1(2r\omega _1)\hfill \end{array}\}mod\omega _2,\omega _3,$$
which are identical to the corresponding equations in §3.2, then show that the leaves of the curve foliation defined by $`\omega _2=\omega _3=0`$ are congruent to the leaves of the corresponding foliation by the $`𝐞_1`$-curves in §3.2.
Finally, note that, setting $`\theta =0`$ (i.e., $`t_1=0`$ and $`r=1`$ in the above structure equations on $`(\frac{\pi }{6},\frac{\pi }{6})\times D`$ gives an immersion of $`D`$ into $`S^5^{\mathrm{\hspace{0.17em}3}}`$ with the property that the cone on the image $`\mathrm{\Sigma }`$ is a special Lagrangian $`3`$-fold. Because the $`\theta `$-curves meet this surface orthogonally, it follows easily that the image of $`(\frac{\pi }{6},\frac{\pi }{6})\times D`$ is exactly $`L_\mathrm{\Sigma }`$ as described in Example 6. Further details are left to the reader. ∎
### 3.7. The ruled family
In this last subsection, I am going to consider the generality of the set of ruled special Lagrangian $`3`$-manifolds.
Examples of ruled special Lagrangian $`3`$-folds in $`^{\mathrm{\hspace{0.17em}3}}`$ were constructed in Harvey and Lawson’s original paper . These included products, special Lagrangian cones, and conormal bundles of minimal surfaces in $`^{\mathrm{\hspace{0.17em}3}}`$. All of these families depend on two functions of one variable in the sense of exterior differential systems.
Harvey and Lawson also showed in \[12, Theorems 4.9, 4.13\] how one could deform the conormal bundle of a minimal surface in $`^{\mathrm{\hspace{0.17em}3}}`$ according to the data of a harmonic function on such a surface and obtain more general ruled special Lagrangian $`3`$-folds. (Borisenko later gave a somewhat different description of the same family.) These examples depend on four functions of one variable in the sense of exterior differential systems.
On the other hand, the construction in Example 4 of twisted special Lagrangian cones provides another family of examples of ruled special Lagrangian $`3`$-folds, again depending on four functions of one variable in the sense of exterior differential systems. It is easy to see that this family is distinct from the family described in \[12, Theorems 4.9\].
In this section, I am going to show that the ruled special Lagrangian $`3`$-folds depend on *six* functions of one variable in the sense of exterior differential systems. Thus, the ‘explicit’ families that have been constructed so far are only a small part of the complete family. For a different description of ruled special Lagrangian $`3`$-folds, one should consult Joyce’s recent article .
#### 3.7.1. Almost CR-structures and Levi-flatness
For the description I plan to give of the ruled special Lagrangian submanifold of $`^{\mathrm{\hspace{0.17em}3}}`$, I will need some facts about a generalized notion of ‘pseudo-holomorphic curves’.
Recall that an *almost CR-structure* on a manifold $`M`$ is a subbundle $`ETM`$ of even dimension equipped with a complex structure map $`J:EE`$. The *rank* of the CR-structure is the rank of $`E`$ as a complex bundle and the *codimension* of the CR-structure is the rank of the quotient bundle $`TM/E`$. A (real) curve $`CM`$ is said to be an *$`E`$-curve* if its tangent line at each point lies in $`E`$. A (real) surface $`SM`$ is said to be *$`E`$-holomorphic* if its tangent plane at each point is a complex line in $`E`$. (In order to avoid confusion, I will not adopt the standard practice of calling these surfaces ‘pseudo-holomorphic curves’, or, indeed, curves of any kind.)
An almost CR-structure $`(E,J)`$ will be said to be *Levi-flat* if, for any $`1`$-form $`\rho `$ on $`M`$ that vanishes on $`E`$, the $`2`$-form $`\mathrm{d}\rho `$ vanishes on all the $`2`$-planes that are complex lines in $`E`$. Note that Levi-flatness is automatic when the codimension of the CR-structure is zero and that Levi-flatness generally has no implications about the ‘integrability’ of the almost CR-structure to a CR-structure, which is a different condition altogether.
###### Proposition 5.
Let $`(E,J)`$ be a real-analytic, Levi-flat, almost CR-structure on $`M`$ and let $`CM`$ be a real-analytic $`E`$-curve. Then there is an $`E`$-holomorphic surface $`SM`$ that contains $`C`$. This surface is locally unique in the sense that, for any two such surfaces $`S_1`$ and $`S_2`$, the intersection $`S_1S_2`$ is also an $`E`$-holomorphic surface that contains $`C`$.
###### Proof.
This is a straightforward application of the Cartan-Kähler Theorem \[2, Chapter III\] so I will only give the barest details. This is a local result, so it suffices to give a local proof.
Let $`r`$ be the rank of $`(E,J)`$ and let $`q`$ be its codimension. For any point $`xM`$, there is an open $`x`$-neighborhood $`UM`$ on which there exist real-analytic $`1`$-forms $`\theta _1,\mathrm{},\theta _q`$ with real values and $`\omega _1,\mathrm{},\omega _r`$ with complex values with the property that the equations $`\theta _1=\mathrm{}=\theta _q=0`$ define the restriction of $`E`$ to $`U`$ and with the property that $`\omega _1,\mathrm{},\omega _r`$ are complex linear on $`E`$ and are linearly independent over $``$ at each point of $`U`$. There are identities of the form
$$\mathrm{d}\theta _\alpha K_{\alpha ij}\omega _i\omega _j+L_{\alpha ij}\omega _i\overline{\omega _j}+\overline{K_{\alpha ij}}\overline{\omega _i}\overline{\omega _j}mod\theta _1,\mathrm{},\theta _q.$$
The hypothesis of Levi-flatness is simply that the functions $`L_{\alpha ij}`$ all vanish identically. Under this hypothesis, the real-analytic exterior differential system $``$ generated algebraically by the $`\theta _\alpha `$ and the real and imaginary parts of the $`2`$-forms $`\omega _i\omega _j`$ is involutive and each of the $`1`$-dimensional integral elements is regular and lies in a unique $`2`$-dimensional integral element. Now apply the Cartan-Kähler theorem. ∎
#### 3.7.2. Oriented lines
Since a ruled $`3`$-manifold in $`^{\mathrm{\hspace{0.17em}3}}`$ can be regarded as a surface in the space of lines in $`^{\mathrm{\hspace{0.17em}3}}`$, it is useful to consider the geometry of this space. It is slightly more convenient to consider the space $`\mathrm{\Lambda }`$ of oriented lines in $`^{\mathrm{\hspace{0.17em}3}}`$, so I will do this.
The space $`\mathrm{\Lambda }`$ is naturally diffeomorphic to the tangent bundle of $`S^5`$. Explicitly, the pair $`(𝐮,𝐯)TS^5`$ consisting of a unit vector $`𝐮S^5`$ and a vector $`𝐯𝐮^{}`$ corresponds to the oriented line with oriented direction $`𝐮`$ that passes through $`𝐯`$. Naturally, I will regard $`𝐮:\mathrm{\Lambda }S^5`$ and $`𝐯:\mathrm{\Lambda }^{\mathrm{\hspace{0.17em}3}}`$ as vector-valued functions on $`\mathrm{\Lambda }`$.
Thus, a curve $`\gamma :(a,b)\mathrm{\Lambda }`$ can be written as $`\gamma (s)=(𝐮(s),𝐯(s))`$ where the curve $`𝐮:(a,b)S^5`$ and the curve $`𝐯:(a,b)^{\mathrm{\hspace{0.17em}3}}`$ satisfy $`𝐮(s)𝐯(s)=0`$ for all $`s(a,b)`$. Such a curve gives rise to a mapping $`\mathrm{\Gamma }:(a,b)\times ^{\mathrm{\hspace{0.17em}3}}`$ by the formula
$$\mathrm{\Gamma }(s,t)=𝐯(s)+t𝐮(s).$$
Assuming that $`\gamma `$ is smooth (resp., real-analytic) then $`\mathrm{\Gamma }`$ is also smooth (resp., real-analytic) and $`\mathrm{\Gamma }`$ will be an immersion except on the locus consisting of those $`(s,t)(a,b)\times `$ where $`\left(𝐯^{}(s)+t𝐮^{}(s)\right)𝐮(s)=0`$. On the locus where it is an immersion, the image of $`\mathrm{\Gamma }`$ is then a ruled surface in $`^{\mathrm{\hspace{0.17em}3}}`$.
More generally, given any smooth (resp., real-analytic) map $`\gamma :P\mathrm{\Lambda }`$ where $`P`$ is a smooth (resp., real-analytic) manifold there is an induced smooth (resp., real-analytic) map $`\mathrm{\Gamma }:P\times ^{\mathrm{\hspace{0.17em}3}}`$ defined by the same formula as above. With the appropriate ‘generic’ assumptions on $`\gamma `$, the mapping $`\mathrm{\Gamma }`$ will be an immersion on some open subset of $`P\times `$ and its image will be a ruled immersion.
There are two natural differential forms on $`\mathrm{\Lambda }`$ that are invariant under the complex isometries of $`^{\mathrm{\hspace{0.17em}3}}`$. These are the pair of $`1`$-forms
$$\theta =J𝐮\mathrm{d}𝐮,\text{and}\tau =J𝐮\mathrm{d}𝐯.$$
It is easy to see that $`\theta `$ and $`\tau `$ are linearly independent, so their common kernel $`ET\mathrm{\Lambda }`$ is a bundle of rank $`8`$. The significance of these two 1-forms is revealed in the following result.
###### Proposition 6.
A curve $`\gamma :(a,b)\mathrm{\Lambda }`$ is tangent to $`E`$ everywhere if and only if the corresponding ruled ‘surface’ $`\mathrm{\Gamma }:(a,b)\times ^{\mathrm{\hspace{0.17em}3}}`$ is $`\omega `$-isotropic.
###### Proof.
This is immediate from the formulae for $`\mathrm{\Gamma }`$ and $`\omega `$. ∎
###### Theorem 6.
There is a complex structure $`J`$ on $`ET\mathrm{\Lambda }`$ with the properties
1. $`(E,J)`$ is a real-analytic, Levi-flat almost CR-structure on $`\mathrm{\Lambda }`$ that is invariant under the complex isometries of $`^{\mathrm{\hspace{0.17em}3}}`$.
2. Any ruled special Lagrangian $`3`$-fold $`L`$ is locally the image of the $`\mathrm{\Gamma }`$ associated to an $`E`$-holomorphic surface $`\gamma :S\mathrm{\Lambda }`$. When $`L`$ is not a $`3`$-plane, this local representation is either unique or admits at most one other such representation.
3. For each $`E`$-holomorphic surface $`\gamma :S\mathrm{\Lambda }`$, the corresponding map $`\mathrm{\Gamma }:S\times ^{\mathrm{\hspace{0.17em}3}}`$ is ruled and a special Lagrangian immersion on a dense open subset of $`S\times `$.
4. Any non-planar special Lagrangian $`3`$-fold $`L`$ that has two distinct rulings is a Lawlor-Harvey-Joyce example for which the $`2`$-dimensional $`3`$-plane sections are quadrics that are doubly ruled.
Before going on to the proof of this result, let me state some immediate corollaries:
###### Corollary 2.
A connected special Lagrangian $`3`$-fold $`L^{\mathrm{\hspace{0.17em}3}}`$ is ruled if and only if the set $`\mathrm{\Lambda }_L`$ of lines that intersect $`L`$ in nontrivial open intervals (which is an analytic subset of $`\mathrm{\Lambda }`$) has dimension at least $`1`$.
###### Proof.
I will only sketch the proof, since the details are straightforward. First, the easy direction: If $`L`$ is ruled, then the analytic set $`\mathrm{\Lambda }_L`$ must have dimension $`2`$ at least.
Conversely, if the dimension of $`\mathrm{\Lambda }_L`$ is at least $`1`$, then it contains an immersed analytic arc $`\gamma :(a,b)\mathrm{\Lambda }`$, which generates a ruled surface $`\mathrm{\Gamma }(D)L`$ for some appropriate domain $`D(a,b)\times `$. The surface $`\mathrm{\Gamma }(D)`$ must be $`\omega `$-isotropic since $`L`$ is Lagrangian. Thus, the arc $`\gamma `$ must be an $`E`$-curve. By Item $`1`$ of Theorem 6 and Proposition 5, this arc lies in an $`E`$-holomorphic surface $`\psi :S\mathrm{\Lambda }`$. By Item $`3`$ of Theorem 6, there is a dense open region $`RS\times `$ so that $`\mathrm{\Psi }(R)`$ is an immersed ruled special Lagrangian $`3`$-fold. It is not hard to see that this $`\mathrm{\Psi }(R)`$ contains at least an open subset of $`\mathrm{\Gamma }(D)`$. Since by Harvey and Lawson’s Theorem 5.5, the real-analytic $`\omega `$-isotropic surface $`\mathrm{\Gamma }(D)`$ lies in a locally unique special Lagrangian $`3`$-fold, it follows that $`\mathrm{\Psi }(R)`$ and $`L`$ must intersect in an open set. Thus $`L`$ is ruled on an open set. By real-analyticity and connectedness, it must be ruled everywhere. ∎
###### Corollary 3.
The ruled special Lagrangian $`3`$-folds in $`^{\mathrm{\hspace{0.17em}3}}`$ depend on six functions of one variable.
###### Proof.
Combine Theorem 6 and Proposition 5. ∎
###### Remark 10 (The characteristic variety).
The characteristic variety of this system turns out to be a pair of complex conjugate points, each of multiplicity $`3`$. This is particularly interesting for the following reason: The condition that the fundamental cubic at each point be singular is a single equation of second order on the special Lagrangian $`3`$-fold. Now, as usual, a Lagrangian manifold can be written as a gradient graph of a potential function, in which case, the special Lagrangian condition is a single second order elliptic equation for the potential. Then the condition that the fundamental cubic be singular is a single *third* order equation for the potential. By the general theory, the characteristic variety of a system consisting of a single elliptic second order equation and a single third order equation consists of at most six points by Bezout’s Theorem. Remarkably, the ‘singular cubic’ system turns out to have such a ‘maximal’ characteristic variety and to be involutive.
This must be quite rare. In fact, so far, I have been unable to find another example of a single pointwise equation on the second fundamental form that is involutive and has six points in its characteristic variety.
Now for the proof of Theorem 6.
###### Proof.
First, I will define the almost CR-structure on $`\mathrm{\Lambda }`$ and show that it is Levi-flat. Consider the mapping $`\lambda :F\mathrm{\Lambda }`$ that sends the coframe $`u:T_x^{\mathrm{\hspace{0.17em}3}}`$ to the oriented line spanned by $`𝐞_1(u)`$ that passes through $`x`$. Since the structure equations give
(3.57)
$$\begin{array}{cc}\hfill \mathrm{d}𝐱_{}& 𝐞_2\omega _2+𝐞_3\omega _3+J𝐞_1\eta _1+J𝐞_2\eta _2+J𝐞_3\eta _3\hfill \\ \hfill \mathrm{d}𝐞_1& 𝐞_2\alpha _{21}+𝐞_3\alpha _{31}+J𝐞_1\beta _{11}+J𝐞_2\beta _{21}+J𝐞_3\beta _{31}\hfill \end{array}\}mod𝐞_1,$$
it follows that the ten $`1`$-forms that appear on the right-hand side of this equation are $`\lambda `$-semibasic and it is evident that $`\lambda ^{}(\theta )=\beta _{11}`$ while $`\lambda ^{}(\tau )=\eta _1`$. The fibers of $`\lambda `$ are cosets of the subgroup of the motion group that fixes an oriented line in $`^{\mathrm{\hspace{0.17em}3}}`$ and hence are diffeomorphic to $`\times \mathrm{SU}(2)`$. In particular, they are connected.
Define complex-valued $`1`$-forms on $`F`$ by
$$\zeta _1=\omega _2+\mathrm{i}\omega _3,\zeta _2=\eta _2\mathrm{i}\eta _3,\zeta _3=\alpha _{21}+\mathrm{i}\alpha _{31},\zeta _4=\beta _{21}\mathrm{i}\beta _{31}.$$
These forms are $`\lambda `$-semibasic and satisfy the equations
$$\mathrm{d}\zeta _1\mathrm{}\mathrm{d}\zeta _40mod\beta _{11},\eta _1,\zeta _1,\mathrm{},\zeta _4,$$
while
(3.58)
$$\begin{array}{cc}\hfill \mathrm{d}\beta _{11}& \zeta _3\zeta _4+\overline{\zeta _3}\overline{\zeta _4}\hfill \\ \hfill 2\mathrm{d}\eta _1& \zeta _1\zeta _4\zeta _2\zeta _3+\overline{\zeta _1}\overline{\zeta _4}\overline{\zeta _2}\overline{\zeta _3}\hfill \end{array}\}mod\beta _{11},\eta _1.$$
Since the fibers of $`\lambda `$ are connected, it follows that there is a (unique) complex structure $`J:EE`$ so that the complex-valued $`1`$-forms on $`\mathrm{\Lambda }`$ that are $``$-linear on $`E`$ pull back to be linear combinations of $`\beta _{11},\eta _1,\zeta _1,\mathrm{},\zeta _4`$. Moreover, the equations (3.58) imply that the almost CR-structure $`(E,J)`$ is Levi-flat, as promised. This structure is clearly real-analytic since it is homogeneous under the action of the complex isometry group on $`\mathrm{\Lambda }`$. (Note also, by the way, that the equations (3.58) also imply that this almost CR-structure is not integrable.) This completes the proof of Item $`1`$.
Now suppose that $`L^{\mathrm{\hspace{0.17em}3}}`$ is a ruled special Lagrangian $`3`$-fold that is not a $`3`$-plane. Then, on a dense open set, this ruling can be chosen to be real-analytic and smooth. Consider the subbundle $`F_L`$ of the adapted frame bundle over $`L`$ that has $`𝐞_1`$ tangent to the ruling direction. Thus, the curves in $`L`$ defined by the differential equations $`\omega _2=\omega _3=0`$ are straight lines and, of course, $`𝐞_1`$ is tangent to these straight lines. It follows that $`\mathrm{d}𝐞_10mod\omega _2,\omega _3`$. (In fact, this is necessary and sufficient that the $`𝐞_1`$-integral curves be straight lines in $`^{\mathrm{\hspace{0.17em}3}}`$.) Since
$$\mathrm{d}𝐞_1=𝐞_2\alpha _{21}+𝐞_3\alpha _{31}+J𝐞_1\beta _{11}+J𝐞_2\beta _{21}+J𝐞_3\beta _{31},$$
it follows, in particular, that $`\beta _{11}\beta _{21}\beta _{31}0mod\omega _2,\omega _3`$. Since $`\beta _{ij}=h_{ijk}\omega _k`$, it follows from this that $`h_{11j}=0`$ for $`j=1`$, $`2`$, and $`3`$. In particular, the fundamental cubic
$$C=h_{ijk}\omega _i\omega _j\omega _k$$
is linear in the direction $`\omega _1`$. Of course, by Proposition 3, it follows that, at points where $`C`$ is non-zero, it is linear in at most three directions. Moreover, by Proposition 3 and Theorem 2, there is no non-planar special Lagrangian $`3`$-fold whose cubic is linear in three directions. Thus, either $`C`$ is linear in exactly two directions on a dense open set, or else it is linear in exactly one direction on a dense open set.
If $`C`$ is linear in exactly two directions on a dense open set, then, again by Proposition 3, it follows that $`C`$ is reducible at every point and, on a dense open set, cannot have an $`\mathrm{SO}(3)`$-stabilizer isomorphic to $`S_3`$, since these are not linear in two distinct variables. It follows that the $`\mathrm{SO}(3)`$-stabilizer at a generic point is $`_2`$, so that, by Theorem 4, $`L`$ must be one of the Lawlor-Harvey-Joyce examples. Moreover, the two linearizing directions, since they represent singular points of the projectivized cubic curve, must lie on the linear factor of $`C`$. Thus, the two possible rulings must lie in the $`2`$-dimensional slices by $`3`$-planes. Of course, this can only happen if the quadrics that are these slices are doubly ruled. Conversely, if the quadrics that are these slices are doubly ruled, then, obviously, $`L`$ must be doubly ruled as well. This establishes Item $`4`$ (as well as the fact that a non-planar special Lagrangian $`3`$-fold cannot be triply ruled).
At any rate, note that $`\beta _{11}=h_{11j}\omega _j=0`$ and that
$`\beta _{21}`$ $`=h_{212}\omega _2+h_{213}\omega _3,`$
$`\beta _{31}`$ $`=h_{312}\omega _2+h_{313}\omega _3=h_{213}\omega _2h_{212}\omega _3,`$
where I have used the symmetry and trace conditions on $`h_{ijk}`$ together with the condition $`h_{111}=0`$. It follows that
$$\left(\beta _{21}\mathrm{i}\beta _{31}\right)\left(\omega _2+\mathrm{i}\omega _3\right)=0.$$
There are now two cases to deal with. Either $`\beta _{21}`$ and $`\beta _{31}`$ vanish identically or they do not.
Suppose first that $`\beta _{21}\beta _{31}0`$. In this case, one can, after restricting to a dense open set, adapt frames so that the fundamental cubic has the form
$$C=h_{222}\left(\omega _{2}^{}{}_{}{}^{2}3\omega _2\omega _{3}^{}{}_{}{}^{2}\right),$$
where $`h_{222}>0`$. In particular, the $`\mathrm{SO}(3)`$-stabilizer of $`C`$ at the generic point is $`S_3`$. Set $`s=h_{222}`$, so that the notation agrees with the notation established in §3.4. Looking back at the structure equations from that section, one sees that
$$\alpha _{21}+\mathrm{i}\alpha _{31}=3(r_3+\mathrm{i}r_2)(\omega _2+\mathrm{i}\omega _3),$$
In particular, $`(\alpha _{21}+\mathrm{i}\alpha _{31})(\omega _2+\mathrm{i}\omega _3)=0`$. Since it has already been established that, in this case,
$$\beta _{11}=\eta _1=\beta _{21}+\mathrm{i}\beta _{31}=\eta _2+\mathrm{i}\eta _3=0,$$
it follows immediately that the natural map from the frame bundle to $`\mathrm{\Lambda }`$ that sends a coframe $`uF_L`$ to $`(𝐱(u),𝐞_1(u))`$ maps the coframe bundle into an $`E`$-holomorphic surface and that this surface is simply the space of lines of the ruling.
Now suppose that $`\beta _{21}`$ and $`\beta _{31}`$ do not vanish identically. Then, by restricting to the dense open set where they are not simultaneously zero, we can reduce frames to arrange that $`h_{212}=0`$, but that $`h_{312}0`$. In fact, there will exist functions $`r0`$, $`s`$, and $`t`$ so that
$$C=6r\omega _1\omega _2\omega _3+s\left(\omega _{2}^{}{}_{}{}^{2}3\omega _2\omega _{3}^{}{}_{}{}^{2}\right)+t\left(3\omega _{2}^{}{}_{}{}^{2}\omega _3\omega _{3}^{}{}_{}{}^{2}\right).$$
This reduces the frames to a finite ambiguity, but I will not worry about this, since it does not impose any essential difficulty. Of course, $`s`$ and $`t`$ cannot vanish identically by Theorem 2. In particular, on this adapted bundle, the following formulae hold:
$$\left(\begin{array}{ccc}\beta _{11}& \beta _{12}& \beta _{13}\\ \beta _{21}& \beta _{22}& \beta _{23}\\ \beta _{31}& \beta _{32}& \beta _{33}\end{array}\right)=\left(\begin{array}{ccc}0& r\omega _3& r\omega _2\\ r\omega _3& s\omega _2+t\omega _3& r\omega _1s\omega _3+t\omega _2\\ r\omega _2& r\omega _1s\omega _3+t\omega _2& s\omega _2t\omega _3\end{array}\right).$$
Now, there are functions $`p_{ij}`$ and $`r_i`$, $`s_i`$, and $`t_i`$ so that
$`dr`$ $`=r_1\omega _1+r_2\omega _2+r_3\omega _3,`$
$`ds`$ $`=s_1\omega _1+s_2\omega _2+s_3\omega _3,`$
$`dt`$ $`=t_1\omega _1+t_2\omega _2+t_3\omega _3,`$
$`\alpha _{32}`$ $`=p_{11}\omega _1+p_{12}\omega _2+p_{13}\omega _3,`$
$`\alpha _{13}`$ $`=p_{21}\omega _1+p_{22}\omega _2+p_{23}\omega _3,`$
$`\alpha _{21}`$ $`=p_{31}\omega _1+p_{32}\omega _2+p_{33}\omega _3.`$
Just as in previous cases of moving frame analyses, substituting these equations into the structure equations for $`\mathrm{d}\beta _{ij}`$ yields 15 equations on these 18 quantities. I will not give the whole solution, since that is not needed for this argument, but will merely note that these equations imply $`p_{21}=p_{31}=0`$ and that $`p_{22}=p_{33}`$ while $`p_{23}+p_{32}=0`$. In particular, this implies
$$(\alpha _{21}+\mathrm{i}\alpha _{31})(\omega _2+\mathrm{i}\omega _3)=0,$$
just as in the first case. Moreover, since $`\beta _{21}=r\omega _3`$ and $`\beta _{31}=r\omega _2`$, it also follows that
$$(\beta _{21}\mathrm{i}\beta _{31})(\omega _2+\mathrm{i}\omega _3)=0.$$
Since it has already been shown that
$$\beta _{11}=\eta _1=\eta _2+\mathrm{i}\eta _3=0,$$
it follows, once again, that the natural map from the frame bundle to $`\mathrm{\Lambda }`$ that sends a coframe $`uF_L`$ to $`(𝐱(u),𝐞_1(u))`$ maps the coframe bundle into an $`E`$-holomorphic surface and that this surface is simply the space of lines of the ruling.
Thus, it has been shown that any ruled special Lagrangian $`3`$-fold is locally the $`3`$-fold generated by an $`E`$-holomorphic surface in $`\mathrm{\Lambda }`$.
The only thing left to check is that every $`E`$-holomorphic surface in $`\mathrm{\Lambda }`$ generates a special Lagrangian $`3`$-fold in $`^{\mathrm{\hspace{0.17em}3}}`$. However, given the analysis already done, this is an elementary exercise in the moving frame and can be safely left to the reader. ∎
###### Remark 11 (The relation with special Lagrangian cones).
It is not difficult to see that there is a Levi-flat almost CR-structure of codimension $`1`$ on $`S^5`$ with the property that its holomorphic surfaces are exactly the links of special Lagrangian cones.
In fact, the mapping $`𝐞_1:\mathrm{\Lambda }S^5`$ is an almost CR-mapping in the obvious sense when $`S^5`$ is given this almost CR-structure. In particular, it follows that any ruled special Lagrangian $`3`$-fold is associated to a special Lagrangian cone that one gets by simply translating all of the ruling lines so that they pass through one fixed point. It is in this sense that all of the ruled special Lagrangian $`3`$-folds in $`^{\mathrm{\hspace{0.17em}3}}`$ are ‘twisted cones’ in some sense.
In light of this fact, it may be that there is a formula for ruled special Lagrangian $`3`$-folds that is analogous to the formula for austere $`3`$-folds given in . I have not yet tried to find this.
On the other hand, this relationship shows that there cannot be a ‘Weierstrass formula’ for the general ruled special Lagrangian like the formula given by Borisenko for the family that he discovered. The reason is that such a formula would, at the very least, imply a Weierstrass formula for the links of special Lagrangian cones. However, it is easy to show that this exterior differential system is equivalent to a Monge-Ampere system in $`5`$-dimensions that, by a theorem of Lie, does not admit a Weierstrass formula. Thus, the best that one can hope for is Weierstrass formulae for special cases.
###### Remark 12 (The generalization to the associative case).
As the reader may know, special Lagrangian $`3`$-folds in $`^{\mathrm{\hspace{0.17em}3}}`$ are special cases of a more general family of calibrated $`3`$-folds in $`^{\mathrm{\hspace{0.17em}7}}`$, namely, the *associative* $`3`$-folds as described §IV of .
Regarding $`^{\mathrm{\hspace{0.17em}7}}`$ as $`\times ^{\mathrm{\hspace{0.17em}3}}`$ and using $`x_0`$ as the standard linear coordinate on the $``$-factor, the $`3`$-form
(3.59)
$$\varphi =dx_0\left(\frac{\mathrm{i}}{2}(\mathrm{d}z_1\mathrm{d}\overline{z_1}+\mathrm{d}z_2\mathrm{d}\overline{z_2}+\mathrm{d}z_3\mathrm{d}\overline{z_3})\right)+\mathrm{Re}(\mathrm{d}z_1\mathrm{d}z_2\mathrm{d}z_3)$$
is a calibration on $`^{\mathrm{\hspace{0.17em}7}}`$, called the *associative calibration*. The $`3`$-folds that it calibrates are said to be *associative*. The associative $`3`$-folds that lie in the hyperplane $`\{0\}\times ^{\mathrm{\hspace{0.17em}3}}`$ are exactly the special Lagrangian $`3`$-folds. However, there are many more associative $`3`$-folds than special Lagrangian $`3`$-folds, since, in particular, Harvey and Lawson prove that every connected real-analytic surface $`S^{\mathrm{\hspace{0.17em}7}}`$ lies in an (essentially) unique associative $`3`$-fold \[12, §IV.4, Theorem 4.1\]
The subgroup of $`\mathrm{GL}(7,)`$ that stabilizes $`\varphi `$ is the compact exceptional group $`\mathrm{G}_2`$. It acts transitively on the oriented lines in $`^{\mathrm{\hspace{0.17em}7}}`$ through the origin, and the $`\mathrm{G}_2`$-stabilizer of an oriented line is $`\mathrm{SU}(3)`$. In particular, the group $`\mathrm{\Gamma }`$ generated by the translations in $`^{\mathrm{\hspace{0.17em}7}}`$ and the rotations in $`\mathrm{G}_2`$ acts transitively on the space $`\mathrm{\Lambda }`$ of oriented lines in $`^{\mathrm{\hspace{0.17em}7}}`$.
It can be shown that there is a unique $`\mathrm{\Gamma }`$-invariant almost complex structure on $`\mathrm{\Lambda }`$ with the property that any pseudo-holomorphic surface $`SL`$ defines a ruled associative $`3`$-fold $`\mathrm{\Sigma }^{\mathrm{\hspace{0.17em}7}}`$ (i.e., the oriented union of the oriented lines in $`^{\mathrm{\hspace{0.17em}7}}`$ that the points of $`S`$ represent) and, conversely, that if $`\mathrm{\Sigma }^{\mathrm{\hspace{0.17em}7}}`$ is a ruled associative $`3`$-fold that is not a $`3`$-plane, then the set of oriented lines $`S\mathrm{\Lambda }`$ that meet $`\mathrm{\Sigma }`$ in at least an interval is a pseudo-holomorphic curve in $`\mathrm{\Lambda }`$.
Further development of this description allows one to give a description of the ruled associative $`3`$-folds of $`^{\mathrm{\hspace{0.17em}7}}`$ that directly generalizes Joyce’s description in of the ruled special Lagrangian $`3`$-folds in $`^{\mathrm{\hspace{0.17em}3}}`$.
The details of these results will be reported on elsewhere.
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# Experimental simulation of a stellar photon bath by bremsstrahlung: the astrophysical 𝛾-process
## Abstract
The nucleosynthesis of heavy proton-rich nuclei in a stellar photon bath at temperatures of the astrophysical $`\gamma `$-process was investigated where the photon bath was simulated by the superposition of bremsstrahlung spectra with different endpoint energies. The method was applied to derive ($`\gamma `$,n) cross sections and reaction rates for several platinum isotopes.
The trans-iron nuclei have been synthesized by neutron capture in the $`s`$\- and $`r`$-processes, except the $`p`$-nuclei ($`p`$ for proton-rich), with relative abundances of the order of 0.01 to 1% . The main production mechanism of the $`p`$-nuclei is assumed to be photodisintegration in the $`\gamma `$-process, i.e. by ($`\gamma `$,n), ($`\gamma `$,p), and ($`\gamma `$,$`\alpha `$) reactions induced on heavier seed nuclei synthesized in the $`s`$\- and $`r`$-processes. Typical parameters for the $`\gamma `$-process are temperatures of $`2T_93`$ ($`T_9`$ is the temperature in units of $`10^9`$ K), densities $`\rho 10^6`$ g/cm<sup>3</sup>, and time scales $`\tau `$ in the order of seconds. Several astrophysical sites for the $`\gamma `$-process have been proposed, whereby the oxygen- and neon-rich layers of type II supernovae seem to be good candidates. However, no definite conclusions have been reached yet , predominantly due to the lack of experimental data for the cross sections and reaction rates of the $`\gamma `$-induced reactions at astrophysically relevant energies. All reaction rates have been derived theoretically using statistical model calculations .
The energy distribution of a thermal photon bath at a temperature $`T`$ is given by the Planck distribution
$$n_\gamma (E,T)=\left(\frac{1}{\pi }\right)^2\left(\frac{1}{\mathrm{}c}\right)^3\frac{E^2}{\mathrm{exp}(E/kT)1}$$
(1)
where $`n_\gamma (E,T)`$ is the number of $`\gamma `$-rays at energy $`E`$ per unit of volume and energy interval. In a photon-induced reaction B($`\gamma `$,x)A the distribution leads to a temperature dependent decay rate $`\lambda (T)`$ of the initial nucleus B
$$\lambda (T)=_0^{\mathrm{}}cn_\gamma (E,T)\sigma _{(\gamma ,\mathrm{x})}(E)𝑑E$$
(2)
with the speed of light $`c`$ and the cross section of the $`\gamma `$-induced reaction $`\sigma _{(\gamma ,\mathrm{x})}(E)`$. Obviously, $`\lambda `$ is also the production rate of the residual nucleus A. In the following we will focus on photodisintegration by the ($`\gamma `$,n) reactions.
A large number of ($`\gamma `$,n) cross sections has been measured over the years . However, most of the data have been obtained around the giant dipole resonance (GDR), i.e. at energies much higher than those in stars, and practically no data exist for the $`p`$-nuclei. The integrand in Eq. (2) is given by the product of the $`\gamma `$ flux $`cn_\gamma (E,T)`$, which decreases steeply with increasing energy $`E`$, and the cross section $`\sigma _{(\gamma ,\mathrm{x})}(E)`$, which increases with $`E`$ approaching the GDR region. The product leads then to a window at an effective energy $`E_{\mathrm{eff}}`$ with a width $`\mathrm{\Delta }`$ similar to the Gamow window for charged-particle-induced reactions. If one assumes a typical threshold behavior of the ($`\gamma `$,n) cross section close to the threshold energy $`E_{\mathrm{thr}}`$, the effective energy is approximately given by $`E_{\mathrm{eff}}=E_{\mathrm{thr}}+\frac{1}{2}kT`$, and the typical width $`\mathrm{\Delta }`$ is in the order of 1 MeV (Fig. 1).
The energy distribution of bremsstrahlung is approximately described by the Schiff formula . Close to the endpoint energy $`E_0`$ (which is also the energy of the incoming electron beam) the bremsstrahlung spectrum decreases steeply with increasing energy. We found that over a narrow energy region the bremsstrahlung spectrum has a similar shape as the Planck spectrum. However, a superposition of several bremsstrahlung spectra $`\mathrm{\Phi }_{\mathrm{brems}}(E_{0,i})`$ with different endpoint energies $`E_{0,i}`$ leads to a quasi-thermal spectrum $`\mathrm{\Phi }_{\mathrm{brems}}^{\mathrm{qt}}(T)`$ which has nearly the same shape as the Planck spectrum over a relatively broad energy range (Fig. 2):
$$cn_\gamma (E,T)\mathrm{\Phi }_{\mathrm{brems}}^{\mathrm{qt}}(T)=\underset{i}{}a_i(T)\mathrm{\Phi }_{\mathrm{brems}}(E_{0,i})$$
(3)
where the strength coefficients $`a_i(T)`$ must be adjusted for each temperature $`T`$. The example shown in Fig. 2 demonstrates that a reasonable agreement between quasi-thermal and thermal spectrum in the astrophysically relevant energy range from 6 to 10 MeV is already obtained by the superposition of six bremsstrahlung spectra. The bremsstrahlung spectra have been calculated from GEANT simulations , which compare well with the observed photon flux of the <sup>11</sup>B($`\gamma `$,$`\gamma ^{}`$) reaction (see below). However, close to the endpoint energy the simulated spectra had to be slightly reduced. Details of the calculated bremsstrahlung spectra will be given elsewhere .
For an application of the quasi-thermal bremsstrahlung spectrum we have chosen to measure the ($`\gamma `$,n) cross sections of several platinum isotopes using the activation technique. The high sensitivity of this technique allows for a concurrent measurement of the ($`\gamma `$,n) cross sections for <sup>190</sup>Pt (natural abundance = 0.014%), <sup>192</sup>Pt (0.782%), and <sup>198</sup>Pt (7.163%) due to similar half-lives of the residual nuclei <sup>189</sup>Pt, <sup>191</sup>Pt, and <sup>197</sup>Pt \[e.g. $`T_{1/2}(^{191}\mathrm{Pt})=2.862\mathrm{d}`$\].
The experiment was performed at the real photon facility of the superconducting Darmstadt linear electron accelerator S–DALINAC . The electron beam with a typical beam current up to 40 $`\mu `$A was stopped in a massive copper disk leading to a photon flux of about $`10^5/`$(keV cm<sup>2</sup> s) at the irradiation position. For the irradiation platinum disks (20 mm diameter, 0.125 mm thickness) of natural isotopic composition were mounted at the target position of the ($`\gamma `$,$`\gamma ^{}`$) setup . The disks were sandwiched between two boron layers with masses of about 650 mg each. For normalization of the photon flux, spectra of resonantly scattered photons from nuclear levels of <sup>11</sup>B were obtained during the activation with two high purity germanium (HPGe) detectors (100% relative efficiency) placed at $`90^{}`$ and $`130^{}`$ relative to the incoming photon beam (for details, see also ). The platinum disks were irradiated for about one day, and then they were mounted in front of a third HPGe detector (30% relative efficiency), where the $`\gamma `$-activity was observed for one day. A typical spectrum is shown in Fig. 3: $`\gamma `$-ray lines from the decay of the platinum isotopes <sup>189</sup>Pt, <sup>191</sup>Pt, and <sup>197</sup>Pt can clearly be identified. Additionally, two lines from the decay of <sup>195m</sup>Pt were detected; this isomer is mainly populated by the ($`\gamma `$,$`\gamma ^{}`$) reaction. In another experiment , the decay curves of several lines were measured and found to be in excellent agreement with the recommended half-lives .
Irradiations were performed using seven endpoint energies from $`E_{0,\mathrm{min}}=7200`$ keV to $`E_{0,\mathrm{max}}=9900`$ keV in steps of 450 keV. In the conventional analysis of the data the shape of the ($`\gamma `$,n) cross section was assumed to exhibit a typical threshold behavior:
$$\sigma =\sigma _0\sqrt{(EE_{\mathrm{thr}})/E_{\mathrm{thr}}}.$$
(4)
The constant $`\sigma _0`$ was derived from the yields measured at different endpoint energies which leads then to the astrophysical decay rate $`\lambda `$ using Eq. (2). The results are summarized in Table I. The disadvantage of this analysis, i.e. the assumed shape of the ($`\gamma `$,n) cross section, can be avoided if one uses the quasi-thermal spectrum $`\mathrm{\Phi }_{\mathrm{brems}}^{\mathrm{qt}}(T)`$ (Fig. 2) for a direct determination of the decay rate. In this case the integral in Eq. (2) is measured directly because the experimental yield per target nucleus $`Y_i`$ in the $`i`$-th irradiation is given by
$$Y_i=\mathrm{\Phi }_{\mathrm{brems}}(E_{0,i})\sigma _{(\gamma ,x)}(E)𝑑E.$$
(5)
A comparison of Eq. (5) with Eqs. (2) and (3) relates the decay rate to the experimental yields $`Y_i`$ by
$$\lambda (T)=\underset{i}{}a_i(T)Y_i.$$
(6)
The average deviation between the Planck distribution and the quasi-thermal distribution is about 10% in the relevant energy region around $`E_{\mathrm{eff}}=E_{\mathrm{thr}}+\frac{1}{2}kT`$. The results of the analysis with the quasi-thermal spectrum are also presented in Table I together with results from statistical model calculations using the code NON-SMOKER .
For the future study of ($`\gamma `$,n) reactions on $`p`$-nuclei with their low abundances the sensitivity of the experiments can be improved with enriched samples. These ($`\gamma `$,n) data might restrict the relevant parameters of $`\gamma `$ process models leading eventually to definite conclusions for the astrophysical site of the $`\gamma `$ process.
Finally, under stellar conditions the target nucleus can be excited in the thermal photon bath reducing the effective neutron threshold by the respective excitation energy. This two-step process must be taken into account in model calculations , because in our experiment the target nucleus is always in the ground state.
###### Acknowledgements.
We thank the S–DALINAC group around H.-D. Gräf for the reliable beam during the photoactivation and A. Richter and U. Kneissl for valuable discussions. This work was supported by the Deutsche Forschungsgemeinschaft (contracts Zi 510/2-1 and Ri 242/12-2). T.R. is supported by a PROFIL fellowship from the Swiss National Science Foundation (grant 2124-055832.98) and by the NSF (grant NSF-AST-97-31569).
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# Gravitationally induced transition to classical behavior above 10¹⁰ proton masses.
## Abstract
The localization length for the center of mass motion of a matter lump, induced by gravitation, is obtained, without using any phenomenological constants. Its dependence from mass and volume is consistent both with unitary evolution of microscopic particles and with the classical behavior of macroscopic bodies required to account for wave function collapse in quantum measurements, the transition between the two regimes being rather sharp. The gravitational interaction of nonrelativistic matter is modelled by a Yukawa Hamiltonian with vanishing pion mass, no gravitational background is needed and the only hypothesis consists in assuming unentanglement between matter and the Yukawa field.
03.65.-w, 03.65.Bz
Reconciling gravitation and microphysics has been one of the fundamental open theoretical problems since the birth of Quantum Mechanics (QM), even more challenging after the renormalization of electrodynamics by Feynman, Tomonaga, and Schwinger. On the other hand, according to several authors, the conventional interpretation of QM is not completely satisfactory due to the dualistic description it gives for measurement processes and for time evolution of isolated microscopic systems.
A possible link between these two issues has been suggested on several grounds. However, if we accept that the solution of the measurement problem in QM, which is central to its experimental relevance, should be derived from a consistent unified theory including QM and gravity, which on the other hand, apart from possible astrophysical and cosmological implications, is hardly expected to have any direct experimental evidence, we run the risk of abandoning a well established and rewarding scientific paradigm. We may be lured to admit that, in order to build a consistent atomic, molecular and condensed matter physics, we need a detailed knowledge of the physics at the Planck scale, while one would expect that the remnants of such faraway lengths and energies in everyday physics could be collected in some constants, which, in a phenomenological theory at low energies, should be fitted experimentally.
Within a phenomenological approach some authors introduced classical stochastic external elements in order to produce the wave function collapse. In that context, in order to fit the experimental evidence for unitary evolution of sufficiently small systems and for localization of macroscopic bodies, a localization length and a time constant are usually fixed as phenomenological constants.
As to the link between collapse models and gravitation, two alternative attitudes are conceivable.
One can just use gravitation as little more than an alibi, supported for instance by the consideration that internal excitations of atomic systems induced by the interaction with an external stochastic field are minimized, and made largely compatible with the existing experimental upper bounds, if the coupling is proportional to the particle mass. This is intuitively evident, since in such a case for long enough wavelengths of the stochastic field the main coupling is to the center of mass of the atomic system.
An alternative option consists in starting from gravitation in order to look for quantitative characterizations of the localization process. From this viewpoint, even if one maintains that the Copenhagen formulation of QM, at most supplemented by a multiworld interpretation, is essentially complete, and appeals to the environment-induced decoherence program only, to account for the localization of macroscopic bodies, one should at least consider the gravitational field as included in the environment.
In this paper we choose the latter option, without addressing the search for a theory of everything, but on the contrary limiting our ambitions and looking for quantitative results on gravitationally induced localization of nonrelativistic matter. To be specific we are going to focus on two fundamental questions:
1) Is the gravitational field vacuum able to localize matter lumps without the need for a gravitational background of cosmological origin?
2) Is the transition to classical behavior, induced by localization, smooth or sharp?
As a result we will find an expression for the mass and volume dependent localization length in the absence of such a background and corresponding to a rather sharp transition.
A peculiarity of our result consists, at variance with current collapse models, in a substantial ineffectiveness of the gravitational vacuum in localizing microscopic particles independently from how long we wait. This avoids the need of fixing a large time constant like other authors are forced to do, to limit unwanted localization of microscopic particles. On the contrary we get a rather sharp transition to an extremely effective localization of macroscopic lumps of condensed matter at around $`10^{10}`$ proton masses, this conferring a classical character to the center of mass motion of macroscopic bodies.
Our main hypothesis stems from the need of reconciling two apparently conflicting requirements. On one side the assumption of a wave function collapse induced by the gravitational field requires its classical character. On the other hand our intention of treating it dynamically and not just as an external noise forces us, for consistency with the quantum nature of matter, to quantize it. On different grounds such paradoxical conditions are reflected in Feynman’s words ”…maybe nature is trying to tell us something here, maybe we should not try to quantize gravity” within his lectures on gravitation. Our proposed way out of this puzzle consists in assuming that the hidden physics of a consistent unified theory is continuously collapsing the entanglement between the matter and the pion field.
To be specific consider nonrelativistic particles of mass $`m`$ whose interaction Hamiltonian with a scalar gravitational potential $`\widehat{\varphi }(x)`$is given by
$$H_I\left[\widehat{n}(𝐤)\right]m\sqrt{4\pi c\mathrm{}G}\underset{R^3}{}𝑑𝐤\frac{\widehat{n}(𝐤)}{\sqrt{k}}\left[\widehat{a}(𝐤)+a^{}(𝐤)\right]=mc\sqrt{4\pi G}\underset{R^3}{}𝑑𝐱\widehat{\psi }^{}(𝐱)\widehat{\psi }(𝐱)\widehat{\varphi }(𝐱),$$
(1)
where $`\widehat{n}(𝐤)`$ is the Fourier transform of the product $`\widehat{\psi }^{}(𝐱)\widehat{\psi }(𝐱)`$ (where sum over spin indices is implied) of the matter field creation and annihilation operators, $`\widehat{a}^{}(𝐤)`$ and $`\widehat{a}(𝐤)`$ are the creation and annihilation operators of the corresponding mode of the scalar potential, $`\mathrm{}`$ and $`G`$ the Planck and the gravitational constants. The total Hamiltonian is taken to be the sum of $`H_I`$, of the matter Hamiltonian, and of
$$H_G=c\mathrm{}\underset{R^3}{}𝑑𝐤k\widehat{a}^{}(𝐤)\widehat{a}(𝐤),$$
(2)
this making our model the analog, for vanishing pion mass, of the Yukawa low energy theory of nuclear interactions, by which the scalar particles associated with the $`\widehat{\varphi }`$ field will be referred to as (gravitational) pions.
Of course one could obtain this nonrelativistic model starting from general relativity, keeping just the conformal excitations of a flat vacuum and then linearizing and quantizing the corresponding Hamiltonian. We prefer to keep the analysis completely independent from Einstein equations of the gravitational field, which presumably are only the large scale manifestation of a fundamental theory that may be well out of reach. Furthermore, whereas we are not able to quantize the Einstein version of the gravitational field, this quantum model reproduces the $`1/r`$ interaction potential as the limit of the Yukawa potential for vanishing pion mass and can then be considered as a low energy effective model for gravitation. As a consequence we consider this model more likely of having a direct quantum relevance.
While our procedure reminds nonrelativistic electrodynamics, a substantial difference is to be kept in mind. For electrodynamics a quantization of the Coulomb field alone is not viable, since a scalar Yukawa-like theory can only produce attraction between like particles. Then, when in nonrelativistic classical electrodynamics one fixes the Coulomb gauge, the Coulomb interaction is separated from the radiation field and inserted into the matter Hamiltonian so that, if one confines his attention to Coulomb interaction, only the matter degrees of freedom are left. On the contrary here pion degrees of freedom are explicitly responsible for the $`1/r`$ law.
To be specific, if one for simplicity replaces $`\widehat{n}(𝐤)`$ in Eq. (1) with the expression, $`1+\mathrm{exp}[i𝐤𝐫]`$, corresponding to two classical point sources respectively in the origin and in the point $`𝐫`$, one finds that, while of course the ground state energy of the $`\widehat{\varphi }`$ field in the presence of these sources diverges due to their singular character, its gradient is well defined and gives the expected $`1/r^2`$ force law. If the two masses are left free to move, then the kinetic energy they acquire in falling towards the common center of mass is balanced by the decrease of the ground state energy of the coherent pion cloud.
This simple exercise gives a relevant clue to answer the first aforementioned fundamental question. At first sight the answer is in the negative since obviously some energy is needed to localize matter and it is precisely this energetic effect to be considered one of the possible experimental evidences of fundamental localization processes. On the other hand, if one considers that the $`\widehat{\varphi }`$ vacuum is not the ground state in the presence of matter, one has to assume that the coherent pion cloud corresponding to the $`\widehat{\varphi }`$ ground state depends on the matter state. Then the localization process is accompanied by the rearrangement of the $`\widehat{\varphi }`$ ground state, which is here too the source of the incremental particle kinetic energy.
Let’s begin by evaluating the gravitational ground state energy in the presence of the $`N`$-particle matter state $`|\mathrm{\Psi }`$ corresponding for computational simplicity to an isotropic Gaussian mass density of dispersion $`\lambda `$:
$$\rho (𝐱)=m\mathrm{\Psi }\left|\widehat{\psi }^{}(𝐱)\widehat{\psi }(𝐱)\right|\mathrm{\Psi }=\frac{Nme^{\frac{x^2}{2\lambda ^2}}}{(2\pi )^{3/2}\lambda ^3}n_\lambda (𝐤)\frac{1}{\left(2\pi \right)^{3/2}}𝑑𝐱\frac{\rho (𝐱)}{m}e^{i𝐤𝐱}=\frac{Ne^{\frac{\lambda ^2k^2}{2}}}{(2\pi )^{3/2}}.$$
(3)
In order to do that we observe that the product state $`|\mathrm{\Psi }|n_\lambda _\varphi `$ is the dressed matter state with the pion cloud in its ground state, namely it is the minimum energy state for $`H_G+H_I[n_\lambda (𝐤)]`$ with a fixed particle factor $`|\mathrm{\Psi },`$ if $`|n_\lambda _\varphi `$ is the coherent state
$$|n_\lambda _\varphi =\mathrm{exp}\left[m^2\frac{4\pi G}{c\mathrm{}}𝑑𝐤k^3n_\lambda (𝐤)^2\right]\mathrm{exp}\left[m\sqrt{\frac{4\pi G}{c\mathrm{}}}𝑑𝐤k^{3/2}n_\lambda (𝐤)\widehat{a}^{}(𝐤)\right]|0_\varphi ,$$
(4)
where $`|0_\varphi `$ denotes the pion vacuum. This immediately follows from the expression
$$H_G+H_I[n_\lambda (𝐤)]=c\mathrm{}\underset{R^3}{}𝑑𝐤k\widehat{b}^{}(𝐤)\widehat{b}(𝐤)+E_0(\lambda ),$$
(5)
where
$$\widehat{b}(𝐤)\widehat{a}(𝐤)+m\sqrt{\frac{4\pi G}{c\mathrm{}}}\frac{n_\lambda (𝐤)}{k^{3/2}}$$
(6)
and the pion ground state energy is given by
$$E_0(\lambda )=4\pi m^2G\frac{d𝐤}{k^2}\left|n_\lambda (𝐤)\right|^2=4\pi G\frac{M^2}{\lambda }.$$
(7)
Assume now that we have an isolated matter lump and that its wave function is the product of a wave function of the center of mass $`𝐱`$
$${}_{CM}{}^{}\mathrm{\Psi }_{\lambda ^{}}^{}(𝐱)=\frac{\mathrm{exp}\left[\frac{x^2}{4\lambda ^{}_{}{}^{}2}\right]}{(2\pi )^{3/4}\lambda ^{3/2}}$$
(8)
with Gaussian squared modulus of dispersion $`\lambda ^{}`$, and the wave function of some stationary inner state, whose particle density, for computational simplicity, is approximated by a Gaussian of dispersion $`\lambda _0`$. Then, ignoring correlations, the corresponding mass density has dispersion $`\lambda `$=$`\sqrt{\lambda _0^2+\lambda ^2}`$.
The total energy of the system including the matter lump and its pion cloud is the sum of the pion ground state energy and the matter kinetic energy
$$E_T(\lambda ^{})=E_0(\sqrt{\lambda _0^2+\lambda ^2})\frac{\mathrm{}^2}{2M}\underset{R^3}{}{}_{CM}{}^{}\mathrm{\Psi }_{\lambda ^{}}^{}(𝐱)_{CM}^2\mathrm{\Psi }_\lambda ^{}(𝐱)𝑑𝐱=\frac{4\pi GM^2}{\sqrt{\lambda _0^2+\lambda ^2}}+\frac{3\mathrm{}^2}{8M\lambda ^2},MNm.$$
(9)
The $`\lambda ^{}`$ value for which $`E_T(\lambda ^{})`$ attains its minimum is the maximum localization length, since less localized states are unstable with respect to conversion of pion field energy into kinetic energy, just like the two point masses mentioned above.
The crucial equation is then
$$\frac{d}{d\lambda ^{}}E_T(\lambda ^{})=0\left(\lambda _0^2+\lambda ^2\right)^3=\frac{64}{9}\pi ^2G^2\mathrm{}^4M^6\lambda ^8,$$
(10)
and, to be specific, we can chose a typical mass density value corresponding to condensed matter $`M/\lambda _0^3=10^{24}m_p/cm^3`$, where $`m_p`$ denotes the proton mass, so that Eq. (10) becomes
$$\left(10^{16}\mu ^{2/3}cm^2+\lambda ^2\right)^3=\frac{64}{9}\pi ^2G^2\mathrm{}^4m_p^6\mu ^6\lambda ^8,\mu M/m_p.$$
(11)
This equation gives
$$\lambda ^{}4.2310^{23}\mu ^3cmif\mu 10^9,\lambda ^{}.81\mu ^{1/2}cmif\mu 10^{10},$$
(12)
which shows that around $`M10^{10}m_p`$ there is a rather sharp transition to classical behavior.
This result implies that the localization due to the gravitational vacuum is microscopically irrelevant even at molecular level, while it accounts for localization of macroscopic bodies and in particular for the collapse of the pointer states in quantum measurement theory. Apart from accounting for this elusive feature of QM, this localization length could have relevant physical consequences at mesoscopic level, where a localization length $`\lambda ^{}10^5cm`$ is comparable with the dimensions of a corresponding lump of condensed matter with $`m10^{10}m_p`$.
We can conclude that, within the present nonrelativistic setting and with the only hypothesis of unentanglement between matter and the pion field, it is proven that gravitation accounts for wave function collapse without assuming a gravitational background of cosmological origin. Of course this does not rule out that the effects of such a background, which may depend on the particular space-time region, should be added to the more fundamental ones proposed here. Furthermore one should remark that, while environment induced decoherence, without including gravitational effects, if not properly controlled, may overshadow gravitational localization, it can only account for the entanglement between measured systems and pointer states, but not for the final collapse of the entangled state.
Finally one should observe that the expression for $`E_0(\lambda )`$ could have been obtained more naively even without quantizing the pion field. We chose to put the classical agent outside the considered model. Of course the final result only shifts the ”Heisenberg cut” to a consistent unified theory outside the scope of the present letter. However we are encouraged in this attempt by Feynman’s words ”People say to me, ‘Are you looking for the ultimate laws of physics?’ No, I’m not.”
Acknowledgments - Financial support from M.U.R.S.T., Italy and I.N.F.M., Salerno is acknowledged
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# Superconductivity from Undressing. II. Single Particle Green’s Function and Photoemission in Cuprates
## I Introduction
Photoemission and optical experiments indicate that in high $`T_c`$ cuprates a transition from an incoherent to a partially coherent state occurs both as the hole doping increases in the normal state and as the system goes superconducting. Basov and coworkers have observed a lowering of c-axis kinetic energy as the transition to the superconducting state occurs in several cuprates, especially in the underdoped situation. It has been established however that the $`magnitude`$ of the c-axis kinetic energy lowering detected is far too small to account for the superconducting condensation energy at least in some cuprates. Ding, Campuzano, Norman and coworkers as well as Feng and coworkers have reported observations of sharp quasiparticle peaks in the superconducting state in angle resolved photoemission emerging from a highly incoherent normal state background along the $`(\pi ,0)`$ direction, close to the $`(\pi /a,0)`$ point. Ding and coworkers have interpreted the photoemission peak in terms of an enhanced quasiparticle weight $`Z`$ in the superconducting state, and Feng and coworkers have suggested that the peak in photoemission is a signature of the superfluid density. Norman and coworkers have analyzed the photoemission observations in terms of a ’mode model’ and emphasized the close connection between their observations and Basov’s observation of kinetic energy lowering. Furthermore, Basov and coworkers have emphasized that kinetic energy lowering seems to occur only when there is a high degree of incoherence in the normal state, and appears to vanish as the normal state becomes more coherent (overdoped regime). They have furthermore proposed that the photoemission experiments suggest that kinetic energy lowering may occur also $`inplane`$ in the cuprates albeit only along the $`(\pi ,0)`$ direction, and for that reason may be difficult to observe directly.
The model of hole superconductivity $`predicted`$, before the experimental observations, that the superconducting condensation energy originates in in-plane kinetic energy lowering and arises from a process of $`undressing`$ $`of`$ $`hole`$ $`carriers`$ as the pairing state develops. Thus it describes both the kinetic energy lowering, arising from the low energy effective Hamiltonian, as well as the high energy optical spectral weight transfer, that has also been observed experimentally. In the first paper of this series (hereafter referred to as I) we formulated more generally the principles of superconductivity through hole undressing, and pointed out that this physics would show up both in the one and two-particle Green’s functions, in qualitative agreement with the observations reported above. Here we report calculation of the single particle Green’s function and spectral weight in the superconducting state and discuss implications for the understanding of photoemission experiments.
## II General principles
In the class of models discussed in I, the wave function renormalization of quasiparticles is a function of the site occupation in a local representation. The wavefunction renormalization arises from coupling to a local boson degree of freedom. Three examples of specific microscopic Hamiltonians describing this physics were discussed in I. The ’coherent’ part of the electron creation operator at site $`i`$ is defined by the following transformation:
$$d_{i\sigma }^{}=[T(TS)\stackrel{~}{n}_{di,\sigma }]\stackrel{~}{d_{i\sigma }}^{}$$
(1)
with $`0S<T1`$. The $`\stackrel{~}{d}`$ operators in Eq. (1) are quasiparticle operators, and $`\stackrel{~}{n}_d`$ is the electron site occupation. Eq. (1) expresses the fact that the electron becomes less coherent as more electrons are added to the band. It should be kept in mind that the ’coherent part’ of the electron operator on the left side of Eq. (1) is not the full electron creation operator, as it does not contain terms that give rise to excited states of the boson degree of freedom.
It will be more useful to use hole operators rather than electron operators throughout this paper; we stress however that the discussion can be consistently carried out in electron as well as in hole representation. In terms of hole operators, the coherent part of the hole creation operator is
$$c_{i\sigma }^{}=[S+(TS)\stackrel{~}{n}_{i,\sigma }]\stackrel{~}{c_{i\sigma }}^{}S[1+\mathrm{{\rm Y}}\stackrel{~}{n}_{i,\sigma }]\stackrel{~}{c_{i\sigma }}^{}$$
(2)
Eq. (2) expresses the fact that the hole quasiparticle weight will increase with the local hole concentration, from $`S`$ in the regime of low hole concentration to $`T`$ for high hole concentration. The high degree of incoherence observed in high $`T_c`$ cuprates for low hole doping implies $`S<<1`$, and the fact that coherence is achieved for relatively small values of hole doping implies that the ’undressing parameter’
$$\mathrm{{\rm Y}}=\frac{T}{S}1$$
(3)
is very large. $`\mathrm{{\rm Y}}`$ is the parameter that drives the transition to the superconducting state. Note that a large $`\mathrm{{\rm Y}}`$ necessarily implies $`S<<1`$, due to the constraint $`T1`$. For the normal state, Eq. (2) implies for the hole operator
$$c_{i\sigma }^{}=S[1+\frac{n}{2}\mathrm{{\rm Y}}]\stackrel{~}{c_{i\sigma }}^{}$$
(4)
with $`n`$ the hole concentration per site, and for the hole number operator
$$n_{i\sigma }=S^2[1+\frac{n}{2}\mathrm{{\rm Y}}]^2\stackrel{~}{n}_{i\sigma }Z(n)\stackrel{~}{n}_{i\sigma }$$
(5)
with $`Z(n)`$ the hole quasiparticle weight. Eq. (5) implies that hole quasiparticles in the normal state become more coherent as the hole concentration increases. In the limit $`S0`$, quasiparticles become completely incoherent in the normal state for low hole concentration and Fermi liquid theory breaks down. That limit is also described by the theory; in that limit, the transition to the superconducting state is a superconductor-insulator transition. Even though for that particular limiting situation Fermi liquid theory does not describe the normal state, we stress that our approach is $`not`$ a ’non-Fermi-liquid’ approach, but instead is deeply rooted in Fermi liquid theory.
Consider the bare kinetic energy in a tight binding model in terms of hole operators
$$H_{kin}=\underset{i,j,\sigma }{}t_{ij}^0(c_{i\sigma }^{}c_{j\sigma }+h.c.)$$
(6)
Replacement of the bare hole operators by the quasiparticle operators using eq. (2) yields
$$H_{kin}=\underset{i,j,\sigma }{}t_{ij}^\sigma (\stackrel{~}{c}_{i\sigma }^{}\stackrel{~}{c}_{j\sigma }+h.c.)$$
(8)
$$t_{ij}^\sigma =t_{ij}^0S^2[1+\mathrm{{\rm Y}}(\stackrel{~}{n}_{i,\sigma }+\stackrel{~}{n}_{j,\sigma })+\mathrm{{\rm Y}}^2\stackrel{~}{n}_{i,\sigma }\stackrel{~}{n}_{j,\sigma }]$$
(9)
Eq. (7) expresses the fact that the hopping amplitude of a hole quasiparticle will be increased, and as a consequence its kinetic energy will be lowered, as the local hole concentration increases; it is a direct consequence of the fact that the quasiparticle coherence increases with local hole concentration as described by Eq. (2). For low hole concentration we can ignore the last term in Eq. (7b) and obtain
$$H_{kin}=\underset{<i,j>,\sigma >}{}[t_{ij}+\mathrm{\Delta }t_{ij}(\stackrel{~}{n}_{i,\sigma }+\stackrel{~}{n}_{j,\sigma })](\stackrel{~}{c}_{i\sigma }^{}\stackrel{~}{c}_{j\sigma }+h.c.)$$
(11)
$$t_{ij}=S^2t_{ij}^0$$
(12)
$$\mathrm{\Delta }t_{ij}=\mathrm{{\rm Y}}t_{ij}$$
(13)
The kinetic energy of the form Eq. (8) is used in the model of hole superconductivity, and leads to pairing and superconductivity for low hole concentration in the presence of appreciable on-site and nearest neighbor Coulomb repulsion. The condition for superconductivity to occur is
$$\mathrm{{\rm Y}}>\sqrt{(1+u)(1+w)}1$$
(14)
where $`u`$ and $`w`$ are dimensionless on-site and nearest neighbor Coulomb repulsions. Hence within this class of models superconductivity is intimately tied to increased quasiparticle coherence. Note that in a model with anisotropy Eq. (8) still implies
$$\frac{\mathrm{\Delta }t_{ij}}{t_{ij}}=\mathrm{{\rm Y}}$$
(15)
$`independent`$ of direction. This assumption was used in our studies with the model of hole superconductivity, and can be seen to be a necessary consequence of the fact that the $`\mathrm{\Delta }t`$ term in the Hamiltonian arises from quasiparticle undressing . A necessary consequence of Eq. (10) is that the superconducting energy gap function has the form
$$\mathrm{\Delta }_k=\mathrm{\Delta }(ϵ_k)$$
(16)
and hence is constant over the Fermi surface $`(ϵ_k=ϵ_F)`$, even for an anisotropic band structure. Thus, Eq. (10) can be understood as a direct consequence of the undressing physics. Finally, Eq. (8) leads to superconductivity through kinetic energy lowering. Hence, within the undressing scenario considered here, kinetic energy lowering as the system goes superconducting is intimately tied to $`swave`$ $`symmetry`$ of the superconducting order parameter as described by Eq. (11).
## III Green’s function: coherent part
The single particle Green’s function is given by a sum of coherent and incoherent parts
$$G_{rs}(\tau )=<Tc_r(\tau )c_s^{}(0)>G_{rs}^{coh}(\tau )+G_{rs}^{incoh}(\tau )$$
(17)
with $`T`$ the time ordering operator. The coherent part of the Green’s function is obtained by replacing the bare fermion operators in Eq. (12) by its coherent parts, given by Eq. (2) in terms of the quasiparticle operators:
$$G_{rs}^{coh}(\tau )=S^2<T[1+\mathrm{{\rm Y}}\stackrel{~}{n}_{r,}(\tau )]\stackrel{~}{c}_r(\tau )][1+\mathrm{{\rm Y}}\stackrel{~}{n}_{s,}(0)]\stackrel{~}{c}_s^{}(0)]>$$
(18)
The normal and anomalous Green’s functions for the quasiparticle operators
$$\stackrel{~}{G}_{rs}(\tau )=<T\stackrel{~}{c}_r(\tau )\stackrel{~}{c}_s^{}(0)>$$
(20)
$$\stackrel{~}{F}_{rs}(\tau )=<T\stackrel{~}{c}_r(\tau )\stackrel{~}{c}_s(0)>$$
(21)
are given by the usual form
$$\stackrel{~}{G}(k,i\omega _n)=\frac{u_k^2}{i\omega _nE_k}+\frac{v_k^2}{i\omega _n+E_k}$$
(23)
$$\stackrel{~}{F}(k,i\omega _n)=u_kv_k[\frac{1}{i\omega _nE_k}\frac{1}{i\omega _n+E_k}]$$
(24)
where the coherence factors $`u_k`$, $`v_k`$ and quasiparticle energies $`E_k`$ are given by the usual BCS expressions
$$u_k^2=\frac{1}{2}(1+\frac{ϵ_k\mu }{E_k})$$
(26)
$$v_k^2=\frac{1}{2}(1\frac{ϵ_k\mu }{E_k})$$
(27)
$$u_kv_k=\frac{\mathrm{\Delta }_k}{2E_k}$$
(28)
$$E_k=\sqrt{(ϵ_k\mu )^2+\mathrm{\Delta }_k^2}$$
(29)
and the gap function $`\mathrm{\Delta }_k`$ is obtained from the BCS solution of the model of hole superconductivity, i.e. the kinetic energy Eq. (8) supplemented with on-site and nearest neighbor Coulomb repulsion. The single particle energy $`ϵ_k`$ in these equations is given by $`ϵ_k=Z(n)ϵ_k^0S^2(1+n\mathrm{{\rm Y}})ϵ_k^0`$, with $`ϵ_k^0`$ the bare kinetic energy given by the Fourier transform of $`(t_{ij}^0)`$.
It can be seen that the extra density operators in Eq. (13) will modify the normal Green’s function introducing anomalous terms, similar to the anomalous terms that occur when calculating the expectation value of the kinetic energy Eq. (7) that lead to the optical sum rule violation. We expand Eq. (13), keeping only linear terms in the density as appropriate to the low hole concentration regime, and use mean field decoupling for the averages to obtain
$$G^{coh}(k,i\omega _n)=S^2[(1+n\mathrm{{\rm Y}})\stackrel{~}{G}(k,i\omega _n)+2f_0\mathrm{{\rm Y}}\stackrel{~}{F}(k,i\omega _n)]$$
(30)
with $`f_0=<\stackrel{~}{c}_i\stackrel{~}{c}_i>`$ the on-site pair amplitude in the superconducting state. We have also performed a space and time Fourier transform. It can be seen that the normal Green’s function has acquired a contribution from the anomalous Green’s function due to the density-dependent dressing.
However the quasiparticle spectral weights derived from Eq. (17) are not positive definite and in fact can become negative in extreme parameter regimes. To remedy this we need to include higher order terms obtained from Eq. (13) by keeping terms with 6 fermion operators. Performing a similar mean field decoupling for these we finally obtain for the Green’s function
$$G^{coh}(k,i\omega _n)=\frac{Z_h}{i\omega _nE_k}+\frac{Z_e}{i\omega _n+E_k}$$
(32)
with
$$Z_h=S^2[[1+n\mathrm{{\rm Y}}]u_kf_0\mathrm{{\rm Y}}v_k]^2$$
(33)
$$Z_e=S^2[[1+n\mathrm{{\rm Y}}]v_k+f_0\mathrm{{\rm Y}}u_k]^2$$
(34)
and
$$f_0=<\stackrel{~}{c}_i\stackrel{~}{c}_i>=\frac{1}{N}\underset{k}{}\frac{\mathrm{\Delta }_k}{2E_k}(12f(E_k)).$$
(35)
The quasiparticle weights $`Z_h`$ and $`Z_e`$ are clearly positive definite. Their sum is not conserved as function of density or temperature because of contributions from the incoherent part of the Green’s function not contained in Eq. (13).
In the absence of undressing ($`\mathrm{{\rm Y}}=0`$) the coherent Green’s function Eq. (18) reduces to the usual BCS form except for the overall factor $`S^2`$. In the presence of undressing ($`\mathrm{{\rm Y}}>0`$) Eq. (18) shows that the coherent part of the Green’s function and spectral function will increase with hole density $`n`$, both for positive and negative energies. Furthermore, as the system goes superconducting the on-site pair amplitude $`f_0`$ develops a positive expectation value. From Eq. (18) this implies that the coherent spectral weight will $`decrease`$ for $`positive`$ energies (hole injection) and $`increase`$ for $`negative`$ energies (hole extraction). The effect in the superconducting state will be largest for parameters where the on-site pair amplitude is large, which corresponds to short coherence length achieved in the strong coupling underdoped regime. The magnitude of these effects both in the normal and superconducting state depend on the magnitude of the undressing parameter $`\mathrm{{\rm Y}}`$. We discuss the implications of these results in subsequent sections.
## IV Results for quasiparticle weights
To illustrate the behavior emerging from the results of the previous section we consider now a specific example. The quasiparticle Hamiltonian is given by the kinetic energy Eq. (8) supplemented by on-site and nearest neighbor Coulomb repulsion
$$H_{Coul}=U\underset{i}{}\stackrel{~}{n}_i\stackrel{~}{n}_i+V\underset{<ij>}{}\stackrel{~}{n}_i\stackrel{~}{n}_j$$
(36)
The BCS solution of this Hamiltonian yields the quasiparticle energies
$$E_k=\sqrt{(ϵ_k\mu )^2+\mathrm{\Delta }_k^2}=\sqrt{a^2(ϵ_k\mu \nu )^2+\mathrm{\Delta }_0^2}$$
(38)
$$\mathrm{\Delta }_k\mathrm{\Delta }(ϵ_k)=\mathrm{\Delta }_m(\frac{ϵ_k}{D/2}+c)$$
(39)
$$\mathrm{\Delta }_0=\frac{1}{a}\mathrm{\Delta }(\mu )$$
(40)
$$\nu =\frac{1}{a}\frac{\mathrm{\Delta }_m}{D/2}\mathrm{\Delta }_0$$
(41)
$$a=\frac{1}{\sqrt{1+(\frac{\mathrm{\Delta }_m}{D/2})^2}}$$
(42)
with $`\mathrm{\Delta }_m`$ and $`c`$ parameters that depend on temperature and doping. The bandwidth $`D`$ in these equations is given by
$$D=D_h(1+n\mathrm{{\rm Y}})$$
(43)
with $`D_h`$ the bandwidth in the limit of zero hole concentration. The quasiparticle gap, i.e. the minimum quasiparticle excitation energy, is given by
$$E_g=\mathrm{\Delta }_0$$
(44)
and occurs at momentum defined by
$$ϵ_k^{[0]}=\mu +\nu $$
(45)
However, if $`ϵ_k^{[0]}`$ is below the bottom of the band, which occurs when the chemical potential is sufficiently below the bottom of the band at low hole concentration, Eq. (22) is not valid, and instead
$$E_g=\sqrt{(\frac{D}{2}\mu )^2+\mathrm{\Delta }(\frac{D}{2})}$$
(46)
We consider a two-dimensional square lattice with only nearest neighbor hopping and $`t_{ij}=t_h`$ in Eq. (8). The quasiparticle bandwidth as the hole concentration goes to zero is $`D_h=2zt_h`$, with $`z=4`$ the number of nearest neighbors to a site. We choose parameters
$`D_h=0.2eV`$ (47)
$`U=5eV`$ (48)
$`V=0.65eV`$ (49)
$`\mathrm{{\rm Y}}=19.2`$ (50)
which imply $`\mathrm{\Delta }t=\mathrm{{\rm Y}}D_h/2z=0.48eV`$. For the present purposes we need not specify the magnitude of the parameter $`S^2`$, which determines the relative weight of coherent and incoherent contributions to the spectral function.
These parameters yield a maximum $`T_c`$ versus hole concentration of $`T_c^{max}=94K`$, as shown in Fig. 1(a), for optimal doping $`n0.045`$. The minimum quasiparticle excitation energy at low temperatures is shown in Fig. 1(b). At low hole concentration it does not go to zero as $`T_c`$ does because the chemical potential falls below the bottom of the band and $`E_g`$ is determined by Eq. (24) rather than by Eq. (22). The behavior of the chemical potential and the bottom of the band versus hole concentration is shown in Fig. 2. The chemical potential crosses the bottom of the band at $`n0.038`$, and $`ϵ_k^0`$ (Eq. (23)) crosses the bottom of the band at $`n0.034`$.
The on-site pair amplitude $`f_0`$ that enters in the expressions for the quasiparticle weights is shown in Fig. 3. As function of doping it follows approximately the behavior of the critical temperature and of the gap parameter $`\mathrm{\Delta }_0`$ (not shown). At low hole doping it goes to zero because the carrier concentration goes to zero, at high hole doping it goes to zero because the coherence length is diverging. As function of temperature, $`f_0`$ behaves approximately like the gap, going to zero at $`T_c`$ as $`(T_cT)^{1/2}`$.
Next we consider the behavior of the quasiparticle weights $`Z_e`$ and $`Z_h`$ as function of temperature. Figure 4 shows the results at the (normal state) Fermi energy, $`ϵ_k=\mu `$, for the optimally doped case ($`n=0.045`$). The values are normalized so that $`Z_e`$ and $`Z_h`$ would be $`0.5`$ for $`\mathrm{{\rm Y}}=0`$. The dashed line shows the value the weights would have for $`f_0=0`$: it is temperature independent and larger than $`0.5`$ because of the undressing due to the average carrier concentration $`n`$. The effect of onset of superconductivity is to increase $`Z_e`$ as the temperature is lowered and to decrease $`Z_h`$. This indicates that there is extra amplitude for electron creation, and less amplitude for hole creation. It may thus be interpreted as a shift of the chemical potential as superconductivity sets in, giving increased hole occupation as the temperature is lowered, or equivalently a shrinking of the electron Fermi sea. This is a surprising result of this calculation, and its implications will be discussed in subsequent sections. Note that the weight $`Z_e`$ increases by almost a factor of $`2`$ between $`T=T_c`$ and $`T=0`$. The magnitude of the increase of course depends on the magnitude of the undressing parameter $`\mathrm{{\rm Y}}`$, and would be larger or smaller for larger or smaller values of $`\mathrm{{\rm Y}}`$ respectively. By adjusting the values of on-site and nearest neighbor Coulomb repulsion in the model it would be possible to obtain the same maximum $`T_c`$ with different values of $`\mathrm{{\rm Y}}`$, as discussed in previous work. Nevertheless we believe that the parameters chosen for this example may be representative of the situation in high $`T_c`$ materials.
Note also that the total weight of the spectral function $`Z_{tot}=Z_e+Z_h`$ increases as the temperature is lowered below $`T_c`$. This indicates that overall there is more coherence in the superconducting than in the normal state, in accordance with qualitative expectations, and this extra spectral weight is transfered from the high energy incoherent part of the spectral function as will be discussed in the following section. However part of the enhancement of $`Z_e`$ at low temperatures relative to its value at $`T_c`$ can be attributed to spectral weight being transfered from negative to positive energies (i.e. corresponding depletion of $`Z_h`$) in addition to spectral weight transfer from the incoherent part of the spectral function.
Similarly Fig. 5 shows the results for an overdoped case, $`n=0.1`$, with $`T_c=68K`$. The behavior is qualitatively similar as the optimally doped case, however the effect of the onset of superconductivity on the spectral weights is considerably smaller because the system is already more coherent in the normal state. This is indicated by the larger values of all the spectral weights relative to the values of the case shown in Fig. 4, due to the enhanced coherence arising from the increased hole concentration. For a much higher hole concentration, as $`T_c`$ approaches zero, the ’gap’ between $`Z_e`$ and $`Z_h`$ in the superconducting state closes, as shown in Fig. 6. It remains always nonzero however as long as $`T_c`$ is nonzero, and there is always some spectral weight transfer from the incoherent region as long as $`T_c`$ is nonzero.
Next we consider the spectral weights for other values of momentum. Figure 7 shows results for $`ϵ_k\mu >0`$. Recall that we are using hole representation, so $`ϵ_k\mu >0`$ means inside the filled Fermi sea for electrons. In the normal state $`Z_e=0`$: since the electron state is full, no new electron can be created in it. Just as in the conventional BCS case, as superconductivity sets in the state becomes partially occupied and $`Z_e0`$ , and $`Z_h`$ correspondingly decreases. However, unlike conventional BCS, $`Z_e`$ and $`Z_h`$ cross in our case and at low temperatures the weight for creating an electron is larger than that for creating a hole, even though we are inside the filled normal state Fermi sea. Clearly this implies that the chemical potential in the superconducting state has changed. Figure 7(b) shows that for $`ϵ_k\mu =16.1meV`$ the weights for electrons and holes coincide at low temperatures; this momentum then corresponds to the new Fermi momentum, $`k_F^{}`$, in the superconducting state. For even larger $`ϵ_k`$, $`Z_e`$ becomes smaller than $`Z_h`$ as in the conventional case, as shown in Fig. 7(c).
The behavior for negative energy (outside of the electron Fermi sea) is shown in Fig. 8. Here, $`ϵ_k\mu `$ was chosen to be at the bottom of the hole band in the optimally doped case. The weight for electron creation is much larger than in the conventional case. Note also that $`Z_e`$ first decreases and then increases as the temperature is lowered, and as $`T0`$ it becomes even larger than its value in the normal state. Such a situation, which is never seen in the conventional case, is possible here due to the non-conservation of $`Z_{tot}`$ because of the transfer of spectral weight from the incoherent part to the coherent part of the spectral function as the temperature is lowered.
Figure 9 shows the spectral weights for an overdoped case, $`n=0.1`$, for values of the momentum above the electron Fermi surface, at the Fermi surface and below the Fermi surface. The behavior is qualitatively similar to that for the optimally doped case, although the differences between the conventional and our case are less pronounced here because this is a weaker coupling regime. In figure 10 we show the spectral weight for an underdoped case, $`n=0.02`$. Here the chemical potential is below the bottom of the band, so the situation comparable to Fig. 4 cannot be attained. Fig. 10 shows the behavior of the spectral weight for $`ϵ_k`$ at its lowest possible value, the bottom of the band, which is qualitatively similar to other cases where $`ϵ_k`$ is above $`\mu `$ such as Fig. 7(a).
Next we consider the behavior of the quasiparticle weights at the chemical potential versus doping in Fig. 11. The upper dot-dashed line is the total spectral weight in the superconducting state, and the dotted line below it is the total spectral weight in the normal state. The difference between the two is the spectral weight transfered from high energy incoherent processes as the system goes superconducting; this difference approaches zero in the overdoped regime. The full lines denote the quasiparticle weights in our case, the dashed lines the usual BCS results ($`u_k^2=Z_h,v_k^2=Z_e`$), which increase approximately linearly with $`n`$ due to the normal state increased coherence with doping and are equal for $`ϵ_k=\mu `$. For low dopings however the chemical potential falls below the bottom of the band and hence we take $`ϵ_k`$ at the bottom of the band rather than at $`\mu `$, this is why the two dashed lines diverge at low dopings. In our case, the weight for electron creation (solid curve labeled $`Z_e`$) is seen to increase rapidly with doping and then taper off for high doping; this latter effect is due to the reduction of the on-site pair amplitude $`f_0`$ for high doping as the coherence length becomes large. The quasiparticle weights $`Z_e`$ and $`Z_h`$ approach each other and the BCS value for high doping, as expected. Note also that there is a narrow doping regime where the electron weight $`Z_e`$ is even larger than the total weight in the normal state (dotted line). This situation can never occur in the conventional BCS case.
We believe the behavior exhibited by $`Z_e`$ in Fig. 11 is relevant to the understanding of the angle resolved photoemission results discussed by Ding et al. In their work, the quasiparticle weight in the superconducting state extracted from photoemission spectra shows similar qualitative behavior to the behavior exhibited by $`Z_e`$ in Fig. 11. We will discuss the relation between $`Z_e`$ and the experimental quantity in a subsequent section. Ding et al also plot $`Z\mathrm{\Delta }`$, the product of their extracted quasiparticle weight and the gap inferred from the photoemission spectra, and point out that its behavior rougly follows the bell-shaped curve of $`T_c`$. Our calculation shows similar behavior, as shown in Fig. 12. Note that the quasiparticle gap itself remains finite in our calculation as the hole concentration goes to zero (Fig. 1 (b)), and also experimentally.
## V Green’s function: incoherent part
To calculate the incoherent part of the Green’s function we now consider a specific model, the generalized Holstein model discussed in I. Our calculation follows closely the calculation of Alexandrov and Ranninger for the conventional Holstein model, and we refer the reader to their seminal work for details which are common to both situations. The site Hamiltonian for our case is given by
$$H=\mathrm{}\omega _0a^{}a+g\mathrm{}\omega _0(a^{}+a)(n_{}+n_{}\gamma n_{}n_{})+Un_{}n_{}$$
(51)
The case of Alexandrov and Ranninger corresponds to $`\gamma =0`$. Using a generalized Lang-Firsov transformation the quasiparticle (polaron) operators $`\stackrel{~}{c}_{i\sigma }`$ are related to the bare fermion (hole) operators by
$$c_{i\sigma }=\stackrel{~}{c}_{i\sigma }X_{i\sigma }$$
(53)
$$X_{i\sigma }=e^{g(a_i^{}a_i)(1\gamma \stackrel{~}{n}_{i\sigma })}$$
(54)
In contrast to Eq. (2), the operator $`c_{i\sigma }`$ here is the full hole destruction operator, including coherent and incoherent parts. The coherent part results from the expectation value of the $`X`$-operators in the zero boson subspace,
$$<X_{i\sigma }>=e^{\frac{g^2}{2}(1\gamma \stackrel{~}{n}_{i\sigma })}$$
(55)
and in particular
$$X_{i\sigma }(\stackrel{~}{n}_{i\sigma }=0)=e^{\frac{g^2}{2}}=S$$
(57)
$$X_{i\sigma }(\stackrel{~}{n}_{i\sigma }=1)=e^{\frac{g^2}{2}(1\gamma )^2}=T$$
(58)
in accordance with Eq. (2).
We wish to calculate the Green’s function
$`G(m,\tau )`$ $`=`$ $`<Tc_0(\tau )c_m^{}(0)>=`$ (59)
$`=`$ $`<Te^{g(a_0^{}(\tau )a_0(\tau ))(1\gamma \stackrel{~}{n}_0(\tau ))}\stackrel{~}{c}_0(\tau )\stackrel{~}{c}_m^{}(0)e^{g(a_m^{}(0)a_m(0))(1\gamma \stackrel{~}{n}_m(0))}>`$ (60)
We expand the exponentials in Eq. (30) using the operator relation
$$e^{g(a^{}a)(1\gamma \stackrel{~}{n}_{})}=e^{g(a^{}a)}+\stackrel{~}{n}_{}(e^{g(a^{}a)(1\gamma )}e^{g(a^{}a)})$$
(61)
and decouple averages over bosons and fermions, following Alexandrov and Ranninger. This leads to
$`G(m,\tau )`$ $`=`$ $`\sigma _0(m,\tau )<T\stackrel{~}{c}_0(\tau )\stackrel{~}{c}_m^{}(0)>`$ (62)
$`+`$ $`[(\sigma _1(m,\tau )\sigma _0(m,\tau )][<T\stackrel{~}{n}_0(\tau )\stackrel{~}{c}_0(\tau )\stackrel{~}{c}_m^{}(0)>+<T\stackrel{~}{c}_0(\tau )\stackrel{~}{c}_m^{}(0)\stackrel{~}{n}_m(0)>]`$ (63)
$`+`$ $`[\sigma _2(m,\tau )2\sigma _1(m,\tau )+\sigma _0(m,\tau )]<T\stackrel{~}{n}_0(\tau )\stackrel{~}{c}_0(\tau )\stackrel{~}{c}_m^{}(0)\stackrel{~}{n}_m(0)>`$ (64)
The boson Green’s functions are defined as
$$\sigma _{i+j}(m,\tau )=<e^{g(1\gamma )^i(a_0^{}(\tau )a_0(\tau ))}e^{g(1\gamma )^j(a_m^{}(0)a_m(0))}>$$
(65)
with $`i,j=0,1`$. At low temperatures they are given by
$$\sigma _\alpha (m,\tau )=S^{2\alpha }T^\alpha [1\delta _{m,0}+\delta _{m,0}e^{g^2(1\gamma )^\alpha D(\tau )}]$$
(67)
$$D(\tau )=[e^{\omega _0|\tau |}+2\frac{cosh\omega _0\tau }{e^{\beta \omega _0}1}]$$
(68)
and in frequency space by
$$\sigma _\alpha (m,i\omega _n)=S^{2\alpha }T^\alpha [\delta _{n,0}\beta +\delta _{m,0}\underset{l=1}{\overset{\mathrm{}}{}}\frac{2l\omega _0g^{2l}(1\gamma )^{\alpha l}}{l!(\omega _n^2+l^2\omega _0^2)}]$$
(69)
We next decouple the fermion averages with the same mean field procedure used to calculate the coherent part of the Green’s function , and calculate the Fourier transform
$$G(k,i\omega _n)=_0^\beta 𝑑\tau e^{i\omega _n\tau }\underset{m}{}e^{ikm}G_{0m}(\tau )$$
(70)
The complete Green’s function is
$$G(k,i\omega _n)=G_{coh}(k,i\omega _n)+G_{inc}(k,i\omega _n)$$
(71)
For each term of the coherent Green’s function, Eq. (18), there is a corresponding term in the incoherent Green’s function. All terms in the coherent Green’s function are of the form
$$G_{coh}^{\alpha ,s}(k,i\omega _n)=cS^{2\alpha }T^\alpha \frac{a_k}{i\omega _nsE_k}$$
(72)
with $`0\alpha 2`$ and $`s=+/1`$. The corresponding term in the incoherent Green’s function is
$$G_{inc}^{\alpha ,s}(k,i\omega _n)=cS^{2\alpha }T^\alpha \underset{l=1}{\overset{\mathrm{}}{}}\frac{g^{2l}(1\gamma )^{\alpha l}}{l!}\frac{1}{N}\underset{k^{}}{}a_k^{}[\frac{n_k^{}}{i\omega _ns(E_k^{}l\omega _0)}+\frac{1n_k^{}}{i\omega _ns(E_k^{}+l\omega _0)}]$$
(73)
The spectral function is obtained as usual from
$$A(k,\omega )=ImG(k,i\omega _n\omega +i\delta )$$
(74)
and results from Eqs. (18), (37)-(39). In particular, the lowest order normal part of the incoherent spectral function is given by
$`A_{inc}^n(k,\omega )`$ $`=`$ $`S^2\times {\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{g^{2l}}{l!}}[1+n[(1\gamma )^le^{\gamma g^2}1]]`$ (75)
$`\times `$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{k^{}}{}}[u_k^{}^2[(1n_k^{})\delta (\omega l\omega _0E_k^{})+n_k^{}\delta (\omega +l\omega _0E_k^{})]`$ (76)
$`+`$ $`v_k^{}^2[n_k^{}\delta (\omega l\omega _0+E_k^{})+(1n_k^{})\delta (\omega +l\omega _0+E_k^{})]]`$ (77)
and the lowest order anomalous contribution by
$`A_{inc}^a(k,\omega )`$ $`=`$ $`S^2\times 2f_0\times {\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{g^{2l}}{l!}}[1+n[(1\gamma )^le^{\gamma g^2}1]]`$ (78)
$`\times `$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{k^{}}{}}(u_k^{}v_k^{})\times [(1n_k^{})\delta (\omega l\omega _0E_k^{})+n_k^{}\delta (\omega +l\omega _0E_k^{})`$ (79)
$``$ $`n_k^{}\delta (\omega l\omega _0+E_k^{})(1n_k^{})\delta (\omega +l\omega _0+E_k^{})].`$ (80)
## VI Results for the spectral function
The spectral function for the models considered here is of the form
$$A(k,\omega )=A_{coh}(k,\omega )+A_{inc}(k,\omega )$$
(82)
$$A_{coh}(k,\omega )=Z_h\delta (\omega E_k)+Z_e\delta (\omega +E_k)$$
(83)
$$A_{inc}(k,\omega )=ImG_{inc}(k,\omega +i\delta )$$
(84)
where the quasiparticle weights $`Z_h`$ and $`Z_e`$, and the incoherent Green’s function $`G_{inc}`$, were discussed in the previous two sections. As seen in Sect. IV, the quasiparticle weight $`Z_e`$ acquires a $`positive`$ contribution from the onset of superconductivity. In a spectroscopic experiment usually only one side of the spectral function is sampled, as the other side is suppressed by the Fermi function. The quantity that will display the enhanced coherence due to undressing exhibited by $`Z_e`$ is
$$I_0(k,\omega )=A(k,\omega )f(\omega )$$
(85)
with $`f`$ the Fermi function. In an experiment there will typically be broadening from experimental resolution, which results in
$$I(k,\omega )=𝑑\omega ^{}F(\omega \omega ^{})I_0(k,\omega ^{})$$
(86)
being measured, with $`F(\omega )`$ a Gaussian with width $`\sigma _\omega `$. There could also be other sources of broadening of the $`\delta `$-functions in the expressions Eq. (43) from lifetime effects. Just as the spectral function, the measured spectrum Eq. (45) will have coherent and incoherent contributions
$$I(k,\omega )=I_{coh}(k,\omega )+I_{inc}(k,\omega )$$
(87)
arising from the coherent and incoherent parts of the spectral function respectively.
Figure 13 shows results for the coherent spectra at the Fermi energy for an underdoped ($`n=0.02`$, labeled $`ud`$), optimally doped ($`n=0.045`$, labeled $`op`$) and overdoped ($`n=0.1`$, labeled $`od`$) case for the parameter values used in Sect. V. For the underdoped case with the chemical potential below the bottom of the band the value of $`ϵ_k`$ at the bottom of the band was used. The dashed lines show the spectra in the normal state at $`T_c`$, and the full lines in the superconducting state at $`T=0.1T_c`$. For each doping, as superconductivity onsets the peaks shift to the left due to the opening of the superconducting gap. Furthermore, the peaks $`grow`$ in magnitude due to the behavior of $`Z_e`$ discussed in Sect. IV. As function of doping the peaks grow in magnitude both in the normal and superconducting state due to the enhanced coherence with increased number of carriers. The superconducting peak in the $`od`$ case is shifted to the right with respect to the $`op`$ case because the superconducting gap is smaller in that region (Fig. 1(b)).
When we include the incoherent part of the spectra the smaller normal state peak can become almost invisible. The results will depend of course on the specific parameters chosen to describe the incoherent background, and we are not suggesting that we are in a position to determine them from first principles. In Figure 14 we show results for a particular set of parameters for the generalized Holstein model. In addition to the parameters already discussed in Sect. IV, including the value of $`\mathrm{{\rm Y}}`$, the new parameters needed are $`S^2`$, $`\omega _0`$ and a broadening factor, given in the figure caption. Note that in the underdoped case (a) the peak in the normal state has become almost invisible, while a sharp peak and a dip are seen in the superconducting spectrum. The dip arises because the background term arising from the second term in Eq. (39) for $`a_k^{}=v_k^{}^2`$, $`s=1`$, is pushed to more negative energies as the superconducting gap opens. As the doping increases the normal state peak becomes more visible, and the overdoped case shows more conventional behavior. Note that the scale in the figures changes with doping and the magnitude of the peaks increases with doping.
Figure 15 shows the temperature dependence of the spectra for the overdoped case. The normal state peak is pushed back continuously as the superconducting gap opens up. In addition, for our case (a) the peak grows in magnitude. To highlight the difference with conventional BCS theory we show in Fig. 15(b) what is obtained with the same parameters in the absence of the term $`f_0`$ in Eq. (18). The peak here first becomes lower and then increases again as the temperature is lowered, but it is always lower or equal to the normal state peak. It is easily seen from the BCS formula that this property is generally true also for other values of the momenta.
Similarly figure 16 shows the temperature dependence of the spectra in the underdoped case. Here, rather than the peak moving continuously, a new peak grows in the superconducting state. The presence of two peaks has not been seen experimentally in photoemission to our knowledge, possibly because of experimental resolution. For the BCS case (b) the peak in the superconducting state is much smaller than in the normal state, while for our case (a) the opposite is true. Results for the temperature dependence in the optimally doped case show behavior intermediate between the overdoped and underdoped cases shown.
## VII The cuprates
We have seen in the previous sections that in systems where superconductivity arises from undressing there is a signature of the formation of the condensate in the single particle spectral function. Specifically, it arises in Eq. (17) from the term involving $`f_0`$, the on-site pair amplitude. Ding et al and Feng et al, discussing experimental results of angle resolved photoemission in cuprates have recently emphasized precisely that feature of the observed spectra, and correlated the growth of the peak in photoemission to quantities related to the superconducting condensate such as the superfluid density and the condensation energy. The spectra calculated within the present theory in the previous section resemble in several aspects the experimental observations in photoemission along the $`(\pi ,0)`$ direction.
Unfortunately, as the alert reader has undoubtedly noticed, the results presented in the previous section with negative $`\omega `$ in a hole representation correspond to $`hole`$ $`destruction`$, or $`electron`$ $`creation`$, that is, $`inverse`$ $`photoemission`$. It is for that case that the experiment would sample $`Z_e`$, the quasiparticle weight for electron creation. Instead, if we calculate spectra for direct photoemission we would find that quasiparticle peaks are $`suppressed`$ by the onset of superconductivity due to the behavior of $`Z_h`$ discussed in Sect. IV.
The present theory does not allow for a switch in the role of the weights $`Z_e`$ and $`Z_h`$: electron-hole asymmetry, of the sign assumed here, is central to the theory. Does this then imply that the theory is irrelevant for description of the cuprates?
We believe this is not the case. We propose that in fact, the photoemission experiments along the $`(\pi ,0)`$ direction close to the $`(\pi /a,0)`$ point sample the part of the spectral function discussed in the previous section, corresponding to hole destruction, or electron creation.
How can photoemission sample electron creation? Recall that in the theory of hole superconductivity the relevant orbitals are oxygen $`p\pi `$ orbitals in the planes. There are however also the oxygen $`p\sigma `$ orbitals, strongly hybridized with the Cu $`d_{x^2y^2}`$ orbitals. Suppose that in a photoemission experiment along the $`(\pi ,0)`$ direction the largest matrix element couples to destruction of a $`d_{x^2y^2}`$ electron. This would not directly couple to the band responsible for superconductivity; however, that process could induce the destruction of an oxygen hole in the $`p\pi `$ orbitals. The proposed situation is schematically depicted in Fig. 17, in an electron representation. Before the photon comes in there is one electron in each $`Cu^{++}`$ atom neighbor to a given $`O`$ atom, and two holes in the $`p\pi `$ orbital on that $`O`$ atom. We assume the energy level structure shown in Fig. 17: an electron from $`Cu^{++}`$ cannot ’fall’ onto the $`O`$ $`p\pi `$ orbital because it is Coulomb-repelled by the electron in the other $`Cu`$ atom. When a photon comes in and knocks out one of the electrons in a $`Cu`$, the other electron can ’fall’ onto the $`O`$ $`p\pi `$ orbital, thus destroying an $`O`$ hole and sampling the quasiparticle weight $`Z_e`$.
It is clear that this qualitative explanation needs further elaboration and experimental confirmation to be convincing. Nevertheless we also point out that it suggests an explanation for why the sharp peaks seen in photoemission along the $`(\pi ,0)`$ direction are not seen along the $`(\pi ,\pi )`$ direction: since Cu $`d_{x^2y^2}`$ orbitals point along the principal axis in the planar square lattice, the coupling to the photon along the $`(\pi ,\pi )`$ direction is likely to be much smaller. For that direction the larger coupling may be to the $`O`$ $`p\pi `$ band itself, in which case inverse rather than direct photoemission would show the enhanced coherence. It is possible that some indication of this effect may have already been seen in tunneling experiments.
## VIII Conclusions
In this paper we have continued exploring the consequences of the physical principle proposed in I: that, in at least some electronic materials in nature, the dressing of quasiparticle carriers is a function of the local carrier concentration, and becomes smaller as the local carrier concentration increases. This physical principle leads to superconductivity occuring in these systems because of lowering of the carriers kinetic energy upon pairing. The superconducting transition, and many of the features of the superconducting state, have already been discussed earlier within the theory of hole superconductivity.
In this paper we explored the consequences of this principle for the single particle Green’s function in the superconducting state. The central result of this paper, Eq. (18), demonstrates that formation of the superfluid condensate will influence the behavior of the single particle spectral function. Eq. (18) is thus the generalization of the BCS spectral function for systems where superconductivity is driven by undressing. Surprisingly, the results show that the enhanced coherence in the superconducting state is displayed in the quasiparticle weight for electron creation but not for electron destruction. We calculated the behavior of the quasiparticle weights as function of temperature, doping dependence and momentum, and highlighted their differences with conventional BCS theory.
Furthermore, we discussed the calculation of the full spectral function including the incoherent contribution for one particular model where superconductivity occurs through undressing, a generalized Holstein model. Our calculation was performed within the Lang-Firsov approximation, and it should be interesting to see whether the qualitative results survive a more exact treatment.
Results for the full spectral function showed several features that resemble experimental observations in photoemission experiments in high $`T_c`$ cuprates, in particular enhanced coherence, as displayed by the quasiparticle peak in the spectra, when the system enters the superconducting state and as the carrier concentration increases both in the normal and in the superconducting state.
This study was strongly motivated by the beautiful experimental results and insightful analysis of the photoemission experiments. Thus it is perhaps disappointing that at the end of the day our calculation predicts these effects, in the simplest one-band model, arising in $`inverse`$ rather than in direct photoemission. Thus some readers may conclude that our calculation is not more than an academic exercise. However, as discussed in Sect. VII, we believe there is a plausible scenario by which the spectral weight for electron creation would be sampled in the photoemission experiments in the cuprates.
While the theory discussed here predicts s-wave rather than d-wave superconductivity we believe it is remarkable how many of the features that appear to be part of the phenomenology of high $`T_c`$ cuprates it exhibits, as a consequence of the $`single`$ $`assumption`$ of a large value of the undressing parameter $`\mathrm{{\rm Y}}`$: (1) incoherence in the normal state at low hole concentration; (2) increased coherence with doping in the normal state; (3) transition to superconductivity for low doping, dissappearing for high doping, and bell-shaped $`T_c`$ versus hole concentration; (4) increased coherence as the system goes superconducting; (5) superconducting transition driven by kinetic energy lowering, optical sum rule violation; (6) non-decrease of the quasiparticle gap at low hole density when $`T_c`$ is going to zero. This latter feature arises in our model from the fact that as the hole concentration decreases and the band becomes narrower the chemical potential falls below the bottom of the band; we believe that many of the unusual properties of underdoped cuprates follow from this simple fact, and in particular that the observed pseudogap is simply the energy difference between the bottom of the band and the chemical potential.
If the theory of hole undressing discussed here describes the cuprates it is likely that it is more generally applicable, because it is based on very general principles. In this regard we note that one of the paradoxes of the conventional explanation of superconductivity is that it is thought to originate in an electron-boson (the electron phonon) coupling that $`opposes`$ $`conductivity`$, i.e. gives rise to resistivity, in the normal state. In a sense the present theory eliminates this paradox. Coupling to a boson is certainly necessary, and that coupling gives rise to enhanced resistivity in the normal state due to enhanced effective mass, but superconductivity arises from a process whereby the coupling to that boson is $`reduced`$ as carriers pair and the system becomes superconducting. The old paradox is however replaced by a new one, that in order for heavily dressed ’confined’ carriers to become less dressed, or ’freer’, it is necessary for them to $`bind`$ in Cooper pairs.
We also note that the principle on which the present theory is based, that an increase in the local hole occupation causes undressing, is likely to be more general than as expressed by Eq. (2): rather than just by the same site occupation, undressing may also be enhanced by hole occupation of neighboring sites, and also by neighboring bond occupation. The possible implications of this for superconductivity and other instabilities of metals will be discussed in future work.
If indeed the essential physics of high $`T_c`$ is hole undressing, what makes a material a high $`T_c`$ superconductor? Presumably, the fact that quasiparticles are heavily dressed in the normal state together with the fact that the undressing process that occurs when the local carrier concentration increases is particularly efficient. Both of those facts are necessary conditions for high $`T_c`$ superconductivity by giving rise to a large $`\mathrm{{\rm Y}}`$ parameter. We will not discuss here what aspects of the chemistry of the cuprates would favor this situation. However, conversely, we may conclude that the reason for a material $`not`$ being a high $`T_c`$ superconductor would be a small value of the parameter $`\mathrm{{\rm Y}}`$, either because quasiparticles are $`not`$ heavily dressed in the normal state (e.g. the case of Aluminum), or, because the quasiparticle dressing in the normal state may not be strongly dependent on the local carrier concentration (e.g. the case of ’heavy fermion’ systems).
## ACKNOWLEDGMENTS
The author is grateful to F. Driscoll for the donation of a computer where the calculations reported here were performed.
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# Finding Z' Bosons Coupled Preferentially to the Third Family at CERN LEP and the Fermilab Tevatron
## 1 Introduction
The Standard Model of particle physics gives an excellent description of physics at the energy scales probed to date. Nonetheless, it does not explain the origins of the masses of the electroweak gauge bosons and the elementary fermions, and must be regarded as a low-energy effective field theory. For a description of the dynamics underlying the generation of mass, we must turn to physics beyond the Standard Model.
Much recent theoretical work on the question of why the top quark is so heavy has suggested that the cause could be additional gauge interactions that single out the third generation fermions. A number of interesting models along these lines extend one (or more) of the Standard Model’s $`\mathrm{SU}(N)`$ gauge groups into an $`\mathrm{SU}(N)\times \mathrm{SU}(N)`$ gauge structure . In general, fermions of the third generation transform under one $`\mathrm{SU}(N)`$ group and those of the first and second generations transform under the other one. When the $`\mathrm{SU}(N)\times \mathrm{SU}(N)`$ spontaneously breaks to its diagonal subgroup, the broken generators correspond to a set of massive $`\mathrm{SU}(N)`$ gauge bosons that couple to fermions of different generations with different strengths.
Many of these models predict the presence of massive $`\mathrm{Z}^{}`$ bosons that couple preferentially to the third-generation fermions. Some theories include an extended $`\mathrm{SU}(2)\times \mathrm{SU}(2)`$ structure for the weak interactions; generally, the first two generations of fermions are charged under the weaker $`\mathrm{SU}(2)`$ and the third generation feels the other, stronger, $`\mathrm{SU}(2)`$ gauge force. Examples include non-commuting extended technicolor (NCETC) models and topflavor models . For many of these models, precision data suggest that the $`\mathrm{Z}^{}`$ must be relatively heavy; but in some non-commuting extended technicolor models, a mass as low as $`400\mathrm{GeV}`$ is not precluded . There are also theories in which a $`\mathrm{Z}^{}`$ boson arises due to an extra $`\mathrm{U}(1)`$ group coupled preferentially to the third-generation quarks (and possibly leptons). Examples are the topcolor-assisted technicolor and flavor-universal topcolor-assisted-technicolor models. In these models, depending on the charge assignments of the ordinary and technifermions, constraints from FCNC and precision electroweak corrections can allow the $`\mathrm{Z}^{}`$ boson to be as light as $`290\mathrm{GeV}`$ .
More generally, electroweak scale $`\mathrm{Z}^{}`$ bosons are also present in string theories , and string-inspired models often yield non-universal couplings for the $`\mathrm{Z}^{}`$. A recent analysis of electroweak precision data actually gives a strong indication of the presence of an extra $`\mathrm{Z}^{}`$ boson with the fits favoring non-universal couplings to the third family.
The literature already contains a number of suggestions about how experiment can set stronger limits on these $`\mathrm{Z}^{}`$ bosons. Note that bounds on $`\mathrm{Z}^{}`$ bosons which do not couple preferentially to the third generation are not directly applicable. For example, in models with extended weak interactions, the presence of $`\mathrm{W}^{}`$ bosons with mass below about 1.5 TeV would cause an enhancement of single top quark production large enough to be visible at the Tevatron’s Run II experiments ; this would provide indirect evidence of similarly light $`\mathrm{Z}^{}`$ bosons. The Run I Tevatron experiments have searched for topcolor $`\mathrm{Z}^{}`$ bosons in $`\mathrm{b}\overline{\mathrm{b}}`$ and $`\mathrm{t}\overline{\mathrm{t}}`$ final states. In these processes, the backgrounds are of QCD strength. As a result, no limit has been set from the $`\mathrm{b}\overline{\mathrm{b}}`$ channel and the recent limit in the $`t\overline{t}`$ channel ($`M_\mathrm{Z}^{}>650\mathrm{GeV}`$) is for a $`\mathrm{Z}^{}`$ that couples only to hadrons and is quite narrow, $`\mathrm{\Gamma }_\mathrm{Z}^{}=.012M_\mathrm{Z}^{}`$. These searches should have greater reach in Run II, due to the higher luminosity and improved detectors. Flavor-changing neutral current effects can also yield constraints on $`\mathrm{Z}^{}`$ bosons with non-universal couplings .
This paper discusses two additional methods of searching for $`\mathrm{Z}^{}`$ bosons that couple primarily to the third family fermions. We first show how existing LEP and Tevatron bounds on the scale of quark-lepton compositeness can be adapted to provide lower bounds on the mass of these $`\mathrm{Z}^{}`$ bosons. We then analyze, the possibility of searching at Tevatron Run II for $`\mathrm{Z}^{}`$ bosons in the channel $`\mathrm{p}\overline{\mathrm{p}}\mathrm{Z}^{}\tau \tau \mathrm{e}\mu `$ in order to exploit the strong $`\mathrm{Z}^{}\tau \tau `$ coupling and the low backgrounds for $`\mathrm{e}\mu `$ final states.
In Section 2, we will review the properties of the $`\mathrm{Z}^{}`$ boson arising in models with an extended weak gauge group and display the existing limits from electroweak precision data. In Section 3 we extract limits on these $`\mathrm{Z}^{}`$ bosons from the LEP and Tevatron compositeness bounds. Section 4 focuses on the Run II search in leptonically-decaying pair-produced taus. We then in Section 5 show how our results are modified for $`\mathrm{Z}^{}`$ bosons in models with an extended hypercharge group, and we mention a few additional search channels which may help improve the reach of Run II in Section 6. Our conclusions are presented in Section 7.
## 2 $`\mathrm{Z}^{}`$ Bosons From Extended Weak Interactions
### 2.1 General properties of $`\mathrm{SU}(2)`$ $`\mathrm{Z}^{}`$ bosons
The models of interest to us include the usual complement of quarks and leptons, along with standard strong and hypercharge interactions. The new physics lies in the weak interactions, which are governed by a pair of $`\mathrm{SU}(2)`$ gauge groups:
$$\mathrm{SU}(2)_h\times \mathrm{SU}(2)_{\mathrm{}}.$$
(2.1)
The $`\mathrm{SU}(2)_h`$ group governs the weak interactions for the third generation (heavy) fermions; the left-handed fermions transform as doublets and the right-handed ones, as singlets under this group. Similarly, the $`\mathrm{SU}(2)_{\mathrm{}}`$ group couples to the first and second generation (light) fermions, whose charges under $`\mathrm{SU}(2)_{\mathrm{}}`$ are as in the standard model. The extended weak group (Equation 2.1) is broken to its diagonal subgroup, $`\mathrm{SU}(2)_L`$ at energy scale $`u`$ by a (composite) scalar field $`\sigma `$, charged under $`\mathrm{SU}(2)_h\times \mathrm{SU}(2)_{\mathrm{}}\times \mathrm{U}(1)_Y`$ as:
$`\sigma `$ $`(2,2)_0,`$ $`\sigma `$ $`=\left(\begin{array}{cc}u& 0\\ 0& u\end{array}\right).`$ (2.4)
The final step in electroweak symmetry breaking could, in principle, proceed through a condensate charged under $`\mathrm{SU}(2)_{\mathrm{}}`$ or one charged under $`\mathrm{SU}(2)_h`$ – or one of each . The first option allows the $`\mathrm{Z}^{}`$ boson to be the lightest , making it the option of greatest phenomenological interest. Hence, we assume that the symmetry breaking $`\mathrm{SU}(2)_L\times \mathrm{U}(1)_Y\mathrm{U}(1)_{\mathrm{em}}`$ is due to a (composite) scalar
$`\mathrm{\Phi }`$ $`(1,2)_{1/2},`$ $`\mathrm{\Phi }`$ $`=\left(\begin{array}{c}0\\ v/\sqrt{2}\end{array}\right).`$ (2.7)
The generator of the $`\mathrm{U}(1)_{\mathrm{em}}`$ group is the electric charge operator:
$$Q=T_{3h}+T_3\mathrm{}+Y,$$
(2.8)
and the corresponding photon eigenstate is:
$$\mathrm{A}^\mu =\mathrm{sin}\theta \left(\mathrm{cos}\varphi \mathrm{W}_{3h}^\mu +\mathrm{sin}\varphi \mathrm{W}_3\mathrm{}^\mu \right)+\mathrm{cos}\theta \mathrm{X}^\mu $$
(2.9)
where $`\theta `$ is the usual weak mixing angle and $`\varphi `$ is an additional mixing angle occasioned by the presence of two weak gauge groups. We can therefore relate the gauge couplings and mixing angles as follows:
$$\begin{array}{c}g_h=\frac{e}{\mathrm{cos}\varphi \mathrm{sin}\theta }=\frac{g}{\mathrm{cos}\varphi }\\ g_{\mathrm{}}=\frac{e}{\mathrm{sin}\varphi \mathrm{sin}\theta }=\frac{g}{\mathrm{sin}\varphi }\\ g_Y=\frac{e}{\mathrm{cos}\theta }.\end{array}$$
(2.10)
For brevity we will write $`s_\varphi \mathrm{sin}\varphi `$ and $`c_\varphi \mathrm{cos}\varphi `$.
In diagonalizing the mass matrix for the neutral gauge bosons, it is convenient to first transform to an intermediate basis ,
$$\mathrm{Z}_1^\mu =\mathrm{cos}\theta \left(c_\varphi \mathrm{W}_{3h}^\mu +s_\varphi \mathrm{W}_3\mathrm{}^\mu \right)\mathrm{sin}\theta \mathrm{X}^\mu $$
(2.11)
$$\mathrm{Z}_2^\mu =s_\varphi \mathrm{W}_{3h}^\mu +c_\varphi \mathrm{W}_3\mathrm{}^\mu ,$$
(2.12)
where the covariant derivative neatly separates into standard and non-standard contributions:
$$D^\mu =^\mu i\frac{g}{\mathrm{cos}\theta }\mathrm{Z}_1^\mu \left(T_{3h}+T_3\mathrm{}\mathrm{sin}^2\theta Q\right)ig\mathrm{Z}_2^\mu \left(\frac{s_\varphi }{c_\varphi }T_{3h}+\frac{c_\varphi }{s_\varphi }T_3\mathrm{}\right).$$
(2.13)
In terms of these states, the neutral mass eigenstates (the $`\mathrm{Z}^0`$ and $`\mathrm{Z}^{}`$ states) are given, to leading order in $`1/x=v^2/u^2`$, by the superpositions:
$$\left(\begin{array}{c}\mathrm{Z}^0\\ \mathrm{Z}^{}\end{array}\right)\left(\begin{array}{cc}1& \frac{c_\varphi ^3s_\varphi }{x\mathrm{cos}\theta }\\ \frac{c_\varphi ^3s_\varphi }{x\mathrm{cos}\theta }& 1\end{array}\right)\left(\begin{array}{c}\mathrm{Z}_1\\ \mathrm{Z}_2\end{array}\right).$$
(2.14)
Expanding the covariant derivative in Equation 2.13 in terms of the mass eigenstates $`\mathrm{Z}^0`$ and $`\mathrm{Z}^{}`$, we find that, to order $`1/x`$:
$$\begin{array}{c}D_\mu =_\mu \frac{ig}{\mathrm{cos}\theta }\mathrm{Z}_\mu ^0\left[\left(1\frac{c_\varphi ^4}{x}\right)T_3\mathrm{}+\left(1+\frac{c_\varphi ^2s_\varphi ^2}{x}\right)T_{3h}\mathrm{sin}^2\theta Q\right]\hfill \\ \hfill ig\mathrm{Z}_{}^{}{}_{\mu }{}^{}\left[\left(\frac{c_\varphi }{s_\varphi }+\frac{c_\varphi ^3s_\varphi }{x\mathrm{cos}^2\theta }\right)T_3\mathrm{}+\left(\frac{s_\varphi }{c_\varphi }+\frac{c_\varphi ^3s_\varphi }{x\mathrm{cos}^2\theta }\right)T_{3h}\mathrm{sin}^2\theta \left(\frac{c_\varphi ^3s_\varphi }{x\mathrm{cos}^2\theta }\right)Q\right].\end{array}$$
(2.15)
For large $`s_\varphi `$ the $`\mathrm{Z}^0`$ boson can maintain a nearly standard model coupling to all fermion species, while the $`\mathrm{Z}^{}`$ boson has a greatly enhanced coupling to the third generation fermions. Moreover, we will see that the $`1/x`$ corrections are small in the phenomenologically interesting region of parameter space, so that the $`\mathrm{Z}^{}`$ boson essentially couples only to left-handed fermions.
To leading order, the mass of the $`\mathrm{Z}^{}`$ boson in the region where $`s_\varphi `$ exceeds $`c_\varphi `$ is
$$M_\mathrm{Z}^{}^2=\left(\frac{ev}{2\mathrm{sin}\theta }\right)^2\frac{x}{s_\varphi ^2c_\varphi ^2}=M_{\mathrm{W}_{\text{SM, tree}}}^2\frac{x}{s_\varphi ^2c_\varphi ^2}.$$
(2.16)
and, to this order, the mass of the $`\mathrm{W}^{}`$ boson is the same. The masses of the $`\mathrm{Z}^0`$ and $`W^\pm `$ bosons are shifted from their tree-level standard model values by identical multiplicative factors, so that there is no change in the predicted value of the $`\rho `$ parameter at order $`1/x`$ .
The width of the $`\mathrm{Z}^{}`$, to leading (here, zeroth) order in $`1/x`$ and in the region where $`s_\varphi >c_\varphi `$ is<sup>1</sup><sup>1</sup>1A fermionic species, $`f`$, contributes to the width of this gauge boson as
$$\mathrm{\Gamma }_{f,\mathrm{Z}^{}}=C_f\frac{\left(M_\mathrm{Z}^{}^24m_f^2\right)^{1/2}}{48\pi }\left[\left(g_{f,R}+g_{f,L}\right)^2\left(1+\frac{2m_f^2}{M_\mathrm{Z}^{}^2}\right)+\left(g_{f,R}g_{f,L}\right)^2\left(1\frac{4m_f^2}{M_\mathrm{Z}^{}^2}\right)\right],$$
(2.17) where $`C_f`$ is a color factor (1 for leptons, 3 for quarks), $`m_f`$ is the fermion mass, and $`g_{f,R}`$, $`g_{f,L}`$ are the right and left handed couplings of the fermion.
$$\frac{\mathrm{\Gamma }_\mathrm{Z}^{}}{\alpha _2M_\mathrm{Z}^{}}=\frac{2}{3}\left(\frac{c_\varphi }{s_\varphi }\right)^2+\left[\frac{5}{24}+\frac{1}{8}\left(1\frac{2m_\mathrm{t}^2}{M_\mathrm{Z}^{}^2}\right)\left(1\frac{4m_\mathrm{t}^2}{M_\mathrm{Z}^{}^2}\right)^{1/2}\mathrm{\Theta }(M_\mathrm{Z}^{}2m_\mathrm{t})\right]\left(\frac{s_\varphi }{c_\varphi }\right)^2,$$
(2.18)
where we have included the effects of the $`\mathrm{t}`$-quark mass, while taking the other fermion masses to be zero. The step function ensures that we only include the top contribution when the $`\mathrm{Z}^{}`$ mass is above the top threshold. To this order, the $`\mathrm{Z}^{}`$ boson couples neither to right-handed fermions nor to $`\mathrm{Z}^0`$$`\mathrm{Z}^0`$ nor to $`\mathrm{W}^+`$$`\mathrm{W}^{}`$. Effects from the composite scalars are not substantial. Figure 1 shows the $`\mathrm{Z}^{}`$ width for $`\mathrm{Z}^{}`$ masses between 350 GeV and 750 GeV as a function of the mixing angle $`s_\varphi `$.
The width of the $`\mathrm{Z}^{}`$ falls to a minimum at approximately $`s_\varphi =0.8`$. The width then grows rapidly as $`s_\varphi `$ becomes larger.
### 2.2 Light $`\mathrm{SU}(2)`$ $`\mathrm{Z}^{}`$ in models of dynamical symmetry breaking
A gauge structure of this kind has been proposed within the context of models of dynamical electroweak symmetry breaking and models in which the vacuum expectation value of a weakly coupled scalar boson breaks the electroweak symmetry . A class of models in which the precision electroweak data allow the $`\mathrm{Z}^{}`$ and $`\mathrm{W}^{}`$ bosons to be particularly light are the non-commuting extended technicolor models . The electroweak symmetry breaking pattern of non-commuting extended technicolor is characterized by a three-stage breakdown from the unbroken, high energy theory to the low energy electromagnetic gauge structure:
$$\begin{array}{cc}\hfill \mathrm{G}_{ETC}\times \mathrm{SU}(2)_{\mathrm{}}\times \mathrm{U}(1)^{}& \stackrel{f}{}\hfill \\ \hfill \mathrm{G}_{TC}\times \mathrm{SU}(2)_h\times \mathrm{SU}(2)_{\mathrm{}}\times \mathrm{U}(1)_Y& \stackrel{u}{}\hfill \\ \hfill \mathrm{G}_{TC}\times \mathrm{SU}(2)_L\times \mathrm{U}(1)_Y& \stackrel{v}{}\hfill \\ \hfill \mathrm{G}_{TC}\times \mathrm{U}(1)_{\mathrm{em}}& .\hfill \end{array}$$
(2.19)
At the scale $`f`$, the extended technicolor gauge group $`\mathrm{G}_{ETC}`$ breaks to the technicolor gauge group and the $`\mathrm{SU}(2)_h`$ (heavy) group. The two $`\mathrm{SU}(2)`$ groups mix and break to their diagonal subgroup at the scale $`u`$. The final breaking of the remaining electroweak symmetry is accomplished at scale $`v`$. If the condensate $`\mathrm{\Phi }`$ responsible for electroweak symmetry breaking at scale $`v`$ is charged under $`\mathrm{SU}(2)_{\mathrm{}}`$ (rather than $`\mathrm{SU}(2)_h`$), the resulting masses of the $`\mathrm{Z}^{}`$ and $`\mathrm{W}^{}`$ bosons can be as low as 400 GeV, as illustrated in Figure 2.
New gauge bosons with such small masses are of great phenomenological interest, as they are within the kinematic reach of Tevatron Run II experiments and their indirect effects may be apparent at LEP 2. In other models with the extended $`\mathrm{SU}(2)_h\times \mathrm{SU}(2)_{\mathrm{}}`$ gauge structure, existing lower limits on the gauge boson masses tend to be of order 1 - 1.5 TeV .
Note that in the context of non-commuting extended technicolor models, the coupling $`g_h`$ is essentially the value of the technicolor coupling at scale $`f`$. We therefore expect $`g_h`$ to be large compared to the weak coupling $`g`$, so that the value of $`c_\varphi ^2`$ (from Equation 2.10) should be relatively small. However, if $`c_\varphi ^2`$ is too small, $`g_h`$ will be above the critical value at which the chiral symmetries of the technifermions break. Thus, as discussed in , we must restrict $`s_\varphi ^2`$ to be smaller than about 0.95, hence $`s_\varphi 0.975`$ (vertical dashed line in Figure 2).
## 3 Limits on an $`\mathrm{SU}(2)`$ $`\mathrm{Z}^{}`$ from Compositeness Searches
For experiments done at energies below the mass of a $`\mathrm{Z}^{}`$ boson, one can approximate the contribution of the $`\mathrm{Z}^{}`$ to fermion-fermion scattering as a contact interaction whose scale is set by the mass of the $`\mathrm{Z}^{}`$ boson. Thus, published experimental limits on compositeness can set a lower bound on $`M_\mathrm{Z}^{}`$.
### 3.1 LEP Data
The LEP experiments ALEPH and OPAL and have recently published limits on contact interactions. Following the notation of , they write the effective Lagrangian for the four-fermion contact interaction in the process $`\mathrm{e}^+\mathrm{e}^{}f\overline{f}`$ as
$$_{\text{contact}}=\frac{g^2}{\mathrm{\Lambda }^2(1+\delta )}\underset{i,j=L,R}{}\eta _{ij}\left(\overline{\mathrm{e}}_i\gamma _\mu \mathrm{e}_i\right)\left(\overline{f}_j\gamma ^\mu f_j\right)$$
(3.1)
where $`\delta =1`$ if $`f`$ is an electron and $`\delta =0`$ otherwise. The values of the coefficients $`\eta _{ij}`$ set the chirality structure of the interaction being studied; OPAL and ALEPH study a number of cases where one of the $`\eta _{ij}`$ is equal to $`\pm 1`$ and the others are zero. Following the convention of taking $`g^2/4\pi =1`$, they determine a lower bound on the scale $`\mathrm{\Lambda }`$ associated with each type of new physics. Of particular interest to us are their limits on contact interactions where the final-state fermions $`f`$ belong to the third generation: $`\mathrm{e}^+\mathrm{e}^{}\mathrm{b}\overline{\mathrm{b}}`$ and $`\mathrm{e}^+\mathrm{e}^{}\tau ^+\tau ^{}`$. Among the limits published by ALEPH and OPAL , those of interest to us are
$`\mathrm{\Lambda }(f=\tau ,\eta _{LL}=+1)`$ $`>\{\begin{array}{cc}3.9\mathrm{TeV}\hfill & \text{ALEPH}\hfill \\ 3.8\mathrm{TeV}\hfill & \text{OPAL}\hfill \end{array}`$ (3.2)
$`\mathrm{\Lambda }(f=\mathrm{b},\eta _{LL}=+1)`$ $`>\{\begin{array}{cc}5.6\mathrm{TeV}\hfill & \text{ALEPH}\hfill \\ 4.0\mathrm{TeV}\hfill & \text{OPAL}.\hfill \end{array}`$ (3.3)
At energies well below the mass of the $`\mathrm{Z}^{}`$ boson, its exchange in the process $`\mathrm{e}^+\mathrm{e}^{}f\overline{f}`$ where $`f`$ is a $`\tau `$ lepton or $`\mathrm{b}`$ quark may be approximated by the contact interaction
$$_{\mathrm{NC}}\frac{e^2}{\mathrm{sin}^2\theta M_\mathrm{Z}^{}^2}\left(\frac{c_\varphi }{2s_\varphi }\left(\overline{\mathrm{e}}_L\gamma _\mu \mathrm{e}_L\right)\right)\left(\frac{s_\varphi }{2c_\varphi }\left(\overline{f}_L\gamma ^\mu f_L\right)\right),$$
(3.4)
based on the $`\mathrm{Z}^{}`$-fermion couplings in Equation 2.15. Comparing this with the contact interactions studied by LEP, we find
$$M_\mathrm{Z}^{}=\mathrm{\Lambda }\sqrt{\frac{\alpha _{\mathrm{em}}}{4\mathrm{sin}^2\theta }}.$$
(3.5)
The limits from $`\tau `$-pair production are, then,
$$M_\mathrm{Z}^{}>\{\begin{array}{cc}365\mathrm{GeV}\hfill & \text{ALEPH}\hfill \\ 355\mathrm{GeV}\hfill & \text{OPAL},\hfill \end{array}$$
(3.6)
and those from $`\mathrm{b}\overline{\mathrm{b}}`$ production are
$$M_\mathrm{Z}^{}>\{\begin{array}{cc}523\mathrm{GeV}\hfill & \text{ALEPH}\hfill \\ 375\mathrm{GeV}\hfill & \text{OPAL}.\hfill \end{array}$$
(3.7)
As Figure 2 illustrates, the limits are comparable to the previous lower bound on $`M_\mathrm{Z}^{}`$ from precision electroweak data in the case where $`s_\varphi `$ is large; for small $`s_\varphi `$ the earlier limits remain stronger. As additional data from the other experiments or higher energies becomes available, the lower bound on $`M_\mathrm{Z}^{}`$ can be updated by using the new lower bound on $`\mathrm{\Lambda }`$ in Equation 3.5.
### 3.2 FNAL Data
The CDF and DØ Collaborations have each searched for the low energy effects of quark-lepton contact interactions on dilepton production in $`110\mathrm{pb}^1`$ of data taken at $`\sqrt{s}=1.8\mathrm{TeV}`$ . Since this process is dominated by first and second generation fermions, the limits on our $`\mathrm{Z}^{}`$ bosons tend to be weaker than those derived from LEP data.
In their analysis, the CDF Collaboration described the effective four-fermi interactions of the first generation fermions due to new physics by an effective Lagrangian including the terms:
$$\begin{array}{c}_{\mathrm{EQ}}\xi _{LL}^0\left(\overline{\mathrm{E}}_L\gamma _\mu \mathrm{E}_L\right)\left(\overline{\mathrm{Q}}_L\gamma ^\mu \mathrm{Q}_L\right)+\xi _{LL}^1\left(\overline{\mathrm{E}}_L\gamma _\mu \tau _a\mathrm{E}_L\right)\left(\overline{\mathrm{Q}}_L\gamma ^\mu \tau _a\mathrm{Q}_L\right)\hfill \\ \hfill +\xi _{LR}^u\left(\overline{\mathrm{E}}_L\gamma _\mu \mathrm{E}_L\right)\left(\overline{\mathrm{u}}_R\gamma ^\mu \mathrm{u}_R\right)+\xi _{LR}^d\left(\overline{\mathrm{E}}_L\gamma _\mu \mathrm{E}_L\right)\left(\overline{\mathrm{d}}_R\gamma ^\mu \mathrm{d}_R\right)+\xi _{RL}^e\left(\overline{\mathrm{e}}_R\gamma _\mu \mathrm{e}_R\right)\left(\overline{\mathrm{Q}}_L\gamma ^\mu \mathrm{Q}_L\right)\\ \hfill +\xi _{RR}^u\left(\overline{\mathrm{e}}_R\gamma _\mu \mathrm{e}_R\right)\left(\overline{\mathrm{u}}_R\gamma ^\mu \mathrm{u}_R\right)+\xi _{RR}^d\left(\overline{\mathrm{e}}_R\gamma _\mu \mathrm{e}_R\right)\left(\overline{\mathrm{d}}_R\gamma ^\mu \mathrm{d}_R\right),\end{array}$$
(3.8)
where $`\mathrm{Q}_L(\mathrm{u},\mathrm{d})_L`$, $`\mathrm{E}_L(\nu _\mathrm{e},\mathrm{e})`$, and the subscripts $`L`$ and $`R`$ denote the left and right helicity projections. The coefficients $`\xi _{ij}`$ are related to the scale of new physics, $`\mathrm{\Lambda }_{ij}`$, as $`\xi _{ij}=g_0^2/\mathrm{\Lambda }_{ij}`$, where $`g_0^2`$ is an effective coupling which grows strong at the compositeness scale: $`g_0^2(\mathrm{\Lambda })/4\pi =1`$. The analysis searched for deviations in the dilepton spectrum from the standard model prediction; the absence of such deviations enabled them to set a lower bound on the scale of the new interactions.
The CDF analysis included fermions beyond the first generation by assuming a kind of universality: electrons and muons have identical contact interactions, all up-type quarks behave alike, all down-type quarks behave alike. They derived separate limits on contact interactions involving different combinations of fermions; for example, assuming that the only contact interaction was one between left-handed muons (electrons) and up-type quarks, they found at $`95\%`$ confidence
$`\mathrm{\Lambda }(\mu _L;\mathrm{u}_L,\mathrm{c}_L,\mathrm{t}_L)`$ $`>4.1\mathrm{TeV}`$ (3.9)
$`\mathrm{\Lambda }(\mathrm{e}_L;\mathrm{u}_L,\mathrm{c}_L,\mathrm{t}_L)`$ $`>3.7\mathrm{TeV}.`$ (3.10)
The presence of a massive $`\mathrm{Z}^{}`$ boson in our model gives rise to four-fermion contact interactions that include the terms
$$_{\mathrm{NC}}\frac{e^2}{\mathrm{sin}^2\theta M_\mathrm{Z}^{}^2}\left(\frac{c_\varphi }{2s_\varphi }\right)^2\left(\overline{\mathrm{e}}_L\gamma _\mu \mathrm{e}_L+\overline{\mu }_L\gamma _\mu \mu _L\right)\left(\overline{\mathrm{u}}_L\gamma ^\mu \mathrm{u}_L+\overline{\mathrm{d}}_L\gamma ^\mu \mathrm{d}_L+\overline{\mathrm{c}}_L\gamma ^\mu \mathrm{c}_L+\overline{\mathrm{s}}_L\gamma ^\mu \mathrm{s}_L\right).$$
(3.11)
In other words, the contact interactions among left-handed fermions in the first two generations all have the same coefficient. The interaction strength for third-generation fermions is different, as seen from Equation 2.15. Thus, the CDF analysis applies to our $`\mathrm{Z}^{}`$ boson only to the extent that initial-state third-generation quarks do not contribute to dilepton production. Since the top quark’s parton distribution function is approximately zero, we can reasonably apply the CDF limits for lepton/up-type-quark contact interactions to our model.
Comparing the interactions 3.8 and 3.11, we find the relationship between $`M_\mathrm{Z}^{}`$ and $`\mathrm{\Lambda }`$ is
$$M_\mathrm{Z}^{}=\mathrm{\Lambda }\left(\frac{c_\varphi }{s_\varphi }\right)\sqrt{\frac{\alpha _{\mathrm{em}}}{4\mathrm{sin}^2\theta }}.$$
(3.12)
The CDF bounds in Equation 3.10 imply
$$M_\mathrm{Z}^{}>\{\begin{array}{cc}\left(\frac{c_\varphi }{s_\varphi }\right)\times 380\mathrm{GeV}\hfill & \text{from dimuons}\hfill \\ & \\ \left(\frac{c_\varphi }{s_\varphi }\right)\times 345\mathrm{GeV}\hfill & \text{from dielectrons}.\hfill \end{array}$$
(3.13)
These limits are comparable to those from the LEP data for $`s_\varphi c_\varphi `$, but become significantly weaker at large $`s_\varphi `$.
The DØ Collaboration has performed a similar analysis for high energy dielectron production , but assuming only first-generation fermions participate in the contact interactions (i.e. the terms explicitly written in Equation 3.8. Since they include only first-generation fermions, their limit
$$\mathrm{\Lambda }(\mathrm{e}_L;\mathrm{u}_L,\mathrm{d}_L)>4.2\mathrm{TeV}$$
(3.14)
applies directly to our $`\mathrm{Z}^{}`$ boson, yielding the constraint
$$M_\mathrm{Z}^{}>\left(\frac{c_\varphi }{s_\varphi }\right)\times 390\mathrm{GeV}$$
(3.15)
which is comparable to the result obtained by CDF.
The TeV 2000 Group Report projects that the limits on the scale of quark-lepton compositeness will be increased to $`\mathrm{\Lambda }67\mathrm{TeV}`$. This would raise the corresponding limits on the mass of these $`\mathrm{SU}(2)`$ $`\mathrm{Z}^{}`$ bosons to $`M_\mathrm{Z}^{}(c_\varphi /s_\varphi )\times 550650\mathrm{GeV}`$.
## 4 Direct Searches for an $`\mathrm{SU}(2)`$ $`\mathrm{Z}^{}`$ in $`\mathrm{p}\overline{\mathrm{p}}\mathrm{Z}^{}\tau ^+\tau ^{}`$
In studying direct production of a $`\mathrm{Z}^{}`$ boson from an extended electroweak gauge structure, we must be aware of several competing issues. The couplings of third generation fermions to the extended gauge sector are enhanced relative to their Standard Model values, while those of the first and second generation particles are reduced. Since the current lower bounds on the $`\mathrm{Z}^{}`$ mass are on the order of $`400\mathrm{GeV}`$, the only machine presently available to perform a direct search is the Fermilab Tevatron. Thus, we are led to searching for a clean signal in third generation final states in a hadronic environment. Given the high mass of the top quark, the large QCD backgrounds for bottom production, and the difficulty of seeing the $`\tau `$-neutrino final state in photon plus missing energy or monojet events, the most promising channel is
$$\mathrm{p}\overline{\mathrm{p}}\mathrm{Z}^{}\tau ^+\tau ^{}+\mathrm{X}.$$
(4.1)
Each $`\tau `$ decays to a final state including either hadrons or one charged lepton and neutrinos. Our analysis concentrates on fully leptonic decays; we discuss possibilities with hadronic final states in Section 6. The three fully leptonic final states are characterized by opposite sign leptons with relatively large missing energy:
$$\tau ^+\tau ^{}\{\begin{array}{cc}\mu ^+\mu ^{}+\text{neutrinos}\hfill & \\ \mathrm{e}^+\mathrm{e}^{}+\text{neutrinos}\hfill & \\ \mathrm{e}^+\mu ^{}(\text{or }\mathrm{e}^{}\mu ^+)+\text{neutrinos}\hfill & \end{array}$$
(4.2)
Dimuon and dielectron final states are also characteristic of a number of Standard Model processes (e.g., Drell-Yan), making it difficult to separate our $`\mathrm{Z}^{}`$ signal from the backgrounds. Thus, we focus on the last channel above, namely oppositely charged, high-$`p_T`$ electron-muon pairs.
We have employed Pythia version 6.127 with our own simple model of the DØ detector to generate events and used our own code to analyze the generated data. Our idealized version of the Run II DØ detector defines the fiducial volume and smears event tracks. The central calorimeters are taken to have a pseudorapidity coverage, $`\eta `$, of $`|\eta |1.1`$, while the end-cap calorimeters are taken to have a coverage of $`1.5|\eta |4.0`$ for jets and $`1.5|\eta |2.5`$ for leptons. The dilepton events may initially be selected by a single or double lepton trigger. Our analysis assumes a trigger times offline selection efficiency of 95%, per lepton. We choose to trigger on the transverse momentum of the leptons, and set the triggering threshhold for each lepton at $`15\mathrm{GeV}`$.<sup>2</sup><sup>2</sup>2We have also considered a combination of this hard $`p_T`$ trigger for one of the leptons and a softer $`p_T`$ trigger for the other. While doing so does increase the number events to analyze, our later analysis cuts end up eliminating these extra events. Jet reconstruction efficiency is taken to be 100% for jets with transverse momenta in excess of $`8\mathrm{GeV}`$.<sup>3</sup><sup>3</sup>3We do not consider basing event triggers on jets, and so do not consider the jet triggering efficiency, which would be expected to be much lower than the efficiency for reconstruction given a previous trigger. This is, however, an issue of some interest, and we will return to it below, in Section 6. We identify a jet as a cluster of hadronic energy in excess of $`8\mathrm{GeV}`$ contained within a cone of base radius $`R=\sqrt{\eta ^2+\varphi ^2}<1`$. Our jet reconstruction code is based on the Pythia cluster finding algorithm. Based on these assumptions, we chose as an event trigger the presence of an electron and a muon of opposite electric charge both with transverse momenta in excess of $`15\mathrm{GeV}`$, with tracks lying within the fiducial volume of the detector.
We generated several sets of signal events corresponding to different $`\mathrm{Z}^{}`$ boson masses and different values of the mixing angle $`\varphi `$ in the region least constrained by precision electroweak data (large $`s_\varphi `$, see Figure 2). We also generated $`\mathrm{e}`$$`\mu `$ events from the four significant sources of Standard Model backgrounds
$$\mathrm{p}\overline{\mathrm{p}}\mathrm{Z}^0/\gamma ^{}\tau ^+\tau ^{}\mathrm{e}\mu +\text{neutrinos}$$
(4.3)
$$\mathrm{p}\overline{\mathrm{p}}\mathrm{W}^+\mathrm{W}^{}\mathrm{e}\mu +\text{neutrinos}$$
(4.4)
$$\mathrm{p}\overline{\mathrm{p}}\mathrm{t}\overline{\mathrm{t}}\mathrm{W}^+\mathrm{W}^{}\mathrm{b}\overline{\mathrm{b}}\mathrm{e}\mu \mathrm{b}\overline{\mathrm{b}}+\text{neutrinos}$$
(4.5)
$$\mathrm{p}\overline{\mathrm{p}}\mathrm{b}\overline{\mathrm{b}}\mathrm{e}\mu \mathrm{q}\overline{\mathrm{q}}+\text{neutrinos},$$
(4.6)
which we will call the $`\mathrm{Z}^0`$, $`\mathrm{W}`$-pair, top, and bottom backgrounds, respectively. The superficially similar backgrounds from charm production were eliminated by the event selection cuts we discuss below. When generating prompt tau leptons from the signal and background processes in Pythia, we have ignored polarization correlations between the tau pairs. This is a reasonable approximation for our study, since the correlations are diluted when considering only the leptonic decays of the tau. A correct accounting of the correlations will most likely improve the separation of signal and background.
For each signal and background process, we generated a minimum of $`3\times 10^4`$ events matching the $`\mathrm{e}\mu `$ final state at $`\sqrt{s}=2\mathrm{TeV}`$. For each event we verified the presence of the opposite charge $`\mathrm{e}\mu `$ pair, and then reconstructed the jets in the event. To the four-momenta of these leptons and jets, we applied smearing functions appropriate to the detector.<sup>4</sup><sup>4</sup>4We have used the following algorithms for track smearing for our Run II detector: for electrons, we performed a Gaussian smearing based on the electron energy, with $`\mathrm{\Delta }E/E=15\%/\sqrt{E}`$; for muons, we performed a Gaussian smearing based on the transverse momentum, with standard deviation given by $`\sigma _{p_T^\mu }=1.5\times 10^3p_{T}^{\mu }{}_{}{}^{2}`$, where the transverse momentum is measured in $`\mathrm{GeV}`$; finally, for jets we performed a Gaussian smearing based on the transverse jet energy, with standard deviation quadratic in the transverse energy with $`\eta `$ dependent coefficients. We accepted events in which the smeared tracks of both the electron and the muon lay in the fiducial volume of our idealized Run II DØ detector. Smeared jets falling outside the detector were dropped from the events. We eliminated events in which both smeared leptons no longer passed the trigger cuts. Finally, the four momenta of both leptons and the surviving jets in the remaining events were stored for later “offline” analysis.
Properly normalized $`p_T`$ distributions of background and signal events following the trigger stage are shown in Figure 3.
Examination of the smeared trigger distributions in this figure suggests offline analysis cuts that will eliminate the majority of the pure electroweak and $`\mathrm{b}`$ backgrounds, while preserving sufficient signal to permit analysis. These $`p_T`$ distributions are symmetric in the “radial” direction in the $`p_T^\mathrm{e}`$-$`p_T^\mu `$ plane, where $`p_T^\mathrm{e}`$ ($`p_T^\mu `$) is the transverse momentum of the electron (muon). For our primary analysis cut, we define the leptonic transverse momentum, $`p_T^L`$, of the event,
$$p_T^L=\sqrt{p_{T}^{\mathrm{e}}{}_{}{}^{2}+p_{T}^{\mu }{}_{}{}^{2}},$$
(4.7)
where we require $`p_T^L`$ to exceed some threshold, $`P_T^{\text{cut}}`$. Our choice of $`P_T^{\text{cut}}`$ was dictated by the requirements of enhancing the signal-to-background ratio while maintaining a sufficiently high absolute signal event rate, and was chosen in conjunction with the other cuts to be described below. We display typical signal-to-background rates before and after the $`p_T^L`$ cut in Figure 4.
Based on the calculated signal-to-background ratio and the absolute signal rates, we placed our $`p_T^L`$ cut at
$$p_T^L60\mathrm{GeV}.$$
(4.8)
This value for $`P_T^{\text{cut}}`$ should effectively reduce the $`\mathrm{W}^+`$$`\mathrm{W}^{}`$ background, since some of the leptons from $`\mathrm{W}`$-pair decay exhibit a characteristic Jacobian peak near $`M_W/2`$.
To further improve signal purity, we consider the presence of hadronic jet activity. For our signal events, we would expect to see no hadronic jet activity originating at the partonic event level. Similarly, we expect no jet activity for the $`\mathrm{Z}^0`$ and $`\mathrm{W}`$-pair backgrounds. However, we always expect activity associated with the top and bottom backgrounds. In particular, we expect two $`\mathrm{b}`$ jets associated with the top decays to $`\mathrm{W}`$-b, and two $`\mathrm{c}`$ jets associated with the bottom decays. Näively then, a cut on jet multiplicity will preferentially remove the top and bottom backgrounds. We have analyzed the expected jet distributions of events surviving the $`p_T^L`$ cut, as measured in our simulation for each type of event considered. This includes the extra jet activity generated by parton showering. We display these distributions in Figure 5; again, by jets we mean here clusters of hadronic activity with energy in excess of $`8\mathrm{GeV}`$.
Note that by rejecting events with jet multiplicity greater than one, we can remove a large majority of the top and bottom backgrounds while minimally impacting the strength of the signal. Comparing the signal-to-background ratio for our model before and after a jet multiplicity cut, we find a significant increase in signal purity, Figure 4.<sup>5</sup><sup>5</sup>5In our simulations we do not account for effects due to pileup and multiple interactions on jet reconstruction. we expect these issues will have minor impact on the final efficiency and the signal-to-background ratio.
Next, we apply a topological cut based on the opening angle between the electron and muon, which we will label $`\theta _{\mathrm{e}\mu }`$. Given the large mass of the $`\mathrm{Z}^{}`$, we expect it to be produced nearly at rest in the detector, and expect the tau pair to be produced back-to-back and highly boosted. Because of this boost, the electron and muon should be travelling nearly collinearly with their respective parent taus, making them approximately back-to-back with one another in the signal events. The background events would not be expected to have an opening angle distribution that is so highly peaked back-to-back. This is borne out by the Monte Carlo simulation. In Figure 6, we display the distribution of opening angles angles for events which have already passed the $`p_T^L`$ and jet multiplicity cuts. We choose to eliminate those events with $`\mathrm{cos}\theta _{\mathrm{e}\mu }>0.5`$, that is, those where the electron and muon are not strongly back-to-back. Note that the vast majority of the signal events pass the topological cut, while the background events are more likely to fail it.
We have displayed the impact of this topological cut on the signal-to-background ratios for various $`p_T^L`$ cuts in Figure 4.
At this point, the remaining background is almost purely $`\mathrm{W}`$-pair. We apply a final topological that that eliminates much of the remaining $`\mathrm{W}`$-pair background, based on the opening angle between the lowest $`p_T`$ lepton and the transverse missing energy in the event, which we label $`\theta _\mathrm{}\mathit{}_T`$. In order for the event to conserve momentum overall, we would expect the missing transverse energy vector to point along the direction of the softer decay lepton in the $`\mathrm{Z}^0`$ and signal events. Hence, we expect this opening angle to be peaked near $`\theta _\mathrm{}\mathit{}_T=0`$ for the signal events, while for the other backgrounds, we would not expect this. Based on the opening angle distributions shown in Figure 7, we eliminate events where $`\mathrm{cos}\theta _\mathrm{}\mathit{}_T<0.9`$, which is particularly effective at eliminating the $`\mathrm{W}`$-pair background, and particularly ineffective at eliminating signal events.
To summarize the effectiveness of our cuts, we present in Table 1 the fraction of each type of event which survives each cut, along with the expected number of events that survive; overall, roughly 40% of the signal events will survive all four cuts, while substantially less than 1% of all background events will similarly survive.
After performing these cuts on our data, we determined normalized signal and background cross sections from our Monte Carlo data, from which we can obtain luminosity bounds for 90% and 95% exclusion, as well as three and five standard deviation discovery bounds. We will explore first the exclusion reach of the Tevatron, followed by the discovery reach.
Exclusion bounds are obtained by calculating the following Poisson test statistic ,
$$r(\sigma _S,\sigma _B,)=1\frac{_{i=0}^N\left(S+B\right)^ie^{(S+B)}/i!}{_{i=0}^NB^ie^B/i!}$$
(4.9)
where $`\sigma _B`$ is the calculated background cross section, $`\sigma _S`$ is the calculated signal cross section, $``$ is the integrated luminosity, $`B=\sigma _B`$ is the expected number of background events, $`S=\sigma _S`$ the expected number of signal events, and $`N`$ is the largest integer smaller than the upper limit on the expected number of background events, that is $`N=\sigma _B=B`$. For each value of $`M_\mathrm{Z}^{}`$ and $`s_\varphi `$, taking the calculated signal and background cross sections, we varied the integrated luminosity and determined the statistic $`r`$. The minimum integrated luminosity required to exclude the model at a given confidence level, $`\mathrm{C}`$, is the luminosity where $`r=\mathrm{C}`$. The algorithm determines the ratio of the total probability of $`N`$ or fewer events occurring in an experiment, for the model of new physics, compared to the probability for standard model physics. At a given confidence level, $`\mathrm{C}`$, the area of overlap between the two probability distributions will be given by $`1\mathrm{C}`$.
We plot the exclusion limits in Figure 8 for a number of $`\mathrm{Z}^{}`$ masses and a range of mixing parameters, $`s_\varphi `$.
For a $`\mathrm{Z}^{}`$ boson of given mass, the luminosity required for exclusion is lowest when the mixing angle is near $`s_\varphi =0.80`$, with an approximately quadratic increase on either side of the minimum. The shape of the exclusion curve reflects the dependence of the $`\mathrm{Z}^{}`$ width on $`s_\varphi `$ (cf. Figure 1): narrow $`\mathrm{Z}^{}`$ bosons are easier to detect. With a few inverse femtobarns of integrated luminosity, the $`\mathrm{Z}^{}\tau \tau \mathrm{e}\mu `$ channel will begin to explore portions of the model parameter space that are not excluded by the precision electroweak data.
Discovery limits are obtained by applying the following algorithm. For a given integrated luminosity, $``$, we expect to observe $`B=\sigma _B`$ background events, and $`S=\sigma _S`$ signal events, for a total of $`S+B`$ expected events. We calculate the Poisson probability, $`P(\mu ,x)`$, that an expected background of $`\mu =B`$ events could fluctuate to give us a false signal of $`x=S+B`$ events total, that is $`P(B,S+B)`$. If the probability of such a fluctuation is smaller than a given confidence level, we are justified in declaring discovery of a new phenomenon at that level. We choose to determine the luminosity, $``$, for three gaussian standard deviation discovery (which we denote as $`3\sigma `$), where $`P(B,S+B)_{3\sigma }1.35\times 10^3`$, and for $`5\sigma `$, where $`P(B,S+B)2.7\times 10^7`$. We plot discovery bounds for a number of $`\mathrm{Z}^{}`$ masses and a range of mixing parameter, $`s_\varphi `$, in Figure 9
As expected from the previously determined exclusion bounds, for a $`\mathrm{Z}^{}`$ boson of a given mass, the luminosity for discovery is lowest when the mixing angle is near $`s_\varphi =0.80`$. As with exclusion, only a few inverse femtobarns of integrated luminosity will be required to discover a $`\mathrm{Z}^{}`$ boson with mass just above the current exclusion bounds.
By studying the $`\mathrm{e}`$$`\mu `$ invariant mass distribution ($`M_{\mathrm{e}\mu }=|p_\mu +p_\mathrm{e}|`$) it may be possible to detect the presence of a $`\mathrm{Z}^{}`$ boson and to determine its mass. As shown in Figure 10,
the invariant mass distribution of the background events which pass all of our cuts peak at around $`100\mathrm{GeV}`$. The centroid of the distribution of signal events is shifted toward higher invariant mass, with the amount of the shift depending on the value of $`M_\mathrm{Z}^{}`$, but almost independent of $`s_\varphi `$, Figure 11.
However, as the background rate is independent of the mixing angle while the signal rate is not, the centroid of the background plus signal distribution is not sufficient on its own to determine the mass. With a sufficiently high data rate, available at the LHC for example, it may be possible to choose more aggressive cuts that would lessen the dependence of the centroid on the background; it may alternatively be possible to perform a background subtraction on the invariant mass distribution, although this approach is more difficult to perform with confidence. With enough data, then, it should be possible to determine $`M_\mathrm{Z}^{}`$ from the measured invariant mass distribution $`M_{\mathrm{e}\mu }`$.
Determining the value of $`s_\varphi `$ will be more challenging. As indicated by the shape of the exclusion curves in Figure 8, the relationship between event-rate and $`s_\varphi `$ is double-valued: a given event rate corresponds to two values of $`s_\varphi `$, one above and one below $`s_\varphi =0.80`$ (approx.). Finding evidence of $`\mathrm{Z}^{}\mathrm{e}^+\mathrm{e}^{},\mu ^+\mu ^{}`$ may allow one to differentiate between the two possible values of $`s_\varphi `$ due to the different forms of the couplings: the first or second generation lepton couplings vary as $`c_\varphi /s_\varphi `$, while to the third generation couplings vary as $`s_\varphi /c_\varphi `$.
## 5 $`\mathrm{Z}^{}`$ Bosons from Extended Hypercharge Interactions
Models with an extended hypercharge gauge group can also produce heavy $`\mathrm{Z}^{}`$ bosons that couple more strongly to the third generation than to the lighter generations . In these theories, the electroweak gauge group is
$$\mathrm{SU}(2)_W\times \mathrm{U}(1)_h\times \mathrm{U}(1)_{\mathrm{}}$$
(5.1)
where third-generation fermions couple to $`\mathrm{U}(1)_h`$ with standard hypercharge values and the other fermions are carry standard hypercharges under $`\mathrm{U}(1)_{\mathrm{}}`$. At a scale above the weak scale, the two hypercharge groups break to their diagonal subgroup, identified as $`\mathrm{U}(1)_Y`$. As a result, a $`\mathrm{Z}^{}`$ boson that is a linear combination of the original two hypercharge bosons becomes massive. This heavy $`\mathrm{Z}^{}`$ boson couples to fermions as
$$i\frac{e}{\mathrm{cos}\theta }\left(\frac{s_\chi }{c_\chi }Y_{\mathrm{}}\frac{c_\chi }{s_\chi }Y_h\right)$$
(5.2)
where $`\chi `$ is the mixing angle between the two original hypercharge sectors
$$\mathrm{cot}\chi =\left(\frac{g_h}{g_{\mathrm{}}}\right)^2.$$
(5.3)
Comparing Equation 5.2 with the covariant derivative for the $`\mathrm{Z}^{}`$ boson from an extended weak group, Equation 2.15, we find three key differences. Two are physically relevant: the overall coupling is of hypercharge rather than weak strength, and the $`\mathrm{Z}^{}`$ couples to both left-handed and right-handed fermions at leading order. One is a matter of convention: mixing angle $`\chi `$ is equivalent to $`\pi /2\varphi `$.
At energies well below the mass of the $`\mathrm{Z}^{}`$ boson, its exchange in the process $`\mathrm{e}^+\mathrm{e}^{}f\overline{f}`$ where $`f`$ is a $`\tau `$ lepton or $`\mathrm{b}`$ quark may be approximated by the contact interaction
$$_{\mathrm{NC}}\frac{e^2}{\mathrm{cos}^2\theta M_\mathrm{Z}^{}^2}\left(\frac{s_\chi }{c_\chi }\left[\overline{\mathrm{e}}\gamma _\mu Y_{\mathrm{}}\mathrm{e}\right]\right)\left(\frac{c_\chi }{s_\chi }\left[\overline{f}\gamma ^\mu Y_hf\right]\right)$$
(5.4)
Comparing this with the contact interactions studied by LEP (see Section 3.1), we find that the LEP data sets its strongest limit through the process $`\mathrm{e}_R^+\mathrm{e}_R^{}\tau _R^+\tau _R^{}`$ ,
$$\mathrm{\Lambda }(f=\tau ,\eta _{RR}=+1)>\{\begin{array}{cc}3.7\mathrm{TeV}\hfill & \text{ALEPH}\hfill \\ 3.7\mathrm{TeV}\hfill & \text{OPAL},\hfill \end{array}$$
(5.5)
which gives a limit on the $`\mathrm{Z}^{}`$ mass of
$$M_\mathrm{Z}^{}=\mathrm{\Lambda }\sqrt{\frac{\alpha _{\mathrm{em}}}{\mathrm{cos}^2\theta }}>370\mathrm{GeV}.$$
(5.6)
This is stronger than the previous limits from precision electroweak data .
We have also used the techniques described in Section 4 to analyze the process $`\mathrm{p}\overline{\mathrm{p}}\mathrm{Z}^{}\tau \tau \mathrm{e}\mu `$ for a $`\mathrm{U}(1)`$ $`\mathrm{Z}^{}`$ boson. Due to the similar form of couplings of the $`\mathrm{Z}^{}`$ bosons to fermions, we obtain results for exclusion and discovery bounds that can be expected from the Tevatron that depend on mixing angle in a similar fashion. We display exclusion bounds in Figure 12
and discovery bounds in Figure 13.
The luminosity required to exclude or discover a $`\mathrm{U}(1)`$ $`\mathrm{Z}^{}`$ boson is a bit greater than for an $`\mathrm{SU}(2)`$ $`\mathrm{Z}^{}`$ boson of the same mass. This difference reflects the fact that the $`\mathrm{U}(1)`$ boson’s coupling to fermions is of hypercharge rather than weak strength.
## 6 Future Searches
Detecting even relatively light $`\mathrm{Z}^{}`$ bosons that couple preferentially to third-generation fermions is clearly a challenge for Tevatron and LEP experiments. Even in the $`p\overline{p}\mathrm{Z}^{}\tau \tau \mathrm{e}\mu `$ process where the signal-to-background ratio can be made quite large, the absolute number of signal events is kept low by the size of the $`\mathrm{Z}^{}`$ boson’s coupling to the light fermions from which it is produced. In the long term, the LHC’s higher center-of-mass energy will allow its experiments to search for these $`\mathrm{Z}^{}`$ bosons without being hampered by low signal event rates. In the meantime, we suggest that a few additional search channels may prove useful.
Obviously, the reach in $`\tau ^+`$$`\tau ^{}`$ final states will be extended beyond that shown in this analysis if use can be made of one or more hadronic tau decays <sup>6</sup><sup>6</sup>6A parton-level study in ref. estimates that a 500 GeV X boson coupling to $`B3L_\tau `$ could be visible via $`X\tau ^+\tau ^{}\text{jet}+\mathrm{}`$ in 2 $`fb^1`$ of data at Run II.. The single-prong decays of the tau, which constitutes about 85% of all decays, may have sufficiently small background since QCD should rarely produce isolated, high-$`p_T`$ tracks. It will be difficult to use these final states to search for $`\mathrm{Z}^{}`$ bosons: it is extremely hard, in the hadronic environment, to trigger on jets, and flavor tagging with high precision is an unresolved problem. Nonetheless, since the branching ratio of $`\tau `$ to hadrons ($`\mathrm{BR}(\tau \text{hadrons})65\%`$) is higher than to leptons ($`\mathrm{BR}(\tau \text{leptons})35\%`$) , even modest jet trigger and flavor tagging efficiencies could prove extremely valuable in searches or measurement of parameters. The ability to use these additional channels with their higher event rates should yield significantly improved mass limits for a given integrated luminosity.
In the semi-leptonic decay scenario, where we have $`\tau ^+\tau ^{}\text{jet}+\mathrm{}`$, the event trigger could be a high-$`p_T`$ electron or muon with, for example, $`p_T>15\mathrm{GeV}`$. In offline processing, one would then reconstruct the jets, and attempt to perform flavor tagging. If the corrected $`\tau `$ tagging efficiency<sup>7</sup><sup>7</sup>7By which we mean the tagging efficiency, after corrections for other objects faking $`\tau `$ jets. can be raised to approximately 15%, then the semi-leptonic events will provide the same event rate for analysis as the fully leptonic events previously considered. Further study of these channels are clearly warranted.
$`\mathrm{Z}^{}`$ bosons arising from extended weak interactions will also be accompanied by $`\mathrm{W}^{}`$ bosons of very similar mass (to leading order, $`M_\mathrm{Z}^{}=M_\mathrm{W}^{}`$). These bosons could be searched for in the process $`\mathrm{p}\overline{\mathrm{p}}\mathrm{W}^{}\tau \nu _\tau `$. Standard model backgrounds would include $`\mathrm{p}\overline{\mathrm{p}}W\mathrm{}\nu _{\mathrm{}}`$ and $`\mathrm{p}\overline{\mathrm{p}}\mathrm{WZ}^0\mathrm{}\nu _{\mathrm{}}\nu \nu `$, both of which should have softer lepton spectra, as well as $`\mathrm{p}\overline{\mathrm{p}}\mathrm{Z}^0+\text{jet}\nu \nu +\text{fake lepton}`$, where the jet is misidentified.
The methods of analysis pursued in Section 4 could productively be applied to models with scalars that have large branching ratios to tau pairs. While a heavy Standard Model Higgs boson does not have a high enough branching ratio for these analyses to provide useful limits, a pseudoscalar Higgs with large branching ratio would be an interesting candidate for study.
## 7 Conclusions
We have discussed two methods of searching for $`\mathrm{Z}^{}`$ bosons that couple primarily to third generation fermions. Bounds on the scale of quark-lepton compositeness derived from data taken at LEP and the Tevatron now imply that $`\mathrm{Z}^{}`$ bosons derived from extended $`\mathrm{SU}(2)_h\times \mathrm{SU}(2)_{\mathrm{}}`$ or $`\mathrm{U}(1)_h\times \mathrm{U}(1)_{\mathrm{}}`$ interactions must have a mass greater than about $`375\mathrm{GeV}`$. The reach of these limits will improve as additional data is taken. As the Tevatron Run II begins, it will become possible to search for $`\mathrm{Z}^{}`$ bosons using the process $`p\overline{p}\mathrm{Z}^{}\tau \tau \mathrm{e}\mu \mathrm{X}`$. We have shown that a combination of cuts based on lepton transverse momenta, jet multiplicity, and event topology, can overcome the standard model backgrounds. With $`30\mathrm{fb}^1`$ of data, the Run II experiments will be able to exclude $`\mathrm{Z}^{}`$ bosons with masses up to $`750\mathrm{GeV}`$. Were a $`\mathrm{Z}^{}`$ boson, instead, discovered, the shape of the $`\mathrm{e}`$$`\mu `$ invariant mass distribution and the relative branching fractions to taus and to muons could reveal the $`\mathrm{Z}^{}`$ mass and coupling strength.
Acknowledgments
The authors thank C. Hoelbling, F. Paige, and M. Popovic for useful conversations, and K. Lane and P. Kalyniak for comments on the manuscript. E.H.S. acknowledges the support of the NSF Faculty Early Career Development (CAREER) program and the DOE Outstanding Junior Investigator program. M.N. acknowledges the support of the NSF Professional Opportunities for Women in Research and Education (POWRE) program. This work was supported in part by the National Science Foundation under grants PHY-9501249 and PHY-9870552, by the Department of Energy under grant DE-FG02-91ER40676, and by the Davis Institute for High Energy Physics.
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# Decoherence Caused by Topology in a Time-Machine Spacetime11footnote 1Intern. J. Mod. Phys. D5, 1-27 (1996)
## 1 Introduction
It is known that General Relativity predicts the possible existence of spacetimes with non-trivial topology. The most intriguing among them are spacetimes with a region (or regions) containing closed time-like curves (CTC’s). These time-machine spacetimes will be considered in the present paper. It is not clear whether the laws of physics permit the development of CTC’s in the course of the evolution of a spacetime from some reasonable initial conditions (see - and references therein). We will not concern this problem here supposing that this is possible. Instead of this, taking spacetime geometry ad hoc, we shall investigate evolution of quantum particles in this geometry.
We shall discuss only non-relativistic processes. In this case the spacetime can be divided into chronal regions (regions without CTC’s) and dischronal regions (containing CTC’s) by slices $`t=\text{const}`$. For simplicity we shall discuss here the spacetime with only one dischronal region. This dischronal region (located between time moments $`t_1`$ and $`t_2`$) is preceded, at $`t<t_1`$, and followed, at $`t>t_2`$, by chronal regions. We shall call the first chronal region (at $`t<t_1`$) the initial region, the second chronal region (at $`t>t_2`$) the final region, and the dischronal region (at $`t_1<t<t_2`$) the time machine.
In a series of rather recent papers (see for example -) the question was analyzed whether the standard laws of physics (classical and quantum) can accommodate, in a reasonable manner, a spacetime with CTC’s. In this paper we shall focus our attention on the problem of propagating non-relativistic non-interacting quantum particles in the spacetime with CTC’s (the time-machine spacetime).
The main conclusion made in the paper is that, in the case when closed time-like lines exist, a sort of superselection (decoherence) may arise in certain conditions, i.e. the superposition principle may be partially violated. In this case the conventional unitary description is not applicable in the time-machine spacetime. Instead, description by a series of partial evolution operators, obeying the generalized unitarity condition, must be applied. The description of this type has been earlier worked out in quantum theory of continuous measurements.
It is well known now that the superposition principle of quantum mechanics is restricted when a measurement is performed. One may say that a sort of superselection exists in the situation of measurement, forbidding to superpose the states corresponding to different superselection sectors. In the case of quantum measurements the superselection sectors correspond to different alternative measurement outputs. The term ‘decoherence’ is often applied to describe this situation -.
The superselection (decoherence) takes a specific form if a continuous (prolonged in time) measurement is performed (see - and references therein). In the procedure of path integration, those paths that are connected with different alternative outputs $`\alpha `$ of the continuous measurement, cannot be summed up. Instead of this, the amplitude $`A_\alpha `$ of each of these alternatives must be calculated by summation of the paths corresponding to the alternative $`\alpha `$. After this, the probability of each alternative $`\alpha `$ can be found as a square modulus of the corresponding amplitude $`A_\alpha `$. The sum of all these probabilities should be equal to unity.
This is valid not only in the situation when the measurement is performed on purpose, but also when a measurement-like interaction of the system of interest with its environment takes place. The latter means that information is recorded in the environment about the state of the system or its evolution. This information may be described with the help of alternatives. Each alternative $`\alpha `$ is a class of states or a class of paths of the system. Different classes should be considered as classical (decohering) alternatives. Superposition of amplitudes corresponding to different alternatives $`\alpha `$ is forbidden.
In practice the situation of the type of a measurement arises each time when “macroscopically distinct” states (or ways of evolution) of the system exist. What states must be considered to be macroscopically distinct, depends on what environment the system has.
This dependence on environment may be illustrated by the well-known two-slit experiment. Propagation of the particle through the first or the second slit suggests two alternatives. Usually these are considered as quantum (‘coherent’, ‘interfering’) alternatives. One cannot know what of this alternative is actually realized. In this case superposition of paths describing propagating through different slits is possible, leading finally to the well-known interference effect.
However this consideration is valid only when the slits are in vacuum or in such a medium which cannot distinguish between two alternatives. If the interaction with the medium (environment) records information in this medium about what slit the particle propagates through, the situation is quite different and the preceding consideration is incorrect. It is possible of course to consider the environment and its interaction with the particle explicitly. However the influence of the environment onto the particle may be taken into account implicitly if one consider two alternatives (corresponding to the slits) to be classical (‘non-coherent’, ‘incompatible’). Superposing the paths passing through different slits is in this case forbidden. This is the basic idea of quantum theory of measurements. More complicated situations of continuous quantum measurements may be found in .
Propagating in a spacetime containing CTC’s, a particle may travel backward in time. This means that it can pass through some time interval (dischronal region) twice or more. Let the number of times the particle returns to its past be $`n`$. Then we have different ways of propagation characterized by the number $`n`$ (in fact, nothing else than a winding number arising from non-trivial topology of the time-machine spacetime).
The situations characterized by different numbers $`n`$ are quite different from the physical point of view. For the given $`n`$ the particle passes through the dischronal region (within the time machine) $`n+1`$ times. An observer in the dischronal region (within the time machine) will see $`n+1`$ particles existing simultaneously. Different values of $`n`$ correspond to different numbers of particles in the dischronal region.
However in non-relativistic quantum mechanics the number of particles is conserved, and a superselection corresponding to this number usually arises. If it really arises, the states corresponding to different numbers of particles may not be superposed.
The superselection (decoherence) connected with the number of particles is not absolute. It takes place for a massive particle possessing a charge of some type (for example electric charge) or spin (spinor charge). Even for such particles the decoherence exists only “in normal conditions”. The latter means that the environment (consisting for example of photons) distinguishes between the states with different numbers of particles. Then any superposition of such states will decohere very quickly so that the superpositions may be considered to be forbidden, not existing.
This may be invalid in special conditions. Superpositions of the states with different numbers of particles exist for example in the superfluid (liquid helium at low temperature). In this case the environment does not distinguish between different numbers of particles, and no superselection connected with the particle number exists.
In all considerations of the present paper we shall suppose that the conditions in the dischronal region of the time machine are “normal” in this sense, i.e. they lead to superselection. Then the number of particles (of the considered type) existing simultaneously in this region presents a classical alternative. If this number is equal to $`n+1`$, this means that the particle $`n`$ times returned from the future to the past. The situations corresponding to different $`n`$ should be treated as non-coherent.
This presumption means that for the environment (medium) within the time machine alternative ways of propagation of the particle corresponding to different $`n`$ are macroscopically distinct. One may say that a sort of measurement is performed in this case, and the result of the measurement is characterized by the number $`n`$. According to what has been said above, the paths characterized by different winding numbers $`n`$ should be considered as non-coherent (‘decohering’). Their summation is forbidden. Only the paths corresponding to the same $`n`$ may be summed up forming corresponding propagators and evolution operators $`U_n`$.<sup>2</sup><sup>2</sup>2For evolution into or from the dischronal region summing over the $`n`$th class of paths is insufficient, and multiplication by a certain projector is needed to obtain $`U_n`$, see Sect. 5. Propagators and evolution operators with different $`n`$ should not be summed up.
Hence, in the case when the superselection of different $`n`$ takes place, the “general” evolution operator $`U=_nU_n`$ (equal to the sum over all paths) makes no sense, and the evolution of the particle in a time-machine spacetime must be described in terms of the ‘partial evolution operators’ $`U_n`$. We shall develop this way of description including proof of generalized unitarity $`_nU_{n}^{}{}_{}{}^{}U_n=\mathrm{𝟏}`$ and generalized multiplicativity $`U_m(t^{\prime \prime },t^{})U_n(t^{},t)=U_{m+n}(t^{\prime \prime },t)`$. The evolution in a chronal region, between two chronal regions and within the dischronal region (within the time machine) will be considered. The resulting theory will be compared with the conventional (coherent) description of evolution as given in .
The paper is organized in the following way. In Sect. 2 general information about propagators and evolution operators in conventional theory is given, and an analogue of this formalism is presented (taken from quantum theory of measurement) for the case when the superposition principle is restricted by superselection. In Sect. 3 the superselection sectors (coherency sectors) are defined corresponding to the topological classes of particle trajectories in the time-machine spacetime. On the basis of this definition in Sects. 45 the generalized unitarity and in Sect. 6 the generalized multiplicativity of evolution in the time-machine spacetime is proven. In Sect. 7 the results obtained are compared with the conventional unitary description of evolution in the time-machine spacetime. Special kind of states ‘trapped’ within the time machine are described in Sect. 8. Sect. 9 contains short concluding remarks.
## 2 Propagators, Paths and Measurements
In the conventional non-relativistic quantum mechanics the propagator $`K(t^{\prime \prime },x^{\prime \prime }|t^{},x^{})`$ is a probability amplitude for the particle to transit from the point $`x^{}`$ at time moment $`t^{}`$ to the point $`x^{\prime \prime }`$ at time moment $`t^{\prime \prime }`$. The propagator may be considered as a kernel of the evolution operator $`U=U(t^{\prime \prime },t^{})`$:
$$\psi _{t^{\prime \prime }}=U(t^{\prime \prime },t^{})\psi _t^{},\psi _{t^{\prime \prime }}(x^{\prime \prime })=K(t^{\prime \prime },x^{\prime \prime }|t^{},x^{})\psi _t^{}(x^{})𝑑x^{}.$$
(1)
An explicit expression for the propagator is presented by the path integral,
$$K(t^{\prime \prime },x^{\prime \prime }|t^{},x^{})=d[x]\mathrm{exp}\left(\frac{i}{\mathrm{}}S[x]\right),$$
(2)
where integration is performed over all paths $`[x]=\{x(t)|t^{}tt^{\prime \prime }\}`$ connecting the points $`(t^{},x^{})`$ and $`(t^{\prime \prime },x^{\prime \prime })`$, and $`S[x]`$ is the action functional calculated along the path $`[x]`$.
The evolution operator (propagator) satisfies the equation
$$U(t^{\prime \prime },t^{})U(t^{},t)=U(t^{\prime \prime },t),K(t^{\prime \prime },x^{\prime \prime }|t^{},x^{})K(t^{},x^{}|t,x)𝑑x^{}=K(t^{\prime \prime },x^{\prime \prime }|t,x)$$
(3)
(for $`tt^{}t^{\prime \prime }`$) which we shall call the property of multiplicativity. Besides, the evolution should conserve scalar products of the states, thus the evolution operator should be unitary that may be written in terms of propagators as follows:
$$K^{}(t^{\prime \prime },x^{\prime \prime }|t^{},x^{})K(t^{\prime \prime },x^{\prime \prime }|t,x)𝑑x^{\prime \prime }=\delta (x^{},x).$$
The formalism of propagators and evolution operators must however be modified if a measurement or an observation of the system is performed. Such a modification is suggested by quantum theory of measurements and particularly by quantum theory of continuous measurements (see - and references therein).
The main feature of the resulting formalism is that a set of classical alternatives $`\alpha `$ should be considered. These alternatives arise as different measurement outputs. However, for emerging the situation of this type, it is not necessary that a measurement or an observation be arranged on purpose. The only condition necessary for this is that some information about quantum system be recorded in classical form in its environment. In the case of the time machine, the topological classes of trajectories described in Sect. 3 will play the role of classical alternatives.
The alternatives $`\alpha `$ are classical in the sense that they are incompatible. Correspondingly, quantum amplitudes corresponding to different alternatives cannot be summed up: they are non-coherent. Particularly, instead of a single propagator $`K(t^{\prime \prime },x^{\prime \prime }|t^{},x^{})`$ or a single evolution operator $`U(t^{\prime \prime },t^{})`$ a series of partial propagators $`K_\alpha (t^{\prime \prime },x^{\prime \prime }|t^{},x^{})`$ or partial evolution operators $`U_\alpha (t^{\prime \prime },t^{})`$ are necessary for describing evolution of the system. The partial propagators have not to be summed up.
How the partial evolution operators may be used to describe the evolution of the measured system? This should be made in different ways in the case of selective situation (for example selective measurement) and non-selective one. Selective situation means that it is known what alternative $`\alpha `$ is realized. The evolution is described in this case by one partial propagator or evolution operator $`U_\alpha (t^{\prime \prime },t^{})`$. The evolution law, in terms of the wave function (state vector) or the density matrix, is following:
$$\psi _t^{}^\alpha =U_\alpha (t^{},t)\psi _t,\rho _t^{}^\alpha =U_\alpha (t^{},t)\rho _t\left(U_\alpha (t^{},t)\right)^{}.$$
(4)
The non-selective situation means that it is not known what concrete alternative is realized. In this case one should sum up over all possible alternatives. However, dealing with classical alternatives, one must sum up probabilities, not amplitudes. This means that summing must be performed in the second of the formulas (4) resulting in the following evolution law:
$$\rho _t^{}=\underset{\alpha }{}\rho _t^{}^\alpha =\underset{\alpha }{}U_\alpha (t^{},t)\rho _t\left(U_\alpha (t^{},t)\right)^{}.$$
(5)
The final density matrix $`\rho _t^{}`$ has to be normalized ($`\mathrm{Tr}\rho _t^{}=1`$) for any normalized initial matrix $`\rho _t`$. This takes place if and only if the following condition is fulfilled:
$$\underset{\alpha }{}\left(U_\alpha (t^{},t)\right)^{}U_\alpha (t^{},t)=\mathrm{𝟏}.$$
(6)
It may be called the generalized unitarity condition. This condition means conservation of probability provided the probability of the alternative $`\alpha `$ to belong to the set $`𝒜`$ of alternatives is defined as follows (see ):
$$\mathrm{Prob}(\alpha 𝒜)=\mathrm{Tr}\underset{\alpha 𝒜}{}\rho _t^{}^\alpha =\mathrm{Tr}\underset{\alpha 𝒜}{}U_\alpha (t^{},t)\rho _t\left(U_\alpha (t^{},t)\right)^{}.$$
(7)
If the alternatives (measurements) $`\alpha _t^{}^{t^{\prime \prime }}`$ corresponding to different time intervals $`[t^{},t^{\prime \prime }]`$ may be considered, their multiplication may be introduced,
$$\beta _t^{}^{t^{\prime \prime }}\alpha _t^t^{}=\gamma _t^{t^{\prime \prime }}.$$
Then the following multiplicative law should be valid for the corresponding evolution operators (propagators):
$$U(\beta _t^{}^{t^{\prime \prime }})U(\alpha _t^t^{})=U(\gamma _t^{t^{\prime \prime }}),tt^{}t^{\prime \prime }.$$
(8)
All these concepts, elaborated previously for continuous quantum measurements, may now be applied to the alternatives corresponding to different topological numbers $`n`$ in the time-machine spacetime.
In the paper the formalism of propagators in the form of path integrals has been applied to describe the evolution of a non-relativistic particle in the time-machine spacetime. The propagator for such a particle was defined as an integral over all paths. It was shown that, in a simple model of the time-machine spacetime, unitarity and multiplicativity for the propagator are violated if at least one of the time moments is in the dischronal time interval.
We argued in Sect. 1 that a sort of superselection may arise in the time-machine spacetime for topologically different paths. If this is the case, then the “general” propagator determined by integrating over all paths makes no sense. Instead, one must use partial propagators just as in theory of continuous quantum measurements. Each of partial propagators should be a sum over topologically equivalent paths.<sup>3</sup><sup>3</sup>3This is valid for the evolution from the past of the time machine to its future. For evolution into or from the time machine the sum over topologically equivalent paths should be multiplied by a certain projector, see Sect. 5. Let us come over to description of the corresponding classes of paths and partial propagators.
## 3 Topological Classes of Paths
Let us consider the concrete model of a time-machine spacetime and define the topological classes of trajectories of a non-relativistic particle in such a spacetime. We shall discuss hereafter the model of a time-machine spacetime that has been used in the paper . It will be evident though that at least some of the results have more general validity.
The spacetime under consideration may be constructed (see Fig. 1) from a usual, ‘chronal’ one by adding two temporary wormholes. The original chronal spacetime will be called the background spacetime. The wormholes added to the background spacetime connect two space regions $`S_1`$ and $`S_2`$ belonging to the time slices $`t_1`$ and $`t_2`$ correspondingly, with $`t_1t_2`$. One of the wormholes $`W_1`$ leads from the past region $`S_1`$ to the future region $`S_2`$. This wormhole is entered if the region $`S_1`$ is approached from its ‘past side’ (but not through the wormhole $`W_2`$). Another wormhole $`W_2`$ leads from the future region $`S_2`$ to the past region $`S_1`$ if the region $`S_2`$ is approached from its ‘past side’ (but not through the wormhole $`W_1`$). About the properties of walls of the wormholes see the paper .
Consider all paths of a non-relativistic particle in this spacetime and divide them in the topological classes. The classes will then be connected with non-coherent sectors (classical alternatives) in description of the particle evolution in the spacetime.
The first class is formed by the paths going through the ‘future-directed’ wormhole $`W_1`$, entering $`S_1`$ and going out from $`S_2`$. Let us denote this class by $`𝒫_{00}`$. The scheme of the path belonging to this class is following:
$$𝒫_{00}:\text{ initial point}S_1W_1\text{final point}.$$
One more class $`𝒫_0`$ includes the paths lying completely in the background spacetime. These paths bypass both entrances $`S_1`$ and $`S_2`$ to the wormholes. This is a class of topologically simple paths. The scheme of the paths of this class is following:
$$𝒫_0:\text{ initial point}\text{final point}.$$
The class $`𝒫_1`$ contains the paths that pass through the wormhole $`W_2`$ once. The path of this type enters the mouth $`S_2`$ of the wormhole at time $`t_2`$ and go out of the mouth $`S_1`$ at time $`t_1`$. The scheme of these paths is following:
$$𝒫_1:\text{ initial point}S_2W_2\text{final point}.$$
The path of the class $`𝒫_2`$ proceeds through the wormhole $`W_2`$ twice, going into $`S_2`$ at time $`t_2`$, going out of $`S_1`$ at $`t_1`$, proceeding through the background spacetime again to time $`t_2`$, once more entering $`S_2`$, again going out of $`S_1`$, and ultimately approaching the final point through the background spacetime. This corresponds to the following scheme:
$$𝒫_2:\text{ initial point}S_2W_2S_2W_2\text{final point}.$$
Generally, the paths of the class $`𝒫_n`$ with $`n0`$ pass through the wormhole $`W_2`$ precisely $`n`$ times, according to the scheme
$$𝒫_n:\text{ initial point}(S_2W_2)^n\text{final point}.$$
It was argued in Introduction that the types of evolution corresponding to different numbers of returns of the particle to its past are usually (in “normal conditions”) macroscopically distinct and therefore cannot be coherently superposed. These types of evolution correspond to the paths belonging to different topological classes. If the decoherence takes place, the propagator defined by integrating over all paths (as in Eq. (2)) makes no sense. It seems reasonable to consider instead the propagators obtained by integrating over the topological classes:
$$K_n(t^{\prime \prime },x^{\prime \prime }|t^{},x^{})=_{𝒫_n}d[x]\mathrm{exp}\left(\frac{i}{\mathrm{}}S[x]\right).$$
(9)
The operators corresponding to these two-point functions will be denoted by $`K_n`$:
$$\psi _{t^{\prime \prime }}(x^{\prime \prime })=\left(K_n(t^{\prime \prime },t^{})\psi _t^{}\right)(x^{\prime \prime })=K_n(t^{\prime \prime },x^{\prime \prime }|t^{},x^{})\psi _t^{}(x^{})𝑑x^{}.$$
If both the initial and final time moments $`t^{},t^{\prime \prime }`$ are in the same chronal region (before or after the time machine), then only one topologically trivial class $`n=0`$ exists, and only one operator $`K_0`$ is defined by Eq. (9). It is evident that, according to general theory (Sect. 2), this operator $`K_0=U_0=U`$ is the evolution operator of the particle in the chronal region. This operator is unitary.
If the time moments $`t^{},t^{\prime \prime }`$ are in different chronal regions or some of them (or both) are in the dischronal region, then there are operators $`K_n`$ corresponding to arbitrary $`n=00,0,1,2,\mathrm{}`$. They seem to be good candidates for describing evolution in the $`n`$th superselection sector ($`n`$th classical alternative). We shall see that this is valid in the case when the time $`t^{}`$ is before the time machine emergence, $`t^{}<t_1`$ and the time $`t^{\prime \prime }`$ is after it disappearing, $`t^{\prime \prime }>t_2`$. In this case the operators $`K_n=U_n`$ play the role of ‘partial evolution operators’ satisfying the generalized unitarity condition, $`_nU_{n}^{}{}_{}{}^{}U_n=1`$. However when the time moment $`t^{\prime \prime }`$ is in the dischronal region (within the time machine), then the partial evolution operators $`U_n`$ will be shown to differ from the operators $`K_n`$ by certain projectors (see Sect. 5).
This is only natural because of the following.
The coherency sectors for evolution of a particle in the time-machine spacetime may be defined as a number of times the particle returns to its past. The set of such superselection sectors is in one-to-one correspondence with the topological classes $`𝒫_n`$ of paths between the time $`t^{}<t_1`$ and the time $`t^{\prime \prime }>t_2`$. Thus the above-defined operators $`K_n`$ must coincide in the case $`t^{}<t_1<t_2<t^{\prime \prime }`$ with what we call ‘partial evolution operators’ $`U_n`$. However, if the final time $`t^{\prime \prime }`$ of the evolution is within the time machine, $`t_1<t^{\prime \prime }<t_2`$, the class $`𝒫_n`$ of paths leading to this time does not correspond to the $`n`$th coherency sector.
Indeed, after the time moment $`t^{\prime \prime }`$ the particle can either escape the time machine or enter the wormhole $`W_2`$ and return to its past one or more times. Therefore, the operator $`K_n`$ (defined by summing up over the class $`𝒫_n`$) does not in fact coincide with the partial evolution operator for the $`n`$th alternative. If however we provide, multiplying $`K_n`$ by the corresponding projector, that the particle escape the time machine after $`t^{\prime \prime }`$, then we have the correct partial evolution operator $`U_n`$.
## 4 Generalized Unitarity I
We considered in Sect. 2 the general situation when non-coherent sectors exist in the evolution of a quantum system. In Sect. 3 coherency sectors connected with topological classes of paths were defined in the time-machine spacetime. In the present section and the next one we shall formulate and prove the generalized unitarity condition for this case.
As has already been argued, the classical alternatives may be identified (at least for the case $`t^{}<t_1<t_2<t^{\prime \prime }`$) with the classes of paths $`𝒫_n`$ introduced in Sect. 3. Operators $`K_n`$ were defined in Sect. 3 by summation over classes $`𝒫_n`$. It is natural to try and identify these operators with the partial evolution operators $`U_n`$ corresponding to the alternatives. This turns out to be correct for the most important case $`t^{}<t_1<t_2<t^{\prime \prime }`$ (evolution from the past of the time machine to its future). This may be accepted also in the cases when both $`t^{},t^{\prime \prime }`$ are in the same chronal region. A more complicated situation arises if $`t^{\prime \prime }`$ is within the time machine (in the dischronal region). The corrections for this case will be introduced in the next section.
Consider first the trivial situations of the evolution within one of two chronal regions. In the case $`t^{}<t^{\prime \prime }<t_1`$ (as well as for $`t_2<t^{}<t^{\prime \prime }`$) there is only one (trivial) class of paths with $`n=0`$ and only one operator $`K_0`$. This operator coincides in fact with the evolution operator in the background spacetime. Evolution in this case does not differ from usual coherent (unitary) evolution and is described by a single evolution operator $`U=U_0=K_0`$ (we omit hereafter the explicit specification of the initial and final time moments). The generalized unitarity condition coincides in this case with the conventional unitarity, $`U^{}U=\mathrm{𝟏}`$.
If the evolution not restricted by a single chronal region is considered, all topological classes $`n=00,0,1,2,\mathrm{}`$ (and corresponding coherency sectors) exist. According to the general formula (6), the generalized unitarity condition must then have the form
$$U_{00}^{}{}_{}{}^{}U_{00}+\underset{n=0}{\overset{\mathrm{}}{}}U_{n}^{}{}_{}{}^{}U_n=\mathrm{𝟏}$$
(10)
where we omitted specification of the time interval. We shall see that this condition is actually fulfilled for the evolution operators $`K_n`$ corresponding to the propagators (9) provided the final time of the evolution, $`t^{\prime \prime }>t_2`$. This is why the operators $`K_n`$ may actually be identified with the partial evolution operators $`U_n`$ in this case.
The generalized unitarity condition (10), as usual for quantum mechanics, expresses conservation of probability. If an initial state is described by the density matrix $`\rho `$, then the probability of realization of $`n`$th alternative is, according to quantum theory of measurements (see Eq (7)),
$$P_n=\mathrm{Tr}\rho _n=\mathrm{Tr}\left(U_n\rho U_{n}^{}{}_{}{}^{}\right).$$
(11)
The condition (10) provides then that
$$P_{00}+\underset{n=0}{\overset{\mathrm{}}{}}P_n=1$$
for an arbitrary initial state $`\rho `$. The same probability interpretation of the generalized unitarity is valid for all other choices of $`t^{}`$, $`t^{\prime \prime }`$ that will be considered below.
To demonstrate that the condition (10) is valid, we shall express the operators $`K_n`$ in a more explicit form, through the evolution operators in the background chronal spacetime (we shall denote them by $`V_i`$) and the evolution operators in the wormholes (denote these operators $`W_i`$). Besides this, two projectors in the background spacetime ($`P_i`$) will be used describing entering the mouths of the wormholes. Expressions for the operators $`K_n`$ may be readily constructed with the help of the description of classes $`𝒫_n`$ (see Sect. 3).
Consider first the case $`t^{}<t_1<t_2<t^{\prime \prime }`$, i.e. the evolution starting before the time machine emergence and finishing after it disappearance.
The paths of the class $`𝒫_{00}`$ go from the starting point to the point at time moment $`t_1`$, then enter the wormhole $`W_1`$, proceed through it to the moment $`t_2`$ and then go from $`t_2`$ to the final point. Summation over all paths of this type gives us the evolution operator $`K_{00}`$ in the form of the product $`V_2W_1P_1V_1`$ where $`V_1`$ and $`V_2`$ are the evolution operators outside the time machine (first of them between starting time $`t^{}`$ and time $`t_1`$, and the second one between $`t_2`$ and the final time $`t^{\prime \prime }`$), and $`W_1`$ is the evolution operator along the wormhole $`W_1`$. The projector $`P_1`$ onto the region $`S_1`$ provides entering the particle into the wormhole.
Analogously the expressions for all operators $`K_n`$ can be found in the case $`t^{}<t_1<t_2<t^{\prime \prime }`$. We shall see below that these operators satisfy the generalized unitarity condition and may be identified with the partial evolution operators. This is why we shall denote them by $`U_n`$:
$$U_{00}=K_{00}=V_2W_1P_1V_1,U_n=K_n=V_2(1P_2)(V_{21}W_2P_2)^nV_{21}(1P_1)V_1,n0.$$
(12)
Here $`V_{21}`$ is the evolution operator in the background spacetime between time moments $`t_1`$ and $`t_2`$. The projector $`P_2`$ provides entering the region $`S_2`$ and therefore the wormhole $`W_2`$ while the complementary projector $`1P_2`$ provides bypassing the wormhole $`W_2`$. Analogously the projector $`1P_1`$ provides bypassing $`W_1`$.<sup>4</sup><sup>4</sup>4We denoted here the unity operator by $`1`$ instead of $`\mathrm{𝟏}`$.
Now we can easily prove the generalized unitarity (10) for the operators $`U_n=K_n`$ defined by Eq. (12). In the proof we shall use unitarity of all $`V_i`$ and $`W_i`$ as well as the properties of projectors, $`P_i^2=P_i`$ and $`P_{i}^{}{}_{}{}^{}=P_i`$. Let us present the expression $`U_{n}^{}{}_{}{}^{}U_n`$ in the form
$`U_{n}^{}{}_{}{}^{}U_n`$ $`=`$ $`A^{}\left(B^{}\right)^nV_{21}^{}{}_{}{}^{}(1P_2)V_{21}B^nA`$ (13)
$`=`$ $`A^{}\left(B^{}\right)^{n1}CB^{n1}AA^{}\left(B^{}\right)^nCB^nA`$
where
$$A=(1P_1)V_1,B=W_2P_2V_{21},C=V_{21}^{}{}_{}{}^{}P_2V_{21}.$$
Summing (13) in $`n`$, we shall obtain
$$\underset{n=1}{\overset{\mathrm{}}{}}U_{n}^{}{}_{}{}^{}U_n=A^{}CA=V_{1}^{}{}_{}{}^{}(1P_1)V_{21}^{}{}_{}{}^{}P_2V_{21}(1P_1)V_1.$$
(14)
Adding two other terms $`U_{0}^{}{}_{}{}^{}U_0`$ and $`U_{00}^{}{}_{}{}^{}U_{00}`$, we shall prove finally the generalized unitarity condition (10).
###### Remark 1
In the proof given above (and in the analogous proofs hereafter) it is supposed implicitly that the terms of the sum (10) tend to zero when $`n`$ tends to infinity. This is not valid in the subspace of some ‘finely tuned’ states, and the formula (10) is not applicable in this subspace. More precisely, this formula, as an equality of two operators, is valid in the space of states $`|\psi `$ for which
$$\underset{n\mathrm{}}{lim}\psi |A^{}\left(B^{}\right)^nCB^nA|\psi =0$$
(in the notations introduced above). From the physical point of view, this means that the probability for the particle (in the considered state) to return $`n`$ times into its past infinitely decreases when $`n`$ tends to infinity. This is right for ‘almost all’ states. We shall see in Sect. 8 however that this is not valid for some finely tuned states trapped within the time machine.
Consider now the case when the initial time moment $`t^{}`$ is within the interval $`[t_1,t_2]`$, but the final time $`t^{\prime \prime }`$ is after $`t_2`$. In this case the paths cannot belong to the class $`𝒫_{00}`$, thus the generalized unitarity takes the form
$$\underset{n=0}{\overset{\mathrm{}}{}}U_{n}^{}{}_{}{}^{}U_n=\mathrm{𝟏}.$$
(15)
Let us prove it.
In the case $`t_1<t^{}<t_2<t^{\prime \prime }`$ we have for the operators $`K_n`$ (that again will prove to coincide with the partial evolution operators $`U_n`$) the following formulas:
$$U_n=K_n=V_2(1P_2)(V_{21}W_2P_2)^nV_{20},n0,$$
(16)
where $`V_{20}`$ denotes the propagator in the background (chronal) spacetime from the initial time $`t^{}`$ (between $`t_1`$ and $`t_2`$) to the time $`t_2`$. Summing up the expression for $`U_{n}^{}{}_{}{}^{}U_n`$ in the same way as above gives the generalized unitarity in the form (15) provided we accept the definition (16).
Let now the final time $`t^{\prime \prime }`$ of the evolution be between $`t_1`$ and $`t_2`$. We shall see that the (generalized) unitarity does not take place for the operators $`K_n`$ in this case, so that these operators cannot be identified with partial evolution operators. Right form of the partial evolution operators will be given in Sect. 5.
Constructing the operators $`K_n`$ in the same way as above (i.e. as integrals over the topological classes of paths), we are led to the following formulas for the case $`t^{}<t_1<t^{\prime \prime }<t_2`$ (when the initial time is earlier than the time machine emergence):
$$K_n=V_{01}(W_2P_2V_{21})^n(1P_1)V_1,n0$$
(17)
($`V_{01}`$ describes evolution between $`t_1`$ and the final time $`t^{\prime \prime }`$). If both initial and final times are within the time machine ($`t_1<t^{}<t^{\prime \prime }<t_2`$), we have
$$K_0=V_{0^{}0},K_n=V_{0^{}1}W_2P_2(V_{21}W_2P_2)^{n1}V_{20},n1.$$
(18)
In both cases (17), (18) the generalized unitarity does not take place.
One can readily see that the algebraic operations used earlier to prove the generalized unitarity, are now impossible because there is no projector of the form $`1P_2`$ in the expressions (17), (18). This gives a hint that the operator $`U_n`$ must differ from $`K_n`$ by the factor of the form $`1Q_2`$ where $`Q_2`$ is a projector. We shall see in Sect. 5 that this may be justified physically.
## 5 Generalized Unitarity II
It has been shown in the preceding section that the operators $`K_n`$ defined with the help of the classes of paths $`𝒫_n`$ and expressed by the formulas (17), (18) do not satisfy the generalized unitarity condition if the evolution is considered to the time moment $`t^{\prime \prime }`$ between $`t_1`$ and $`t_2`$.The physical reason of this is that the classes of paths do not correspond in this case to the coherency sectors, i.e. to macroscopically distinguishable (distinct) physical situations. Therefore, the operators $`K_n`$ are not in this case partial evolution operators.
Indeed, the situations are macroscopically distinct if they differ by the number of particles in the time interval $`[t_1,t_2]`$. If the path belongs to the class $`𝒫_n`$ and ends after $`t_2`$, it describes a quite definite number of particles, $`n+1`$ (because returning to the past is impossible after time $`t_2`$). However, if the path ends in the point $`t^{\prime \prime }<t_2`$, then the number of particles in the interval $`[t_1,t_2]`$ is indefinite, since it is not known whether the particle will escape the time machine after $`t^{\prime \prime }`$ or enter the back wormhole $`W_2`$, return to its past and propagate through the time machine once more. To fix the sector of coherency (classical alternative), we have to introduce projectors $`1Q_2`$ and $`Q_2`$ corresponding to the two possibilities described above: escaping the time machine and entering it once more.
It is evident how this may be done. If the operator $`V_{20}`$ describes the evolution between the time moments $`t^{\prime \prime }`$ (the final point of evolution) and $`t_2`$, then the operator
$$Q_2=V_{20}^1P_2V_{20}=V_{20}^{}{}_{}{}^{}P_2V_{20}$$
(19)
projects on those states at time $`t^{\prime \prime }`$ which after the evolution to the time $`t_2`$ will enter the region $`S_2`$ and return to the past. Accordingly, the operator $`1Q_2`$ distinguishes the states corresponding to the particle escaping from the time machine.
Having the projector $`1Q_2`$ and using Eq. (17), one can find the partial evolution operators for the case $`t^{}<t_1<t^{\prime \prime }<t_2`$:
$$U_n=(1Q_2)K_n=(1Q_2)V_{01}(W_2P_2V_{21})^n(1P_1)V_1,n0.$$
(20)
Let us remark that, in a sense, these operators, not (17), correspond to the topological classes of paths $`𝒫_n`$. Indeed, they actually fix the number of times the particle returns to its past, not only before time $`t^{\prime \prime }`$, but also after this. The operators (20) provide that the paths from the only one class $`𝒫_n`$ contribute the evolution of the particle after $`t^{\prime \prime }`$ as well as before this time.
Expressing $`Q_2`$ in (20) through $`P_2`$ (due to (19)) and using the equation $`V_{21}=V_{20}V_{01}`$ (multiplicativity for $`V_{21}=V(t_2,t_1)`$), we have finally the following formula for the partial evolution operators in the case $`t^{}<t_1<t^{\prime \prime }<t_2`$:
$$U_n=V_{20}^{}{}_{}{}^{}(1P_2)(V_{21}W_2P_2)^nV_{21}(1P_1)V_1,n0.$$
(21)
Let us try to prove the generalized unitarity for these partial evolution operators. One may for example substitute $`P_2`$ in the formula (20) by its expression through $`Q_2`$ to get
$$U_n=(1Q_2)V_{01}(W_2V_{20}Q_2V_{01})^n(1P_1)V_1,n0.$$
(22)
Then, acting just as in Sect. 4, we have
$$\underset{n=0}{\overset{\mathrm{}}{}}U_{n}^{}{}_{}{}^{}U_n=V_{1}^{}{}_{}{}^{}(1P_1)V_1.$$
(23)
The generalized unitarity is not yet fulfilled. The reason is that one more partial evolution operator must be considered, describing one more channel of evolution. Indeed, if the initial time is before $`t_1`$, then the particle have the possibility to enter the forward-directed wormhole $`W_1`$. The probability of this alternative is to be taken into account.
To take it into account, we must interpret the concept of a ‘final time moment’ in another way. So far we supposed that the ‘final time’ is nothing else than a space-like surface ($`t^{\prime \prime }=\text{const}`$) in the background spacetime (the spacetime without wormholes). However when the final time is within the time machine (i.e. later than $`t_1`$ but earlier than $`t_2`$) this ‘time’ may be thought of as a sum of the space-like surface in the background spacetime plus a space-like slice of the wormhole $`W_1`$.
With this wider concept of the final moment, the particle may arrive at this moment not only through the background spacetime but also through the wormhole $`W_1`$. The first possibility has been taken into account by the operators (22). The latter possibility is described by the operator
$$U_{00}=W_{01}P_1V_1.$$
(24)
Adding the term $`U_{00}^{}{}_{}{}^{}U_{00}`$ to the sum (23) gives, in the case $`t^{}<t_1<t^{\prime \prime }<t_2`$, the generalized unitarity condition in the form of Eq. (10).
At last, consider the case $`t_1<t^{}<t^{\prime \prime }<t_2`$ when both initial $`t^{}`$ and final $`t^{\prime \prime }`$ time moments are within the interval $`[t_1,t_2]`$. Using the partial evolution operators obtained from (18) by inclusion the projectors $`1Q_2=V_{20}^{}{}_{}{}^{}(1P_2)V_{20}`$, we have
$$U_n=V_{20^{}}^{}{}_{}{}^{}(1P_2)(V_{21}W_2P_2)^nV_{20},n0,$$
(25)
we shall easily prove the generalized unitarity in this case too:
$$\underset{n=0}{\overset{\mathrm{}}{}}U_{n}^{}{}_{}{}^{}U_n=\mathrm{𝟏}.$$
(26)
One more remark may be made in connection with the argument preceding Eq. (24). If the concept of time moment within the time machine is changed (as was discussed in the above-mentioned argument), the same wider concept must be applied not only to ‘final’, but also to the ‘initial time moment’. Then, besides the operators (25), one more operator should be added,
$$U_{00}=W_{0^{}0}$$
(27)
describing evolution in the wormhole $`W_1`$ between two slices of this wormhole.
It seems at the first glance that this additional operator violates the (generalized) unitarity. One can see however that this is not the case. Indeed, by the unity operator $`\mathrm{𝟏}`$ in Eq. (26), just as in all preceding formulas, we meant the unity operator in the space of states localized on the space-like surface of the background spacetime. If we express this in the notation explicitly, Eq. (26) should read
$$\underset{n=0}{\overset{\mathrm{}}{}}U_{n}^{}{}_{}{}^{}U_n=\mathrm{𝟏}_{\text{background}}.$$
(28)
The operator (27) satisfy the relation
$$U_{00}^{}{}_{}{}^{}U_{00}=\mathrm{𝟏}_{\text{wormhole}}$$
(29)
where the r.h.s. is the unity operator for the states localized on the space-like slice of the wormhole $`W_1`$. Summing up both preceding formulas, one has the generalized unitarity in the form
$$U_{00}^{}{}_{}{}^{}U_{00}+\underset{n=0}{\overset{\mathrm{}}{}}U_{n}^{}{}_{}{}^{}U_n=\mathrm{𝟏}_{\text{complete}}.$$
(30)
Quite analogously, the operator of the form
$$U_{00}=V_2W_{20}$$
(31)
may be added to the series of operators (16). The remarks analogous to those of the preceding paragraph may be made in this case. The generalized unitarity takes then the form (30) (instead of (15)) for the case $`t_1<t^{}<t_2<t^{\prime \prime }`$ too.
Thus, the generalized unitarity takes place for an arbitrary time interval. Therefore, the concept of probability may be used correctly even within the time machine.
Let us make one more remark. We saw that an additional projecting factor $`1Q_2`$ should be including in the partial evolution operators $`U_n`$ corresponding to the final time within the time machine, $`t_1<t^{\prime \prime }<t_2`$. The same logic seems to require inclusion of an analogous factor $`1Q_1`$ in the expressions for $`U_n`$ in the case when the initial time is within the time machine, $`t_1<t^{}<t_2`$. This actually may be done. Then we guarantee that the operators $`U_n`$ describe propagation of only those states that have resulted from the evolution starting earlier than $`t_1`$. This however is not necessary because the operators $`U_n`$ may be quite reasonably applied to arbitrary states prepared at time $`t^{}`$ within the time machine. This is why the projector $`1Q_1`$ turned out to be not necessary for providing the generalized unitarity.
Preparation of the state within the time machine could be restricted by the condition of self-consistency if interaction of the particle with its duplicates were taken into account. We however consider non-interacting particles (i.e. we suppose that the interaction is negligible). Therefore, the state of one particle at time $`t^{}`$ within the time machine may be prepared arbitrarily.<sup>5</sup><sup>5</sup>5Notice that the complete description of the states of all particles at this moment (including the original particle and its duplicates) is not arbitrary. Indeed, if we are preparing, at time $`t^{}`$ between moments $`t_1`$ and $`t_2`$, the state that will $`n`$ times return to its past, then it is already known at $`t^{}`$ that $`n`$ duplicates of the particle exist coming from the region $`S_1`$ at time $`t_1`$. Further evolution of this arbitrary state will be described by one of the partial evolution operators $`U_n`$. Any number $`n`$ may be realized in this evolution, the probability of each of them being determined by the formula (11).
## 6 Multiplicativity for the Time Machine
Let us address now the question of multiplicativity for the evolution in the time-machine spacetime. Multiplicativity of the propagators or evolution operators (3) means that the evolution during some time interval may be considered in two stages as evolution during two subintervals. With the decoherence caused by measurement, the multiplicativity is described by the formula (8). We should now try and interpret this formula for the evolution in the time-machine spacetime with the specific decoherence (superselection) arising in this case.
The winding number $`n`$ introduced in Sect. 3 through the classes of paths $`𝒫_n`$ hints how multiplicativity might be defined and proved in the case of the time machine. Indeed, this number is equal to the number of times the particle returns to its past travelling through the wormhole $`W_2`$. If we have two alternatives characterized by the numbers $`n_1`$ and $`n_2`$ correspondingly, then the product of these alternatives corresponds to the number $`n=n_1+n_2`$. This means simply that if the particle returns $`n_1`$ times and then once more returns $`n_2`$ times to its past, then ultimately it returns $`n_1+n_2`$ times. It is evident that the (generalized) multiplicativity (8) might have, in the case of a time machine, the following form that can be readily proved for the operators $`K_n`$:
$$K_{00}(t^{\prime \prime },t^{})K_{00}(t^{},t)=K_{00}(t^{\prime \prime },t),K_m(t^{\prime \prime },t^{})K_n(t^{},t)=K_{m+n}(t^{\prime \prime },t),n0.$$
(32)
The relation (32) may be readily proved for the operators $`K_n`$ with the help of the path-integral representation (9). The proof is based on the fact that each path $`p𝒫_{m+n}`$ may be presented as a product of the paths $`p_1𝒫_m`$ and $`p_2𝒫_n`$, and vice versa, the product of any paths $`p_1𝒫_m`$, $`p_2𝒫_n`$ gives a path belonging to $`𝒫_{m+n}`$. By ‘product’ we mean passing through one of the paths and then through another one.The product of two paths is defined only if the first path ends in the point where the second one starts (see , Chapter 10 for details of this algebra of paths).
Instead of this, the relation (32) for operators $`K_n`$ may easily be proved by using explicit expressions for these operators derived in Sect. 4. The latter way of consideration may be applied also to the partial evolution operators $`U_n`$ which, as we know, do not always coincide with the operators $`K_n`$. We shall see that for partial evolution operators $`U_n`$ the second of the relations (32) is valid only when at least one of the numbers $`m`$, $`n`$ is zero.
It is evident that the generalized multiplicativity (32) is valid in the trivial cases $`t^{}<t^{\prime \prime }<t_1`$ and $`t_2<t^{}<t^{\prime \prime }`$, when it reduces to the conventional multiplicativity because both $`m=n=0`$:
$$U_0(t^{\prime \prime },t^{})U_0(t^{},t)=U_0(t^{\prime \prime },t).$$
Trivial is also the case when one of the intervals $`[t,t^{}]`$ or $`[t^{},t^{\prime \prime }]`$ lies in one of the chronal regions. Then the multiplicativity takes one the following forms that can be easily proven:
$$U_n(t^{\prime \prime },t)=U_0(t^{\prime \prime },t^{})U_n(t^{},t)\text{or}U_n(t^{\prime \prime },t)=U_n(t^{\prime \prime },t^{})U_0(t^{},t).$$
Consider now non-trivial situations.
Summing up the results of Sects. 45, we have the following formulas for the partial evolution operators $`U_{00}(t^{\prime \prime },t^{})`$, $`U_n(t^{\prime \prime },t^{})`$ (where $`n0`$):
$`t^{}<t_1<t_2<t^{\prime \prime }:`$
$`U_{00}`$ $`=`$ $`V_2W_1P_1V_1,U_n=V_2(1P_2)(V_{21}W_2P_2)^nV_{21}(1P_1)V_1.`$
$`t_1<t^{}<t_2<t^{\prime \prime }:`$
$`U_{00}`$ $`=`$ $`V_2W_{20},U_n=V_2(1P_2)(V_{21}W_2P_2)^nV_{20}.`$
$`t^{}<t_1<t^{\prime \prime }<t_2:`$
$`U_{00}`$ $`=`$ $`W_{01}P_1V_1,U_n=V_{20}^{}{}_{}{}^{}(1P_2)(V_{21}W_2P_2)^nV_{21}(1P_1)V_1.`$
$`t_1<t^{}<t^{\prime \prime }<t_2:`$
$`U_{00}`$ $`=`$ $`W_{0^{}0},U_n=V_{20^{}}^{}{}_{}{}^{}(1P_2)(V_{21}W_2P_2)^nV_{20}.`$
One can verify straightforwardly that, of all multiplicativity relations (32), the following are valid also for the partial evolution operators:
$$U_{00}(t^{\prime \prime },t^{})U_{00}(t^{},t)=U_{00}(t^{\prime \prime },t),U_0(t^{\prime \prime },t^{})U_n(t^{},t)=U_n(t^{\prime \prime },t),n0.$$
(33)
If $`m0`$, the second of the relations (32) is not valid for $`U_n`$. Moreover, if we defined the partial evolution operators with the time argument within the time machine in such a way as it was supposed in the remark at the end of Sect. 5, then even the products of the form $`U_0U_n`$ with $`n0`$ would also give zero.
We see therefore that a topologically non-trivial evolution (described by the evolution operator $`U_n`$ with $`n1`$) cannot be presented as the product of two topologically non-trivial evolutions. Topologically non-trivial evolution is ‘integral in time’, it cannot be followed step by step: first $`m0`$ returns to the past, then again $`n0`$ returns that finally gives $`m+n`$ returns.
This is not astonishing. If we describe the evolution within the time machine, not achieving an exit from it, then this part of evolution cannot be characterized by a definite number $`n`$ of returns of the particle to the past (this number depends on the subsequent stage of evolution). In a sense, evolution within the time machine is not local in time, influence of the future cannot be excluded.
We succeeded in constructing operators describing evolution to the time moment $`t^{\prime \prime }`$ within the time machine, for example $`U_n(\text{in,past})`$. However a special projector in this operator provides escaping from the time machine after time $`t^{\prime \prime }`$. This operator guarantees that the future stage of evolution (after $`t^{\prime \prime }`$) will not influence the past evolution (until $`t^{\prime \prime }`$). Thus, non-locality of evolution within the time machine is evident even from the structure of partial evolution operators, not only from the form of multiplicativity for them.
On the other side, the operators $`K_n`$ (but not $`U_n`$) possess the property of multiplicativity (32) for arbitrary time arguments.
One may object that expressing multiplicativity in terms of the operators $`K_n`$, in fact non-physical within the time machine, makes no sense (if the superselection in $`n`$ takes place). In fact, each of these operators describes propagation of non-physical states (forbidden superpositions). However, the product of such operators,
$$K_{n_N}(t^{\prime \prime },\tau _{N1})\mathrm{}K_{n_2}(\tau _2,\tau _1)K_{n_1}(\tau _1,t^{})=K_{n_1+n_2+\mathrm{}+n_N}(t^{\prime \prime },t^{}),$$
with $`t^{}<t_1<t_2<t^{\prime \prime }`$, is a correct partial evolution operator. No non-physical state is propagated due to this operator. Thus, even though each of the operators $`K_n`$ describes propagation of non-physical states within the time machine, the result of the propagation from the past (in respect to the time machine) to the future of the time machine is presented by these operators correctly.
One may say that the operators $`K_n`$ give multiplicative description of evolution even within the time machine, but at the price of introducing non-physical states in intermediate stages.
One more interpretation of the results obtained is complementarity between (generalized) unitarity and multiplicativity. One may describe the evolution within the time machine by the operators $`U_n`$ (unitary in the generalized sense but satisfying only trivial multiplicativity relations) or by $`K_n`$ (satisfying the generalized multiplicativity but not unitarity).
## 7 Comparison with the Unitary Theory
We showed in Sects. 45 that the description of the particle evolution in the time-machine spacetime with the help of the partial propagators (partial evolution operators) is correct from the point of view of probabilities. This statement means that the partial evolution operators $`U_n`$ satisfy the condition of generalized unitarity (10).
In the conventional theory, when no decoherence (superselection) is supposed, the evolution is described by a single evolution operator $`U`$ which must be unitary, $`U^{}U=1`$. In a number of papers evolution of particles in a time-machine spacetime was described in this way (see and references therein). Let us compare two types of theories. This is not difficult because we used here the same model of spacetime as in .
It was shown in that a single evolution operator $`U=U(t^{\prime \prime },t^{})`$ is unitary in the case $`t^{}<t_1<t_2<t^{\prime \prime }`$ when the evolution from the past of the time machine to its future is considered. At the first glance, this contradicts to our conclusion about generalized unitarity of the partial evolution operators $`U_n`$ in this case.
Indeed, the evolution operator $`U`$ was defined in by summation over all paths. In our terms, this means that this operator is equal to
$$U=U_{00}+\underset{n=0}{\overset{\mathrm{}}{}}U_n$$
(34)
where $`U_n`$ are defined by Eq. (12). Unitarity for this operator means that
$$U_{00}^{}{}_{}{}^{}U_{00}+\underset{n=0}{\overset{\mathrm{}}{}}U_{00}^{}{}_{}{}^{}U_n+\underset{n=0}{\overset{\mathrm{}}{}}U_{n}^{}{}_{}{}^{}U_{00}+\underset{m,n=0}{\overset{\mathrm{}}{}}U_{m}^{}{}_{}{}^{}U_n=\mathrm{𝟏}.$$
(35)
The formulas (10) and (35) differ by the non-diagonal terms (they are absent in the former case). Both formulas may be valid only if a sum of non-diagonal terms is zero. It turns out however that this actually takes place in the considered case ($`t^{}<t_1<t_2<t^{\prime \prime }`$), so that both forms of unitarity turn out to be valid. Let us show this with the evolution operators (12).
Consider an off-diagonal term in (35), $`U_{n+k}^{}{}_{}{}^{}U_n`$, $`k1`$. Using Eq. (12), one may present this term in the form
$$U_{n+k}^{}{}_{}{}^{}U_n=A^{}(B^{})^{n+k1}CB^{n1}AA^{}(B^{})^{n+k}CB^nA$$
(36)
where it is denoted
$$A=(1P_1)V_1,B=W_2P_2V_{21},C=V_{21}^{}{}_{}{}^{}P_2V_{21}.$$
The formula (36) is valid for all $`n1`$. Summing up this expression in $`n`$ from 1 to $`\mathrm{}`$ and adding the terms $`U_{k1}^{}{}_{}{}^{}U_{00}`$ and $`U_{k}^{}{}_{}{}^{}U_0`$, we have
$`U_{k1}^{}{}_{}{}^{}U_{00}+{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}U_{n+k}^{}{}_{}{}^{}U_n`$
$`=`$ $`V_{1}^{}{}_{}{}^{}(1P_1)(V_{21}^{}{}_{}{}^{}P_2W_{2}^{}{}_{}{}^{})^{k1}V_{21}^{}{}_{}{}^{}(1P_2)W_1P_1V_1`$
$`+`$ $`V_{1}^{}{}_{}{}^{}(1P_1)(V_{21}^{}{}_{}{}^{}P_2W_{2}^{}{}_{}{}^{})^kV_{21}^{}{}_{}{}^{}(1P_2)V_{21}(1P_1)V_1`$
$`+`$ $`V_{1}^{}{}_{}{}^{}(1P_1)(V_{21}^{}{}_{}{}^{}P_2W_{2}^{}{}_{}{}^{})^kV_{21}^{}{}_{}{}^{}P_2V_{21}(1P_1)V_1`$
Two last terms here can be summed up to give a more simple expression:
$`U_{k1}^{}{}_{}{}^{}U_{00}+{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}U_{n+k}^{}{}_{}{}^{}U_n`$ (37)
$`=`$ $`V_{1}^{}{}_{}{}^{}(1P_1)(V_{21}^{}{}_{}{}^{}P_2W_{2}^{}{}_{}{}^{})^{k1}V_{21}^{}{}_{}{}^{}(1P_2)W_1P_1V_1`$
$`+`$ $`V_{1}^{}{}_{}{}^{}(1P_1)(V_{21}^{}{}_{}{}^{}P_2W_{2}^{}{}_{}{}^{})^k(1P_1)V_1`$
In the r.h.s. of the latter expression however both terms are equal to zero. The reason is the following equations:
$$(1P_2)W_1=0,(1P_1)W_2=0.$$
(38)
The first of them is a consequence of the fact that after propagation in the wormhole $`W_1`$ the particle is in the space region $`S_2`$. The operator $`P_2`$ projects on this region, and $`1P_2`$ projects on the complementary region (all the space-like surface $`t_2=\text{const}`$ of the background spacetime excluding the region $`S_2`$). Going out of the wormhole $`W_1`$ the particle cannot be in this complementary region. Therefore, the first equation in (38) is valid. Analogously, the second equation in (38) follows from the fact that after travelling in the wormhole $`W_2`$ the particle turns out to be in the space region $`S_1`$, and projecting, by $`1P_1`$, on the complementary region gives zero.
Using the first equation from (38) and a conjugate to the second equation we see that the expression (37) as well as its conjugate are equal to zero. Therefore,
$$U_{k1}^{}{}_{}{}^{}U_{00}+\underset{n=0}{\overset{\mathrm{}}{}}U_{n+k}^{}{}_{}{}^{}U_n=U_{00}^{}{}_{}{}^{}U_{k1}+\underset{n=0}{\overset{\mathrm{}}{}}U_{n}^{}{}_{}{}^{}U_{n+k}=0.$$
(39)
Thus, the off-diagonal terms do not contribute the sum (35), and both unitarity and generalized unitarity are valid in the considered case: evolution from the initial chronal region (earlier than the time machine) to the final chronal region (after it).
This does not mean of course that both descriptions, one with the superselection and one without superselection, are equivalent. According to the coherent description, the evolution law is following:
$$\rho ^{}=U\rho U^{}=U_{00}\rho U_{00}^{}{}_{}{}^{}+\underset{n=0}{\overset{\mathrm{}}{}}U_{00}\rho U_{n}^{}{}_{}{}^{}+\underset{n=0}{\overset{\mathrm{}}{}}U_n\rho U_{00}^{}{}_{}{}^{}+\underset{m,n=0}{\overset{\mathrm{}}{}}U_m\rho U_{n}^{}{}_{}{}^{}.$$
According to the theory with decoherence (superselection), the evolution should be described by partial evolution operators according to a general formula (5) that takes the following form in our case:
$$\rho ^{}=U_{00}\rho U_{00}^{}{}_{}{}^{}+\underset{n=0}{\overset{\mathrm{}}{}}U_n\rho U_{n}^{}{}_{}{}^{}.$$
(40)
This leads of course to different predictions. Thus, the two theories have different physical contents.<sup>6</sup><sup>6</sup>6Let us repeat once more that decoherence of different $`n`$ leading to generalized unitarity must take place “in ordinary conditions” while coherent description may be valid “in special conditions” when the environment does not distinguish between different numbers of particles in the dischronal region.
Consider different locations of times $`t^{}`$, $`t^{\prime \prime }`$ (the start and the end of the evolution) in respect to the time machine. In all cases when at least one of the time moments $`t^{}`$, $`t^{\prime \prime }`$ is within the time machine, no unitarity is found in the paper . However, using the definition (20), we can prove Eq. (39) in just the same way as above for the case $`t^{}<t_1<t^{\prime \prime }<t_2`$. Thus, in this case operator (34) is also unitary. There is again no contradiction with the result of because in that paper not this operator, but $`_nK_n`$ was considered as an evolution operator (both definitions coincide for $`t^{}<t_1<t_2<t^{\prime \prime }`$).
Formally unitarity for $`t^{}<t_1<t^{\prime \prime }<t_2`$ is caused by the projector $`1Q_2`$ in the definition of $`U_n`$. Physically the reason of this is that this projector provides escaping the particle from the time machine (see Sect. 5). Thus, though we deal with $`t^{\prime \prime }<t_2`$, actually the situation is, in a sense, equivalent to the case $`t^{\prime \prime }>t_2`$.
For the cases $`t_1<t^{}<t_2<t^{\prime \prime }`$ and $`t_1<t^{}<t^{\prime \prime }<t_2`$ the relation $`U^{}U=1`$ is invalid, and we are left only with the generalized unitarity. We could introduce in these cases too a sort of projector providing unitarity of the operator $`_nU_n`$. However there is no physical ground for this (see the remark at the end of Sect. 5).
## 8 States Trapped in the Time Machine
We considered the evolution of a particle entering the time machine from the past chronal region, circulating several times within the time machine and then escaping from it into the future chronal region. There are however the states that exist only within the time machine but not in the chronal regions. These states may be described at arbitrary time moment between $`t_1`$ and $`t_2`$. We shall describe them at time $`t_1`$.
Consider the state $`|\psi _1`$ at time $`t_1`$ possessing the properties
$$(1P_2)V_{21}|\psi _1=0,(1W_2V_{21})|\psi _1=0$$
It is evident that the particle in this state “bites its own tail”: after evolution to $`t_2`$ it enters the wormhole $`W_2`$ and after going out of the wormhole it turns out to be in the state identical with the initial state. After this the same cycle of evolution is repeated. The particle in such a state travels through the interval $`[t_1,t_2]`$, but it never passes any time moment before $`t_1`$ or any time moment after $`t_2`$.
There is evident generalization of this state, providing passing the dischronal region twice, three times or generally $`n`$ times before reaching the initial state. Such a state satisfies the following conditions at time $`t_1`$:
$`(1P_2)(V_{21}W_2)^n^{}V_{21}|\psi _n`$ $`=`$ $`0\text{ for }n^{}<n,`$
$`(1(W_2V_{21})^n)|\psi _n`$ $`=`$ $`0.`$ (41)
Then the particle “bites its own tail” after winding $`n`$ causal loops: after $`n`$ cycles, each of which contains travelling through the interval $`[t_1,t_2]`$, entering the wormhole $`W_2`$ and passing this wormhole, the state is identical to the initial one.
The states thus described are nothing else than a quantum version of the classical “Jinnee” discussed in , the classical bodies moving along closed time-like trajectories within a time machine.
One can construct a state in such a way as to provide passing the dischronal region infinite number of times never escaping to the future, but with the initial state never being repeated. The condition for this is (at time $`t_1`$)
$$(1P_2)(V_{21}W_2)^n^{}V_{21}|\psi _{\mathrm{}}=0\text{ (all }n^{}).$$
(42)
The particle in such a state cannot go out of the time machine though probably its state never becomes identical to the initial state. The states (8) may be considered to be special cases of (42).
One more generalization can be considered: the state existing earlier than the time machine emerged (before $`t_1`$) but trapped in TM so that it cannot be found at times later than $`t_2`$. Let us characterize this state at time $`t^{}<t_1`$:
$`P_1V_1|\psi _n`$ $`=`$ $`0\text{ (the particle does not enter }W_1\text{)},`$
$`(1P_2)(V_{21}W_2)^n^{}V_{21}(1P_1)V_1|\psi _n`$ $`=`$ $`0,n^{}1\text{ (it does not escape TM)}.`$
It may be shown that the state of the particle in this case cannot be repeated after a finite number $`n`$ of cycles. The condition for such a repetition,
$$(1(V_{21}W_2)^n)V_{21}(1P_1)V_1|\psi _n=0\text{ (the state is repeated after }n\text{ cycles)}.$$
would be inconsistent due to the relation $`W_2=P_1W_2`$ (expressing that the mouth leading out from the wormhole $`W_1`$ is in the region $`S_1`$).
Therefore, it is possible that a particle enters the time machine from the past and stays within it, infinitely repeating evolution from $`t_1`$ to $`t_2`$ and backward. Moreover, the states obtained by time reversal from those already considered, may also be defined. Then the particle escapes from the time machine into the future after infinite number of cycles within the time machine, never being in the past chronal region.
Existence of the states trapped in the time machine contradicts to the generalized unitarity condition in the form given earlier. Indeed, this condition provides that any state evolves finally to the future chronal region through one of the channels enumerated by $`n`$. The trapped states do not at all escape from the time machine. These states belong to the class of ‘finely tuned’ states which the generalized unitarity condition is not valid for (see Remark 1 at page 1).
For example, application of the (generalized) unitarity condition (26) for evolution from the time $`t_1`$ (considered as a time within the time machine) to the same time $`t_1`$ should seemingly forbid the trapped states. However, analysing attentively the proof of this condition, we see that it is supposed in this proof that the operator
$$I_N=1\underset{n=0}{\overset{N}{}}U_{n}^{}{}_{}{}^{}U_n=\left(V_{21}^{}{}_{}{}^{}P_2W_{2}^{}{}_{}{}^{}\right)^NV_{21}^{}{}_{}{}^{}P_2V_{21}\left(W_2P_2V_{21}\right)^N$$
tends to zero when $`N\mathrm{}`$. The expectation value of this operator for almost all states (as well as the trace of this operator multiplied by almost all density matrices) decreases with increasing $`N`$. However, for the states considered in the present section (the trapped states) we have, as a consequence of Eq. (42),
$$\psi |\mathrm{\hspace{0.17em}1}\underset{n=0}{\overset{N}{}}U_{n}^{}{}_{}{}^{}U_n|\psi =\psi |\psi =1.$$
Hence the operator $`I_N`$ does not decrease and the generalized unitarity condition is not fulfilled in the space of the trapped states.
The physical interpretation of this fact in terms of probabilities is following. For this very special class of states (not escaping to the future chronal region) no part of probability goes into channels enumerated by $`n`$. It might be reasonable to introduce one more channel corresponding to $`n=\mathrm{}`$ and claim that all probability, for these states, goes into this channel. This leads to the following remark.
###### Remark 2
In a general case (for arbitrary states) the generalized unitarity condition should be written in the form
$$U_{00}^{}{}_{}{}^{}U_{00}+\underset{n=0}{\overset{N}{}}U_{n}^{}{}_{}{}^{}U_n+I_N=\mathrm{𝟏}$$
(where the operator $`I_N`$ is given above for the case $`t^{}=t^{\prime \prime }=t_1`$ and can be readily written for all other choices of $`t^{}`$ and $`t^{\prime \prime }`$). The expectation value of this formula for an arbitrary state $`|\psi `$,
$$\psi |U_{00}^{}{}_{}{}^{}U_{00}|\psi +\underset{n=0}{\overset{N}{}}\psi |U_{n}^{}{}_{}{}^{}U_n|\psi +\psi |I_N|\psi =1,$$
gives the probability distribution for different channels of evolution (the last term in the l.h.s. gives the probability for all channels with numbers more than $`N`$). If the limit
$$\psi |I_{\mathrm{}}|\psi =\underset{N\mathrm{}}{lim}\psi |I_N|\psi $$
exists, the formula
$$\psi |U_{00}^{}{}_{}{}^{}U_{00}|\psi +\underset{n=0}{\overset{\mathrm{}}{}}\psi |U_{n}^{}{}_{}{}^{}U_n|\psi +\psi |I_{\mathrm{}}|\psi =1$$
gives the probability distribution for all channels with finite $`n`$ and for the channel with $`n=\mathrm{}`$ corresponding to the trapped states.
## 9 Conclusion
We considered in this paper how evolution of a non-relativistic non-interacting quantum particle in the spacetime with closed time-like curves (a time-machine spacetime) should be described. A simple non-relativistic model of such a spacetime was used corresponding to emergence of the time machine at some time moment and its disappearance at another moment.
In the paper evolution of a particle in this spacetime was described with the help of the evolution operator $`U`$ found by path integration. It was shown that such an operator is multiplicative and unitary only when the evolution is considered between the time moments belonging to the chronal regions (including the case when one of the time moments is in the past chronal region and the other is in the future chronal region).
In our paper the arguments were given that superselection (decoherence) may arise for the evolution in this spacetime (in “normal conditions”). The superselection sectors are enumerated in this case by the number of times the particle returns to its past. Therefore, evolution of the particle is described by a family of partial evolution operators $`U_n`$ (or partial propagators) instead of a single evolution operator $`U`$.
In the case when the evolution is considered into the future chronal region (i.e. to the time moment after the time machine disappearance), the partial propagators may be correctly defined as integrals over the corresponding topological classes of paths. If the final time of evolution is within the time machine, the path-integral expressions for the partial propagators (evolution operators) should be corrected by the projectors providing escaping the particle from the time machine after a certain number of returns to the past.
With these definitions, the family of partial evolution operators satisfy the generalized unitarity condition $`_nU_{n}^{}{}_{}{}^{}U_n=1`$ and multiplicativity condition $`U_nU_m=U_{n+m}`$ (the latter is fulfilled only for at least one of the numbers $`n`$, $`m`$ equal to zero). It is shown that the generalized unitarity is compatible with the unitarity of the operator $`U=_nU_n`$ in the case of evolution from the past of the time machine to its future. Nevertheless, the operator $`U`$ cannot describe correctly the evolution if the superselection exists.
A special class of states of the particle is described emerging and disappearing simultaneously with the time machine, i.e. existing only within it. These states present a quantum analogue of the classical “Jinnee” states investigated earlier . Besides this, the states are considered that either 1) exist in the past of the time machine but then are trapped in it and never escape into the future of the time machine, or 2) never exist in the past of the time machine, have infinite number of cycles within the time machine and then finally escape into its future. For all these states the generalized unitarity condition should be corrected to take into account the channel of evolution with $`n=\mathrm{}`$.
From conceptual point of view, the superselection of the considered type may be thought of as the influence of the future on the past, i.e. a sort of “consistency conditions”.
The coherent (unitary) evolution in the time-machine spacetime is also possible, but only in quite special conditions when the environment of the particle in the dischronal region does not distinguish between different numbers of particles.<sup>7</sup><sup>7</sup>7This question needs special investigation with a concrete model of an environment. A very good discussion of the role of environment in different situations may be found in . It is interesting that in the case of coherent evolution no information from the future to the past can be carried by the particle. If the particle leaves some information about the future in its past, then the interaction responsible for the information transfer leads to decoherence so that the winding number $`n`$ becomes definite.
ACKNOWLEDGEMENTS
One of the authors (M.M.) is grateful to J. Audretsch, V.A. Namiot and H.-D. Zeh for discussing the question of decoherence. The work was supported in part by the Danish National Research Foundation through the establishment of the Theoretical Astrophysics Center, the Danish Natural Science Research Council through grant 11-9640-1, and by the Deutsche Forschungsgemeinschaft.
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# Quasar Variability in the Framework of Poissonian Models
## 1 Introduction
Variability is one of the defining properties of Active Galactic Nuclei (AGN) and a potentially powerful discriminant among different scenarios for the physics of the central engine. Optical–UV variability studies in the past decade have put a lot of effort into investigating the effects of continuum variations upon the emission lines, leading to substantial progress in the diagnostics of the properties of the Broad Line Region (Netzer & Peterson 1997). Yet, no comparable progress has been achieved in understanding the origin of the continuum fluctuations.
In this paper we discuss quasar variability in the context of Poissonian models. By Poissonian we mean any physical system in which the variations are due to the stochastic superposition of independent flares, occurring at a given mean rate but randomly distributed in time, a process also called “christmas-tree”, “shot-noise”, “discrete events” or “subunits” model. This approach has two key advantages: (1) It provides a simple mathematical framework which can be applied to observed light curves to constrain the most relevant flare properties, such as energy, time-scale and rate (e.g., Cid Fernandes, Aretxaga & Terlevich 1996). (2) It encompasses a large class of physical models for AGN. For instance, scenarios as diverse as accretion-disk instabilities of several kinds (Haardt, Maraschi & Ghisellini 1994; Kawaguchi et al. 1998), disruption of stars in the gravitational field of super-massive black-hole (Peterson & Ferland 1986; Ayal, Livio & Piran 2000), stellar collisions (Keenan 1978; Courvoisier, Paltani & Walter 1996), supernovae (Aretxaga & Terlevich 1994; Aretxaga, Cid Fernandes & Terlevich 1997), and even extrinsic variability models such as micro-lensing (Hawkings 2000) all share a common Poissonian nature. It is the combination of these two points that makes the Poissonian interpretation of AGN variability attractive, as its application to observed data may provide feasibility tests on different theories.
The early realization that most luminous sources tend to vary less (Uomoto, Wills & Wils 1976; Pica & Smith 1983) qualitatively favored a Poissonian interpretation. However, these and subsequent quasar monitoring studies found that the slope ($`\beta `$) of the fractional variability ($`\delta \sigma (L)/\overline{L}`$) versus mean luminosity relation, $`\delta \overline{L}^\beta `$, is shallower than the $`\beta =1/2`$ value predicted in a simple Poissonian model (Cristiani, Vio & Andreani 1990; Trevese et al. 1994; Hook et al. 1994, Cristiani et al. 1996, Paltani & Courvoisier 1997), which lead some of these studies to rule out such models. Some studies, however, find slopes consistent with $`\beta =1/2`$ (Cid Fernandes et al. 1996; Garcia et al. 1999), while others even question the very existence of a correlation (Bonoli et al. 1979; Netzer & Sheffer 1983; Giallongo, Trevese & Vagnetti 1991; Lloyd 1984; Cimatti, Zamorani & Marano 1993; Netzer et al. 1996).
The slope of the variability-luminosity relation has thus been a controversial issue, with conflicting results reported in the literature. Selection effects and different coverages of the luminosity-redshift plane are likely causes for these discrepancies, as discussed in Hook et al. (1994). Garcia et al. (1999) argued that, although the observed variability-luminosity relation in their sample is shallower than $`\beta =1/2`$, good agreement between data and model is achieved after taking into account the increase in variability with frequency. Moreover, Cid Fernandes et al. (1996) argued that the photometric (at least for photographic data) and sampling uncertainties, along with wavelength/redshift effects, propagate to a poorly defined variability-luminosity relation. Furthermore, they showed that a slope of -1/2 is only expected in the simplest of Poissonian models, in which the flare energy, time-scale and background contribution are held fixed as universal constants for all objects. Discarding Poissonian models on the basis of $`\beta 1/2`$ is thus both risky and an oversimplified interpretation of Poissonian processes.
Quasar variability studies are reaching a new level of quality due to the efforts of several groups which engaged into long term, differential CCD photometric and spectroscopic monitoring programs (Borgeest & Schramm 1994; Netzer et al. 1996; Sirola et al. 1998; Giveon et al. 1999—hereinafter G99; Garcia et al. 1999; Kaspi et al. 2000). These studies represent an enormous improvement over pre-existing variability databases, both in sampling and photometric quality. While most previous studies (by necessity) based their analysis on properties of the ensemble of objects, these new data allow the study of quasar variability properties on an object-by-object basis. With the continuation and steady improvement of these programs, a clearer picture of the phenomenology of quasar variability will eventually emerge.
In this paper we review the formal relations between observable variability properties and the parameters in a general Poissonian model (§2). This approach is here seen as a valid step towards a physical understanding of the nature of AGN variability, which is particularly relevant given the current lack of an accepted paradigm for this ubiquitous phenomenon. Emphasis is put on the use of multi-wavelength data to constrain the basic model parameters. In §3 we apply the formalism to the photometric (B and R) quasar monitoring data collected over the past decade in the Wise Observatory (G99), by means of basic light curve statistics and an analysis of the individual Structure Functions. Simulations are used to verify the consistency of the results. A discussion on the use of higher moments of the data (skewness and kurtosis) is also presented (Appendix A). In Section 4 we discuss our results and present a correlation analysis of the variability with other properties, as well as possible interpretations in the context of AGN models. Finally, in Section 5 we summarize our main results and outline prospects of future work.
## 2 Formalism
Poissonian models for AGN variability have been studied both qualitatively and quantitatively, and for wavelengths across the electromagnetic spectrum, from radio (e.g., Dent 1972), to optical–UV (e.g., Cid Fernandes et al. 1996; Garcia et al. 1999) and X-rays (e.g., Lehto 1989; Almaini et al. 2000). In this section we follow the detailed formulation of this problem developed by Cid Fernandes (1995). The resulting formulae rest upon basic probability and time-series theory (e.g., Papoulis 1965).
In a Poissonian scenario the instantaneous monochromatic luminosity $`L_\lambda (t)`$ is due to the superposition of a variable component, $`V_\lambda (t)`$, and an underlying background component $`C_\lambda `$:
$$L_\lambda (t)=V_\lambda (t)+C_\lambda $$
(1)
where $`V_\lambda (t)`$ is made by the superpositions of flares $`l_\lambda (t)`$ with random “birth-dates” $`t_i`$:
$$V_\lambda (t)=\underset{i}{}l_\lambda (tt_i)$$
(2)
The first two moments (mean luminosity and relative variability) of the variable component are
$$\overline{V_\lambda }=\nu _\lambda E_\lambda $$
(3)
$$\frac{\sigma (V_\lambda )}{\overline{V_\lambda }}=\frac{1}{(\nu _\lambda \tau _\lambda )^{1/2}}$$
(4)
In these expressions $`\nu _\lambda `$ is the mean rate of flares, $`E_\lambda `$ is the monochromatic energy of individual flares, and $`\tau _\lambda `$ is the flare “life-time”, defined by
$$\tau _\lambda \frac{E_\lambda ^2}{l_\lambda ^2(t)𝑑t}$$
(5)
It can be shown that these relations also hold when one allows for the (likely) possibility that the flares within a given object are not all identical, i.e., when $`l_\lambda (t)=l_\lambda (t,𝐱)`$, where $`𝐱`$ denotes a general set of parameters (size, density, cooling time, …). The only modifications in this more general case is that $`E_\lambda `$ above means the average $`E_\lambda (𝐱)`$ over the probability distribution of $`𝐱`$, and similarly for $`\tau _\lambda `$.<sup>1</sup><sup>1</sup>1To be precise, if $`p(𝐱)`$ is the probability density of $`𝐱`$, one finds
$$\tau _\lambda \overline{E_\lambda }^2/l_\lambda ^2(t,𝐱)p(𝐱)𝑑𝐱𝑑t$$
The total mean luminosity and relative variability are both affected by the underlying constant component, turning eqs. (3) and (4) into
$$\overline{L_\lambda }=\nu _\lambda E_\lambda +C_\lambda $$
(6)
$$\delta _\lambda \frac{\sigma (L_\lambda )}{\overline{L_\lambda }}=v_\lambda \frac{1}{(\nu _\lambda \tau _\lambda )^{1/2}}$$
(7)
where $`v_\lambda \overline{V_\lambda }/\overline{L_\lambda }`$ is the fraction of the mean luminosity which is actually due to the variable component. Allowance for the possible contribution of a non-variable component is essential, though it has often been forgotten when discussing Poissonian models. Physically, $`C_\lambda `$ can be associated with several sources. These can be either around the nucleus and extrinsic to the variability generation process, as the host galaxy stellar contribution, or, more interestingly, a part of the variable continuum source which remains stable over long time-scales, such as the non-flaring part of an accretion disk.
The Poissonian model thus involves four basic parameters: the rate ($`\nu _\lambda `$), energy ($`E_\lambda `$) and life-time ($`\tau _\lambda `$) of the flares, plus $`C_\lambda `$. The shape of the flares may be regarded as a further degree of freedom, but it has little effect upon the analysis presented here. Equations (6) and (7) relate these parameters to just two observables. Even estimating $`\tau _\lambda `$ through a Structure Function (SF) analysis, one has a non-closed system, with three observables and four variables. Higher moments of the light curve could in principle be used as further constraints (see Appendix A), but these are so badly affected by sampling uncertainties that presently they do not provide useful constraints. Another piece of information that can help constraining the model parameters is that the constant component has to be at most as strong as the observed minimum in the light curve:
$$0C_\lambda L_{\lambda ,\mathrm{min}}$$
(8)
It is reasonable to presume that the same flares are seen across a narrow spectral band, such as the optical–UV. Empirical support for this hypothesis comes from the high similarity between the continuum light curves in different bands (e.g., Cutri et al. 1985; Krolik et al. 1991; G99). Indeed, the whole method of reverberation mapping relies on an equivalent hypothesis, namely, that the fluctuations in the ionizing continuum can be mapped by those of the optical–UV continuum (Peterson 1993). We therefore may write $`\nu _\lambda =\nu `$ and $`\tau _\lambda =\tau `$. The developments below could also be made allowing for $`\lambda `$ dependent time-scales, but we shall adopt the simpler constant $`\tau _\lambda `$ scenario.
A first consequence of this assumption is that one can isolate the spectral shape of flare energy directly from the spectral behavior of the standard deviation $`\sigma _\lambda `$:
$$E_\lambda =\left(\frac{\tau }{\nu }\right)^{1/2}\sigma _\lambda \sigma _\lambda $$
(9)
It is well established that the amplitude of the variations increase towards shorter wavelengths (Cutri et al. 1985; Edelson, Krolik & Pike 1990; Kinney et al. 1991; Cristiani et al. 1997; Di Clemente et al. 1996), with most sources becoming bluer as they brighten, which immediately tells us that the variable component is blue. This fact also has the very interesting consequence (see eq. ) that $`v_\lambda `$ must vary with wavelength, which can only be understood invoking an underlying component.
Another way of seeing this is to rewrite (7) as $`\delta _\lambda =v_\lambda 𝒩^{1/2}`$, where $`𝒩=\nu \tau `$ is the mean number of living flares at any time. Since we are assuming that the same flares are seen in different wavebands, the fact that $`\delta _\lambda `$ decreases towards the red can only be made consistent with a Poissonian scenario if the relative contribution of the background,
$$c_\lambda \frac{C_\lambda }{\overline{L_\lambda }}=1v_\lambda ,$$
increases with $`\lambda `$. Of course, the decomposition of optical–UV spectra of AGN into a variable plus a constant component has been proposed before, both on observational and on theoretical grounds. Here we proved that the spectral behavior of $`\delta _\lambda `$ implies the existence of a constant source if one is to keep within the framework of Poissonian models.
One can construct families of possible spectra for the constant component by writing
$$C_\lambda =\overline{L_\lambda }𝒩^{1/2}\sigma _\lambda $$
(10)
The simultaneous analysis of the light curve statistics in different wavebands thus has interesting consequences, but we still do not have a closed system. Except for the life-time $`\tau `$, which can be estimated through a structure function analysis, there is no way to determine absolute values for $`\nu `$, $`E_\lambda `$ or $`C_\lambda `$. Nonetheless, the wavelength information, coupled with the condition that $`C_\lambda `$ must satisfy (8), can yield improved constraints upon $`𝒩`$. One may combine the above relations to obtain
$$\left(\frac{1\mu _\lambda }{\delta _\lambda }\right)^2𝒩\left(\frac{1}{\delta _\lambda }\right)^2$$
(11)
where $`\mu _\lambda L_{\lambda ,\mathrm{min}}/\overline{L_\lambda }`$. Both lower and upper limits can be made more stringent by considering the whole wavelength information available. We can thus define $`𝒩_{\mathrm{min}}`$ by using the maximum value of the lower limit above, and conversely for $`𝒩_{\mathrm{max}}`$. This range of allowed $`𝒩`$ translates into corresponding ranges for $`\nu `$, $`E_\lambda `$ and $`C_\lambda `$. Exactly how stringent $`N_{\mathrm{min}}`$ is depends on how close $`L_{\lambda ,\mathrm{min}}`$ gets to $`C_\lambda `$. As the chances of the light curve reaching the background level decrease with the increasing superposition of events, one expects the $`N_{\mathrm{min}}`$ limit to get progressively less stringent as $`𝒩`$ increases. The quadratic dependence on the observed quantities also conspires to broaden the range of $`𝒩`$.
This very general and straightforward formalism can be applied to several variability data sets. In the next section, we apply it to one of the best sets of quasar optical light curves presently available.
## 3 Analysis of the Wise Observatory quasar light curves
### 3.1 Description of the data
Giveon et al. (1999) have presented the results of a long term B and R photometric monitoring of 42 optically selected nearby quasars from the Palomar Green (PG) sample. The observations were collected with the Wise Observatory 1m telescope over the 1991 to 1998 period. Spectrophotometric measurements at Wise and Steward Observatories were used to complement the light curves of 13 of the objects. Previous results on this and related monitoring campaigns have been published in Maoz et al. (1994), Netzer et al. (1996), and a spectroscopic study of a subsample of the objects has been recently concluded (Kaspi et al. 2000). The reader is referred to G99 for details of the observations. The following quantities are medians over the sample: $`n_{obs}=32`$ observations per object, rest-frame sampling interval of 33 days, rest-frame light curve span of 5.9 years, photometric uncertainty $`=0.015`$ mag (in B), $`z=0.16`$; $`M_B=23.3`$.
The apparent magnitude light curves were converted to monochromatic luminosities $`L_\lambda `$ at 4400 (B) and 6400 Å (R) using the same cosmological parameters as G99 ($`H_0=70`$ km s$`^1`$Mpc<sup>-1</sup>, $`q_0=0.2`$, $`\mathrm{\Lambda }=0`$). The K-correction was also performed as in G99, assuming a power-law spectrum, whose slope is defined by the median B $``$ R color. This was a minor correction because of the low redshifts in this sample. Galactic extinction corrections were made with the values of $`A_B`$ extracted from NED<sup>2</sup><sup>2</sup>2The NASA/IPAC Extragalactic Database (NED) is operated by the Jet Propulsion Laboratory, California Institute of Technology, under contract with the National Aeronautics and Space Administration., and the extinction law of Cardelli, Clayton & Mathis (1989, with $`R_V=3.1`$), but were mostly negligible ($`A_B0.17`$, median $`=0.03`$) because of the high galactic latitude of the objects (Schmidt & Green 1983). Emission lines contribute $`10\%`$ to the B and R band luminosities of PG quasars (G99); their effect upon the present analysis is also negligible.
The analysis presented below was carried out with the resulting $`L_B(t)`$ and $`L_R(t)`$ light curves. We note that the estimates of $`𝒩`$, $`\nu `$, $`\tau `$ and the background fractional contribution are cosmology, extinction and K-correction independent. Only the absolute values of $`E_\lambda `$ and $`C_\lambda `$ depend on such factors.
### 3.2 Structure Function Analysis: Estimating time-scales and the asymptotic variance
The most commonly employed tool to extract information on the variability time-scales out of quasar light curves is the Structure Function (Simonetti, Cordes & Heeschen 1985; Hook et al. 1994; Cristiani et al. 1996). The SF is defined by the mean of $`[L(t+\mathrm{\Delta }t)L(t)]^2`$ over a light curve. For a Poissonian sequence of flares the SF is known to be simply proportional to the SF of a single isolated flare (Appendix B), and therefore only has structure for time-scales shorter than the flare duration. For large $`\mathrm{\Delta }t`$ the SF converges to twice the intrinsic variance of the process: $`SF_\lambda (\mathrm{\Delta }t)2\sigma _\lambda ^2=2E_\lambda ^2\nu /\tau `$ (eq. ). This happens because variations on such long time-scales effectively correspond to independent samples of the $`L_\lambda (t)`$ process, as the light curve loses memory of its past.
The SF provides estimates for both the variability time scale and amplitude. We use the notation $`\sigma _{\lambda ,SF}^2`$ to distinguish the asymptotic variance derived from the SF from that computed directly from the light curve $`\sigma _\lambda ^2=\overline{L_\lambda ^2}\overline{L_\lambda }^2ϵ_\lambda ^2`$ (where the last term corrects for the small effects of photometric error), which may be somewhat underestimated for light curves not much longer than the correlation time-scale since they do not sample the whole power contained in the fluctuation power spectrum. Likewise, $`\delta _{\lambda ,SF}\sigma _{\lambda ,SF}/\overline{L_\lambda }`$ denotes the long term net variability.
The SFs were fit with theoretical functions corresponding to different flare evolution models. In all fits we are interested in just two quantities: The flare life-time $`\tau `$, defined by (5), and $`\sigma _{\lambda ,\mathrm{SF}}^2`$. Four models for $`l(t)`$ were explored: (1) square; (2) exponential; (3) symmetric triangle and (4) asymmetric triangle. Equations and plots for the corresponding SFs are displayed in Appendix B. These are clearly toy-models for the radiative evolution of physical flares, but as shown in Appendix B (Fig. 10) the SF is not very sensitive to $`l(t)`$. While it is unfortunate that little information about $`l(t)`$ can be retrieved with this technique, this does make the estimates of $`\tau `$ and $`\sigma _{\lambda ,\mathrm{SF}}^2`$ more robust. This is confirmed by the fits, which showed that all $`l(t)`$ models produce very similar parameters (with the exception of exponential flares, which in some cases do not converge and are generally poorer). For this reason, we present only the results for square flares.
Let $`L_i`$ ($`i=1,\mathrm{}N`$) be the luminosity of a quasar at an epoch $`t_i`$ in a given spectral band. The SF of this light curve may be defined as follows. If $`t_i`$ and $`t_j`$ are two distinct epochs of the light curve, and $`t_i>t_j`$, then
$$s_{ij}=s(\mathrm{\Delta }t_{ij})=(L_iL_j)^2(ϵ_i+ϵ_j)^2$$
(12)
where $`\mathrm{\Delta }t_{ij}=(t_it_j)/(1+z)`$ is the interval between these two epochs in the quasar rest-frame. In this expression, the observational errors $`ϵ`$ are subtracted in quadrature from the luminosity difference in order to “remove” the observational noise from the SF. The data to be fitted, hence, are the set of pairs $`(\mathrm{\Delta }t_{ij},s_{ij})`$. The SF parameters were estimated by direct minimization of the sum of the square residuals between data and model. Note that no binning was applied to data before the fitting. All but the fourth model in Appendix B have two parameters: a time scale $`\tau `$ and the asymptotic limit $`SF(\mathrm{})=2\sigma _{\lambda ,SF}^2`$ (the fourth model in Appendix B has two time-scales). The errors in the parameters were obtained by bootstrap. In this technique an observed distribution with $`n_{obs}`$ data points is resampled many times, each time picking a random set of $`n_{obs}`$ data points (allowing for repetition) out of the original ones. The underlying hypothesis is that the actually observed values trace the distribution of the measured quantities. For each resampled data set the parameters are fitted, and the parameter errors are estimated from the variance of the values estimated in all resampled data sets. Initially the B and R SFs were fitted separately. A global estimate of $`\tau `$ combining the B and R data was then performed, fitting both SFs with the same value of $`\tau `$ but different $`\sigma _{B,SF}^2`$ and $`\sigma _{R,SF}^2`$.
B and R SFs for five objects are illustrated in Fig. 1 along with their individual and global fits. The SFs in this figure were computed binning the squared luminosity differences in $`\mathrm{\Delta }t`$ such that each bin contained 25 points. This was done just for plotting purposes, as the fitting procedure does not involve binning. The error bars in Fig. 1 were computed bootstrapping the light curves 1000 times.
Though the fits were satisfactory for most objects (e.g., PG 0838 and 1307 in Fig. 1), very often the SFs exhibit complex shapes with ups and downs (e.g., PG 1613 and 2251). Such a non monotonic behavior does not rule out the hypothesis that the light curve is made up by the Poissonian superposition of flares with simple profiles, as large departures from the predicted SF also occur for simulated light curves when these are sampled a finite number of times, as will be explained in Fig. 4 and §3.4.1. Theoretical SFs are derived under the assumption of infinitely long light curves, and it would be naive to expect observed SFs to exhibit the smooth, well behaved shapes predicted by theory. This point has to be kept in mind when evaluating the fits performed in this section. One must therefore exert caution when interpreting the fit parameters, particularly $`\tau `$.
Results of the fits are presented for all 42 sources in Table 1, along with other useful light curve statistics. Though in many cases the life-times inferred from both wavelengths agreed well, we find a large scatter around the $`\tau _B=\tau _R`$ line (see also G99). For the reasons discussed above, we believe this is more likely an artifact of the finite span of the light curves than a real effect, and in what follows we shall only refer to the values of $`\tau `$ obtained with the combined B and R fits. The asymptotic net standard deviations $`\delta _{B,SF}`$ and $`\delta _{R,SF}`$ in Table 1 also correspond to those obtained in the combined fits, but we remark that these were very similar for the individual and global fits.
For four objects, PG 0026, 1100, 1427 and 1626, we find that the SF does not converge over the time span of the observations. All these quasars present rising or decreasing trends over all the length of the observations (Fig. 2 of G99). This suggests either long ($`5`$ yr) time scales, an underlying slowly varying component, or else may be indicative of a non-statistically stationary process. In addition, for PG 1114 and 1512 we find time-scales close to the length of the light curve. These six objects were discarded from further analysis.
Fig. 2 shows how the net variability measured directly from the light curve compares with the SF-determined long term net variability. One finds that $`\delta _{\lambda ,SF}`$ is some 17% larger than $`\delta _\lambda `$ for both B and R, confirming that the latter slightly underestimates the long term variability. Also shown in Fig. 2c is a comparison of the net variabilities in the B and R bands, to illustrate the fact that $`\delta _{B,SF}>\delta _{R,SF}`$ by typically $`30\%`$.
The diversity of time-scales obtained in this analysis (from $`0.5`$ to 3 yr) is consistent with the results of G99, who found a wide range of zero-crossing auto-correlation time-scales. In fact, our values of $`\tau `$ are similar to theirs, but with a substantial scatter. In the context of Poissonian models, this diversity reinforces the suspicion (§1) that quasar flares are not all the same, and that the $`\delta \overline{L}^{1/2}`$ may not be universal. Indeed, in this more general scenario, the “generalized Poissonian model” of Cid Fernandes et al. (1996), every quasar lies in one of a family of $`\delta \overline{L}^{1/2}`$ laws (see Fig. 6). Members of each such family have different rates, but same energy, life-time and background contribution (eqs. and ). Also, this intrinsic diversity casts doubts upon the meaning and usefulness of ensemble SFs.
### 3.3 Constraints on the energies, rate and background contribution
The asymptotic net variabilities $`\delta _{B,SF}`$ and $`\delta _{R,SF}`$ obtained above, along with the values of $`\overline{L}_\lambda `$ and $`\mu _\lambda `$ (listed in Table 1), can be directly applied on eq. (11) to constrain the allowed range of $`𝒩`$, $`c_B`$ and $`c_R`$. We remark that this is nearly independent of the SF analysis, since $`\delta _{\lambda ,SF}`$ is not too different from $`\delta _\lambda `$ (Fig. 2) and $`\tau `$ is not employed in these constraints. The results of this calculation are presented in Table 2.
Except for 3 out of 36 quasars, the upper limit $`𝒩_{max}`$ is always given by the B band, since, as already noted by G99, the variations there are larger than in R (Fig. 2c), in agreement with the general tendency of AGN. As anticipated (§2), the ranges of the obtained $`𝒩`$s are wide (factor of $`20`$). Still, they provide useful estimates of a physically meaningful quantity in Poissonian models. From Table 2 one concludes that the typical number of living events in the G99 quasars is of order 5–100.
A corollary of $`𝒩_{max}`$ being defined by the B light curves is that non trivial lower limits for the background component can only be obtained for the R band for the majority of objects (Table 2). In some cases, like for PG 1309 and 1354, one finds lower limits of as much as 40% in R, indicating a substantial contribution from an underlying constant spectral component. The upper limits for $`c_R`$ are given by $`\mu _R`$, tabulated in Table 1, which for these two quasars are 90 and 86% respectively.
#### 3.3.1 Energies and rates
The $`𝒩`$ and $`c_\lambda `$ limits above make no explicit use of the life-times. Estimates for the allowed range of flare rates and energies require knowledge of $`\tau `$. These are listed in Table 2 for the life-times found in §3.2. Rates between 1 and a few hundred yr<sup>-1</sup>, and monochromatic energies from $`3\times 10^{45}`$ to $`10^{49}`$ erg Å<sup>-1</sup> are found for $`E_B`$ and $`E_R`$. Note that the energies for individual quasars are more constrained (typically to within a factor of 4) than either $`\nu `$ or $`𝒩`$, since they only depend on the square root of $`\nu `$. Furthermore, inspection of Table 2 shows that there is not a single value of either $`E_B`$ or $`E_R`$ that is simultaneously compatible with all lower and upper limits for the G99 quasars. This, again, points to a diversity of flare properties in quasars.
### 3.4 Light curve simulations
The estimates of the basic Poissonian parameters presented here rely on just the mean and minimum luminosities plus the SF-based estimates of the flare life-time and asymptotic variance, this latter quantity being little different from the light curve variance. The actual light curves contain much more information, as they are defined by a particular realization of birth-dates of the flares and their detailed radiative evolution. Retrieving this information from these observations is however an impossible task. We have experimented using higher moments (skewness and kurtosis) of the light curves as further constraints, but this proved fruitless in practice (Appendix A).
Since our estimates are based on so little information, there is no guarantee that a superposition of flares with the inferred properties would bear any morphological resemblance to the observed light curves. In order to verify whether Poissonian models within the bounds defined in §3.3 can produce quasar-looking light curves we have performed a series of Monte Carlo light curve simulations. Our goal here is to broadly assess the “morphological compatibility” of models and data in a qualitative way, based on a simple visual inspection. The simulations make use of the value of $`\tau `$ listed in Table 1, and vary $`\nu `$ within its empirically defined limits for each quasar (Table 2). The flare energy and background luminosity are determined self consistently through eqs. (9) and (10). Square shaped flares were used for consistency with the results already presented.
In Fig. 3 we present three illustrative examples of this experiment. The top curves in each panel show the B band light curves for PG 0052, 1309 and 1404, as observed in the Wise Observatory campaign. The thin solid line in the bottom illustrates a randomly chosen realization of the light curve, computed with $`\nu `$ set to half-way between $`\nu _{min}`$ and $`\nu _{max}`$. The thick solid line shows the simulation which, among thousands of runs, best matches the observed light curve in a $`\chi ^2`$ sense, whereas the dashed curve shows the same simulation after sampling it with the same rest-frame pattern as the data and with Gaussian perturbations added mimicking the sequence of photometric errors reported in G99. All model light curves were vertically displaced for clarity. By construction, their mean luminosity is essentially identical to the observed one, as are their variances.
Considering the simplicity of the model and the requirement of consistency with the constraints defined in §3.3, the resemblance between observations and the models is very satisfactory. The “best-matches” are particularly striking. Even better matches would be obtained fitting the light curves, i.e., adjusting both the global parameters and the individual flare birth-dates to optimize the residuals. Of course, the meaning and usefulness of such detailed fits would be questionable given the number of parameters involved. For instance, there are anywhere between 36 and 145 flares for the length of the PG0052 light curve, and up to 1316 in PG1404. Indeed, with so many degrees of freedom, it might be impossible to ever disprove such models. Unlike the flare dates, however, one has much less liberty to play with the global parameters, and our estimates of $`\tau `$, $`\nu `$, $`E_\lambda `$ and $`c_\lambda `$ go a long way towards constraining the simulations to a region of parameter space capable of producing quasar looking light curves. Simulations with much longer or shorter time scales, for instance, do not resemble the observations at all. We therefore conclude that even a simple Poissonian superposition of square flares is capable of producing light curves which are very similar to quasars and simultaneously compatible with their basic light curve statistics.
#### 3.4.1 Structure Function
The SFs of PG 0052 and 1404 are compared to the corresponding simulated light curves in Fig. 4. As anticipated, relatively large deviations from the smooth theoretical curve occur even for model light curves. The bottom panels show how the SF is improved with light curves twice to five times longer than the observed ones. These were constructed simply extending the sampling of the model by patching together 2 (panels e and f) or 5 (g and h) sequences of the actual observing dates for these two quasars, thus maintaining the G99 sampling pattern. Even for such long (13–45 yr) hypothetical light curves the SF oscillates, though one sees it gradually converging to its statistically expected shape and amplitude. Increasing the sampling rate is not as benefic as increasing the length of the data train, as SF oscillations on intervals $`\mathrm{\Delta }t`$ are only averaged out for $`\mathrm{\Delta }t\tau `$, i.e, as independent (separated by $`\tau `$) portions of the light curve are sampled a statistically significant number of times. This explains why even long light curves exhibit long term SF ‘noise’ but are much better defined on short time scales (Fig. 4e–h).
These experiments illustrate the inherent difficulty in modeling SFs of stochastic processes. As a further example, we note that the SF analysis of a longer (11 yr) light curve of PG 1226 (= 3C 273) by Paltani, Courvoisier & Walter (1988) suggests a value of $`\tau `$ of $`0.6`$ yr, a factor of 4 smaller than the 2.3 yr we found for the 5.5 yr G99 monitoring of this quasar. The SF of PG 1226 is similar to that of PG 1229 in the B band (Fig. 1), with an initial peak at small $`\mathrm{\Delta }t`$ ($`0.4`$ yr), followed by oscillations on longer time scales. Whereas for PG 1229 the global fits favored the small $`\tau `$, in PG 1226 the amplitude of the 0.5 yr peak is smaller, and a longer $`\tau `$ was favored, similar to what is seen in the R band SF of PG 0838 (Fig. 1). Even considering the differences in technicalities of SF analysis between different works, one is forced to conclude that the actual uncertainties in the SF parameters is likely to substantially exceed the formal fit-errors (about 5% of the parameter values).
Overall, we conclude that SF analysis for individual objects is just starting to become possible, despite the excellent quality of the G99 light curves. We are however confident that the fits performed in this paper are correct to within better than an order of magnitude, and thoroughly meet our goal of providing rough but useful estimates for the basic parameters of Poissonian models for the optical variability of quasars.
## 4 Discussion
Poissonian models provide a simple and elegant mathematical framework which relates basic parameters to measurable quantities, as demonstrated by the results presented above. The same technique can be applied to other data sets and spectral bands, thus increasing the number of constraints. In particular, it will be interesting to apply it to the spectrophotometric monitoring observations of 28 quasars, all in the G99 sample, reported in Kaspi et al. (2000). This should yield a clearer picture of the (optical) spectral energy distributions of the variable and constant components, as well as better constrained parameters than was possible with only two wavelengths. An illustration of how important it is to properly define $`\overline{L_\lambda }`$ and $`\delta _\lambda `$ is given in Fig. 5.
The $`E_B`$ limits listed in Table 2 encompass both the estimates of Cid Fernandes et al. (1996) and Garcia et al. (1999) for different samples. Their estimates, however, were based on the properties of the ensemble of quasars, and assumed typical variability time-scales, whereas here we attempted to treat each quasar individually. In fact, this is the first time that constraints on Poissonian parameters are computed on an object by object basis for a large number of quasars.
A further comparison can be made with the results of Paltani et al. (1988) for 3C273. We estimate from their plots a monochromatic flare energy of $`2\times 10^{47}`$ erg Å<sup>-1</sup>, about an order of magnitude smaller than $`E_{B,min}`$ in Table 2. This discrepancy is explained by the already mentioned $`5`$ times smaller $`\tau `$ found by those authors, as well as their $`2`$ times smaller asymptotic net variability compared to the $`\delta _{SF,B}=0.13`$ in Table 2, and another factor of $`1.5`$ due to the redshift effect (see below). The discrepancy in the asymptotic SF variance is mostly due to the fact that they allow for a slowly ($`10`$ yr) varying component, attributed to a blazar behavior. Two of the other 6 radio loud objects in G99 sample (PG 1100 and 1512) show evidence for long time-scale variations which could be due to a similar slow component (marginal evidence that the radio loud objects in this sample present longer time-scales can be seen in Fig. 6c). It is in fact adequate to recall that AGN variability has often been described in terms of the superposition of rapid and slowly varying components (e.g., Lyutyi 1977, Pica et al. 1988). A caveat in the present analysis is thus that what we have been calling a “constant component” may well correspond to an underlying slow process.
We note that no attempt was made to correct any of the variability indices nor the inferred parameters for the variability-$`\lambda `$ effect which plague fixed band photometric light curves of quasars, leading to overestimated variability amplitudes (Cristiani et al. 1996; Cid Fernandes et al. 1996; Aretxaga et al. 1997; Garcia et al. 1999). Fortunately, the low redshifts of the G99 quasars minimize this effect. For the median $`z`$ of the sample, and using the parameterization of the $`\delta \times \lambda `$ relation of Garcia et al. (1999), we estimate that the variability indices in Table 1 (columns 4 to 7) would need to be multiplied by typically 0.8 to obtain the rest-frame indices. This would increase $`𝒩_{min}`$, $`𝒩_{max}`$ and the corresponding $`\nu `$ limits by some 56% while reducing the energies by a factor of 0.64. These are relatively small factors considering the allowed ranges for the parameters and were not applied in Table 2. Furthermore, such corrections would soon become obsolete since the spectral data collected by Kaspi et al. (2000) allows the direct computation of rest-frame continuum variability indices. We did, however, experiment with the $`\lambda `$-correction in the correlation analysis presented next.
### 4.1 Correlation Analysis
Perhaps the most puzzling result of G99 was the finding that there are few convincing correlations between variability indices and a large array of other observed properties for the PG quasars. In Table 3 we synthesize the results of a correlation analysis analogous to that performed by G99, but for our SF-based variability indices. The table is similar to their Table 4, from which the multi-wavelength data was borrowed, and lists the percentage probabilities ($`P_r`$) of no correlation in a Spearman’s rank test, small values indicating significant correlations and negative values corresponding to anticorrelations. Entries between parentheses correspond to correlations after correction of the variability amplitudes with the $`\delta `$-$`\lambda `$-$`z`$ relation following Garcia et al. (1999).
The anti-correlation between the variability amplitude (here measured by $`\delta _{SF,B}`$ and $`\delta _{SF,R}`$) and the optical luminosity found in other studies (e.g., Hook et al. 1994) is essentially absent for this sample, as illustrated in Fig. 6a. Somewhat more significant anti-correlations are obtained by applying the $`\lambda `$ correction (Fig. 6b; Table 3), but in both cases the correlations here are even weaker than those found by G99. The dotted lines in Fig. 6a and b mark the prediction of simple Poissonian models in which $`c_\lambda `$, $`E_\lambda `$ and $`\tau `$ are identical for all quasars. This illustrates the often claimed failure of this model (e.g. Hook et al. 1994), and the necessity to allow for a diversity of values for these parameters to account for the observed variability properties in the framework of Poissonian processes.
The variability time-scale seems to increase with luminosity (Fig. 6c), but, as in G99, with a large scatter. Previous evidence for such a correlation was reported by Cristiani et al. (1996), in an analysis of the ensemble SF of 486 objects spanning a much larger range of luminosities and redshifts. It is interesting to note that applying their Model E SF-fits to the typical absolute magnitudes of G99 quasars yield $`\tau `$ between $`0.5`$ and 2.5 yr, compatible with our individual SF fits. We also find $`\tau `$ to correlate positively with the $`\delta _{SF}`$’s, but this may be due to a trade-off effect in the SF fitting (§3.4.1). The weak anticorrelation with $`\alpha _{ox}`$ found by G99 is much stronger here (Fig. 6d).
The good correlations between the variability amplitude and the equivalent widths of H$`\beta `$ and HeII$`\lambda `$ 4686 found by G99 are reproduced at about the same high level of significance (Fig. 7a). Unlike G99, however, we find the optical to X-ray index $`\alpha _{ox}`$ to be significantly anticorrelated with both $`\delta _{SF,B}`$ (Fig. 7d) and $`\delta _{SF,R}`$. All these correlations are improved with the $`\lambda `$-correction.
A correlation was found between the ratio of the $`B`$ to $`R`$ asymptotic standard deviations and $`\overline{L_B}`$ as well as $`z`$. This could be due to the redshift effects already discussed, but in this case a $`\delta `$-$`\lambda `$ law different from that of Garcia et al. (1999) would be implied, since the $`z`$-effects upon $`\sigma _{B,\mathrm{SF}}/\sigma _{R,\mathrm{SF}}`$ cancels out in their parameterization. Alternatively, if this interpretation is proved wrong, one would conclude that the flares are bluer in more luminous sources (eq. 9).
Overall, we largely confirmed G99 results, with a few minor differences (e.g., Fig. 6d) undoubtedly due to the different methodologies employed in the definition of variability amplitudes and time-scales. Despite the few good correlations found, whose meaning we discuss below, a striking result of this analysis is the large scatter present even in the best correlations obtained. In any scenario for AGN, the variability properties must somehow be dictated by the physical conditions prevalent in the active nucleus (accretion rate, black hole mass, orientation, etc.). Since these conditions must also define other aspects of AGN phenomenology, it is natural to expect that properties driven by common causes should exhibit some degree of correlation. Given the controversial history of correlations in quasar variability studies, however, it is perhaps not so surprising that so few good correlation were found.
Regarding Poissonian models, one would obviously like to search for correlations between the basic parameters and other properties. The constraints upon $`𝒩`$, $`\nu `$, $`E_\lambda `$ and $`c_\lambda `$ derived in this paper are however not strong enough to warrant a proper correlation analysis. Taking the middle value between $`𝒩_{min}`$ and $`𝒩_{max}`$ as a measure of $`𝒩`$, we obtain an anticorrelation between $`𝒩`$ and EW(H$`\beta `$) and a positive correlation with $`\alpha _{ox}`$, both with $`P_r`$ at the 1% level. The flare rate follows the same trends. For the reason discussed above, it would be premature to give much emphasis to these correlations.
#### 4.1.1 Interpretation of the correlations with EW(H$`\beta `$)
A simple interpretation of the significant correlation between EW(H$`\beta `$) and the variability amplitude is possible by postulating that the variable component dominates the ionization of the gas. This idea is also compatible with our earlier conclusion (§2) that the variability-wavelength anticorrelation indicates the dilution of a blue variable component by a red underlying background. In this hypothesis both variable ($`V_\lambda `$) and constant ($`C_\lambda `$) components contribute to the continuum under H$`\beta `$, but only the former is proportional to the ionizing luminosity ($`L_{ion}`$), so one predicts
$$EW(H\beta )\frac{\nu E_B}{\nu E_B+C_B}$$
(13)
where we used the B band because of its proximity to H$`\beta `$. This relation may be rewritten as
$$EW(H\beta )𝒩^{1/2}\delta _{SF,B}=v_B=1c_B$$
(14)
which reveals a proportionality between EW(H$`\beta `$) and $`\delta _{SF,B}`$. Obviously, the same prediction applies to EW(HeII).
Using $`𝒩=(𝒩_{min}+𝒩_{max})/2`$, the product $`𝒩^{1/2}\delta _{SF,B}`$ correlates at the $`P_r=0.6\%`$ level with EW(H$`\beta `$), in agreement with the prediction. Since our estimates of $`𝒩`$ are not independent of $`\delta _{SF,B}`$ (eq. 11), it is perhaps more meaningful to test the predicted relation using $`c_B`$, for which we have a robust upper limit imposed by the minimum in the light curve ($`\mu _B`$ in Table 1). In Fig. 7b we see that the expected anticorrelation is confirmed with a high significance (Table 3). Even considering that $`\mu _B`$ is an upper limit to $`c_B`$ and the scatter in the plot, it is interesting to see that the trend is roughly linear, as predicted.
$`\mu _B`$ is also strongly correlated with $`\alpha _{ox}`$ ($`P_r=0.1\%`$), the X-ray spectral index $`\alpha _x`$ ($`P_r=1.8\%`$), and anticorrelated with EW(\[OIII\]) ($`P_r=0.8\%`$), \[OIII\] to H$`\beta `$ peak intensities ratio ($`P_r=0.2\%`$) and EW(HeII) ($`P_r=0.5\%`$). The correlation with $`\alpha _{ox}`$, in particular, gives strength to our working hypothesis, insofar as this index can be interpreted as indicative of the ratio between the constant and variable (ionizing) components.
By admitting that different quasars have different background fractions, and that $`V_\lambda `$ dominates the ionization, we naturally explain at least some of the correlations observed. This argues against $`C_\lambda `$ being part of the ionizing source, such as hot but non-flaring portions of the accretion disk surface. Outer, colder regions of the disk can contribute to $`C_\lambda `$, as do the host galaxy starlight and other circum-nuclear sources.
The alternative hypothesis that $`L_{ion}`$ is proportional to the total $`V_\lambda +C_\lambda `$ components would not induce the correlations above, while the hypothesis that $`L_{ion}`$ is governed by the constant component alone can be straightforwardly rejected, since this would not produce emission line variability.
We note that our eq. (13) is formally and philosophically identical to eq. (10) of Aretxaga & Terlevich (1994), who identify $`C_\lambda `$ with a young nuclear star-cluster and $`V_\lambda `$ with compact SN remnants exploding out of the same cluster (see also Aretxaga et al. 1997). Their analysis of EW(H$`\beta `$) aimed to understand the near universal value of this quantity across the AGN luminosity scale (Binette, Fosbury & Parker 1993). In this model, a range of EW(H$`\beta `$) values could be linked to either different SN explosion energies, or, as is more likely, to different H$`\beta `$ and continuum production efficiencies in the SN remnants (see Cid Fernandes 1997 for a review of the pros and cons of the starburst model). Stochastic effects could also play a role in defining the EW(H$`\beta `$) and variability relation, something which can happen in any Poissonian scenario. Indeed, the tentatively identified anticorrelation between $`𝒩`$ (and $`\nu `$) and EW(H$`\beta `$) may be indicative of stochastic effects in the G99 data, and may be partly responsible for the scatter in Fig. 7a and d.
### 4.2 Relation to physical models
The ultimate goal of the Poissonian description of quasar variability is of course to extract information on the fundamental properties of the variability phenomenon and use them to guide and discriminate between physical theories. A detailed discussion of the implications of the estimates presented in §3 to different scenarios for AGN variability is beyond the scope of the present paper, but we would like to take one particular model as an illustration of the possible associations between the observationally constrained parameters discussed here and physically meaningful quantities.
In the model proposed by Haardt et al. (1994), blobs emerge from the surface of an accretion disk and release magnetically stored energy in the form of rapid X-ray flares. In the X-ray regime, the background fraction in this model would be associated with the fraction of the disk area covered by blobs (see Figure 1 of Galeev, Rosner & Vaiana 1979), which equals the ratio of charge and discharge times. As they postulate the optical-UV variability to be driven by the X-ray flares, $`c_\lambda `$ is related to the fraction of the reprocessed luminosity and the optical-UV emission from the underlying stable component of the disk, but note that comparison with the limits on $`c_\lambda `$ derived in this paper require allowance for the diluting effects of larger scale sources, such as the host galaxy or an extended scattering region. One of the predictions of this model is that at any time there are of order 10 active blobs, independent of the source luminosity, consistent with the limits on $`𝒩`$ established here. Also the reprocessed luminosity, which we would identify with $`(1c)\overline{L}=\nu E`$ in our formalism, is predicted to be similar to the X-ray luminosity. It is not clear, however, if the time smearing of the minutes long flares during the disk illumination which results in the reprocessing of the blobs energies can be made compatible with the $``$ 0.5–3 yr time scales implied by the SF analysis. Surely this must impose some limits upon the geometry of the reprocessing surface. Furthermore, it remains to be established whether the flare energies can be made compatible with the limits laid out in §3.3. A more detailed scrutiny of the Haardt et al. (1994) model must therefore await more theoretical developments. On the observational side, a Poissonian analysis of X-ray light curves of quasars (similar to that done by Cid Fernandes 1995 for EXOSAT light curves of Seyferts) could also set constraints upon this particular model.
Despite their generality and appealing aspects, Poissonian models are by no means the only possible description of AGN variability. For instance, Mineshige & Shields (1990) showed that thermal limit cycles similar to those known to occur in dwarf novae (Warner 1995) can also take place in AGN disks, leading to eruptions followed by quiescent periods. Such a model, in which the degree of variability is “self-regulated”, would invalidate the formalism employed here, which fundamentally rests upon the hypothesis of independence of the flares and statistical stationarity. The same applies to scenarios involving random walk-like or other kinds of state dependent behavior (e.g., Begelman & De Kool 1991; Stern, Svensson & Sikora 1991) or periodic modulations associated to, for instance, precessing jets (e.g., Abraham & Romero 1999). Testing the Poissonity of AGN light curves is however a hard task. Fourier phase coherence studies like that performed by Krolik, Done & Madejski (1993) for X-ray light curves of Seyferts, or intermittency tests (Vio et al. 1992) like the one carried out by Longo et al. (1996) for the historical light curve of NGC 4151, can in principle help discriminating between Poissonian models and scenarios where the variations are due to coherent oscillations of a single entity, but these require more data points than are currently available from quasar optical monitoring studies. Therefore, until proven wrong, the Poissonian description may be regarded as a useful tool to unveil physically meaningful properties of AGN.
## 5 Summary and Conclusions
We have reviewed the Poissonian formalism for quasar variability in an attempt to provide a general framework which allows fundamental parameters to be estimated from good monitoring data. This was applied to the 6 yr long B & R light curves of 42 PG quasars obtained by G99, yielding constraints for the energy, rate and lifetimes of the flares. Our main results can be summarized as follows.
(1) The only reasonable way to account for the fact that quasars vary more at shorter wavelengths within a Poissonian scenario is to include the diluting effects of an underlying “non-variable background” component redder than the spectral energy distribution of the putative flares.
(2) A wide range of flare energies, lifetimes, and/or background fractions has to be invoked to account for the observed variability properties (amplitudes and time-scales). This “stretching” of the simplest Poissonian scenario (in which all parameters are the same for every quasar) is warranted by the model independent result that quasars present a wide range of variability time scales, from $`0.5`$ to more than 3 yr, as inferred from a Structure Function analysis.
(3) The mean number of living flares is constrained to be of order $`𝒩5`$ to 100, and lower limits for the R-band background contribution of typically 25% are established. These estimates are independent of cosmology, K-correction and extinction, and little sensitive to the SF analysis.
(4) Flare rates between $`1`$–100 yr<sup>-1</sup> and monochromatic flare energies in the $`10^{4648}`$ erg Å<sup>-1</sup> range are implied by the data. Overall, the Poissonian parameters for individual quasars are constrained to within about an order of magnitude.
(5) Light Curve simulations were performed and demonstrate the ability to reproduce the observed morphology of quasars light curves extremely well even for a blatantly simplistic “on/off” square-shaped model for the evolution of the individual flares. Indeed, the whole Poissonian analysis is highly insensitive to the flare shape.
(6) Experiments were performed on the use of higher moments of the light curve as further constraints, but found to be of little use at present.
(7) The variability properties of the PG quasars present little correlation with other properties. Even the best correlations identified, like that between the variability amplitude and EW(H$`\beta `$), present a substantial scatter, confirming the results of G99.
(8) The EW(H$`\beta `$) $`\times `$ variability amplitude was interpreted in a scenario where only the variable component participates in the ionization of the line emitting gas, consistent with conclusion (1) above. Correlations with EW(HeII) and the X-ray to optical spectral index further support this interpretation.
Progress on this line of approach to AGN variability will require longer light curves to obtain more accurate estimates of the variability properties, which are important to constrain both Poissonian and non-Poissonian models. Regardless of possible improvements on the time-series analysis techniques, it is important to encourage the continuation of the current CCD monitoring campaigns. Spectral information, some of which is already available (Kaspi et al. 2000), will also be valuable in deriving better constraints for the model parameters and a more detailed picture of the spectral behavior of the flares and background component.
The most appealing aspect of Poissonian models is their generality. Scenarios as diverse as accretion disk instabilities, stellar collisions and supernovae all fall under the large Poissonian “umbrella”. Detailed modeling will be required to bridge the gap between this mathematical formalism and the physics of AGN variability and to fully explore the constraints made possible by the application of this technique.
It is a pleasure to thank the Wise Observatory group for their generous attitude of releasing to the scientific community the results of a laborious and lengthy observational project. We also thank T. Heckman and J. Krolik for their comments on an earlier version of this paper. RCF thanks the hospitality of Johns Hopkins University, where this work was developed. Support for this work was provided by the National Science Foundation through grant # GF-1001-99 from the Association of Universities for Research in Astronomy, Inc., under NSF cooperative agreement AST-9613615. LV acknowledges an MSc fellowship awarded by CAPES. Partial support from CNPq, FAPESP and PRONEX are also acknowledged.
## Appendix A Skewness and Kurtosis in Poissonian models
The Poissonian formalism allows the computation of higher moments of the $`L(t)`$ distribution, though most applications to AGN variability stop on the variance (eq. ). In this appendix we present expressions for the theoretically expected third (skewness) and fourth (kurtosis) moments, and discuss their applicability to the G99 data set. These moments were not used to constrain the flare properties of quasars in the main text, and are presented here only for the purposes of completeness and future reference.
The skewness ($`\gamma `$) and kurtosis ($`\kappa `$) are defined by
$$\gamma =\sigma ^3\overline{[L\overline{L}]^3}$$
(A1)
$$\kappa =\sigma ^4\overline{[L\overline{L}]^4}3$$
(A2)
The $`3`$ term in (A2) makes the kurtosis of a Gaussian distribution 0. Expressions for $`\gamma `$ and $`\kappa `$ can be derived extending the formalism employed by Cid Fernandes (1995), which yields (Vieira 2000)
$$\gamma =\varphi _\gamma (\nu \tau )^{1/2}$$
(A3)
$$\kappa =\varphi _\kappa (\nu \tau )^1$$
(A4)
In these expression $`\varphi _\gamma `$ and $`\varphi _\kappa `$ play the role of “shape factors”:
$$\varphi _\gamma =\frac{\tau ^2}{E^3}l^3(t)𝑑t$$
(A5)
$$\varphi _\kappa =\frac{\tau ^3}{E^4}l^4(t)𝑑t$$
(A6)
Both $`\varphi _\gamma `$ and $`\varphi _\kappa `$ equal 1 for square shape flares. Exponential flares result in $`\varphi _\gamma =4/3`$ and $`\varphi _\kappa =2`$, while for triangular flares $`\varphi _\gamma =9/8`$ and $`\varphi _\kappa =27/20`$. The above relations reveal interconnections between the different moments in a Poissonian model: $`\gamma \delta `$ and $`\kappa \delta ^2`$, where $`\delta =\sigma /\overline{L}`$ (eq. ). Unlike for $`\delta `$, an underlying non-variable component does not affect $`\gamma `$ nor $`\kappa `$. In principle, this allows an estimate of the background fraction from the ratio of $`\gamma `$ (or $`\kappa `$) to $`\delta `$.
The main problem concerning the use of higher light curve moments is that, as can be seen in eqs. (A3) and (A4), the predicted moments rapidly tend to the Gaussian limit ($`\gamma =\kappa =0`$) with increasing $`𝒩=\nu \tau `$, i.e., as the superposition of events grows higher. This is obviously a consequence of the Central Limit Theorem (e.g., Papoulis 1965).
The detection of deviations from Gaussianity in quasar light curves is known to be problematic (Press & Rybicki 1997), and the G99 light curves are no exception. In Fig. 8 we present the skewness and kurtosis for the Wise Observatory data. Corrections for the effects of photometric errors upon the moments were applied, but were negligible, given the excellent accuracy of the G99 photometry. At first sight, the results in Fig. 8 would seem to immediately rule out any Poissonian model, as one finds negative $`\gamma `$ and $`\kappa `$ for about half of the objects, whilst the theory (eqs. \[A3\] and \[A4\]) predict only positive values!
This, however, is likely to be an effect of sampling. To demonstrate this, we have run a series of Poissonian light curve simulations and compared the output 2nd, 3rd and 4th moments with the predicted values as a function of observational parameters such as the length of the light curve ($`T_{obs}`$) and the number of observations ($`n_{obs}`$). The results clearly show that negative $`\gamma `$ and $`\kappa `$ occur very often also in simulated light curves, particularly for $`T_{obs}10\tau `$, as illustrated in Fig. 9. The scatter decreases for increasing $`n_{obs}`$, but the agreement between analytical and simulated moments is only achieved for large $`n_{obs}`$ and $`T_{obs}\tau `$. This bias is essentially insensitive to flare profile or their rate. The conclusion here is that one cannot use the currently available data to strongly constrain the higher moments of the $`L(t)`$ process. On the positive side, these experiments reinforce our conclusion that the actual flare shape is irrelevant for most of the analysis presented in this paper (see also next section).
We finalize by noting that, in analogy with what was done for the second moment, higher moments should be computed from their asymptotic ($`\mathrm{\Delta }t\mathrm{}`$) behavior estimated via SFs of the corresponding order. These, however, are subjected to large uncertainties due to the high powers involved and are not presented here.
## Appendix B Structure Functions for simple flare profiles
The structure function of a Poissonian sequence of flares is:
$$SF(\mathrm{\Delta }t)=SF(\mathrm{})s(\mathrm{\Delta }t)$$
(B1)
where $`SF(\mathrm{})=2\nu E^2/\tau `$ is twice the asymptotic variance of the $`L(t)`$ process, and $`s(\mathrm{\Delta }t)`$, the normalized SF of individual flares, is given by
$$s(\mathrm{\Delta }t)=1\frac{l(t+\mathrm{\Delta }t)l(t)𝑑t}{l^2(t)𝑑t}$$
(B2)
and is sensitive only to the shape of flare time profile $`l(t)`$.
### B.1 Square flares
For $`l(t)=l_0`$ between $`0<t<T`$ and 0 otherwise one obtains a a linear SF up to $`\mathrm{\Delta }t=T`$, with
$$s(\mathrm{\Delta }t)=\frac{\mathrm{\Delta }t}{T}$$
(B3)
and 1 for $`\mathrm{\Delta }t>T`$. The life-time $`\tau `$ (given by eq. 5) of square flares equals $`T`$. The dotted lines in Fig. 10 show the resulting SF.
### B.2 Exponentially decaying flares
For $`l(t)=l_0e^{t/T}`$ flares one finds
$$s(\mathrm{\Delta }t)=1e^{\mathrm{\Delta }t/T}$$
(B4)
while the life-time is $`\tau =2T`$. This SF is shown as a dot-dashed line in Fig. 10.
### B.3 Symmetric Triangular flares
Linear flares with equal rise and decay times
$`l(t)`$ $`=`$ $`l_0\times \{\begin{array}{ccc}(1+t/T)\hfill & ;\mathrm{for}& \hfill Tt0\\ (1t/T)\hfill & ;\mathrm{for}& \hfill 0tT\end{array}`$ (B7)
produce the following SF:
$`s(\mathrm{\Delta }t)=`$ $`\{\begin{array}{ccc}\frac{3}{4}\left(\frac{\mathrm{\Delta }t}{T}\right)^3+\frac{3}{2}\left(\frac{\mathrm{\Delta }t}{T}\right)^2\hfill & ;\mathrm{for}& \hfill 0\mathrm{\Delta }tT\\ +\frac{1}{4}\left(\frac{\mathrm{\Delta }t}{T}\right)^3\frac{3}{2}\left(\frac{\mathrm{\Delta }t}{T}\right)^2+3\frac{\mathrm{\Delta }t}{T}1\hfill & ;\mathrm{for}& \hfill T\mathrm{\Delta }t2T\\ 1\hfill & ;\mathrm{for}& \hfill 2T\mathrm{\Delta }t\end{array}`$ (B11)
The relation between $`\tau `$ and $`T`$ for triangular shots is $`\tau =3T/2`$. The solid line in Fig. 10 shows this SF.
### B.4 Asymmetric Triangular flares
Triangular flares with unequal rise ($`T_r`$) and decay ($`T_d`$) times
$`l(t)`$ $`=`$ $`l_0\times \{\begin{array}{ccc}(1+t/T_r)\hfill & ;\mathrm{for}& \hfill T_rt0\\ (1t/T_d)\hfill & ;\mathrm{for}& \hfill 0tT_d\end{array}`$ (B14)
produce more complicated SFs:
$`s(\mathrm{\Delta }t)=`$ $`\{\begin{array}{ccc}\frac{1}{2}\frac{\mathrm{\Delta }^3}{T}\left(\frac{1}{T_r^2}+\frac{1}{T_rT_d}+\frac{1}{T_d^2}\right)+\frac{3}{2}\frac{\mathrm{\Delta }^2}{T_rT_d}\hfill & ;\mathrm{for}& 0\mathrm{\Delta }tT_r\hfill \\ 1\frac{1}{2}\frac{\mathrm{\Delta }^3}{TT_d^2}+\frac{3}{2}\frac{\mathrm{\Delta }}{T_d}\frac{3}{2T}\left(\frac{T_r^2}{6T_d}+\frac{T_r}{2}+\frac{T_d}{3}\right)\hfill & ;\mathrm{for}& T_r\mathrm{\Delta }tT_d\hfill \\ \frac{1}{2}+\frac{1}{2}\frac{\mathrm{\Delta }^3}{TT_rT_d}+\frac{3}{2}\frac{\mathrm{\Delta }^2}{T_rT_d}\hfill & & \\ +3\frac{\mathrm{\Delta }}{T}\left(1+\frac{T_d}{2T_r}+\frac{T_r}{2T_d}\right)\frac{1}{2T}\left(\frac{T_d^2}{T_r}+\frac{T_r^2}{T_d}\right)\hfill & ;\mathrm{for}& T_d\mathrm{\Delta }tT\hfill \\ 0\hfill & ;\mathrm{for}& T\mathrm{\Delta }t\hfill \end{array}`$ (B20)
where $`T=T_r+T_d`$. The solution above corresponds to $`T_r<T_d`$. The corresponding expression for flares that spend more time rising than decaying is identical, with $`T_r`$ swapped by $`T_d`$. The life-time in either case is given by $`\tau =3T/4`$. G99 finding that quasars spend more time on the rise than fading gives some motivation to using $`T_r>T_d`$, but the SF, being a mean of squared differences, does distinguish between $`T_r`$ and $`T_d`$ (see Kawaguchi et al. 1998 for techniques to explore the time-assymetry of flares). Fig. 10 shows the SF for a ratio of 10 between $`T_r`$ and $`T_d`$ (or vice-versa) as a dashed line.
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# Every diassociative A-loop is Moufang
## 1. Introduction
A *loop* $`(L,)`$ consists of a nonempty set $`L`$ with a binary operation $``$ on $`L`$ such that (i) given $`a,bL`$, the equations $`ax=b`$ and $`ya=b`$ each have unique solutions $`x,yL`$, and (ii) there exists an identity element $`1L`$ satisfying $`1x=x1=x`$ for all $`xL`$. As usual, we abbreviate the binary operation by juxtaposition. Two varieties of loops which have been widely discussed in the literature are the Moufang loops and the A-loops.
A *Moufang loop* is a loop satisfying the identity
$$x(yxz)=(xyx)z.$$
(1.1)
These were introduced by R. Moufang in 1934 , and are discussed in detail in the texts by Bruck and Pflugfelder . By Moufang’s Theorem (, VII.4; , IV.2.9), every Moufang loop is diassociative; that is, the subloop $`x,y`$ generated by any pair of elements is a group.
For $`xL`$, the left and right translations by $`x`$ are defined by $`yL(x)=xy`$ and $`yR(x)=yx`$, respectively. The *multiplication group* of $`L`$ is the permutation group $`\mathrm{Mlt}(L)=R(x),L(x):xL`$ generated by all left and right translations. The *inner mapping group* is the subgroup $`\mathrm{Mlt}_1(L)`$ fixing $`1`$. If $`L`$ is a group, then $`\mathrm{Mlt}_1(L)`$ is the group of inner automorphisms of $`L`$.
In 1956, R.H. Bruck and L.J. Paige defined an *A-loop* to be a loop in which every inner mapping is an automorphism. Many of the basic theorems about A-loops are contained in ; for example, A-loops are always power associative (every $`x`$ is a group), but not necessarily diassociative. In the same paper, Bruck and Paige included a detailed study of the diassociative A-loops, pointing out that these satisfy “many of the properties of Moufang loops”. In hindsight, this is not surprising, since, as we will show:
###### Theorem 1.
Every diassociative A-loop is a Moufang loop.
For commutative loops, this was proved in 1958 by J.M. Osborn . Conversely, every commutative Moufang loop is an A-loop (see Bruck , Lemma VII.3.3). However, not all Moufang loops are A-loops; , together with the results of the present paper, provide a simple description of the diassociative A-loops as a sub-variety of the Moufang loops (see Corollary 2). Further work on A-loops is contained in .
By our Theorem 1, we have:
###### Corollary 1.
For an A-loop, the following are equivalent:
1. $`L`$ *has the* inverse property, *i.e.,* $`x^1(xy)=y`$ *and* $`(xy)y^1=x`$ *for all* $`x,yL`$;
2. $`L`$ *has the* alternative property, *i.e.,* $`x(xy)=x^2y`$ *and* $`(xy)y=xy^2`$ *for all* $`x,yL`$;
3. $`L`$ *is diassociative.*
4. $`L`$ *is a Moufang loop*.
The equivalence of the first three items is from Bruck and Paige , Theorem 3.1. (One may begin with even weaker hypotheses, but we will not pursue this here.)
In any loop, the inner mapping group $`\mathrm{Mlt}_1(L)`$ is generated by the left, right, and middle inner mappings defined, respectively, by:
$`L(x,y)`$ $`=L(x)L(y)L(yx)^1`$
$`R(x,y)`$ $`=R(x)R(y)R(xy)^1`$
$`T(x)`$ $`=R(x)L(x)^1`$
(, IV.1, , I.5.2). Bruck and Paige (, (3.42)) showed that diassociative A-loops satisfy:
$$\left[R(x,y)R(y,x)\right]^1T(x)T(y)=T(xy)$$
(1.2)
Furthermore, they showed (see Corollary on p. 315) that for Moufang A-loops, the map $`T:L\mathrm{Mlt}_1(L)`$ (where $`xT(x)`$) is a homomorphism (i.e., $`T(x)T(y)=T(xy)`$, so $`R(x,y)=R(y,x)^1`$). Not surprisingly, one of our key lemmas will be:
###### Lemma 1.
If $`L`$ is a diassociative A-loop, then $`T:L\mathrm{Mlt}_1(L)`$ is a homomorphism.
The nucleus, $`\mathrm{Nuc}(L)`$, of an inverse property loop $`L`$ is the normal subloop of all elements that associate with all pairs of elements from $`L`$, i.e., $`\mathrm{Nuc}(L)=\{xL:(xy)z=x(yz)`$ for all $`y,zL\}`$. By results already in the literature, we have the following corollary to Theorem 1:
###### Corollary 2.
$`L`$ is a diassociative A-loop if and only if $`L`$ is Moufang and $`L/\mathrm{Nuc}(L)`$ is a commutative loop of exponent three.
###### Proof.
In any Moufang loop, each $`T(x)`$ is a pseudo-automorphism with companion $`x^3`$, and each $`R(x,y)=L(x^1,y^1)`$ is a pseudo-automorphism with companion the commutator $`(x,y)`$ (, Lemma VII.2.2). In general, if $`c`$ is a companion of the pseudo-automorphism $`\phi `$, then $`c`$ is in the nucleus iff $`\phi `$ is an automorphism. Thus all cubes and commutators are in the nucleus iff all inner mappings are automorphisms. ∎
Every Moufang A-loop is an $`M_4`$ loop in the terminology of Pflugfelder ; that is, it satisfies the identity $`(xy)(zx^4)=(xyz)x^4`$ (since cubes are in the nucleus and $`(xy)(zx)=(xyz)x`$ is a Moufang identity). We do not know whether an $`M_4`$ loop $`L`$ must be an A-loop. By , Theorem 1, $`L`$ is Moufang and $`L/\mathrm{Nuc}(L)`$ has exponent three, but it is not clear whether $`L/\mathrm{Nuc}(L)`$ is necessarily commutative. We also do not know whether every loop isotope of a Moufang A-loop is a (Moufang) A-loop. This would be true if every $`M_4`$ loop is an A-loop, since the $`M_k`$ loops are isotopically invariant (, Theorem 2; , IV.4.12)
Our investigations were aided by the automated deduction tool OTTER developed by McCune ; see Section 4 for further discussion.
## 2. Preliminaries
In preparation for the proofs of Lemma 1 and Theorem 1, we now establish some notation and recall some basic results from . Let $`L`$ be a diassociative A-loop. One can then derive many equations relating the the $`L(x,y)`$, $`R(x,y)`$, and $`T(x)`$.
Define the permutation $`J`$ of $`L`$ by: $`xJ=x^1`$. Conjugating by $`J`$, we have $`R(x)^J=JR(x)J=L(x^1)`$; likewise, $`L(x)^J=R(x^1)`$ and $`L(x,y)^J=R(x^1,y^1)`$. Note that $`\phi ^J=\phi `$ for all automorphisms $`\phi `$ of $`L`$; in particular for all $`\phi \mathrm{Mlt}_1(L)`$. Taking $`\phi =L(x,y)`$, we have:
$$L(x,y)=R(x^1,y^1)$$
(2.1)
for $`x,yL`$. Furthermore, from , ((3.31) and (3.32)) we have the following formulas for the inverses of the right and left inner mappings:
$`R(x,y)^1`$ $`=R(y^1,x^1)`$ (2.2)
$`L(x,y)^1`$ $`=L(y^1,x^1)`$ (2.3)
The fact that each $`T(x)`$ is an automorphism implies immediately:
$`R(y)T(x)`$ $`=T(x)R(x^1yx)`$ (2.4)
$`L(y)T(x)`$ $`=T(x)L(x^1yx)`$ (2.5)
Another useful inner mapping is defined by
$$C(x,y)=R(x)L(y)R(x^1)L(y^1).$$
(2.6)
Since $`C(x,y)^J=C(x,y)`$, we also have:
$$C(x,y)=L(x^1)R(y^1)L(x)R(y).$$
(2.7)
Also, by (3.41):
$$C(x,y)=R(x,y)R(y,x)^1.$$
(2.8)
Further equations relating the $`C(x,y),R(x,y),L(x,y)`$ will be proved later (see Corollaries 3 and 4). As pointed out in , in any loop, if $`\phi `$ is an automorphism which fixes an element $`p`$, then $`\phi `$ commutes with $`L(p)`$ and $`R(p)`$. In particular (, Lemma 3.3(i,ii,iii)), if $`p,q,r`$ are contained in any subgroup of $`L`$, then:
$`R(p)R(q,r)`$ $`=R(q,r)R(p)\text{}L(p)R(q,r)=R(q,r)L(p)`$ (2.9)
$`R(p)L(q,r)`$ $`=L(q,r)R(p)\text{}L(p)L(q,r)=L(q,r)L(p)`$ (2.10)
$`R(p)C(q,r)`$ $`=C(q,r)R(p)\text{}L(p)C(q,r)=C(q,r)L(p).`$ (2.11)
One consequence is that the factors in the right and left inner mappings can by cyclically permuted:
$`R(x,y)`$ $`=R(y)R(y^1x^1)R(x)=R(y^1x^1)R(x)R(y)`$ (2.12)
$`L(x,y)`$ $`=L(y)L(x^1y^1)L(x)=L(x^1y^1)L(x)L(y).`$ (2.13)
## 3. Proofs
###### Proof of Lemma 1.
For $`x,y,zL`$, we compute
$`zL(xy)T(x)`$ $`=yL(x)R(z)T(x)`$
$`=yC(x^1,z^1)R(z)L(x)T(x)`$ $`\mathrm{by}(\text{2.7})`$
$`=yC(x^1,z^1)R(z)R(x)`$
$`=yR(z)R(x)C(x^1,z^1)`$ $`\mathrm{by}(\text{2.11})`$
$`=yR(z)L(z^1)R(x)L(z)`$ $`\mathrm{by}(\text{2.6})`$
$`=yT(z)R(x)L(z).`$
By the mirror of this calculation and switching $`x`$ and $`y`$, we obtain:
$$zR(xy)T(y^1)=xT(z^1)L(y)R(z).$$
But by (2.4), we have
$$yT(z)R(x)L(z)=yR(zxz^1)R(z)=xT(z^1)L(y)R(z).$$
Hence, $`L(xy)T(x)=R(xy)T(y^1)`$, so that $`T(x)T(y)=L(xy)^1R(xy)=T(xy)`$. ∎
###### Corollary 3.
$`R(x,y)^1`$ $`=R(y,x)`$ (3.1)
$`R(x,y)`$ $`=R(x^1,y^1)=L(x,y)=L(x^1,y^1)`$ (3.2)
$`C(x,y)`$ $`=C(x^1,y^1)=R(x,y)^2.`$ (3.3)
###### Proof.
(3.1) follows from Lemma 1 and (1.2). To get (3.2), apply (2.2) and (2.1). Then, (3.3) follows by using (2.8). ∎
###### Lemma 2.
For all $`x,y,z`$ in a diassociative A-loop,
$$(yx)C(z,y)=(yx)C(z^1,x).$$
(3.4)
###### Proof.
Let $`a=(yx)z^1`$. Then
$`(yx)C(z,y)`$ $`=(yx)C(z^1,y^1)`$ $`\mathrm{by}(\text{3.3})`$
$`=(yx)R(z^1)L(y^1)R(z)L(y)`$
$`=aL(y^1)R(z)L(y)`$
$`=(y^1a)R(a^1(yx))L(y)`$
$`=(yx)L(a^1)L(y^1a)L(y)`$
$`=(yx)L(y,a^1)`$ $`\mathrm{by}(\text{2.13})`$
$`=(yx)L(y^1,a)`$ $`\mathrm{by}(\text{3.2})`$
$`=(ya^1)(ax)`$
$`=(yx)R(x^1,a^1)`$
$`=(yx)L(x^1,a^1)`$ $`\mathrm{by}(\text{3.2})`$
$`=(yx)L(a^1)L(xa)L(x^1)`$ $`\mathrm{by}(\text{2.13})`$
$`=x^1(xaz)`$
$`=(yx)R(z^1)L(x)R(z)L(x^1)`$
$`=(yx)C(z^1,x).`$
###### Proof of Theorem 1.
For $`x,y,zL`$, we compute
$`x(y(xz))`$ $`=xR(z)L(y)L(x)`$
$`=xC(z,y)L(y)R(z)L(x)`$
$`=(yx)C(z,y)R(z)L(x)`$ $`\mathrm{by}(\text{2.11})`$
$`=(yx)C(z^1,x)R(z)L(x)`$ $`\mathrm{by}(\text{3.4})`$
$`=(yx)R(z)C(z^1,x)L(x)`$ $`\mathrm{by}(\text{2.11})`$
$`=(yx)L(x)R(z)`$
$`=(xyx)z.`$
###### Corollary 4.
$`C(x,z)=L(z,x)=R(z,x)`$, and $`C(x,z)^3=I`$.
###### Proof.
By the Moufang equation, $`R(xz)L(x)=R(x)L(x)R(z)`$. Hence, $`R(x^1)R(xz)R(z^1)=L(x)R(z)L(x^1)R(z^1)`$, so that (by (2.12) and (2.7)) $`R(z^1,x^1)=C(x^1,z^1)`$. Now use Corollary 3. ∎
## 4. Computer-aided Proofs
We comment further on our use of McCunes program OTTER . This is a general-purpose automated reasoning program which will prove theorems from axioms in first-order logic. In comparison with human reasoning, it is strongest in equational reasoning, and weakest in domains such as set theory, where there are many propositional connectives and alternations of quantifiers. Thus, most of the new mathematics to come out of automated reasoning has been in fields close to algebra. The book by Wos and Pieper describes general methods for applying automated reasoning to problems in mathematics and other areas. Many new theorems proved by OTTER occur in the book by McCune and Padmanabhan .
Many authors (as in ) simply use the OTTER output as the proof of a theorem. This is mathematically sound, since although OTTER’s search procedure is rather complex, the program can be made to output a simple proof object, which can be independently verified by a short `lisp` program. However, OTTER’s proofs are often long sequences of complicated equations which carry little intuitive content, and it is useful to re-express them in a form which a human reader can easily understand and verify.
Some discussion of the procedure for “humanizing” proofs occurs in . This was applied in the case of loop theory in , and in the present paper, where much of the argument is cast in the spirit of Bruck and Paige , emphasizing group-theoretic properties of the $`R(x)`$ and $`L(x)`$, rather than equations in the loop product and inverse. For example, in Corollary 4, the statement $`C(x,z)^3=I`$ conveys more information to most human readers than does the equivalent equation,
$$z^1(z((z^1(z((z^1(z(yx)x^1))x)x^1))x)x^1)=y,$$
which might (in its `ascii` form) be a typical line of OTTER output. However, some proofs seem to require direct computations in the loop itself. These proofs, although easy enough to verify by hand, may lack some motivation. The need for such computations probably explains why the results of this paper have not been found before.
###### Acknowledgements.
We wish to thank Tomaš Kepka for suggesting this problem to us.
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# The GHz-Peaked Spectrum radio galaxy 2021+614: Detection of slow motion in a compact symmetric object
## 1 Introduction
Despite many years of study of extragalactic radio sources, it is still unclear how they are formed and evolve. A crucial element in the study of their early evolutionary stages is to identify the young counterparts of old and extended FR I/FR II objects. Good candidates for young radio sources are those with peaked spectra, GHz-Peaked Spectrum (GPS) sources and Compact Steep-Spectrum (CSS) sources, because they are small in angular size as expected for young sources. GPS sources are characterized by a simple convex radio spectrum peaking at a frequency of about 1 GHz and are typically 100 pc in size. CSS sources have peaks in their spectra at lower frequencies and have projected linear sizes of $`<15`$ kpc.
The best direct evidence for very low kinematic ages has now been found for a few GPS radio galaxies. Measurements of hotspot advance speeds, from which an age estimate can be made, have been obtained for 0108+388 (Owsianik et al. Owsianik2 (1998)), 0710+439, 2352+495 (Owsianik & Conway Owsianik1 (1998); Owsianik et al. Owsianik3 (1999)) and 1943+456 (Polatidis Polatidis (1999)). The hotspot advance speeds are typically of order $`0.1h^1c`$, translating into ages of a few hundred to a few thousand years. All these sources belong to the morphological class of Compact Symmetric Objects (CSO), which are characterized by their small size ($`<500`$ pc) and symmetric radio morphology.
Having established that CSOs are young it is interesting to examine them at a number of points along their evolutionary track in the $`P`$$`D`$ (luminosity – linear size) diagram, and to carry out investigations at several radio wavelengths with a range of angular resolution and limiting flux density. We are investigating a sample of 11 bright GPS sources at 5 GHz with VSOP. These observations have been complemented by VLBI observations at 15 GHz to obtain matched-beam spectral index data. The 11 GPS sources in our sample are all those known in November 1995 with declination $`\delta >25\mathrm{°}`$, peak frequency $`\nu _{\mathrm{peak}}5`$ GHz, and peak flux-density $`S_{\mathrm{peak}}`$ $`{}_{}{}^{}{}_{}{}^{>}`$ 0.5 Jy/beam.
Here we report on one such object, the GPS radio galaxy 2021+614, also named OW 637. It is one of the strongest GPS sources with a total flux density of 2.5 Jy at 5 GHz and a radio luminosity of $`6\times 10^{37}h^2`$ W between 100 MHz and 100 GHz. The radio spectrum of 2021+614 has a broad, relatively flat peak centred at about 4 GHz, and falls off at lower and higher frequencies. The flattening of the spectrum at the highest radio frequencies ($`100`$ GHz) indicates the presence of a very compact component (Steppe et al. Steppe (1988)). The flux density above the spectral peak shows variability. Seielstad et al. (Seielstad (1983)) detect a 20% total flux density change on a time scale of 10 years. Aller et al. (Aller (1992)) observe some variability at 15 GHz, but not at 5 GHz, on a time scale of 5 years. High angular-resolution VLBI observations of 2021+614 at 2.3, 5 and 8.4 GHz have been published by Wittels et al. (Wittels (1982)), Bartel et al. (1984a ), Pearson & Readhead (Pearson (1988)) and Conway et al. (Conway (1994)). Based on the radio morphology and decomposition of the radio spectrum into contributes from individual components these authors all prefer a core-jet classification for 2021+614. In addition, Conway et al. (Conway (1994)) investigated structural changes in the source and identified the apparent centroid shift of the two main components as real motion. Cawthorne et al. (Cawthorne (1993)) determined that there is no significant linearly polarised emission from any of the components at 5 GHz, with upper limits of 5 mJy. The source was also observed by Kellermann et al. (Kellermann (1998)) as part of a 2-cm VLBI survey.
The optical counterpart of 2021+614 is an elliptical galaxy at redshift 0.2266. It is a highly reddened Narrow Line Radio Galaxy (NLRG) most probably with a considerable dust component within the optical object (Bartel et al. 1984b ). The shape of the \[OIII\] ($`\lambda _0=5007`$ Å) emission line profile is asymmetric and has a velocity dispersion of 780 km/sec. Deep CCD imaging by O’Dea et al. (O'Dea1 (1990)) shows that the galaxy has a prominent compact nucleus and two possible companions within $`12\mathrm{}`$.
In this paper we report new data from VSOP and global VLBI observations. Combining these data with older VLBI observations we determine the increase in separation between the two strongest components first detected by Conway et al. (Conway (1994)). In Sect. 2 we describe our observations and the data used to quantify this increase. In Sect. 3 we present the morphology and spectral-index distribution for 2021+614, and we calculate the separation rate of the components. We introduce a complementary approach for studies of changes in source structure, transfering the problem from the image coordinate plane into the spatial frequency plane. Further discussions of the results are in Sect. 4 where we propose that the morphological classification for the source is indeed CSO rather than core-jet and deduce its age. In Sect. 5 we summarize our main conclusions. In Appendix A we elaborate on the problems encountered in deducing the separation rate and its uncertainty and develop a method which allows to measure the relative increase in separation occuring between two epochs by means of the amplitude interference patterns.
## 2 Observations and Data Reduction
The VSOP satellite HALCA observed 2021+614 at 5 GHz on November 6, 1997 together with a 15-station ground-based array composed of all 10 VLBA and 5 of the 10 EVN radio telescopes (Effelsberg, Medicina, Noto, Onsala and Torun) plus the VLA in phased-array mode. The on-source time was 9 hours for the ground telescopes and 6 hours for the satellite. The tracking stations used to downlink the data and relay the local oscillator signal to the satellite were located at the Deep Space Network sites at Goldstone in California (USA) and Tidbinbilla in Australia. The VLBI observing run at 15 GHz took place on October 9, 1998 using the VLBA and the 100-m Effelsberg radio telescope for three 30-minute scans over a range of hour angles. Both data sets were correlated at the NRAO Array Operations Center in Socorro, NM, USA.
The 5 and 15 GHz VLBI images are shown in Fig. 1 (left and middle). The VSOP image was obtained following standard procedures for editing, *a-priori* amplitude calibration and fringe-fitting as recommended by the AIPS Cookbook (NRAO AIPS package) for space-VLBI data. Imaging was carried out with the Caltech Difference Mapping program (Difmap, Shepherd et al. Shepherd (1995)) applying uniform weighting to the data points. In order to produce a dynamic-range limited image the data had to be taken through several iterations of phase, and phase and amplitude self-calibration. Because of their low sensitivity, baselines between HALCA and the small ground telescopes needed extensive flagging during imaging. Relatively high-SNR fringes on space baselines can be only seen between HALCA and the 100-m Effelsberg radio telescope, and the phased VLA which has the equivalent sensitivity of a single 115-m antenna. The longest baseline measures 524 M$`\lambda `$ corresponding to a resolution of 0.3 mas (uniform weighting) and was achieved between the VLBA antenna located at Saint Croix in the Virgin Islands and HALCA.
## 3 Results
### 3.1 The sub-mas morphology
From Fig. 1 we see that 2021+614 has a simple symmetric structure at 5 and 15 GHz, dominated by two bright components. These two components are labeled B and D, following Bartel et al. (1984a ). A third component, labeled A, is visible only on the 5 GHz image, indicating that it has a steep spectral index ($`\alpha <2`$). In addition, a central component, E, is visible in both images and a jet-like feature connects this central component to components C and D. The jet-like feature appears bent in the 5 GHz image, but not in the 15 GHz image. Low surface brightness, extended structure can be seen east of component C and D. All components appear to be resolved to some degree, except the central feature. Low level side lobes near component D are due to the sparse sampling of the uv plane in region between Earth-Earth and Earth-HALCA baselines and limit the dynamic range of the image.
### 3.2 The 5 GHz – 15 GHz spectral-index distribution
The distribution of the spectral index $`\alpha `$ ($`S_\nu \nu ^\alpha `$) shown in Fig. 1 (right) is derived from the 5 and 15 GHz images. Both images were restored with a 0.6-mas circular beam and the spectral index calculated within regions whose intensity was higher than $`5`$ $`\times `$ rms noise in both images. Three regions with different spectral index characteristics can be seen: a steep spectrum ($`\alpha 0.7`$) north-eastern component, an inverted spectrum ($`\alpha 0.4`$) south-western complex and a central component with steep spectral index ($`\alpha 0.6`$). The inverted-spectrum border around the south-western complex is most probably artificial and due to the differences in beam position angle of the two combined images and possible noise effects in these low intensity areas.
### 3.3 Modelfitting
As mentioned above we observed 2021+614 at 5 GHz with VSOP and at 15 GHz with the VLBA and Effelsberg; additionally we had three 15-GHz uv-data sets from snapshot observations made by K.I. Kellermann et al. (Kellermann (1998)). We also make extensive use of Conway’s six-component model for 5 GHz data taken in November 1982 (priv. comm., Conway et al. Conway (1994)).
Before modelfitting can be performed, the weights of the single visibility measurements must be determined. This guarantees a meaningful significance for the reduced $`\chi ^2`$ as an indicator for the goodness of fit. For each integration time the weight is the reciprocal amplitude variance for the visibility measurement and is obtained from the internal scatter of the data within the averaging interval. The averaging time is limited by the coherence time, which for earth-based VLBI observations at 5 and 15 GHz is well above the 2-min averaging time adopted. In addition, to avoid complicated models with numerous components for the 5 GHz VSOP observation, the highest spatial frequencies were excluded.
We performed modelfitting using the Difmap `modelfit` program. This program fits the Fourier-transformed image-plane model to the real and imaginary parts of the complex visibilities (in contrast to other programs which fit the visibility amplitudes and closure phases) and tacitly assumes that the visibility phases are well calibrated. Thus, we modelfitted uv-data sets which have been self-calibrated beforehand.
The modelfitting provides a description of the most prominent characteristics of the brightness distribution – the number of components, their position, size, shape and intensity of source components – with as few parameters as possible.
For each uv-data set we followed the same modelfitting procedure. The starting model contained two circular gaussians each with a full-width-at-half-maximum (FWHM) of 0.7 mas located at the position of the highest pixel of the component. After the initial cycle of modelfitting additional components were added to improve the fit. Parameters of all components were allowed to vary during the process. This procedure was repeated until the reduced $`\chi ^2`$ did not decrease any further. In order to keep the number of parameters needed to fit the brightness distribution as small as possible we used elliptical gaussian components as well as components represented by delta functions.
For a model that is a good approximation to the data, the expected value of reduced $`\chi ^2`$ should be about 1. The reduced $`\chi ^2`$ characterizing our best fitting models (e.g. for the 5 GHz observation) was never less than 2.7. This indicates that the fit is poor. The reason for this is most likely that our simple models of a small number of components do not reproduce the extended, low surface brightness emission. However, this is not important in investigating the intensity and motion of the main components.
One of the problems in modelfitting is that more than one minumum in the $`\chi ^2`$ function may exist. We investigated the parameter space for pathological behaviour of the $`\chi ^2`$ function around the fitted global minimum, as a function of the most important parameters. Nearby local minima occur when the intensity or shape of one of the components degenerates. The associated source models can be rejected on the basis of their containing negative or unnaturally elongated components. The minima found for all uv data sets could be identified as being global.
For the 15 GHz observations the modelfitting procedure yielded models with five components, whereas for the 5 GHz observation a seven-component model best met the requirements. The best fitting model parameters for the 5 GHz data set are listed in Table LABEL:tab1. Fig. 2 (left) shows visibilities for Earth baselines sampled during the 1997.8 observation at 5 GHz and used for modelfitting. The contour image shown on the right of the uv-coverage represents the sum of the seven components in the model. The rightmost panel in Fig. 2 shows Conway’s six-component model for the 5 GHz observation from 1982.9.
### 3.4 Separation rate
We are interested in measuring the separation between the two brightest source components and changes of that separation in time. Fig. 3 incorporates all the separation measurements between components B and D, the two brightest components at 5 GHz, available to us as a function of observing year. The linear regression fit to the 5 GHz data points (triangles) shows that the component separation is changing at a rate of $`14.9\pm 0.2`$ $`\mu `$as/yr. In the introduction to Appendix A we discuss in more detail the estimation of the uncertainty for this measurement.
The four separation measurements at 15 GHz (squares) do not show the same progressive increase of separation seen in the 5 GHz data, whereas the 8.3 GHz data points (crosses) are consistent with separating components.
Striking evidence that an increase in separation has occurred in 2021+614 can be seen in Fig. 4 where we plot the visibility amplitudes versus projected uv-distance parallel to the axis defined by the components B and D. The double source structure along this line can be seen clearly, as can the effect of the extended nature of the individual components. Fig. 4 (top) shows the self-calibrated visibility amplitudes from our VSOP observation at epoch 1997.8 out to a projected baseline length of 180 M$`\lambda `$. Fig. 4 (middle) shows projected model visibility amplitudes for the data shown in the upper panel. Fig. 4 (bottom) represents Conway’s best model for the 1982.9 observation at uv-loci identical to those sampled during the 1997.8 observation. Comparing the fringe patterns from the upper and middle panel we recognize the missing flux density at short baselines resulting from unmodelled extended structure - an issue we discussed shortly in section 3.3.
Minima in the uv-plane occur at spatial frequencies of $`u_{\mathrm{min},n}=(2n1)/(2d)`$ for the $`n^{\mathrm{th}}`$ minimum measured along the position angle of the line connecting the two components, where $`d`$ is the distance between the two components in radians (Fomalont & Wright Fomalont (1974)). For the 1997.8 data this implies that the distance from the origin of the uv-plane to the first minimum is $`14.7`$ M$`\lambda `$, whereas for the 1982.9 observation the position of the first minimum occurs at $`15.2`$ M$`\lambda `$. The effect of the increase in separation is more easily recognizable for the higher order minima: it is qualitatively evident that an inward shift in the position of the minima has occurred, as expected for a separating source.
The fact that a small increase in separation between source components translates into an easily measurable change in the interference pattern led us to develop a method which compares two uv-data sets obtained at different epochs directly, parametrizing the time evolution, i.e., the structural change of the source. This method helps to overcome the problems connected with the estimation of errors for distance and separation-rate measurements. These problems are outlined in section *A.1*. The critical points of our method are the feasibility and realization of the direct comparison of the uv-data. In the sections following *A.1* we describe this procedure and apply it to our data. The method calculates the two-dimensional factorized increase in separation directly in the uv-plane and provides error estimates for those numbers.
We find that the visibility-amplitude interference pattern for the 1997.8 observation has to be stretched by $`4.43\pm 0.05`$% in the u-direction and by $`2.12\pm 0.10`$% in the v-direction in order to overlap with the 1982.9 observation.
The fact that the two multiplicative factors differ from each other tells us that the separation is not a simple linear increase along the line defined by the position of the two components at epoch 1982.9. The angular polar coordinate $`\psi _2`$ of component B with respect to D at epoch 1997.8 is related to that at epoch 1982.9, $`\psi _1`$, by $`\mathrm{tan}\psi _2=s_u/s_v\mathrm{tan}\psi _1`$; $`s_u`$ and $`s_v`$ are the stretching factors determined in the appendix. Note that the angular polar coordinate of component B changes from $`32.9\mathrm{°}`$ to $`33.5\mathrm{°}`$ between the two epochs (see Fig. 4).
### 3.5 Variability
Additional qualitative differences in source characteristics between the 1987.9 and 1997.8 observations can be directly established from Fig. 4 and quantified using the component models. Extrapolating the measured visibility amplitudes at low uv-spacing down to zero uv-spacings gives the total flux density of the source. There appears to be a decrease in intensity of about 4% over the 15 yr period. However, this may be due to amplitude calibration errors which can be as high as 5%.
On the other hand, it is immediately evident from Fig. 4 (middle and lower panel) that a change in relative intensity of the two brightest components has occurred. In Conway’s model for the 1982.9 observation the ordinate value of the minima are close to zero, indicating two major components with almost equal intensity, beating against each other. Fifteen years later the components are no longer equal in intensity and consequently the minima lie well above zero. The comparison of the source models provides an explanation in terms of changes in component intensity. The C-D complex increased 20% in intensity solely due to component C, whereas the B component decreased by an equivalent amount, simulating constant intensity, within amplitude calibration errors, for the source as a whole.
The most striking difference between the two models shown in Fig. 2 is the component south-east of the central component detected in the 1982.9 data, but not required in the model for the 1997.8 data. This component, labeled G in Fig. 2 (right), seems to have faded out to a level below the threshold adopted for our models. However, in the clean-component image from the 1997.8 data (Fig. 1, left) there is faint extended structure seen at the position corresponding to component G.
### 3.6 Self-similar evolution
It is of interest to check whether the ongoing source evolution – detected as an increase in component separation – follows the self-similar evolution model, which has been proposed for young radio sources, such as GPS and CSS sources (Snellen et al. Snellen1 (1997), Snellen2 (1999); Snellen & Schilizzi 2000a , 2000b ).
Self-similar evolution of a simple two-component compact symmetric source requires a proportional increase of the component sizes as the source components separate from each other. An increase in the size of the components in the image plane must be accompanied in the uv-plane by a proportional decrease of the FWHM of the upper envelope of the amplitudes which convolves the amplitude variation due to the beating of the two components. These decaying fringes can be seen in Fig. 4, where the upper envelope for the 1997.8 data is traced as a dotted line. However, self-similar growth is not observed when the uv-data from epochs 1982.9 and 1997.8 are compared – instead, we observe a decrease of the component size as the source expands, contrary to that expected in the self-similar growth model. This can be seen from Fig 4 (bottom) where the gaussian-like upper envelope from the 1997.8 fringes lies above the 1982.9 epoch maxima, indicating that source components have shrunk in size. The observed shrinkage of 20% is not due to high spatial frequency information from self-calibrating the complete 5-GHz VSOP uv data set before modelfitting. For the purpose of detecting changes in source structure we flagged the HALCA baselines before performing self-calibration.
## 4 Discussion
### 4.1 Morphological classification & age
Compact radio sources can roughly be classified into two morphological groups whose different appearances are believed to arise from orientation effects. For sources showing *symmetric double* structure, the radio axis, along which the individual components are aligned, lies near the plane of sky, whereas in sources with *core-jet* morphology the radio axis is pointing more towards the observer. In the latter case the observed structure is highly affected by projection and relativistic effects. In which category does 2021+614 belong?
The high resolution VSOP image (Fig. 1, left) reveals many details not seen in earlier observations. It shows that at the higher resolution provided by the space baselines, a compact component is visible between component B and D at the end of a low brightness jet in the direction of component D. Its central position relative to the two most prominent components, B (NE hotspot/lobe) and D (SW hotspot) and the low surface-brightness linear feature (jet) connecting it with component C (SW lobe) suggest its identification as the central engine of activity – the core.
The presence of component A and the extended emission seen at 5 GHz south-east of the central component are possibly not consistent with a classification as a CSO. While component A could be outward moving plasma emitted from the core at an earlier epoch than component B and D, it could be argued that D is the core component, since it is the most compact component and has the most inverted spectrum between 5 and 15 GHz and thus the highest turnover frequency. In addition component D is situated at the end of a linear arrangement of source components C, E, B and A which could be interpreted as regions of high emissivity (knots) along the path of an outward flowing jet. We note that Conway et al. (Conway (1994)) give component G (Fig. 2, right) and the extended emission detected east of component D as a possible counterjet identification. They argue that the misalignment could be explained by projection effects. On the 5 GHz VSOP map, however, this component is resolved out, which is not consistent with identifying it as a compact counterjet/hotspot.
Our rate of separation of components B and D of $`14.9\pm 0.2`$ $`\mu `$as/yr obtained from a linear regression fit, corresponds to an apparent speed of separation of $`0.14h^1c`$ in the source reference frame. This means that the two components were ejected $`380\pm 7`$ yr ago assuming constant velocity. With respect to the weak central component at the nucleus the speed of separation is $`0.07h^1c`$.
The “Interference Pattern Method” of parametrizing the structural change in the source (see Appendix A) shows that the percentage increase in the separation along the position angle defined by the u and v stretching factors $`s`$ is 3.0%. This increase over a timerange $`\mathrm{\Delta }t`$ of 14.9 years between epoch 1982.9 and 1997.8 implies that components B and D were ejected $`t=s^1\mathrm{\Delta }t=500\pm 10`$ yr ago, assuming constant velocity. The apparent inconsistency between this value and the $`380\pm 7`$ yr source age deduced above is caused by underestimation of the uncertainty by the linear regression fit process. A conservative estimate for the source age and its error given by the average of the two age values, is $`440\pm 80`$ yr. The corresponding separation rate and hotspot advance speed are $`=13\pm 3\mu `$as/yr and $`0.06\pm 0.01h^1c`$, respectively.
The subluminal character of the separation speed argues in favour of a CSO classification for 2021+614. And so do the low, total linear polarization of the source and the absence of compact components south-west of component D, the nucleus in the core-jet scenario. An undetectable counter-jet would require strong relativistic beaming effects, which would imply high apparent expansion speeds, on the order of speed of light or higher. Alternatively, the source 2021+614 might be a member of the blazar group seen face-on and observed at a extremely small angle $`\theta 1/\gamma `$ from the line of sight, where $`\gamma `$ is the Lorentz factor – resulting in low apparent pattern speeds. However, it seems very unlikely that 2021+614 is a blazar for a number of reasons, including, the large radio luminosity of the source together with its optical identification as an elliptical galaxy, the stellar component in the optical spectrum and the upper limit of 0.09 on the line flux ratio H$`\beta `$($`\lambda _04865`$)/\[OIII\]($`\lambda _05007`$) found by Bartel et al. (1984a ). All these points argue convincingly for a classification as a radio galaxy. In general, the optical properties – spectral and morphological – of 2021+614 are similar to those of other (radio-loud) NLRG or (radio-quiet) Seyfert 2 galaxies. These objects are seen edge-on, following the Unified Models for AGN.
Therefore, regarding the measured increase in separation as a real expansion, we are confident in assigning 2021+614 to the class of compact symmetric objects.
### 4.2 Self-similar evolution & hotspot advance speeds
Measurements of the hotspot advance speeds and infered kinematic ages for CSOs trace out an interesting evolution scenario for young radio sources. Hotspot advance speeds for CSOs span a range from $`0.06h^1c`$ to $`2.0h^1c`$ (Owsianik et al. Owsianik2 (1998), Owsianik & Conway Owsianik1 (1998), Owsianik et al. Owsianik3 (1999), Polatidis Polatidis (1999)). This indicates that radio galaxies spend a few thousand years in the GPS/CSO evolution stage. During this stage radio galaxies apparently do not grow following a *simple* self similar evolution scheme. This statement is based on the differences in structure between the 1982.9 and 1997.8 models detected in 2021+614. However, the observed decrease of the component sizes in this particular source does not reduce the importance of the self-similar evolution model for young radio sources. The timescales characterizing the GPS phenomenon as a whole are measured in thousands of years. Changing local environmental conditions on sub-parsec scales in the NLR (Narrow Line Region) medium caused by high density clouds could produce shock fronts which increase the compression of the ram-pressure confined and shock-ionized radio emitting plasma in the hotspots during short time scales of tens of years. On longer time scales, self-similar growth is recovered because it is controlled by the average external density of the NLR into which the radio galaxy expands and by the power with which the jet is driven forward.
Calculations carried out by Owsianik et al. (Owsianik et al. Owsianik1 (1998), Owsianik2 (1998) and Owsianik3 (1999)) for CSOs together with age estimates for those sources indicate similar environmental conditions for all of them. Moreover, all objects studied so far are members of the group of bright and therefore powerful GPS radio galaxies with an intrinsic radio power output of $`10^{26.026.6}`$ W/Hz at 5 GHz. Central engines powering GPS radio galaxies of similar radio luminosities create hotspots with similar internal pressure under similar environmental conditions (external density, density profile, magnetic field). In particular, the internal pressure and size of the working surface of the hotspot determine the kinetic power transported by the jet, which is responsible for generating the observed hotspot advance speeds. Therefore, not surprisingly, all hotspot advance speeds are of the same order.
## 5 Conclusions
We provide strong evidence that 2021+614 is one of a small group of compact symmetric sources for which speeds of separation have been measured. All have apparent ages of a few hundred to a few thousand years, indicating their youthfulness.
In the case of 2021+614 we measure a separation of $`16.1h^1`$ pc and a separation rate of $`0.12\pm 0.02h^1c`$ between the two dominant components. These components are associated with lobes and/or hotspots. The hotspot-advance speed is $`0.06\pm 0.01h^1c`$. From separation and separation rate measurements we deduce an apparent age of $`440\pm 80`$ yr. All results are measured in the source reference frame.
We do not observe self-similar growth in 2021+614 over a timerange of 15 yr but we argue that this does not rule out the self similar evolution scheme for young radio sources over longer timescales.
## Appendix A Errors in separation measurements
In order to generate an error estimate for the separation rate of 14.9 $`\mu `$as/yr deduced from 5 GHz data shown in Fig. 3 we need to consider the uncertainties for each individual data point. This is not a simple issue since we do not have the original data from which the two data points at 1982.9 and 1987.7 were deduced.
However, based on the assumption that a linear relation fits the separation measurements well and errors for individual data points are equal, a linear regression fit can provide error estimates for the separation measurement. Adopting this procedure we obtain an uncertainty of 2.1 $`\mu `$as for separation measurements and 0.2 $`\mu `$as/yr for the separation rate at 5 GHz. This method, however, excludes the possibility of an independent estimate for the goodness-of-fit.
Statistically correct treatments, such as elliptical gaussian fits in the image plane (e.g. AIPS task `JMFIT`) yield very small positional uncertainties for high dynamic range images (Fomalont Fomalont2 (1999)). The image in Fig. 2 (middle) allows the separation between component B and D to be determined with an accuracy an order of magnitude better as compared to the linear regression fit.
In order to circumvent this inconsistency between error estimations for the separation measurements and to obtain a second, independent source age estimate we developed a new method which we present below.
### A.1 The “Interference Pattern Method”
Our goal is to determine changes in source structure occured between 1982.9 and 1997.8 directly in the uv plane, using the amplitude interference patterns. These changes are parametrized by a two-dimensional stretching factor whose components along the u and v-axis are $`s_\mathrm{u}`$ and $`s_\mathrm{v}`$, and by an amplitude correction factor $`s_\mathrm{a}`$. Our method provides estimates for these quantities and for their uncertainties. In this way, interpreting the observed shift in the position of the visibilty amplitude minima between the two epochs – best seen in Fig. 4 – as real expansion, we are able to determine the source age, the separation rate and the hotspot advance speed.
### A.2 Parametrizing the time evolution/structural change
Conway’s 1982.9 data set was obtained during a 9 hour global VLBI observing run using 4 antennas in the USA and the 100-m dish in Effelsberg. Consequently, compared to our 1997.8 data set the uv coverage is sampled more sparsely and the two observations measure different spatial frequencies.
In order to obtain two data sets suitable for detecting an increase in separation between the two most prominent components we prepared the two uv data sets as follows: *1982.9* – a simple two-component model, comprising components B and D from Conway’s six-component model for the 1982.9 observations, was Fourier-transformed back into the uv-plane. The model visibility amplitudes were calculated on a $`561\times 561`$ pixel grid with mesh size of 0.5 M$`\lambda `$ ranging from $`140.0`$ to $`140.0`$ M$`\lambda `$ in both the u and v-dimension. This grid spacing was fine enough to detect a decrease in the position of the first minimum of a few percent; *1997.8* – all the model components except B and D were uv-subtracted from the 1997.8 observations. In order to take account of the flux density variability of the components, the relative intensities of B and D were adjusted to match Conway et al’s values. Then, the visibility amplitudes for the 1997.8 data, sampled along the uv-tracks shown in Fig. 2 (left), were gridded onto the $`561\times 561`$ pixel mesh.
The gridding of the visibility data was done by assigning the visibility value to the nearest grid point. Multiple assignments near the uv-plane centre were averaged together, but no special weighting was applied to those data points. Fig. 5 shows a grey-scale plot of visibilty amplitudes measured during the 1997.8 observation, prepared as explained above and gridded onto the $`561\times 561`$ pixel mesh.
After having prepared the data sets we assumed that all remaining differences between 1997.8 “data” and 1987.9 “model” are either due to real expansion or a residual multiplicative amplitude correction factor, $`s_\mathrm{a}`$. The parametrization of the change in separation was done by applying multiplicative factors, $`s_\mathrm{u}`$ and $`s_\mathrm{v}`$, to the sampled uv-points and associating the amplitude value of the original point with the new coordinates. Such an evolution model is, strictly speaking correct, only for unresolved source components which separate from or contract towards each other. For a pair of resolved components, modelled by elliptical gaussians, the resulting fringe pattern is convolved with an elliptical gaussian whose major and minor axis FWHM are inversely proportional to the corresponding parameters in the image plane. We investigated the impact of this simplification on the fitting process and concluded that it is minor and therefore negligible at the levels of accuracy involved.
### A.3 Fitting for the scale factors $`s_\mathrm{u}`$, $`s_\mathrm{v}`$ and $`s_\mathrm{a}`$
The $`\chi ^2`$ parameter measuring the goodness of fit was defined as the sum of the squared differences of the visibility amplitudes:
$$\chi ^2=\underset{i=1}{\overset{N}{}}\frac{[s_\mathrm{a}A_1(s_\mathrm{u}u_i,s_\mathrm{v}v_i)A_2(u_i,v_i)]^2}{\sigma _i^2},$$
where the indices 1 and 2 refer to the 1982.9 model and 1997.8 data, respectively; $`N`$ is the total number of gridded visibility measurements; $`\sigma _i`$ is the error associated with the visibilty amplitude data point. Values $`>1`$ for $`s_\mathrm{u}`$ and $`s_\mathrm{v}`$ indicate inward shifting minima as the source evolves. The result provided by the fitting routine gave a value of a few tens for the minimum reduced $`\chi ^2`$ rather than unity. This fit is too poor to use the $`(\chi ^2+1)`$-contour projections onto the parameter axes to determine the 1-$`\sigma `$ uncertainties for $`s_a`$, $`s_\mathrm{u}`$ and $`s_\mathrm{v}`$. Nevertheless, the procedure outlined above provided us the best-fit values for the overall amplitude correction factor, $`s_\mathrm{a}=0.88`$ and for the scale factors in u and v-direction, $`s_\mathrm{u}=1.0445`$ and $`s_\mathrm{v}=1.0202`$.
The percentage increase in separation $`s`$ along the position angle defined by the u and v stretching factors $`s_u`$ and $`s_v`$ is given by
$$s=[(s_\mathrm{u}1)^2\mathrm{sin}^2\psi _1+(s_\mathrm{v}1)^2\mathrm{cos}^2\psi _1]^{1/2}$$
and is 3.0%, where $`\psi _1`$ is the angular polar coordinate of component B with respect to D at epoch 1987.9.
The most difficult problem to tackle during the fitting process is to ensure that differences in component size and shape between different observations do not influence substantially the best-fit values of the parameters. The relatively high residual correction of 0.88 needed for the amplitude parameter, and the difference between $`s_\mathrm{u}`$ and $`s_\mathrm{v}`$, are due to different component sizes at the two epochs. A more sophisticated method would take these details into account and be more widely applicable. However, for our purpose, which is limited to the determination whether or not structural change in 2021+614 can be explained by subluminal motion of components, the systematic errors introduced are not relevant.
### A.4 Bootstrap method for estimating errors in the scale factors $`s_\mathrm{u}`$ and $`s_\mathrm{v}`$.
The bootstrap method is a very powerful error estimation technique applicable when there is not enough knowledge about the nature of the measurement errors to do a proper Monte Carlo simulation. We use this method to estimate the errors for the parameters $`s_\mathrm{u}`$ and $`s_\mathrm{v}`$.
For the 1997.8 observation we generated one hundred uv-data samples each containing a different 50% of the original data. The visibility measurements in each sample were selected randomly. We fitted the one hundred data sets to Conway’s model visibilites as outlined in the previous section for the whole data set, but this time fitting only the two scale factors $`s_\mathrm{u}`$ and $`s_\mathrm{v}`$, and setting $`s_\mathrm{a}`$ to its best-fit value of 0.88. The standard deviations of the best-fit parameters provide the errors associated with these quantities. We found $`s_\mathrm{u}=1.0443\pm 0.0005`$ and $`s_\mathrm{v}=1.0212\pm 0.0010`$. This means that using 50% of the available data we are able to determine a 4% decrease in separation of the minima along the u-axis and a 2% decrease along the v-axis at an accuracy level of 5% and 10%, respectively – *clear evidence for real motion*. The difference in accuracy is due to the higher percentage increase along the u-axis. The values found agree with the scale factors determined using all the data, reported at the end of the previous section. The errors improve by a factor of $`1/\sqrt{2}`$ if twice as many amplitude measurements are used.
In order to check the results we obtained for the scale-factor errors $`s_\mathrm{u}`$ and $`s_\mathrm{v}`$ we can carry out a simple calculation. 3% is the minimum change in scale factor which moves the uv-points by at least one pixel for uv-loci with $`\sqrt{u^2+v^2}>15`$ M$`\lambda `$. With about 5200 uv-points obeying this constraint, and taking into account that each measured and gridded point contributes significantly to the shift determination, the accuracy of the shift determination is improved by a factor of $`1/\sqrt{5200}`$, i.e. $`0.03/\sqrt{5200}0.0004`$, which in first approximation agrees with the errors reported above.
###### Acknowledgements.
This research was supported by the European Commission’s TMR Programme, “Access to Large-Scale Facilities”, under contract No. ERBFMGECT950012. We gratefully acknowledge the VSOP Project, led by the Japanese Institute of Space and Astronautical Science in cooperation with many organizations and radio telescopes around the world. The National Radio Astronomy Observatory is a facility of the National Science Foundation operated under cooperative agreement by Associated Universities, Inc. We thank J.E. Conway for providing the component model for the 1982 observation and the absolute separation measurements for both of his observations. Special thanks go to R.C. Vermeulen and K.I. Kellermann for donating self-calibrated uv-data sets from their 2 cm observations of 2021+614. We thank the referee Hugh D. Aller for careful reading and helpful suggestions.
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# References
SU(3) Clebsch-Gordan Coefficients<sup>1</sup><sup>1</sup>1 published in Romanian in St. Cerc. Fiz. 36, 3 (1984)
Marius Grigorescu
Department of Applied Mathematics
The University of Western Ontario
London, ON, Canada N6A 5B7
Abstract: The purpose of this paper is to find out a set of general recurrence formulas for the calculus of the $`SU(3)`$ Clebsch-Gordan coefficients. The first six sections are introductory, presenting the notations and general group theoretical methods applied to $`SU(3)`$. The following eight sections are devoted to a detailed treatment of the carrier spaces of the irreducible representations and their direct product.
I. Introduction
The first applications of the SU(3) group theory have occurred in nuclear physics, following the attempts to describe the nuclear collective properties in the frame of the shell model. The complex problem of the states classification and determination of the energy spectrum was simplified by using the dynamical symmetry of the Hamiltonian. This is reflected by certain regularities observed in the spectra, and indicates the existence of a set of operators which commute with the many-body Hamiltonian and generate a Lie algebra. Because the new ”constants of motion” are not related to the already present geometrical symmetries of the physical system, the spectrum shows an additional degeneracy. The subspace of the states with the same energy carries a representation of this extended algebra, integrable to the Lie group of dynamical symmetry. The group $`SU(3)`$ represents the dynamical symmetry of the isotropic 3-dimensional harmonic oscillator, whose degenerate energy eigenstates provide the basis of a space of irreducible representation (irrep). The irreducible spaces and representations will be denoted by $`V(P,Q)`$ and $`D(P,Q)`$, respectively, where $`P`$ and $`Q`$ are positive integers. The classification of the harmonic oscillator states may be obtained by using the complete set of commuting operators, whose eigenvalues are the indices of the basis vectors of the $`SU(3)`$ irreducible representations ch. 13. The classification of the states of the nucleon system in the harmonic oscillator potential was obtained by Elliott , and it was applied to the light nuclei. Bargmann and Moshinsky have considered in addition to the harmonic oscillator potential also the residual two-particle interaction \- . This Hamiltonian is $`SU(3)`$-invariant, and to obtain the eigenfunctions of a system consisting of 3 nucleons it is necessary to know the $`SU(3)`$ Clebsch-Gordan (CG) coefficients for the product $`D(P,Q)D(P_1,0)`$.
The coupling of the oscillator wave functions in the 2s-1d shell was obtained by Hecht by using the CG coefficients for the products
$$D(P,Q)D(P_1,0),D(P,Q)D(P_1,1),D(P,Q)D(0,Q_1).$$
(1)
For even-even nuclei was used with success the interacting boson model . In this model only the symmetry of the valence nucleons is considered, supposing that they are coupled in $`s`$ and $`d`$\- boson pairs. At a fixed number of bosons the most general Hamiltonian of this system has the $`SU(6)`$ symmetry, and the states are labeled by the eigenvalues of the Casimir operators for a subgroup chain of $`SU(6)`$. There are only three possible subgroup chains, one of them starting with $`SU(3)`$. The energy levels calculated by using this chain are in a good agreement with the experiment for the nuclei with half-filled valence shell. In this case the $`SU(3)`$ symmetry describes rotational spectra, and is independent of the degeneracy of the single-particle levels.
The importance of the simple Lie group theory for the elementary particle physics was exposed in detail in ref. . The discovery of the flavor quantum numbers such as the isospin, hypercharge, charm, conserved by the strong interactions has lead to the introduction of the dynamical symmetry models based on $`SU(2)`$, $`SU(3)`$, and $`SU(4)`$, respectively. These symmetries are only approximate, because the particles assigned to the multiplets have different masses. However, the classification is correct, because the symmetry breaking interactions do not change the values of the internal quantum numbers.
In general, a hadronic model with $`n_F`$ constituent quarks (flavors) is related to the $`SU(n_F)`$ symmetry. The $`n_F`$ quarks can be distinguished by the values of a set of $`n_F`$ flavor quantum numbers. If there are no charmed particles, the hadronic systems are correctly described by the octet model, based on the flavor symmetry group $`SU(3)^F`$. A basic reference to this model and its applications is . In this reference several tables of CG coefficients are included, and are used to calculate the matrix elements of the mass operator. This leads to a relationship between the strengths of the coupling between the pseudoscalar meson octet and the barionic currents in the case of the Yukawa coupling. The tables are also used to express the hadronic wave functions in terms of quarks.
The $`SU(3)`$ invariance of the strong interactions leads also to useful equations for the calculus of the multiplet dependent factor of the hadron scattering amplitude . In a many-quark system, these amplitudes can be obtained only by knowing the CG series and coefficients for the direct product of two arbitrary representations.
In the octet model the barions are $`L^\pi =0^+`$ states of a three quark system, contradicting the expected behavior of a system of three particles of spin $`1/2`$, obeying the Fermi-Dirac statistics. This puzzling situation was solved by introducing the color quantum number taking three possible values, in addition to the flavor and Dirac spinor indices of the relativistic quark wave function . The colored quarks are usual spin $`1/2`$ Fermions, but unobservable, because by postulate, only the ”white” systems can be free. The system of three colored quarks is antisymmetric to the permutation of the color indices, and is invariant to their transformation by $`SU(3)`$. However, contrary to $`SU(3)^F`$, the global color transformation group, $`SU(3)^C`$, is an exact symmetry. Assuming that the quark Lagrangean is invariant not only to global, but also to local $`SU(3)^C`$ transformations, it is possible to formulate a Yang-Mills theory of the strong interactions, the quantum chromodynamics . Although based on $`SU(3)`$, this theory does not rely on the $`SU(3)`$ irreducible representations, and the commutation relations between the gauge fields and currents are completely specified by $`su(3)`$, the Lie algebra of $`SU(3)`$.
In certain physical situations the $`SU(3)`$ group is too restrictive, and it is necessary to use larger semisimple Lie groups, as $`SU(n)`$. The general theory of semisimple Lie groups was elaborated in the classical works of Cartan and Weyl . The development of the quantum theory of angular momentum has determined the intensive study of $`SU(2)`$. The results obtained by Wigner and Racah in the theory of the $`SU(2)`$ representation and irreducible tensor operators, have solved practically all the problems of the atomic spectroscopy. The extension of these results to $`SU(n)`$ was given in the series of papers , by Baird abd Biedenharn. In ref. -II the method of boson generators is applied to recover the general formulas of Gel’fand and Zetlin for the matrix elements of the generators. A particular attention is given to $`SU(3)`$, because in its study are encountered features characteristic for $`SU(n)`$ in general, but which are not revealed by the simple structure of $`SU(2)`$. The last papers \- IV,V concern the calculus of the CG coefficients as matrix elements of the irreducible tensor operators in the Gel’fand basis, by using the Wigner-Eckart theorem. The main difficulty, represented by the occurrence of the multiplicities, is solved by the classification of the tensor operators using a conjugation operation.
The computational method used to obtain the $`SU(3)`$ and $`SU(4)`$ CG coefficients was presented in ref. , and , respectively. An alternative procedure, based on the Wigner-Eckart theorem, is presented in ref. . There nonorthonormal isoscalar factors are calculated by using the recurrence formulas satisfyed by the matrix elements of the irreducible tensor operators.
A general method for the calculus of the CG coefficients for every group $`SU(n)`$, based on the properties of the projection operators associated to the matrix elements of the irreducible representations is presented in ref. , ch. 7. The CG coefficients obtained this way are expressed as linear combinations of integrals over the group parameters .
The products of Eq. (1) have been reduced by Hecht, using the recurrence formulas between the CG coefficients. Isoscalar factors of particular interest, as well as formulas for the general ones, have been derived in , by using expansions of $`SU(3)`$ invariant polynomials.
The results presented in these works solve completely the problems related to the $`SU(3)`$ irreducible representations and the related CG coefficients. However, I have considered that a self-contained presentation of some of these results, as well as of a simple and general method to calculate the isoscalar factors, might be useful. The sections II-VI are devoted to the basic definitions. The notations used for the $`SU(3)`$ algebra, its fundamental and general tensor representations, as well as the CG series, are introduced. The sections VII-IX contain a detailed study of the irrep spaces and their direct product, while in sect. X-XIV the formulas obtained before are applied to the explicit calculation of the CG coefficients in particular situations of interest.
II. The SU(3) group
The group $`SU(3)`$ consists of the linear unimodular transformations of the complex space $`C^3`$ which leave invariant the Hermitian bilinear form
$$<a,b>=a_1^{}b_1+a_2^{}b_2+a_3^{}b_3,a,bC^3.$$
(2)
The matrices of the unitary linear operators associated to these transformations in the basis
$$\{x_jC^3,j=1,2,3;<x_j,x_k>=\delta _{jk}\}$$
(3)
are unitary unimodular matrices $`U`$ which give a 3-dimensional representation of the group $`SU(3)`$. An analytic structure on this group can be introduced by considering the real and the imaginary part of the 9 complex elements of the matrix $`U`$ as analytical coordinates of a point in the space $`R^{18}`$. These 18 real parameters are constrained by 10 algebraic relations determined by the conditions
$$detU=1,UU^{}=I,$$
(4)
such that only eight of them are independent. The remaining parameters can be expressed as analytical functions of them, and therefore, $`SU(3)`$ is a real, eight-dimensional, analytical manifold.
If $`U_1,U_2SU(3)`$, then the coordinates of $`U=U_1U_2^1`$ are analytical functions of the coordinates of $`U_1`$ and $`U_2`$. Thus, $`SU(3)`$ is a real Lie group.
By using the continuity of the matrix elements and Eq. (4) it can be shown that $`SU(3)`$ is homeomorphic with a compact set in $`R^8`$ -IV. Moreover, it can be proved by recurrence that $`SU(n)`$ is simply connected if $`SU(n1)`$ is simply connected -(VIII sect. 4). As $`SU(1)=\{I\}`$ is simply connected, then every $`SU(n)`$, and in particular $`SU(3)`$, is simply connected.
The coordinates on $`SU(3)`$, $`\{q_j,j=1,8\}`$, are defined such that the origin of the space $`R^8`$ corresponds to the unit matrix $`I`$. Their choice is simplified by the existence of the subgroups $`SU(2)`$ and $`SUD(3)=D(3)SU(3)`$ ( consisting of the diagonal $`3\times 3`$ unitary unimodular matrices), for which the coordinates are known. The four remaining parameters are chosen usually in the form given by Nelson . Other parametrizations are given in -.
Every $`SU(3)`$ matrix can be obtained from the identity $`I`$ by a series of infinitesimal transformations corresponding to a continuous variation of the group parameters, and can be represented in the form
$$U(q_1,\mathrm{},q_8)=e^{\frac{i}{2}_{k=1}^8q_k\lambda _k}.$$
(5)
The parameters $`q_k`$, $`k=1,8`$ of this representation define a canonical system of coordinates -(IX, sect. 3), and the matrices of the generators of the infinitesimal transformations
$$F_k=\frac{1}{2}\lambda _k,k=1,8$$
(6)
give a representation in $`C^3`$ of the real Lie algebra $`su(3)`$. These matrices are the elements of a basis in the real linear space of the trace 0, $`3\times 3`$ Hermitian matrices. Because $`su(2)su(3)`$, it is convenient to chose these matrices in the form given by Gell-Mann
$$\lambda _1=\left[\begin{array}{ccc}0& 1& 0\\ 1& 0& 0\\ 0& 0& 0\end{array}\right],\lambda _2=\left[\begin{array}{ccc}0& i& 0\\ i& 0& 0\\ 0& 0& 0\end{array}\right],\lambda _3=\left[\begin{array}{ccc}1& 0& 0\\ 0& 1& 0\\ 0& 0& 0\end{array}\right],$$
(7)
$$\lambda _4=\left[\begin{array}{ccc}0& 0& 1\\ 0& 0& 0\\ 1& 0& 0\end{array}\right],\lambda _5=\left[\begin{array}{ccc}0& 0& i\\ 0& 0& 0\\ i& 0& 0\end{array}\right],\lambda _6=\left[\begin{array}{ccc}0& 0& 0\\ 0& 0& 1\\ 0& 1& 0\end{array}\right],$$
(8)
$$\lambda _7=\left[\begin{array}{ccc}0& 0& 0\\ 0& 0& i\\ 0& i& 0\end{array}\right],\lambda _8=\frac{1}{\sqrt{3}}\left[\begin{array}{ccc}1& 0& 0\\ 0& 1& 0\\ 0& 0& 2\end{array}\right],$$
(9)
As there is an $`SU(3)`$ transformation which brings every $`SU(3)`$ matrix to the form $`SUD(3)`$, the elements of $`SU(3)`$ belong to equivalence classes, such that every class contains only one element of $`SUD(3)`$. Consequently, the characters for the elements of $`SU(3)`$ will be analytical functions of the two variables which parameterize $`SUD(3)`$. An invariant measure, normalized in the character space, can be found in .
The matrices $`\lambda `$ satisfy the relationship
$$\lambda _i\lambda _j=\frac{2}{3}\delta _{ij}I+\underset{k=1}{\overset{8}{}}(d_{ijk}+if_{ijk})\lambda _k.$$
(10)
Here $`d_{ijk}`$ is a real, symmetric, tensor of order three, $`f_{ijk}`$ is also real, but antisymmetric, and $`\delta _{ij}`$ is the Kronecker symbol. The non-vanishing components of these tensors are
$$f_{123}=1,f_{458}=f_{678}=\frac{\sqrt{3}}{2}$$
(11)
$$f_{147}=f_{246}=f_{345}=f_{257}=f_{156}=f_{367}=\frac{1}{2},$$
and
$$d_{118}=d_{228}=d_{338}=d_{888}=\frac{1}{\sqrt{3}},$$
(12)
$$d_{448}=d_{558}=d_{668}=d_{778}=\frac{1}{2\sqrt{3}},$$
$$d_{146}=d_{157}=d_{256}=d_{344}=d_{355}=d_{247}=d_{366}=d_{377}=\frac{1}{2}.$$
By using Eq. (10) it is possible to derive the commutation and anti-commutation relationships
$$[F_i,F_j]=i\underset{k=1}{\overset{8}{}}f_{ijk}F_k,\{F_i,F_j\}=\frac{1}{3}\delta _{ij}I+\underset{k=1}{\overset{8}{}}d_{ijk}F_k.$$
(13)
The rank of the $`su(3)`$ algebra, defined as the dimension of the maximal nilpotent Lie subalgebra, it is given by the dimension of the $`sud(3)`$ subalgebra, which equals 2. The two commuting elements of $`su(3)`$ are $`F_3`$ and $`F_8`$, the basis in Eq. (3) being chosen such that their matrices are diagonal. Their diagonal elements are related to the quantum numbers characteristic to the systems with $`SU(3)`$ symmetry, and can be used to label the states. For the hadron classification the states are labeled by the eigenvalues $`i_3`$ and $`y`$ of the operators $`I_3`$ and $`Y`$,
$$I_3=F_3,Y=\frac{2}{\sqrt{3}}F_8.$$
(14)
The structure constants $`c_{jk}^i`$, $`i,j,k=1,8`$, defined by the commutation relations
$$[iF_j,iF_k]=\underset{j=1}{\overset{8}{}}c_{jk}^l(iF_l)$$
(15)
are $`c_{jk}^l=f_{jkl}`$. The Cartan metric tensor is defined by
$$g_{ij}=\underset{k,l=1}{\overset{8}{}}c_{il}^kc_{jk}^l,$$
(16)
and has the simple expression $`g_{ij}=3\delta _{ij}`$. Therefore the Killing form
$$(X,Y)=\underset{i,j=1}{\overset{8}{}}g_{ij}X^iY^j$$
(17)
is non-degenerate, negative definite, and the algebra $`su(3)`$ is semisimple. This means that $`su(3)`$ contains no commutative ideal, and implies that $`SU(3)`$ is also a semisimple Lie group , Ch. XI, sect. 4.
The properties summarized above show that $`SU(3)`$ is a real, compact, simply connected, semisimple Lie group. Such a group has no proper connected invariant subgroup, but it may have a discrete one. For $`SU(3)`$ this subgroup is
$$Z_3=\{I,e^{\frac{2}{3}i\pi }I,e^{\frac{4}{3}i\pi }I\}.$$
(18)
The factor group $`SU(3)/Z_3`$ is obtained by identifying the elements of $`SU(3)`$ which are different by a factor $`e^{\frac{2}{3}i\pi }`$ or $`e^{\frac{4}{3}i\pi }`$, and is triply connected. This factor group is important in the elementary particle physics, because only the hadron states assigned to its irreps have integer charge and hypercharge.
III. The su(3) algebra
The linear representations theory of the semisimple Lie algebras $`L`$, and the corresponding groups, rely on the decomposition of $`L`$ in a direct sum of invariant subspaces with respect to the restriction of the adjoint representation
$$Xad_X=[X,],XL$$
(19)
to the nilpotent subalgebra
$$H=\{X,YL/(ad_X)^kY=0,0<kZ\}.$$
(20)
An important example is provided by the Gauss decomposition, $`L=N^++H+N^{}`$, where $`N^\pm `$ are nilpotent subalgebras, while $`N^\pm +H`$ are solvable , ch. 1, sect. 6.
The $`su(3)`$ algebra and the $`SU(3)`$ group have no Gauss decomposition because are semisimple compact ( , ch. 3, sect. 6), but they represent compact real forms of the algebra $`sl(3,C)`$, respectively the group $`SL(3,C)`$, which are semisimple complex, and have a Gauss decomposition. The representations of $`su(3)`$ and $`SU(3)`$ which appear in applications are completely determined by the representations of $`sl(3,C)`$, respectively $`SL(3,C)`$ by the unitary Weyl trick , XI, XII.
The Lie algebra $`sl(3,C)`$, ($`A_2`$), is the subalgebra of $`gl(3,C)`$ generated by the traceless $`3\times 3`$ matrices. A basis of $`sl(3,C)`$ is represented by the matrices $`F_k`$, $`k=1,8`$ introduced in the previous section, but this choice is not appropriate for the Gauss decomposition. Instead, it is convenient to express the basis elements as linear combinations of the basis of the $`gl(3,C)`$ algebra. The Weyl basis of $`gl(3,C)`$ is represented by nine $`3\times 3`$ matrices $`\{e_{ik},i,k=1,3\}`$ which have a single non-vanishing element, equal to 1,
$$(e_{ik})_{\alpha \beta }=\delta _{i\alpha }\delta _{k\beta }.$$
(21)
Their commutator
$$[e_{ik},e_{jl}]=\delta _{kj}e_{il}\delta _{il}e_{jk},$$
(22)
is a matrix of $`sl(3,C)`$, and therefore the same commutation relations will be satisfied by the nine traceless matrices $`A_k^i=e_{ik}\delta _{ik}I/3`$,
$$[A_k^i,A_l^j]=\delta _{kj}A_l^i\delta _{il}A_k^j.$$
(23)
For applications in nuclear physics it is useful to remark that operators $`e_{ik}`$ which satisfy the commutation relations (22) can be constructed by using boson operators. Thus, if $`a_i^{}`$ and $`a_i`$, $`i=1,2,3`$, denote the creation and annihilation boson operators, ($`[a_i,a_k^{}]=\delta _{ik}`$, $`[a_i,a_k]=[a_i^{},a_k^{}]=0`$), which appear in the Hamiltonian of the isotropic harmonic oscillator,
$$h_0=\mathrm{}\omega \underset{i=1}{\overset{3}{}}(a_i^{}a_i+\frac{1}{2}),$$
(24)
then $`e_{ik}=a_i^{}a_k`$ satisfy Eq. (22), and $`[e_{ik},h_0]=0`$.
The relationship between the two basis sets for $`sl(3,C)`$, $`\{A_k^i;i,k=1,3`$ and $`\{F_k,k=1,8\}`$ takes a simple form in terms of the complex combinations
$$I_\pm =F_1\pm iF_2,K_\pm =F_4\pm iF_5,L_\pm =F_6\pm iF_7.$$
(25)
By using these new operators, the relationship between the two basis sets is given by the equations
$$A_1^1=\frac{1}{\sqrt{3}}F_8+F_3,A_2^2=\frac{1}{\sqrt{3}}F_8F_3,A_3^3=\frac{2}{\sqrt{3}}F_8,$$
(26)
$$A_2^1=I_+,A_3^1=K_+,A_3^2=L_+,$$
$$A_1^2=I_{},A_1^3=K_{},A_2^3=L_{}.$$
(27)
The algebra $`sl(3,C)`$ decomposes in the Cartan subalgebra $`H`$, which is generated by the elements $`F_3`$ and $`F_8`$, and the subalgebras $`N^\pm `$, generated by $`I_\pm `$, $`K_\pm `$ and $`L_\pm `$. By using the expansion
$$F_\rho =\underset{k=1}{\overset{3}{}}\mathrm{\Phi }_k(\rho )e_{kk}$$
(28)
where $`\rho `$ is 3 or 8, the $`sl(3,C)`$ commutation relations in the Cartan-Weyl basis are expressed by
$$[F_\rho ,A_k^j]=\alpha _{jk}(\rho )A_k^j$$
(29)
$$[A_k^i,A_l^j]=\delta _{kj}A_l^i\delta _{il}A_k^j,ik,jl.$$
(30)
The coefficients $`\alpha _{jk}`$ are linear functions on $`H`$, defined by
$$\alpha _{jk}(\rho )=\mathrm{\Phi }_j(\rho )\mathrm{\Phi }_k(\rho ),$$
(31)
and the set
$$\mathrm{\Delta }=\{\alpha _{jk};j,k=1,2,3\}$$
(32)
represents the root system.
The linear combinations of the roots generate the dual space of $`H`$, denoted $`\stackrel{~}{H}`$. If $`\alpha _{12}`$, $`\alpha _{13}`$ and $`\alpha _{23}`$ are chosen positive, then $`\alpha _{12}`$ and $`\alpha _{23}`$ are the simple roots. The new structure constants
$$C_{\rho (jl)}^{(ik)}=\alpha _{jl}(\rho )\delta _{ij}\delta _{kl}$$
(33)
give the restriction of the metric tensor to the Cartan subalgebra
$$g_{\rho \sigma }=(F_\rho ,F_\sigma )=\underset{j,l=1}{\overset{3}{}}\alpha _{jl}(\rho )\alpha _{jl}(\sigma )=3\delta _{\rho \sigma },$$
(34)
with $`\rho ,\sigma =`$ 3 or 8. Thus, the restriction of the Killing form to $`H`$ is non-degenerate, and for every $`\alpha \stackrel{~}{H}`$ exists an unique element $`h_\alpha H`$, denoted in the following also by $`\alpha `$, such that $`hH`$
$$\alpha (h)=(h_\alpha ,h),(\alpha ,\beta )(h_\alpha ,h_\beta ),\alpha ,\beta \stackrel{~}{H}.$$
(35)
Reciprocally, this relationship associates linear functions to the elements of $`H`$. The linear functions $`\alpha _3`$ and $`\alpha _8`$ associated by Eq. (29) to $`F_3`$ and $`F_8`$ are identically zero. However, Eq. (34) and (35) lead to a covariant orthogonal basis $`\widehat{g}_\sigma `$, $`\sigma =`$ 3 or 8, in $`\stackrel{~}{H}`$, defined by $`\widehat{g}_\sigma (\rho )=g_{\sigma \rho }`$. The covariant coordinates of the roots from $`\mathrm{\Delta }`$ in this basis are given by Eq. (35).
Another covariant basis in $`\stackrel{~}{H}`$ is represented by the simple roots. These have the length $`1/\sqrt{3}`$ with respect to the Killing form, and span an angle of 120 degrees. The matrices associated by Eq. (35) to these roots are
$$\alpha _{12}=\underset{\rho ,\sigma }{}g^{\rho \sigma }\alpha _{12}(\rho )F_\sigma =\frac{1}{3}F_3=\frac{1}{6}\left[\begin{array}{ccc}1& 0& 0\\ 0& 1& 0\\ 0& 0& 0\end{array}\right],$$
(36)
$$\alpha _{23}=\underset{\rho ,\sigma }{}g^{\rho \sigma }\alpha _{23}(\rho )F_\sigma =\frac{1}{6}(\sqrt{3}F_8F_3)=\frac{1}{6}\left[\begin{array}{ccc}0& 0& 0\\ 0& 1& 0\\ 0& 0& 1\end{array}\right].$$
(37)
The contravariant vectors $`\alpha ^\mu `$, $`\mu =`$ 12 or 23 associated to the simple roots are defined by the relationship
$$\frac{(\alpha ^\mu ,\alpha _\nu )}{(\alpha _\nu ,\alpha _\nu )}=\frac{1}{2}\delta _{\mu \nu },$$
(38)
and the coordinates of an arbitrary element $`M\stackrel{~}{H}`$ in the contravariant basis are
$$m_\mu =2\frac{(M,\alpha _\mu )}{(\alpha _\mu ,\alpha _\mu )}.$$
(39)
IV. Fundamental representations
The irreducible representations of a Lie group $`G`$ which admits a Gauss decomposition of the form $`G=Z^{}DZ^+`$, with $`Z^\pm `$, $`D`$, generated by $`N^\pm `$, $`H`$, respectively, are induced by the one-dimensional representations (characters) $`\tau (\delta )`$, $`\delta D`$ of the subgroup $`D`$. Let $`R_g`$, $`gG`$, a representation of $`G`$ on the finite dimensional linear space $`V`$. The elements of $`V`$ which are eigenstates of the operators $`R_\delta `$, $`\delta D`$ are called weight vectors. In particular, the weight vectors $`|M>`$,
$$R_\delta |M>=\tau _M(\delta )|M>,\delta D$$
(40)
which remain invariant to the action of the subgroup $`Z^+`$,
$$R_z|M>=|M>,zZ^+,$$
(41)
are called highest weight vectors.
In the linear irreps theory of the group $`GL(n,C)`$ are proved the following important theorems:
Th. I. Every carrier space $`V`$ of a finite dimensional irrep $`R`$ of $`GL(n,C)`$ contains an unique highest weight vector, cyclical for $`V`$.
Th. II. The irrep induced by the character $`\tau _M`$ occurrs in the decomposition of an reducible representation on the space $`V`$ with a multiplicity equal to the number of the highest weight vectors $`|M>`$ contained by $`V`$ ch. VI, sect. 3.1.
Th. III. The analytical complex inductive characters for the representations of the group $`GL(3,C)`$ have the expression
$$\tau (\delta )=\gamma _1^{m_1}\gamma _2^{m_2}\gamma _3^{m_3},\delta =\left[\begin{array}{ccc}\gamma _1& 0& 0\\ 0& \gamma _2& 0\\ 0& 0& \gamma _3\end{array}\right]D,$$
(42)
with $`m_1m_2m_3`$, integers ch. 8 sect. 3.
The finite dimensional irreps of the group $`SU(3)`$ can be realized on the irrep spaces of its complex extension, $`SL(3,C)`$. The restriction of the characters $`\tau `$ of $`DGL(3,C)`$ to $`SUD(3)D`$ defines a character $`\tau ^0`$ which specifies completely the irreps of $`SU(3)`$. This restriction is obtained by considering only matrices $`\delta `$ with elements $`\gamma _i`$, $`i=1,2,3`$ such that
$$|\gamma _i|=1,\gamma _1\gamma _2\gamma _3=1.$$
(43)
If the simple roots are $`\alpha _{12}`$ and $`\alpha _{23}`$, then every matrix $`dSUD(3)`$ can be written as
$$d=e^{ih},h=t^{12}\alpha _{12}+t^{23}\alpha _{23},$$
(44)
where the real parameters $`t^{12}`$ and $`t^{23}`$ are
$$t^{12}=2\frac{(h,\alpha ^{12})}{(\alpha _{12},\alpha _{12})},t^{23}=2\frac{(h,\alpha ^{23})}{(\alpha _{23},\alpha _{23})}.$$
(45)
Replacing $`\alpha _{12}`$ and $`\alpha _{23}`$ from Eq. (44) by the expressions of Eq. (36), (37), the matrix elements of $`d`$ take the form
$$\gamma _1=e^{\frac{i}{6}t^{12}},\gamma _2=e^{\frac{i}{6}(t^{12}t^{23})},\gamma _3=e^{\frac{i}{6}t^{23}}.$$
(46)
These elements satisfy the conditions of Eq. (43), and therefore, by Th. III, the inductive character can be written in the form
$$\tau _{\underset{¯}{M}}^0(d)=e^{i(\underset{¯}{M},h)},$$
(47)
with $`\underset{¯}{M}\stackrel{~}{H}`$ defined by $`\underset{¯}{M}=(m_1m_2)\alpha ^{12}+(m_2m_3)\alpha ^{23}`$. The element $`\underset{¯}{M}`$ of $`\stackrel{~}{H}`$ is called highest weight. Therefore, every irrep of $`SU(3)`$ is completely specified by two non-negative integers, $`P=m_1m_2`$ and $`Q=m_2m_3`$, and is denoted $`D(P,Q)`$. The carrier space of this irrep will be denoted in the following by $`V(P,Q)`$. By using the same notation for the elements of the $`su(3)`$ algebra and the corresponding representation operators in $`V(P,Q)`$, the expansion of Eq. (40) and (41) for $`SU(3)`$ near the identity leads to the equations
$$\alpha _{12}|\underset{¯}{M}>=(\alpha _{12},\alpha ^{12})P|\underset{¯}{M}>,\alpha _{23}|\underset{¯}{M}>=(\alpha _{23},\alpha ^{23})Q|\underset{¯}{M}>,$$
(48)
and, respectively
$$I_+|\underset{¯}{M}>=0,L_+|\underset{¯}{M}>=0.$$
(49)
Thus, according to Eq. (36), (37), $`|\underset{¯}{M}>`$ is an eigenvector of $`F_3`$ and $`F_8`$, and by Eq. (14), of $`I_3`$ and $`Y`$, with the eigenvalues
$$(i_3)_{\underset{¯}{M}}=\frac{P}{2},(y)_{\underset{¯}{M}}=\frac{P+2Q}{3}.$$
(50)
Sometimes it is convenient to choose $`\alpha _{13}`$ and $`\alpha _{32}`$ as simple roots. In this case, the highest weight vector, denoted by $`|M>`$, will satisfy the equations
$$K_+|M>=0,L_{}|M>=0,$$
(51)
and the highest weight $`M=P\alpha ^{13}+Q\alpha ^{32}`$ is the reflection of the weight $`\underset{¯}{M}`$ with respect to the contravariant vector $`\alpha ^{13}=\alpha ^{12}`$. The quantum numbers $`i_3`$ and $`y`$ labeling $`|M>`$ are
$$(i_3)_M=\frac{P+Q}{2},(y)_M=\frac{PQ}{3}.$$
(52)
A basis in $`V(P,Q)`$ is represented by the eigenvectors $`|m>`$ of the operators $`I_3`$ and $`Y`$. The couple $`(i_3,y)`$ of the eigenvalues labels the vectors $`|m>`$, and will be denoted in the first part of this work by $`m`$. These eigenvalues are related to the components $`m_{13}`$ and $`m_{32}`$ of the weight $`m`$ in the contravariant basis by
$$i_3=\frac{m_{13}+m_{32}}{2},y=\frac{m_{13}m_{32}}{3}.$$
(53)
According to Th. I, all basis vectors $`|m>|i_3,y>`$ can be obtained by the application of the operators $`L_+`$, $`I_{}`$ and $`K_{}`$ on the highest weight vector $`|(P+Q)/2,(PQ)/3>`$.
The numbers $`P`$ and $`Q`$ are related to the eigenvalues $`f`$ and $`g`$ of the two Casimir operators $`F`$ and $`G`$ ch. 9 sect. 4 defined by
$$F=\underset{k=1}{\overset{8}{}}F_k^2=\frac{1}{2}\underset{i,k=1}{\overset{3}{}}A_k^iA_i^k,G=\frac{1}{2}\underset{i,k,l=1}{\overset{3}{}}(A_l^iA_i^kA_k^l+A_i^lA_k^iA_l^k),$$
(54)
such that
$$f=\frac{P^2+PQ+Q^2}{3}+P+Q$$
(55)
$$g=\frac{1}{9}(PQ)(2P+Q+3)(P+2Q+3).$$
The representations $`D(1,0)`$ and $`D(0,1)`$ are 3-dimensional, non-equivalent, and are called fundamental representations. These have the highest weights $`\alpha ^{13}`$ and $`\alpha ^{32}`$, respectively, and every representation can be constructed by the decomposition of their multiple direct product.
By applying the commutator $`[F_\rho ,A_k^j]`$, $`\rho =3,8`$, of Eq. (29) to the weight vector $`|m>`$, and then using the eigenvalues equations
$$F_\rho |m>=m_\rho |m>,m_3=i_3,m_8=\frac{\sqrt{3}}{2}y,$$
(56)
one can find the relation
$$F_\rho A_k^j|m>=(m_\rho +\alpha _{jk}(\rho ))A_k^j|m>,$$
(57)
which gives the weight of the vector $`A_k^j|m>`$.
The representation matrices of the $`su(3)`$ generators in the space $`V(1,0)`$ have the form of Eqs. (7),(8),(9), and the basis consists of the highest weight vector $`x_1=|\frac{1}{2},\frac{1}{3}>`$ and
$$x_2=I_{}x_1=|\frac{1}{2},\frac{1}{3}>,x_3=K_{}x_1=|0,\frac{2}{3}>.$$
(58)
The weights $`i_3,y`$ of the vectors $`x_1`$, $`x_2`$, $`x_3`$, can be represented on the root diagram by the vertices 1,2,3, respectively, of the triangle shown in Fig. 1.
The representation matrices of the $`SU(3)`$ group on $`V(1,0)`$ are just the elements of the group of Eq. (5). These act by transforming the vectors $`x_j`$, $`j=1,2,3`$ in $`\underset{¯}{x}_j`$,
$$\underset{¯}{x}_j=\underset{k=1}{\overset{3}{}}U_j^kx_k.$$
(59)
A new three-dimensional irreducible representation is obtained by the complex conjugation of Eq. (59),
$$\underset{¯}{y}^j=\underset{k=1}{\overset{3}{}}(U_j^k)^{}y^k.$$
(60)
Here $`y^j=x_j^{}`$ are contravariant basis vectors which generate the representation space $`V(1,0)^{}`$. The matrices of the $`SU(3)`$ elements in the $`D(1,0)`$ representation are
$$U|_{(1,0)^{}}=(U|_{(1,0)})^{}=e^{\frac{i}{2}_{k=1}^8\lambda _k^{}a_k},$$
(61)
which shows that the matrices of the $`su(3)`$ generators are
$$F_i|_{(1,0)^{}}=(F_i|_{(1,0)})^{}.$$
(62)
The matrices of the operators $`F_\rho |_{(1,0)^{}}`$ and $`A_k^i|_{(1,0)^{}}`$ can be found by using Eq. (25), such that
$$F_\rho |_{(1,0)^{}}=(F_\rho |_{(1,0)})^{},A_k^i|_{(1,0)^{}}=(A_i^k|_{(1,0)})^{}.$$
(63)
Here the matrices $`F|_{(1,0)}`$ and $`A_k^i|_{(1,0)}`$ are real, and therefore
$$F_\rho |_{(1,0)^{}}=F_\rho |_{(1,0)},A_k^i|_{(1,0)^{}}=A_i^k|_{(1,0)}.$$
(64)
The weight diagram of the vectors $`y^k`$ is obtained by reflection with respect to the origin of the weight diagram for $`x_k`$, $`k=1,2,3`$, and coincides with that of the basis in $`V(0,1)`$. If the basis vectors are labeled by their weights, such that $`x_m`$ denotes the basis vectors of Eq. (58), and $`y^m`$ the vectors $`y^k`$, then
$$y^m=(x_m)^{}.$$
(65)
By convention, the matrix elements of the operators $`I_\pm `$ and $`K_\pm `$ in the canonical basis of $`V(P,Q)`$ are chosen to be positive . This convention is violated by the basis $`\{y^i,i=1,2,3\}`$, because, as it follows from Eq. (64),
$$I_+y^1=y^2,I_{}y^2=y^1,K_+y^1=y^3,K_{}y^3=y^1.$$
(66)
However, a canonical basis in $`V(0,1)`$, denoted $`\{\eta _i,i=1,2,3\}`$ can be obtained from $`\{y^i,i=1,2,3\}`$ by the transformation
$$\eta _1=y^1,\eta _2=y^2,\eta _3=y^3,$$
(67)
or
$$\eta _m^{}=\underset{m}{}G_m^{}^mx_m^{}=\underset{m}{}G_m^{}^my^m.$$
(68)
The transformation matrix $`G`$, as well as the basis vectors, is defined up to a phase factor. The choice of Eq. (67) corresponds to
$$G_m^{}^m=(1)^{\frac{1}{3}+e_m}\delta _{m^{},m},$$
(69)
where $`e_m=(i_3+y/2)_m`$ is the electric charge of the state $`m=(i_3,y)`$ .
In the following, the spaces $`V(1,0)`$ and $`V(0,1)`$ will be denoted also by $`V(3)`$ and $`V(3^{})`$, respectively, because they are 3-dimensional.
V. Tensor representations
Def. I. The object $`T_{i_1,\mathrm{},i_P}`$ is called covariant tensor of rank $`P`$ with respect to $`SU(3)`$ if at the action of $`USU(3)`$ has the transformation law
$$\underset{¯}{T}_{i_1,\mathrm{},i_P}=U_{i_1}^{k_1}\mathrm{}U_{i_P}^{k_P}T_{k_1,\mathrm{},k_P}.$$
(70)
Here and in the following the summation convention of the repeated indices is considered.
Def. II. The object $`T^{i_1,\mathrm{},i_Q}`$ is called contravariant tensor of rank $`Q`$ with respect to $`SU(3)`$ if at the action of $`USU(3)`$ has the transformation law
$$\underset{¯}{T}^{i_1,\mathrm{},i_Q}=(U_{i_1}^{k_1})^{}\mathrm{}(U_{i_Q}^{k_Q})^{}T^{k_1,\mathrm{},k_Q}.$$
(71)
The representations determined by such transformation formulas are called tensor representations, and are presented in detail in ch. 10, sect. 2. The components of the tensors I and II can be considered as elements of a $`3^P`$, respectively $`3^Q`$ \- dimensional space. Mixed tensors can be obtained by the direct product of the covariant and contravariant tensors. These spaces carry $`SU(3)`$ representations which are in general reducible. Therefore, of a particular interest are those tensors, called irreducible, whose components are the basis elements in spaces of $`SU(3)`$ irreducible representation.
The action of $`SU(3)`$ defined by Eq. (70),(71) commutes with the action of the permutation group on the set of tensor indices. Therefore, irreducible are only the tensors which, as functions of indices, provide an irreducible representation of the permutation group. Such tensors are obtained by linear combinations of the type - I and II tensors. The resulting tensor should be either symmetric, or antisymmetric with respect to the permutation of well- defined subsets of indices.
Objects with the transformation properties of the tensors I and II are represented by the set of basis elements in the spaces obtained by the multiple direct product of the spaces $`V(3)`$ and $`V(3^{})`$, respectively. The spaces
$$V(P)=V^1(3)\mathrm{}V^P(3),V(Q)=V^1(3^{})\mathrm{}V^P(3^{}),$$
(72)
have as basis elements
$$(a)T_{i_1\mathrm{}i_P}=x_{i_1}^{(1)}\mathrm{}x_{i_P}^{(P)},(b)T^{j_1\mathrm{}j_Q}=y^{(1)j_1}\mathrm{}x^{(Q)j_Q}.$$
(73)
The generators of the infinitesimal transformations in $`V(P)`$ have the form
$$F_k=F_k^{(1)}I^{(2)}\mathrm{}I^{(P)}+\mathrm{}I^{(1)}I^{(2)}\mathrm{}F_k^{(P)},$$
(74)
where $`F_k^{(i)}`$, $`k=1,8`$, and $`I^{(i)}`$ denote the $`su(3)`$ generator, and, respectively, the unit operator, in the space $`V^i(3)`$. Similarly is defined the action on $`V(P)`$ of the $`su(3)`$ operators $`\{A_k^i,i,k=1,2,3\}`$.
The direct product of the highest weight vectors $`x_1^{(k)}`$, $`k=1,\mathrm{},P`$ give the element of the tensor $`T^_P`$
$$T_{11\mathrm{}1}^_P=|\underset{¯}{M}>=x_1^{(1)}\mathrm{}x_1^{(P)},$$
(75)
which is a highest weight vector in $`V(P)`$, with $`(i_3,y)_{\underset{¯}{M}}=(P/2,P/3)`$. The expressions of the representation operators $`I_{}`$, $`K_{}`$, are similar to $`F_k`$ of Eq. (74), and are symmetric with respect to $`I_{}^{(k)}`$, $`K_{}^{(k)}`$, $`k=1,P`$. Therefore, their action on the highest weight component $`T_{11\mathrm{}1}^_P`$ generates the components of a symmetric, irreducible, covariant tensor. The symmetry with respect to the permutation of the lower indices allows to replace the Kronecker product of Eq. (73) by the Young product. In the Young product, the vectors $`x_i^{(k)}`$ and $`x_i^{(j)}`$, from $`V^k(3)`$ and $`V^j(3)`$, with $`kj`$, are supposed to be the same, and are both denoted by $`x_i`$. This procedure leads to a mapping $`T^_PT^P`$, where $`T^P`$ is a polynomial of degree P in three variables having the same transformation properties under $`SU(3)`$ as the basis elements of $`V(3)`$,
$$T_{i_1\mathrm{}i_P}^P=x_{i_1}\mathrm{}x_{i_P}=(x_1)^{p_1}(x_2)^{p_2}(x_3)^{p_3},$$
(76)
where $`p_1`$, $`p_2`$, $`p_3`$ are non-negative integers such that
$$p_1+p_2+p_3=P.$$
(77)
The components of $`T^P`$ are eigenvectors of the operators $`F_3`$ and $`F_8`$, and generate an orthogonal basis in $`V(P,0)`$. Similarly it is possible to obtain irreducible representations equivalent to $`D(0,Q)`$, on spaces generated by contravariant symmetric tensors,
$$T^{Qj_1\mathrm{}j_Q}=y^{j_1}\mathrm{}y^{j_Q}=(y^1)^{q_1}(y^2)^{q_2}(y^3)^{q^3},$$
(78)
where $`q_1`$, $`q_2`$, $`q_3`$ are non-negative integers such that
$$q_1+q_2+q_3=Q.$$
(79)
The $`su(3)`$ generators act on the tensors $`T^P`$ and $`T^Q`$ as differential operators, completely defined by their action in the spaces $`V(3)`$ and $`V(3^{})`$.
The dimension of the space $`V(P,0)`$ equals the number of the components of the tensor $`T^P`$. This is given by the number of partitions of $`P`$ in $`p_1`$, $`p_2`$, $`p_3`$, such that
$$dimV(P,0)=N_P=\underset{p_1=0}{\overset{P}{}}\underset{p_2=0}{\overset{Pp_1}{}}1=\frac{(P+1)(P+2)}{2}.$$
(80)
Similarly can be obtained the dimension of the space generated by the tensor $`T^Q`$,
$$dimV(Q,0)^{}=dimV(0,Q)=N_Q=\frac{(Q+1)(Q+2)}{2}.$$
(81)
Representations which are equivalent to $`D(P,Q)`$ can be obtained by the decomposition of the direct product $`D(P,0)D(Q,0)^{}`$. The basis of the product space has the form
$$T_{i_1\mathrm{}i_P}^{j_1\mathrm{}j_Q}=(x_1)^{p_1}(x_2)^{p_2}(x_3)^{p_3}(y^1)^{q_1}(y^2)^{q_2}(y^3)^{q_3},$$
(82)
with
$$\underset{i=1}{\overset{3}{}}p_i=P,\underset{i=1}{\overset{3}{}}q_i=Q.$$
(83)
This tensor is symmetric with respect to the permutation of the upper, as well as of lower indices. By contraction with the invariant tensor $`\delta _j^i\delta _{ij}`$, one obtains the sequence of tensors $`T^{(k)}`$,
$$T_{i_{k+1}\mathrm{}i_P}^{(k)j_{k+1}\mathrm{}j_Q}=\delta _{j_k}^{i_k}T_{i_k\mathrm{}i_P}^{(k1)j_k\mathrm{}j_Q}.$$
(84)
The spaces $`V^{(k)}`$ generated by them are all $`SU(3)`$ invariant, and satisfy the relationship
$$V(P,0)V(Q,0)^{}=V^{(0)}V^{(1)}\mathrm{}V^{(n)},n=\mathrm{min}(P,Q).$$
(85)
The orthogonal complement of the space $`V^{(k+1)}`$ in $`V^{(k)}`$ with respect to the scalar product defined by the contraction of the indices is represented by the linear combinations of the tensors from $`V^{(k)}`$ which are traceless with respect to any contraction, and is denoted by $`V_0^{(k)}`$. Thus,
$$V^{(k)}=V_0^{(k)}V^{(k+1)},$$
(86)
and
$$V(P,0)V(Q,0)^{}=\underset{k=0}{\overset{n}{}}V_0^{(k)}=\underset{k=0}{\overset{n}{}}V(Pk,Qk).$$
(87)
The representation on the space $`V_0^{(0)}`$ is equivalent to $`D(P,Q)`$, and
$$dimV(P,Q)=dimV^{(0)}dimV^{(1)}=N_PN_QN_{P1}N_{Q1}=$$
(88)
$$=\frac{(P+1)(Q+1)(P+Q+2)}{2}.$$
This is a special case of the general Weil’s formula for the dimension of the irreps for every simple Lie group ch. X sect. 13.4, ch. 8 sect. 8.
The weight diagram of $`V(P,Q)`$ is obtained by substracting from the highest weight $`(i_3,y)_M=((P+Q)/2,(PQ)/3)`$ linear combinations with integer coefficients of the simple roots $`\alpha _{12}`$, $`\alpha _{13}`$, and $`\alpha _{32}`$. This diagram has three symmetry axes, and is bounded by the polygonal line drawn in Fig. 2. The number of weights of this diagram is smaller than $`dimV(P,Q)`$, and to label the states it is necessary to find additional operators, which commute with $`I_3`$ and $`Y`$, but not with all $`A_k^i`$, $`ik`$.
The problem of labeling the states of the $`SU(n)`$ irrep spaces is solved by the canonical factorization $`SU(n)U(1)SU(n1)`$ -II. The additional operators are in this case the Casimir operators of the subgroups $`SU(k)`$, $`k=2,\mathrm{},n1`$. For $`SU(3)`$ the additional operator can be chosen as the Casimir operator $`I^2`$ of the $`su(2)`$ subalgebra $`\{I_{},I_3,I_+\}`$,
$$I^2=F_1^2+F_2^2+F_3^2,$$
(89)
which satisfies
$$[I^2,F_3]=0,[I^2,Y]=0,[I^2,I_\pm ]=0.$$
(90)
The basis vectors of the space $`V(P,Q)`$ can be labeled by the eigenvalues of the operators which form a complete set, in this case $`F`$, $`G`$, $`I^2`$, $`I_3`$ and $`Y`$. However, in applications, instead of the eigenvalues $`f`$ and $`g`$ is more convenient to use the integers $`P`$ and $`Q`$. The eigenvalue equations for the labeling operators have the known form
$$I^2|PQii_3y>=i(i+1)|PQii_3y>,$$
(91)
$$I_3|PQii_3y>=i_3|PQii_3y>,$$
(92)
$$Y|PQii_3y>=y|PQii_3y>,$$
(93)
and the vectors of the canonical basis are subject to the normalization condition
$$<PQii_3y|PQi^{}i_3^{}y^{}>=\delta _{ii^{}}\delta _{i_3i_3^{}}\delta _{yy^{}}.$$
(94)
The orthonormal basis in $`V(P,0)`$ can be constructed by the normalization of Eq. (76) with the factor $`a^P(i,i_3,y)`$, such that
$$|P0ii_3y>=a^P(i,i_3,y)(x_1)^{p_1}(x_2)^{p_2}(x_3)^{p_3}.$$
(95)
By acting with the operators $`I^2`$, $`I_3`$, $`Y`$ on both sides of this equation one obtains
$$p_1=i+i_3,p_2=ii_3,p_3=\frac{P}{3}y,$$
(96)
and then, by using Eq. (77), one obtains a linear relationship between isospin and hypercharge,
$$i=\frac{P}{3}+\frac{y}{2}.$$
(97)
The weights of the basis vectors of the space $`V(P,0)`$ are not degenerate, and the states labeling can be achieved by using only $`(i_3,y)m`$. Therefore, in the following will be used the notation
$$\xi _P^m\xi _{Pi_3}^y|P0ii_3y>=a^P(i_3,y)(x_1)^{p_1}(x_2)^{p_2}(x_3)^{p_3}.$$
(98)
The complex conjugation of this equation leads to the expression of the contravariant basis vectors,
$$\xi _m^Q=(\xi _Q^m)^{}=a^Q(i_3,y)(y^1)^{i+i_3}(y^2)^{ii_3}(y^3)^{\frac{Q}{3}y}.$$
(99)
The vector space generated by the basis $`\xi _m^Q`$ carries an irrep equivalent to $`D(0,Q)`$, but where the matrix elements of the operators $`I_\pm `$ and $`K_\pm `$ are not positive definite. The transition to a basis $`\eta _Q^m`$ in which the phase convention is fulfilled, is generated by the transformation of Eq.(67), such that
$$\eta _Q^m=a^Q(i_3,y)(y^1)^{i+i_3}(y^2)^{ii_3}(y^3)^{\frac{Q}{3}y}=(1)^{i+i_3}(\xi _Q^m)^{}.$$
(100)
By using Eq. (89) and (63) it follows that Eq. (100) represents the canonical basis of the space $`V(0,Q)`$,
$$\eta _Q^m|0Qii_3y>=(1)^{i+i_3}|Q0ii_3y>^{},$$
(101)
or, by using the notation
$$k=\frac{Q}{3}\frac{y}{2},$$
(102)
$$\eta _Q^m\eta _{Qk_3}^y=|0Qkk_3y>=$$
(103)
$$=(1)^{\frac{Q}{3}(k+\frac{y}{2})}a^Q(k_3,y)(y^1)^{kk_3}(y^2)^{k+k_3}(y^3)^{\frac{Q}{3}+y}.$$
VI. The Clebsch-Gordan series
The space obtained by the direct product of two irrep spaces is not irreducible, and it should be decomposed in a direct sum of irreducible spaces. The basis of the product space can be related to the canonical basis of the direct sum of irrep spaces by a unitary transformation $`S`$, such that
$$S(V(P1,Q1)V(P2,Q2))=_k_{\gamma =1}^{m_k}V^\gamma (P_k,Q_k).$$
(104)
This transformation brings the matrices of representation to a block-diagonal form,
$$S(R(P1,Q1)R(P2,Q2)))S^1=_km_kR(P_k,Q_k).$$
(105)
The elements of the matrix $`S`$ are called Clebsch-Gordan coefficients, and the set of irrep spaces $`V^\gamma (P_k,Q_k)`$ which appear in the right-hand side of Eq. (104) represent the Clebsch-Gordan series. The index $`\gamma `$ is necessary to distinguish between the spaces which are isomorphic, and carry the same irrep. The main difficulty which appears with respect to the case of the group $`SU(2)`$ is due to the fact that in general, for certain representations $`D(P_k,Q_k)`$, $`m_k`$ can be greater than 1.
The basis of the direct product space has the form of a mixed tensor, without permutation symmetry. The linear combinations of the components of this tensor obtained by contraction and antisymmetrization generate subspaces which are $`SU(3)`$ invariant. Therefore, the basis of the irrep spaces is generated by those mixed tensors which are symmetric at the permutation of the upper and lower indices, and traceless with respect to any contraction. The tensors which give the decomposition of the direct product can be constructed explicitly below, for each particular case.
(A). $`D(P_1,0)D(P_2,0)`$ and $`D(0,Q_1)D(0,Q_2)`$
The direct product basis is
$$T_{(i_1\mathrm{}i_{P_1})(i_{P_1+1}\mathrm{}i_{P_1+P_2})}=T_{1i_1\mathrm{}i_{P_1}}T_{2i_{P_1+1}\mathrm{}i_{P_1+P_2}},$$
(106)
a tensor symmetric to permutations within each subset of lower indices. By contraction with the antisymmetric invariant contravariant tensor $`ϵ^{ijk}`$ one obtains the sequence of tensors $`T^{(k)}`$,
$$T_{(i_{k+1}\mathrm{}i_{P_1})(i_{P_1+k+1}\mathrm{}i_{P_1+P_2})}^{(k)j_1\mathrm{}j_k}=ϵ^{j_ki_ki_{P_1+k}}T_{(i_k\mathrm{}i_{P_1})(i_{P_1+k}\mathrm{}i_{P_1+P_2})}^{(k1)j_1\mathrm{}j_{k1}},$$
(107)
which are traceless with respect to every contraction, but not symmetric with respect to the permutation of the lower indices. The spaces $`V^{(k)}`$ generated by them satisfy the relationship
$$V(P_1,0)V(P_2,0)=V^{(0)}V^{(1)}\mathrm{}V^{(n)},n=\mathrm{min}(P_1,P_2).$$
(108)
The tensors $`T^{(k)}`$ can be symmetrized with respect to every pair of lower indices placed in different subsets. This procedure leads to the irreducible tensors $`T_S^{(k)}`$,
$$T_{S(i_1\mathrm{}i_{P_1+P_22k})}^{(k)j_1\mathrm{}j_k}=T_{(\mathrm{}i_r\mathrm{})(\mathrm{}i_s\mathrm{})}^{(k)j_1\mathrm{}j_k}+T_{(\mathrm{}i_s\mathrm{})(\mathrm{}i_r\mathrm{})}^{(k)j_1\mathrm{}j_k},$$
(109)
which generate spaces denoted $`V_S^{(k)}`$. The equations (107) and (109) lead to a decomposition of the space $`V^{(k)}`$ of the form
$$V^{(k)}=V_S^{(k)}V^{(k+1)},V_S^{(k)}=V(P_1+P_22k,k).$$
(110)
By using the Eq. (108) and (110) one obtains the Clebsch-Gordan series
$$D(P_1,0)D(P_2,0)=\underset{k=0}{\overset{n}{}}D(P_1+P_22k,k),n=\mathrm{min}(P_1,P_2).$$
(111)
Similarly it is possible to show that
$$D(0,Q_1)D(0,Q_2)=\underset{k=0}{\overset{n}{}}D(k,Q_1+Q_22k),n=\mathrm{min}(Q_1,Q_2).$$
(112)
(B). $`D(P,0)D(0,Q)`$
According to Eq. (87), the decomposition of this direct product has the form
$$D(P,0)D(0,Q)=\underset{k=0}{\overset{n}{}}D(Pk,Qk),n=\mathrm{min}(P,Q)$$
(113)
The basis of the space $`V_0^{(k)}`$ is represented by the traceless tensors $`T_0^{(k)}`$, and can be constructed by using the method of projection operators ch. 7 sect. 3. Thus,
$$T_{0i_{k+1}\mathrm{}i_P}^{(k)j_{k+1}\mathrm{}j_Q}=_{SU(3)}𝑑U[R(Pk,Qk)_U]_{ab}U_{i_{k+1}}^{\beta _{k+1}}\mathrm{}$$
(114)
$$\mathrm{}U_{i_P}^{\beta _P}(U_{j_{k+1}}^{\alpha _{k+1}})^{}\mathrm{}(U_{j_Q}^{\alpha _Q})^{}T_{\beta _{k+1}\mathrm{}\beta _P}^{(k)\alpha _{k+1}\mathrm{}\alpha _Q}.$$
Here $`USU(3)`$ denotes a group element, $`dU`$ is an invariant measure on $`SU(3)`$, and $`[R(Pk,Qk)_U]_{ab}`$ is a matrix element of $`U`$ in the representation $`D(Pk,Qk)`$. The indices $`a,b`$, in the vector notation, are such that $`a(i,i_3,y)`$ it is the same as $`(_{i_{k+1}\mathrm{}i_P}^{j_{k+1}\mathrm{}j_Q})`$ in the tensor notation, and $`b`$ is arbitrary, but fixed.
(C). $`D(P_1,Q_1)D(P_2,Q_2)`$.
The basis of the reducible space obtained by the direct product is represented by the components of the tensor
$$T_{(i_1\mathrm{}i_{P_1})(i_{P_1+1}\mathrm{}i_{P_1+P_2})}^{(j_1\mathrm{}j_{Q_1})(j_{Q_1+1}\mathrm{}j_{Q_1+Q_2})}=T_{1i_1\mathrm{}i_{P_1}}^{j_1\mathrm{}j_{Q_1}}T_{2i_{P_1+1}\mathrm{}i_{P_1+P_2}}^{j_{Q_1+1}\mathrm{}j_{Q_1+Q_2}},$$
(115)
where $`T_1`$ and $`T_2`$ are traceless symmetric tensors. The reduction can be performed in this case by using the procedure suggested by Coleman in ref. , consisting of two steps:
(1) - the product space is decomposed in a direct sum of reducible invariant spaces $`V(P,P^{};Q,Q^{})`$ generated by traceless tensors.
(2) - the spaces $`V(P,P^{};Q,Q^{})`$ are decomposed in a sum of irreducible spaces.
(1) - Following the procedure applied in the case (B), it is possible to construct traceless tensors $`T^{(m,n)}`$
$$T_{(i_{m+1}\mathrm{}i_{P_1})(i_{P_1+n+1}\mathrm{}i_{P_1+P_2})}^{(m,n)(j_{n+1}\mathrm{}j_{Q_1})(j_{Q_1+m+1}\mathrm{}j_{Q_1+Q_2})}=\delta _{j_{Q_1+m}}^{i_m}T_{(i_m\mathrm{}i_{P_1})(i_{P_1+n+1}\mathrm{}i_{P_1+P_2})}^{(m1,n)(j_{n+1}\mathrm{}j_{Q_1})(j_{Q_1+m}\mathrm{}j_{Q_1+Q_2})}=$$
(116)
$$\delta _{j_n}^{i_{P_1+n}}T_{(i_{m+1}\mathrm{}i_{P_1})(i_{P_1+n}\mathrm{}i_{P_1+P_2})}^{(m,n1)(j_n\mathrm{}j_{Q_1})(j_{Q_1+m+1}\mathrm{}j_{Q_1+Q_2})}.$$
The spaces generated by the tensors $`T^{(m,n)}`$ are denoted $`V^{(m,n)}`$, and the orthogonal complements of the spaces $`V^{(m+1,n)}`$ and $`V^{(m,n+1)}`$ in $`V^{(m,n)}`$ are denoted by $`V^{(m_0,n)}`$ and $`V^{(m,n_0)}`$, respectively. The bases of these spaces can be obtained, for instance, by using Eq. (114), and $`V^{(m_0,n_0)}=V(P_1m,P_2n;Q_1n,Q_2m)`$. Because
$$V^{(m,n)}=V^{(m_0,n_0)}V^{(m_0,n+1)}V^{(m+1,n_0)}V^{(m+1,n+1)}$$
(117)
one obtains the decomposition
$$V(P_1,Q_1)V(P_2,Q_2)=\underset{m_0=0}{\overset{\mathrm{min}(P_1,Q_2)}{}}\underset{n_0=0}{\overset{\mathrm{min}(P_2,Q_1)}{}}V^{(m_0,n_0)}=$$
(118)
$$=\underset{m=0}{\overset{\mathrm{min}(P_1,Q_2)}{}}\underset{n=0}{\overset{\mathrm{min}(P_2,Q_1)}{}}V(P_1m,P_2n;Q_1n,Q_2m).$$
(2) - Let $`T_{0(i_1\mathrm{}i_r)(i_{r+1}\mathrm{}i_{r+r^{}})}^{(j_1\mathrm{}j_s)(j_{s+1}\mathrm{}j_{s+s^{}})}T_{0(r)(r^{})}^{(s)(s^{})}`$ a traceless basis tensor in $`V(r,r^{};s,s^{})`$. By symmetrization with respect to any pair of upper indices belonging to different subsets one obtains a tensor $`T_{0(r)(r^{})}^{(s+s^{})}`$. This is symmetric in the upper indices, and represents the basis of a space denoted $`V(r,r^{};s+s^{})`$. The orthogonal complement of the space $`V(r,r^{};s+s^{})`$ in $`V(r,r^{};s,s^{})`$, denoted $`V(r+r^{}+1;s1,s^{}1)`$, is generated by the tensor
$$ϵ_{ij_1j_{s+1}}T_{0(i_1\mathrm{}i_r)(i_{r+1}\mathrm{}i_{r+r^{}})}^{(j_1\mathrm{}j_s)(j_{s+1}\mathrm{}j_{s+s^{}})}$$
(119)
which is traceless and symmetric in all lower indices. Thus,
$$V(r,r^{};s,s^{})=V(r,r^{};s+s^{})V(r+r^{}+1;s1,s^{}1).$$
(120)
The complete decomposition of the subspaces $`V(r,r^{};s+s^{})`$ and $`V(r+r^{}+1;s1,s^{}1)`$ can be obtained by the method used in the case (A), and the corresponding series are similar to Eq. (111) and (112). Therefore, the final result is
$$V(r,r^{};s,s^{})=V(r+r^{};s+s^{})\underset{k=1}{\overset{\mathrm{min}(r,r^{})}{}}V(r+r^{}2k;s+s^{}+k)$$
(121)
$$\underset{k=1}{\overset{\mathrm{min}(s,s^{})}{}}V(r+r^{}+k,s+s^{}2k).$$
The formulas of Eq. (118) and (121) solve the problem of the Clebsch-Gordan series for $`SU(3)`$ in the general case. The same permutation symmetry arguments as above, lead to the rules which give the decomposition of the product between the Young tableaux associated to the irreps ch. 8. The problem of the decomposition of the direct product for an arbitrary semisimple Lie group was solved by Kostant and Steinberg ch. 8 sect. 8. A purely geometrical procedure to obtain the decomposition was given by Speiser , .
VII. The weights multiplicity
The explicit form of the canonical basis of the spaces $`V(P,0)`$ and $`V(0,Q)`$ is given by Eq. (98) and (103), respectively. The canonical basis of $`V(P,Q)`$ can be obtained by decomposing the direct product between $`V(P,0)`$ and $`V(0,Q)`$. By using the notation $`s(P,Q)`$ for the irreps labels, Eq. (104) leads to
$$|sii_3y>_\gamma =\underset{\mu mjk}{}$$
(122)
$$<_{jm\mu ki_3my\mu }^{s_1s_2}|_{ii_3y}^{s\gamma }>|s_1jm\mu >|s_2ki_3my\mu >$$
where $`<_{jm\mu ki_3my\mu }^{s_1s_2}|_{ii_3y}^{s\gamma }>`$ are the $`SU(3)`$ Clebsch-Gordan coefficients, and $`\gamma `$ is a label used to distinguish between the subspaces which all carry the same irrep $`s`$. Because
$$|s_1jm\mu >|s_2ki_3my\mu >$$
(123)
are not eigenstates of the operator $`I^2`$, and $`SU(3)U(1)SU(2)`$, it is convenient to consider the transformation matrix $`S`$ as a product of two unitary matrices, denoted $`\alpha `$ and $`C`$, such that
$$S=\alpha C$$
(124)
and
$$\alpha \alpha ^{}=I,CC^{}=I.$$
(125)
The elements of the matrices $`C`$ and $`\alpha `$ are the $`SU(2)`$ Clebsch-Gordan coefficients, and, respectively, the isoscalar factors. By using this factorization (the Racah lema ), Eq. (122) becomes
$$|sii_3y>_\gamma =\underset{\mu mjk}{}\alpha _{\mu jk}^{iyss_1s_2\gamma }\times $$
(126)
$$C_{mi_3mi_3}^{jki}|s_1jm\mu >|s_2ki_3my\mu >$$
with
$$C_{j_3k_3i_3}^{jki}(<jj_3|<kk_3|)|jkii_3>.$$
(127)
If $`s_1=(P,0)`$, $`s_2=(0,Q)`$ and $`s=(P,Q)`$, then the summation on $`j`$ and $`k`$ in Eq. (126) reduces to the one on $`\mu `$, and the expression of the canonical basis in the space $`V(P,Q)`$ becomes
$$|PQii_3y>=\underset{\mu m}{}\alpha _{\mu \frac{P}{3}+\frac{\mu }{2}\frac{Q}{3}\frac{y\mu }{2}}^{iy(PQ)(P0)(0Q)}C_{mi_3mi_3}^{\frac{P}{3}+\frac{\mu }{2}\frac{Q}{3}\frac{y\mu }{2}i}\xi _{Pm}^\mu \eta _{Qi_3m}^{y\mu }.$$
(128)
The weights multiplicity is given by the number of the values taken by $`i`$ in this expression, at fixed values of $`i_3`$ and $`y`$. By using the properties of the coefficients $`C`$ one obtains
$$|kj|ik+j,$$
(129)
where $`j=P/3+\mu /2`$ and $`k=Q/3(y\mu )/2`$.
The minimum value of $`i`$ is the same for all the terms of the sum in Eq. (128), equal to
$$i_{min}=|kj|=|\frac{QP}{3}\frac{y}{2}|=$$
(130)
$$\frac{PQ}{3}+\frac{y}{2}ify\frac{2}{3}(PQ)(a)$$
and
$$\frac{QP}{3}\frac{y}{2}ify\frac{2}{3}(PQ)(b).$$
The maximum value of $`i`$ for an arbitrary term of the sum is
$$i_{max}=k+j=\frac{P+Q}{3}\frac{y}{2}+\mu ,$$
different from term to term. The highest value of $`i_{max}`$ will be denoted by $`i_M`$. The coefficients $`C_{j_3k_3i_3}^{jki}`$ vanish when $`i>j+k`$, but the sum of Eq. (128) is not zero if it contains at least one term, which means $`ii_M`$. The range of $`i`$ in $`|PQii_3y>`$ is limited by the inequalities
$$|i_3|i,i_{min}ii_M$$
(131)
where
$$i_M=\frac{P+Q}{3}\frac{y}{2}+(\mu )_{max}.$$
(132)
The maximum value of $`\mu `$ which appears in Eq. (128), denoted $`(\mu )_{max}`$, can be found by using Eq. (77), (79), (96) and (103),
$$p_3=\frac{P}{3}\mu ,q_3=\frac{Q}{3}+y\mu ,$$
(133)
$$0p_3P,0q_3Q.$$
(134)
These lead to the inequalities
$$\frac{2}{3}P\mu \frac{P}{3},y\frac{2}{3}Q\mu \mu +\frac{Q}{3},$$
(135)
which give the maximum value
$$(\mu )_{max}=\mathrm{min}(\frac{P}{3},y+\frac{Q}{3})=$$
(136)
$$\frac{P}{3}ify\frac{PQ}{3}$$
and
$$y+\frac{Q}{3}ify\frac{PQ}{3}.$$
This result shows that the value of $`i_M`$ is
$$i_M=\frac{2P+Q}{3}\frac{y}{2}ify\frac{PQ}{3}(a),$$
(137)
and
$$i_M=\frac{P+2Q}{3}+\frac{y}{2}ify\frac{PQ}{3}(b).$$
The equations (130) and (137) can be represented in the orthogonal frame $`(i,y)`$ by straight lines. Thus, one obtains the diagram of Fig. (3), bounded by the lines $`u`$, $`d`$, $`\nu `$ and $`\delta `$ defined by the equations
$$u:i=\frac{y}{2}+\frac{P+2Q}{3},\nu :i=\frac{y}{2}+\frac{PQ}{3},$$
$$d:i=\frac{y}{2}+\frac{2P+Q}{3},\delta :i=\frac{y}{2}\frac{PQ}{3}.$$
These lines cross at the points $`A`$, $`B`$, $`C`$, $`D`$, with coordinates
$$A(\frac{Q}{2},\frac{2P+Q}{3});B(\frac{P+Q}{2},\frac{PQ}{3});$$
$$C(\frac{P}{2},\frac{P+2Q}{3});D(0,\frac{2}{3}(QP)).$$
The quantum numbers $`(i,y)`$ in Eq. (128) may take values only within the parallelogram bounded by these lines. The whole set of quantum numbers $`(i,i_3,y)`$ which label the canonical basis of $`V(P,Q)`$ can be represented by the sites of a 3-dimensional lattice (Fig. 4). The projection of this lattice in the plane $`(i_3,y)`$ is the weight diagram of Fig. 2, and its volume is equal to the dimension of the space $`V(P,Q)`$, given in Eq. (88) . The same result was obtained by Biedenharn -II by using the matrix elements of the operators $`K_\pm `$ and $`L_\pm `$. The formulas which allow the calculus of the weights multiplicities for any semisimple Lie group have been given by Kostant (1959) and Freudenthal (1969) ch. X sect. 13.4.
VIII. The matrix elements of the generators
The $`U(n)`$ irrep spaces and the matrix elements of the generators can be obtained by using the Gel’fand-Zetlin method ch. 10. In this work the matrix elements of the operators $`I_\pm `$, $`K_\pm `$ and $`L_\pm `$ are calculated by extending the procedure applied for $`SU(2)`$, of using the matrix elements of the commutators $`[A_k^i,A_i^k]`$ which are known.
(A). Denoting by $`B_{ii_3}`$ the matrix elements of the operator $`I_{}`$, the commutation relations
$$[Y,I_\pm ]=0,[I_3,I_\pm ]=\pm I_\pm ,[I^2,I_\pm ]=0$$
(138)
show that
$$I_{}|PQii_3y>=B_{ii_3}|PQii_31y>,$$
(139)
$$I_+|PQii_3y>=B_{ii_3+1}|PQii_3+1y>.$$
Thus, the matrix element of the commutator $`[I_+,I_{}]=2I_3`$ in the state $`|PQii_3y>`$ is
$$B_{ii_3}^2B_{ii_3+1}^2=2i_3.$$
(140)
By using the conditions $`B_{ii_3}0`$, $`B_{ii}=0`$, one obtains
$$B_{ii_3}=\sqrt{i(i+1)i_3(i_31)}.$$
(141)
(B). In the adjoint representation (Eq. (19) (29)) the operators $`K_\pm `$ and $`L_\pm `$ are eigenvectors of $`ad_{I_3}`$ and $`ad_Y`$,
$$[I_3,K_\pm ]=\frac{\pm 1}{2}K_\pm ,[Y,K_\pm ]=\pm K_\pm ,$$
(142)
$$[I_3,L_\pm ]=\frac{\pm 1}{2}L_\pm ,[Y,L_\pm ]=\pm L_\pm ,$$
but not for $`ad_{I^2}`$, because
$$[I^2,K_+]=L_+I_++K_+(\frac{3}{4}+I_3),[I^2,L_+]=K_+I_{}+L_+(\frac{3}{4}I_3).$$
(143)
Therefore, the vectors $`K_+|PQii_3y>`$ and $`L_+|PQii_3y>`$ are linear combinations of the canonical basis vectors of $`V(P,Q)`$ which have the weights given by Eq. (57),
$$K_+|PQii_3y>=\underset{i^{}}{}\gamma _{ii_3y}^i^{}|PQi^{}i_3+\frac{1}{2}y+1>,$$
(144)
$$L_+|PQii_3y>=\underset{i^{}}{}\omega _{ii_3y}^i^{}|PQi^{}i_3\frac{1}{2}y+1>.$$
(145)
By acting with the operator $`I^2`$ on both sides of Eq. (144) and (145), and then by using Eq. (141), (143) and (94), we obtain a homogeneous system of equations for the coefficients $`\gamma _{ii_3y}^i`$ and $`\omega _{ii_3y}^i^{}`$. The determinant of this system vanishes only if
$$i^{}=i\pm \frac{1}{2},$$
(146)
and in this case the solution is
$$\gamma _{ii_3y}^{i+\frac{1}{2}}=\sqrt{\frac{i+i_3+1}{ii_3}}\omega _{ii_3+1y}^{i+\frac{1}{2}},$$
(147)
$$\gamma _{ii_3y}^{i\frac{1}{2}}=\sqrt{\frac{ii_3}{i+i_3+1}}\omega _{ii_3+1y}^{i\frac{1}{2}}.$$
According to Eq. (146) the sums in Eq. (144) and (145) contain only two terms, such that
$$K_+|PQii_3y>=\gamma _{ii_3y}^{i+\frac{1}{2}}|PQi+\frac{1}{2}i_3+\frac{1}{2}y+1>+$$
(148)
$$\gamma _{ii_3y}^{i\frac{1}{2}}|PQi\frac{1}{2}i_3+\frac{1}{2}y+1>,$$
$$L_+|PQii_3y>=\omega _{ii_3y}^{i+\frac{1}{2}}|PQi+\frac{1}{2}i_3\frac{1}{2}y+1>+$$
(149)
$$\omega _{ii_3y}^{i\frac{1}{2}}|PQi\frac{1}{2}i_3\frac{1}{2}y+1>.$$
The operator $`K_+`$ in Eq. (148) can be replaced by the commutator $`[I_+,L_+]=K_+`$, and then, by using Eq. (149), (141), and (94), it is possible to find two recurrence relationships in $`i_3`$ for $`\gamma _{ii_3y}^{i+\frac{1}{2}}`$ and $`\gamma _{ii_3y}^{i\frac{1}{2}}`$. Denoting by $`\chi _{iy}\gamma _{iiy}^{i+\frac{1}{2}}`$ and $`\kappa _{iy}\gamma _{iiy}^{i\frac{1}{2}}`$, the solution of these relationships is
$$\gamma _{ii_3y}^{i+\frac{1}{2}}=\sqrt{\frac{i+i_3+1}{2i+1}}\chi _{iy},\gamma _{ii_3y}^{i\frac{1}{2}}=\sqrt{\frac{ii_3}{2i}}\kappa _{iy}.$$
(150)
This result, together with Eq. (147), determine the isospin dependent factor of the coefficients $`\omega `$, as
$$\omega _{ii_3y}^{i+\frac{1}{2}}=\sqrt{\frac{ii_3+1}{2i+1}}\chi _{iy},\omega _{ii_3y}^{i\frac{1}{2}}=\sqrt{\frac{i+i_3}{2i}}\kappa _{iy}.$$
(151)
By convention the matrix elements of the operators $`\{A_k^i;i,k=1,2,3\}`$ are real. Moreover, they should be positive for $`I_\pm `$ and $`K_\pm `$, which means that the coefficients $`\chi `$ and $`\kappa `$ are positive real numbers. This property indicates that the operators $`K_{}=(K_+)^{}`$ and $`L_{}=(L_+)^{}`$ act on the canonical basis vectors according to
$$K_{}|PQii_3y>=\gamma _{i+\frac{1}{2}i_3\frac{1}{2}y1}^i|PQi+\frac{1}{2}i_3\frac{1}{2}y1>+$$
(152)
$$\gamma _{i\frac{1}{2}i_3\frac{1}{2}y1}^i|PQi\frac{1}{2}i_3\frac{1}{2}y1>,$$
$$L_{}|PQii_3y>=\omega _{i+\frac{1}{2}i_3+\frac{1}{2}y1}^i|PQi+\frac{1}{2}i_3+\frac{1}{2}y1>+$$
(153)
$$\omega _{i\frac{1}{2}i_3+\frac{1}{2}y1}^i|PQi\frac{1}{2}i_3+\frac{1}{2}y1>.$$
To determine the factors $`\chi `$ and $`\kappa `$ it is necessary to find two new recurrence formulas, together with the related initialization values. These formulas are provided by the matrix elements of the commutators $`[K_+,K_{}]`$ and $`[L_+,L_{}]`$ in an arbitrary state $`|PQii_3y>`$,
$$<PQii_3y|[K_+,K_{}]|PQii_3y>=i_3+\frac{3}{2}y,$$
(154)
and
$$<PQii_3y|L_+,L_{}]|PQii_3y>=i_3+\frac{3}{2}y.$$
(155)
By using Eq. (148), (149), (152), (153) to express the action of the commutators in terms of $`\gamma `$ and $`\omega `$, and then Eq. (150), (151) to express $`\gamma `$ and $`\omega `$ by $`\chi `$ and $`\kappa `$, we get
$$\frac{1}{2i+1}\kappa _{i+\frac{1}{2}y1}^2+\frac{1}{2(i+1)}\chi _{i\frac{1}{2}y1}^2\frac{1}{2i+1}\chi _{iy}^2$$
(156)
$$\frac{1}{2(i+1)}\kappa _{iy}^2=\frac{3y}{2(i+1)}$$
and
$$\frac{1}{2i+1}\kappa _{i+\frac{1}{2}y1}^2+\frac{1}{2i}\chi _{i\frac{1}{2}y1}^2\frac{1}{2i+1}\chi _{iy}^2+\frac{1}{2i}\kappa _{iy}^2=1.$$
(157)
The equations (156) and (157) relate the factors $`\chi `$ at the points $`4:(i,y)`$ and $`3:(i1/2,y1)`$, to the factors $`\kappa `$ at the points $`4:(i,y)`$ and $`1:(i+1/2,y1)`$ (Fig. 5). By adding these equations $`\kappa _{i+\frac{1}{2}y1}^2`$ is eliminated, and one obtains a relationship between $`\chi _3`$, $`\chi _4`$ and $`\kappa _4`$ (Fig. 5). This gives $`\kappa _{iy}^2`$ as a function of $`\chi _3`$ and $`\chi _4`$. By the transformation $`ii+\frac{1}{2}`$, $`yy1`$, one can further obtain $`\kappa _{i+\frac{1}{2}y1}^2`$ as a function of $`\chi _1`$ and $`\chi _2`$. When the resulting expressions of $`\kappa _{iy}^2`$ and $`\kappa _{i+\frac{1}{2}y1}^2`$ are introduced in Eq. (157), one obtains
$$\frac{i+1}{i+\frac{1}{2}}\chi _{iy1}^2+\chi _{iy+1}^2=\frac{i+\frac{3}{2}}{i+1}\chi _{i+\frac{1}{2}y}^2+\chi _{i\frac{1}{2}y}^2.$$
(158)
Similarly, one can find a homogeneous equation for $`\kappa `$,
$$\kappa _{iy1}^2+\frac{i+\frac{1}{2}}{i}\kappa _{iy+1}^2=\frac{i+1}{i+\frac{1}{2}}\kappa _{i+\frac{1}{2}y}^2+\kappa _{i\frac{1}{2}y}^2.$$
(159)
By introducing the notations
$$\varphi _{iy}=\frac{i+1}{i+\frac{1}{2}}\chi _{iy1}^2\chi _{i\frac{1}{2}y}^2,$$
(160)
and
$$\psi _{iy}=\frac{i+1}{i+\frac{1}{2}}\kappa _{i+\frac{1}{2}y}^2\kappa _{iy1}^2,$$
one can see that Eq. (158) and (159) can be expressed as
$$\varphi _{iy}=\varphi _{i+\frac{1}{2}y+1},$$
(161)
respectively
$$\psi _{iy}=\psi _{i\frac{1}{2}y+1}.$$
(162)
These equations state the invariance of the function $`\varphi _{iy}`$ on lines parallel to $`u`$, and of the function $`\psi _{iy}`$ on lines parallel to $`d`$. The action of the operators $`A_k^i`$ on the boundary lines of the diagram of Fig. 3 cannot lead to weights outside the diagram, and therefore, the values of $`\chi `$ and $`\kappa `$ on the boundary will be used to initiate the recurrence formulas. Denoting by $`(\mathrm{\Delta },\tau )`$ and $`(\mathrm{\Omega },\theta )`$ the coordinates $`(i,y)`$ of the points placed on the line $`d`$, respectively $`\nu `$, the initial conditions are
$$\chi _{\mathrm{\Delta }\tau }=0,\kappa _{\mathrm{\Omega }\theta }=0.$$
(163)
These conditions indicate the fact that $`\varphi _{iy}`$ and $`\psi _{iy}`$ vanish on $`d`$, respectively $`\nu `$, while Eq. (161), (162) show further that they vanish everywhere. This result can be expressed in the form
$$z_{iy1}^2=z_{i\frac{1}{2}y}^2,w_{i+\frac{1}{2}y}^2=w_{iy1}^2,$$
(164)
where the factors $`z_{iy}`$ and $`w_{iy}`$ are related to $`\chi _{iy}`$ and $`\kappa _{iy}`$ by the relationships
$$\chi _{iy}^2=\frac{z_{iy}^2}{2(i+1)},\kappa _{iy}^2=\frac{w_{iy}^2}{2i+1},$$
(165)
and satisfy the conditions
$$z_{\mathrm{\Delta }\tau }=0,w_{\mathrm{\Omega }\theta }=0.$$
(166)
The factors $`z`$ and $`w`$ take the same values in all points of any line parallel to $`d`$, respectively $`u`$. Therefore, to calculate $`\chi `$ and $`\kappa `$ at any point of the diagram of Fig. 3 it is enough to know the values of $`z`$ on $`u`$ or $`\nu `$, and of $`w`$ on $`d`$ or $`\delta `$.
When the expressions of $`\chi `$ and $`\kappa `$ given by Eq. (165) are introduced in the system of Eq. (156), (157), one obtains a system in $`z`$ and $`w`$, which can be solved with respect to $`z_{i\frac{1}{2}y1}^2`$ and $`w_{i+\frac{1}{2}y1}^2`$,
$$\frac{1}{i}z_{i\frac{1}{2}y1}^2=\frac{1}{i+\frac{1}{2}}z_{iy}^2\frac{1}{i(2i+1)}w_{iy}^2+3y+2i+2,$$
(167)
$$\frac{1}{i+1}w_{i+\frac{1}{2}y1}^2=\frac{1}{i+\frac{1}{2}}w_{iy}^2+\frac{1}{(i+1)(2i+1)}z_{iy}^2+3y2i.$$
(168)
By replacing $`(i,y)`$ with $`(\mathrm{\Omega },\theta )`$ in Eq. (167), respectively with $`\mathrm{\Delta },\tau )`$ in Eq. (168), and using Eq. (166), we obtain two simple recurrence relationships between the factors $`z`$ on the line $`\nu `$, and $`w`$ on $`d`$,
$$\frac{1}{\mathrm{\Omega }}z_{\mathrm{\Omega }\frac{1}{2}\theta 1}^2=\frac{1}{\mathrm{\Omega }+\frac{1}{2}}z_{\mathrm{\Omega }\theta }^2+3\theta +2\mathrm{\Omega }+2,$$
(169)
$$\frac{1}{\mathrm{\Delta }+1}w_{\mathrm{\Delta }+\frac{1}{2}\tau 1}^2=\frac{1}{\mathrm{\Delta }+\frac{1}{2}}w_{\mathrm{\Delta }\tau }^2+3\tau 2\mathrm{\Delta }.$$
(170)
These recurrence formulas can be initialized by using the values
$$z_{\frac{P}{2}\frac{P+2Q}{3}}=0,w_{\frac{P}{2}\frac{P+2Q}{3}}=0,$$
(171)
given by Eq. (166), and the solutions are
$$z_{\mathrm{\Omega }\theta }^2=(2\frac{PQ}{3}+\theta +1)(\frac{P+2Q}{3}\theta )(\frac{2P+Q}{3}+\theta +2)$$
(172)
$$w_{\mathrm{\Delta }\tau }^2=(2\frac{2P+Q}{3}\tau +1)(\frac{P+2Q}{3}\tau )(\frac{QP}{3}+\tau +1).$$
(173)
The values of $`z^2`$ and $`w^2`$ at any other point of the diagram of Fig. 3 can be obtained by using Eq. (164), such that
$$z_{iy}^2=(\frac{2P+Q}{3}i\frac{y}{2})(\frac{P+2Q}{3}+i+$$
(174)
$$\frac{y}{2}+2)(\frac{PQ}{3}+i+\frac{y}{2}+1)$$
$$w_{iy}^2=(\frac{QP}{3}+i\frac{y}{2})(\frac{P+2Q}{3}i+$$
(175)
$$\frac{y}{2}+1)(\frac{2P+Q}{3}+i\frac{y}{2}+1).$$
The coefficients $`\chi _{iy}`$ and $`\kappa _{iy}`$ are further determined by Eq. (165), and finally, Eq. (150), (151), complete the calculation by giving the matrix elements of the generators. In the following the action of the generators $`K_\pm `$ and $`L_\pm `$ will be expressed by using the notation of ref. , as
$$K_\pm |PQii_3y>=a_\pm (ii_3y)|PQi+\frac{1}{2}i_3\pm \frac{1}{2}y\pm 1>+$$
$$b_\pm (ii_3y)|PQi\frac{1}{2}i_3\pm \frac{1}{2}y\pm 1>,$$
$$L_\pm |PQii_3y>=c_\pm (ii_3y)|PQi+\frac{1}{2}i_3\pm \frac{1}{2}y\pm 1>+$$
$$d_\pm (ii_3y)|PQi\frac{1}{2}i_3\pm \frac{1}{2}y\pm 1>.$$
The coefficients $`a_\pm `$, $`b_\pm `$, $`c_\pm `$, $`d_\pm `$, can be found by the comparison between these equations and Eq. (148), (149) (152), (153).
In the particular case $`Q=0`$ the diagram of Fig. 3 reduces to a line, given by Eq. (97), and the matrix elements are
$$a_+^P(i_3,y)=\sqrt{(\frac{P}{3}+\frac{y}{2}+i_3+1)(\frac{P}{3}y)},b_+^P=0,$$
(176)
$$c_+^P(i_3,y)=\sqrt{(\frac{P}{3}+\frac{y}{2}i_3+1)(\frac{P}{3}y)},d_+^P=0.$$
(177)
In the case $`P=0`$ the diagram of Fig. 3 reduces to the line of Eq. (102), and the matrix elements are
$$a_+^Q=0,b_+^Q(i_3,y)=\sqrt{(\frac{Q}{3}\frac{y}{2}i_3)(\frac{Q}{3}+y+1)},$$
(178)
$$c_+^Q=0,d_+^Q(i_3,y)=\sqrt{(\frac{Q}{3}\frac{y}{2}+i_3)(\frac{Q}{3}+y+1)}.$$
(179)
IX. The Clebsch-Gordan coefficients
The Clebsch-Gordan series and the explicit form of the basis tensors in the irrep spaces appearing in the decomposition of the direct product $`V(P_1,Q_1)V(P_2,Q_2)`$ have been presented in sect. VI. The expressions obtained for these tensors become cumbersome when the dimension of the factor spaces increases, and therefore, the tensor representation is useful only in simple cases. However, the canonical basis of the spaces $`V^\gamma (P_k,Q_k)`$ in Eq. (104) can be constructed by using the theorems I-III presented in sect. IV. According to these theorems, the basis can be constructed by following a three-step procedure:
1. Find the subspaces $`W(P_k,Q_k)`$ generated by the vectors given in Eq. (122) which satisfy Eq. (48),(49), namely
$$I_+|P_kQ_k\underset{¯}{i}\underset{¯}{i}\underset{¯}{y}>_\gamma =0,$$
(180)
$$K_+|P_kQ_k\underset{¯}{i}\underset{¯}{i}\underset{¯}{y}>_\gamma =0,L_+|P_kQ_k\underset{¯}{i}\underset{¯}{i}\underset{¯}{y}>_\gamma =0,$$
(181)
and
$$I_3|P_kQ_k\underset{¯}{i}\underset{¯}{i}\underset{¯}{y}>_\gamma =\underset{¯}{i}|P_kQ_k\underset{¯}{i}\underset{¯}{i}\underset{¯}{y}>_\gamma ,$$
(182)
$$Y|P_kQ_k\underset{¯}{i}\underset{¯}{i}\underset{¯}{y}>_\gamma =\underset{¯}{y}|P_kQ_k\underset{¯}{i}\underset{¯}{i}\underset{¯}{y}>_\gamma ,$$
where $`\gamma `$ takes values from 1 to $`m_k=dimW(P_k,Q_k)`$, $`\underset{¯}{i}=P_k/2`$, $`\underset{¯}{y}=(P_k+2Q_k)/3`$, and
$$A_k^i=A_k^{(1)i}I^{(2)}+I^{(1)}A_k^{(2)i},i,k=1,2,3.$$
(183)
2. Chose an orthonormal basis in each space $`W(P_k,Q_k)`$, such that
$${}_{\gamma }{}^{}<P_kQ_k\underset{¯}{i}\underset{¯}{i}\underset{¯}{y}|P_kQ_k\underset{¯}{i}\underset{¯}{i}\underset{¯}{y}>_\gamma ^{}=\delta _{\gamma \gamma ^{}},\gamma ,\gamma ^{}=1,m_k.$$
(184)
3. Generate the basis of the subspaces $`V^\gamma (P_k,Q_k)`$ by applying the operators $`I_{}`$, $`K_{}`$, $`L_{}`$ on the highest weight vectors $`|P_kQ_k\underset{¯}{i}\underset{¯}{i}\underset{¯}{y}>_\gamma `$.
If the tensor product is simply reducible, ( $`m_k`$ is 0 or 1), the decomposition is obtained simply by the transition to a basis of $`V(P_1,Q_1)V(P_2,Q_2)`$ in which the Casimir operators constructed with the generators of Eq. (183) are diagonal. Because $`F`$ and $`G`$ are Hermitian, the subspaces corresponding to different eigenvalues $`f`$, $`g`$, and implicitly with different indices $`(P,Q)`$, are orthogonal and irreducible. The matrix of the linear transformation from the eigenvectors of the set of operators
$$(a):\{F_1,F_2,G_1,G_2,I^2,I_1^2,I_2^2,I_3,Y_1,Y_2\}$$
to the eigenvectors of the set
$$(b):\{F,G,F_1,F_2,G_1,G_2,I^2,I_3,Y\}$$
is the matrix of the isoscalar factors, and it can be obtained by finding the common eigenvectors of the set $`(b)`$.
The multiplicities appear when the space obtained by the direct product has a symmetry higher than $`SU(3)`$. This means that it exists a Hermitian operator $`X`$ which commutes with all the operators of the set $`(b)`$, and is independent of them, which should be added to obtain the complete set of compatible observables. The expression given by Moshinsky for this operator is
$$X=\frac{1}{2}\underset{i,j,k=1}{\overset{3}{}}(A_k^{(1)i}A_i^{(1)j}A_j^{(2)k}+A_k^{(1)i}A_j^{(1)k}A_i^{(2)j}).$$
(185)
The subspaces generated by the eigenvectors of the set (b) of operators can be further decomposed by diagonalizing the operator $`X`$. This method is cumbersome, and in practice it is not used, but is equivalent to the three-step procedure presented above. Thus, every vector satisfying the system of Eq. (180)-(182), is also an eigenvector of the operators $`F`$ and $`G`$. This follows by using the commutation relations of Eq. (23) to rearrange the operators appearing in Eq. (54) such that $`I_+`$, $`K_+`$ and $`L_+`$ are placed to the right of $`I_{}`$, $`K_{}`$, $`L_{}`$. When the resulting expression is applied on the vectors which satisfy the Eq. (180)-(182), the only non-vanishing terms will be those dependent only on $`I_3`$ and $`Y`$, which lead to multiplication by a constant. Because $`F`$ and $`G`$ are Hermitian, the solutions of Eq. (180)-(182), denoted by $`|s>`$, which correspond to different sets $`(P,Q)`$, will be orthogonal. The remaining eigenvectors are generated in the step 3. They will be eigenvectors of $`F`$ and $`G`$ with the same values as $`|s>`$, because $`F`$ and $`G`$ commute with all $`A_k^i`$. This shows that the solutions $`|s>`$ corresponding to different values $`P`$ and $`Q`$ are highest weight vectors in orthogonal irreducible subspaces.
When there are multiplicities, there are several independent solutions, $`\{|s>_\gamma ;\gamma =1,dim(W_s)\}`$, with the same indices $`(P,Q)`$. In this case, one of the vectors $`|s>`$ is arbitrarily chosen, and then the other can be obtained by the Gram-Schmidt procedure, such that are orthonormal. This leads to a decomposition of the product space in a direct sum of irreducible orthogonal subspaces, and to an unitary matrix of isoscalar factors. The vectors of the canonical basis in these subspaces are linear combinations of the form of Eq. (126)
$$|sii_3y>=\underset{\mu ,m,j,k}{}\alpha _{\mu jk}^{iyss_1s_2}C_{mi_3mi_3}^{jki}|s_1jm\mu >|s_2ki_3my\mu >.$$
(186)
For given values of $`s_1`$ and $`s_2`$, $`s`$ in Eq. (186) can take all the values appearing in the Clebsch-Gordan series. Fixing a certain value of $`s`$ means to specify the allowed intervals for the variation of $`i`$ and $`y`$. These intervals have been determined in sect. VII, and the result is presented in the diagram of Fig. 3.
The possible values of the summation indices $`\{j,\mu \}`$ and $`\{k,y\mu \}`$ can be represented in similar diagrams, denoted $`A`$ and $`B`$, corresponding to the carrier spaces of the irreps $`s_1`$, respectively $`s_2`$. In the orthogonal frame of axes $`(k,\mu ,j)`$, the possible values of these indices are represented by the points $`\pi `$ for which the projection in the plane $`(j,\mu )`$ is a node of the lattice $`A`$, and in the plane $`(k,\mu )`$ a node of the lattice $`B`$. This condition is expressed by the inequalities
$$\frac{2P_1+Q_1}{3}\mu \frac{P_1+2Q_1}{3}$$
(187)
$$\frac{2P_2+Q_2}{3}y\mu \frac{P_2+2Q_2}{3},$$
(188)
or, in compact form,
$$\mathrm{max}(\frac{2P_1+Q_1}{3},y\frac{P_2+2Q_2}{3})$$
(189)
$$\mu \mathrm{min}(\frac{P_1+2Q_1}{3},y+\frac{2P_2+Q_2}{3}).$$
However, not all the points $`(k,\mu ,j)`$ which satisfy this inequality appear as summation indices in Eq. (186), but only those for which
$$k+ji,|kj|i.$$
(190)
Geometrically, this means that not all the points selected by Eq. (189), represented by the 3D lattice of Fig. 6, appear as indices of the isoscalar factors, but only those which satisfy also Eq. (190). This condition selects the points which are inside and on the border of the space bounded by three planes which are parallel to the $`\mu `$-axis, and cross the $`k`$ and $`j`$ axes at the points with coordinates $`(k,j)=`$ $`(i,0)`$, $`(0,i)`$, and $`(0,i)`$ (Fig. 7). The final lattice obtained by this selection procedure is complicated in general, and it changes at the variation of $`i`$ and $`y`$.
With these considerations, we can proceed by following the three-step procedure presented at the beginning of this section.
1. According to Eq. (186), the highest weight vectors have the expression
$$|s\underset{¯}{i}\underset{¯}{i}\underset{¯}{y}>=\underset{\mu ,m,j,k}{}\alpha _{\mu jk}^sC_{m\underset{¯}{i}m\underset{¯}{i}}^{jk\underset{¯}{i}}|s_1jm\mu >|s_2k\underset{¯}{i}my\mu >,$$
(191)
which satisfies Eq. (180) and (182) ($`\alpha _{\mu jk}^s\alpha _{\mu jk}^{\underset{¯}{i}\underset{¯}{y}ss_1s_2}`$). By using Eq. (181) with the notations of Eq. (150), one obtains the following recurrence relations
$$\sqrt{\frac{\underset{¯}{i}+jk}{2j}}\chi _{j\frac{1}{2}\mu 1}^{s_1}\alpha _{\mu 1j\frac{1}{2}k+\frac{1}{2}}^s$$
(192)
$$\sqrt{\frac{(j+k\underset{¯}{i}+1)(\underset{¯}{i}+kj+1)}{(2j+1)(\underset{¯}{i}+j+k+2)}}\kappa _{j+\frac{1}{2}\mu 1}^{s_1}\alpha _{\mu 1j+\frac{1}{2}k+\frac{1}{2}}^s$$
$$+\sqrt{\frac{\underset{¯}{i}j+k+1}{2k+1}}\chi _{k\underset{¯}{y}\mu }^{s_2}\alpha _{\mu jk}^s$$
$$+\sqrt{\frac{(j+k\underset{¯}{i}+1)(\underset{¯}{i}k+j)}{2(k+1)(\underset{¯}{i}+j+k+2)}}\kappa _{k+1\underset{¯}{y}\mu }^{s_2}\alpha _{\mu jk+1}^s=0$$
and
$$\sqrt{\frac{\underset{¯}{i}+jk}{(2j)^2}}\chi _{j\frac{1}{2}\mu 1}^{s_1}\alpha _{\mu 1j\frac{1}{2}k+\frac{1}{2}}^s$$
(193)
$$+\sqrt{\frac{(\underset{¯}{i}+kj+1)(j+k\underset{¯}{i}+1)}{(2j+1)(\underset{¯}{i}+j+k+2)}}\kappa _{j+\frac{1}{2}\mu 1}^{s_2}\alpha _{\mu 1j+\frac{1}{2}k+\frac{1}{2}}^s$$
$$\sqrt{\frac{(\underset{¯}{i}j+k+1)(j+k\underset{¯}{i})^2}{2j(2k+1)}}\chi _{ky\mu }^{s_2}\alpha _{\mu jk}^s$$
$$+\sqrt{\frac{(\underset{¯}{i}+kj+2)^2(\underset{¯}{i}k+j)(j+k\underset{¯}{i}+1}{4j(k+1)(\underset{¯}{i}+j+k+2)}}\kappa _{k+1\underset{¯}{y}\mu }^{s_2}\alpha _{\mu jk+1}^s=0.$$
These relationships contain the factors $`\alpha `$ at the points $`A(k+1,\mu ,j)`$, $`B(k,\mu ,j)`$, $`C(k+1/2,\mu 1,j+1/2)`$, and $`D(k+1/2,\mu 1,j1/2)`$. By eliminating $`\alpha _C\alpha _{\mu 1j+\frac{1}{2}k+\frac{1}{2}}^s`$ between them one obtains a relationship between $`\alpha _A`$, $`\alpha _B`$ and $`\alpha _D`$,
$$\frac{2j+1}{\sqrt{2j}}\chi _{j\frac{1}{2}\mu 1}^{s_1}\alpha _{\mu 1j\frac{1}{2}k+\frac{1}{2}}^s$$
(194)
$$+\sqrt{\frac{(\underset{¯}{i}+jk)(\underset{¯}{i}+kj+1)}{2k+1}}\chi _{k\underset{¯}{y}\mu }^{s_2}\alpha _{\mu jk}^s$$
$$+\sqrt{\frac{(\underset{¯}{i}+j+k+2)(j+k\underset{¯}{i}+1)}{2(k+1)}}\kappa _{k+1\underset{¯}{y}\mu }^{s_2}\alpha _{\mu jk+1}^s=0,$$
and by eliminating $`\alpha _D\alpha _{\mu 1j1/2k+1/2}^s`$, a relationship between $`\alpha _A`$, $`\alpha _B`$ and $`\alpha _C`$,
$$\sqrt{\frac{2j+1}{\underset{¯}{i}+j+k+2}}\kappa _{j+\frac{1}{2}\mu 1}^{s_1}\alpha _{\mu 1j+\frac{1}{2}k+\frac{1}{2}}^s\sqrt{\frac{j+k\underset{¯}{i}+1}{2k+1}}\chi _{k\underset{¯}{y}\mu }^{s_2}\alpha _{\mu jk}^s$$
(195)
$$+\sqrt{\frac{(\underset{¯}{i}+jk)(\underset{¯}{i}+kj+1)}{2(k+1)(\underset{¯}{i}+j+k+2)}}\kappa _{k+1\underset{¯}{y}\mu }^{s_2}\alpha _{\mu jk+1}^{s_2}=0.$$
These equations, combined with the inequalities
$$\mathrm{max}(\frac{2P_1+Q_1}{3},\frac{PP_2+2(QQ_2)}{3})$$
(196)
$$\mu \mathrm{min}(\frac{P_1+2Q_1}{3},\frac{P+Q_2+2(P_2+Q)}{3}),$$
and
$$|jk|\frac{P}{2}j+k$$
(197)
determine the isoscalar factors for the highest weight vectors. The solution of these recurrence relationships is unique if all $`\alpha _{\mu jk}^s`$ can be expressed in terms of only one of them. In this case, the absolute value of this reference factor can be found by using the normalization relation
$$\underset{\mu jk}{}|\alpha _{\mu jk}^s|^2=1,$$
(198)
while its sign is given by the convention used in : among the isoscalar factors of maximum $`j`$, are chosen as positive those of maximum $`k`$.
The occurrence of the multiplicities is determined by the structure of the recurrence relationships of Eq. (194) and (195), with the constraints of Eq. (196), (197). Let $`\mu _{min}`$ and $`\mu _{max}`$ be the effective extreme values of $`\mu `$, different in general from the limits appearing in Eq. (196). One can see that the recurrence relationships determine completely the factors $`\alpha `$ only in two cases:
I. The lattice of Fig. 6 attached to the highest weight $`(\underset{¯}{i},y)`$ is bounded at $`\mu _{min}`$ by a segment parallel to the axis $`k`$, and at $`\mu _{max}`$ by an arbitrary figure.
II. The lattice of Fig. 6 is bounded at $`\mu _{min}`$ by an arbitrary figure, and at $`\mu _{max}`$ by a segment parallel to the axis $`j`$. In any other case the multiplicity will be greater than 1, equal to the minimum number of factors $`\alpha `$ required to initiate the recurrence relationships.
As an example, in Fig. 6 is represented the diagram associated to the highest weight in the case $`P=P_1=P_2`$, $`Q=Q_1=Q_2`$. It is easy to see that the multiplicity of the representation $`D(P,Q)`$ in the decomposition of the direct product is equal to the number of nodes of the lattice $`A`$ placed on the $`j`$ \- axis. If $`(PQ)/3`$ is an integer, this number represents the multiplicity of the weight $`i_3=0`$, $`y=0`$ from the space $`V(P,Q)`$, which is $`1+\mathrm{min}(P,Q)`$. By this, one recovers the result stated in ref. .
2. The highest weight vector $`|s>|PQ\underset{¯}{i}\underset{¯}{i}y>`$ given by Eq. (194) - (197) depends in general on free parameters which cannot be specified by the $`SU(3)`$ symmetry alone. In principle, it is possible to chose these parameters such that to specify an orthonormal basis in $`W(P,Q)`$. However, the result has a physical meaning only if there is a Hermitian operator associated to an observable, commuting with all the operators of the set (b), which has a diagonal matrix in this orthonormal basis.
3. Every point of the lattice of Fig. 3 is associated by Eq. (186) to a row in the matrix of the isoscalar factors. The calculation of the factors $`\alpha _{\mu jk}^s`$ along the row $`(\underset{¯}{i},\underset{¯}{y})`$ was presented at step 1. The factors $`\alpha _{\mu jk}^{iys}`$ from other rows can be related to $`\alpha _{\mu jk}^s`$ by an additional set of recurrence formulas.
Acting by $`L_{}`$ on both sides of Eq. (186) it is possible to find a recurrence relation between the factors placed on lines parallel to $`d`$,
$$\sqrt{\frac{2i+2}{i+j+k+2}}\kappa _{i+\frac{1}{2}y1}^s\alpha _{\mu jk+\frac{1}{2}}^{i+\frac{1}{2}y1s}=\sqrt{\frac{i+jk}{2j}}\kappa _{j\mu }^{s_1}\alpha _{\mu +1j\frac{1}{2}k+\frac{1}{2}}^{iys}$$
(199)
$$+\sqrt{\frac{(i+kj+1)(j+ki+1)}{(2j+1)(i+j+k+2)}}\chi _{j\mu }^{s_1}\alpha _{\mu +1j+\frac{1}{2}k+\frac{1}{2}}^{iys}$$
$$+\sqrt{\frac{i+kj+1}{2k+1}}\kappa _{k+\frac{1}{2}y\mu 1}^{s_1}\alpha _{\mu jk}^{iys}$$
$$\sqrt{\frac{(ik+j)(j+ki+1)}{2(k+1)(i+j+k+2)}}\chi _{k+\frac{1}{2}y\mu 1}^{s_2}\alpha _{\mu jk+1}^{iys}.$$
By using this formula it is possible to calculate successively the isoscalar factors $`\alpha _{\mu jk}^{iys}`$ for any $`(i,y)`$, if the factors along the lines parallel to $`u`$ or $`\nu `$ are known.
By acting with $`K_{}`$ on both sides of Eq. (186) one arrives at
$$\sqrt{\frac{(i+kj+1)(j+ki+1)(2i+2)}{(i+j+k+2)(2i+1)^2}}\kappa _{i+\frac{1}{2}y1}^s\alpha _{\mu jk+\frac{1}{2}}^{i+\frac{1}{2}y1s}$$
(200)
$$+\sqrt{\frac{i+jk}{2i+1}}\chi _{i\frac{1}{2}y1}^s\alpha _{\mu jk+\frac{1}{2}}^{i\frac{1}{2}y1s}$$
$$=\sqrt{\frac{2j+1}{i+j+k+2}}\chi _{j\mu }^{s_1}\alpha _{\mu +1j+\frac{1}{2}k+\frac{1}{2}}^{iys}+\sqrt{\frac{j+ki+1}{2k+1}}\kappa _{k+\frac{1}{2}y\mu 1}^{s_2}\alpha _{\mu jk}^{iys}$$
$$+\sqrt{\frac{(i+j+k+1)(i+jk)}{2(k+1)(i+j+k+2)}}\chi _{k+\frac{1}{2}y\mu 1}^{s_2}\alpha _{\mu jk+1}^{iys}.$$
By eliminating the factor $`\alpha _{\mu jk+\frac{1}{2}}^{i+\frac{1}{2}y1s}`$ between Eq. (199) and (200) one obtains a recurrence relation on lines parallel to $`u`$,
$$\chi _{i\frac{1}{2}y1}^s\alpha _{\mu jk+\frac{1}{2}}^{i\frac{1}{2}y1s}=\sqrt{\frac{(ij+k+1)(j+ki+1)}{2j(2i+1)}}\chi _{j\mu }^{s_1}\alpha _{\mu +1j\frac{1}{2}k+\frac{1}{2}}^{iys}$$
(201)
$$+\sqrt{\frac{(i+j+k+2)(i+jk)}{(2i+1)(2j+1)}}\chi _{j\mu }^{s_1}\alpha _{\mu +1j+\frac{1}{2}k+\frac{1}{2}}^{iys}$$
$$+\sqrt{\frac{(i+jk)(j+ki+1)}{(2i+1)(2k+1)}}\kappa _{k+\frac{1}{2}y\mu 1}^{s_2}\alpha _{\mu jk}^{iys}$$
$$+\sqrt{\frac{(i+j+k+2)(ij+k+1)}{(2i+1)(2k+2)}}\chi _{k+\frac{1}{2}y\mu 1}^{s_2}\alpha _{\mu jk+1}^{iys}.$$
The equations (199) and (201) determine completely the isoscalar factors $`\alpha _{\mu jk}^{iys}`$ as a function of $`\alpha _{\mu jk}^{\underset{¯}{i}\underset{¯}{y}s}\alpha _{\mu jk}^s`$ (the indices $`s_1`$, $`s_2`$ of $`\alpha `$ appearing in Eq. (186) have been omitted). The range of the indices $`\mu `$, $`j`$, $`k`$, as well as the shape of the lattice associated to each row in the axis frame $`(k,\mu ,j)`$, are given by Eq. (189) and (190). In terms of this lattice, Eq. (199) and (201) can be pictured geometrically by the structure represented in Fig. 8, as expressing an unknown factor $`\alpha `$ at the point $`P:(k+1/2,\mu ,j)`$ as a function of known factors at the points 1,2,3,4.
In certain particular situations it is convenient to use also the recurrence formula obtained by acting on both sides of Eq. (186) with the operator $`F`$, constructed with the generators of Eq. (183). This formula is
$$\rho _{\mu jk}^{s_1s_2}\alpha _{\mu jk}^{siy}=\lambda _{\mu jk}^{s_1s_2}\alpha _{\mu 1j\frac{1}{2}k\frac{1}{2}}^{iys}+\lambda _{\mu +1j+\frac{1}{2}k+\frac{1}{2}}^{s_1s_2}\alpha _{\mu +1j+\frac{1}{2}k+\frac{1}{2}}^{iys}$$
(202)
$$\tau _{\mu jk}^{s_1s_2}\alpha _{\mu 1j+\frac{1}{2}k+\frac{1}{2}}^{iys}\tau _{\mu +1j\frac{1}{2}k\frac{1}{2}}^{s_1s_2}\alpha _{\mu +1j\frac{1}{2}k\frac{1}{2}}^{iys}$$
$$+\xi _{\mu jk}^{s_1s_2}\alpha _{\mu 1j\frac{1}{2}k+\frac{1}{2}}^{iys}+\xi _{\mu +1j+\frac{1}{2}k\frac{1}{2}}^{s_1s_2}\alpha _{\mu +1j+\frac{1}{2}k\frac{1}{2}}^{iys}$$
$$+\omega _{\mu jk}^{s_1s_2}\alpha _{\mu 1j+\frac{1}{2}k\frac{1}{2}}^{iys}+\omega _{\mu +1j\frac{1}{2}k+\frac{1}{2}}^{s_1s_2}\alpha _{\mu +1j\frac{1}{2}k+\frac{1}{2}}^{iys}$$
where
$$\rho _{\mu jk}^{s_1s_2}=f_sf_{s_1}f_{s_2}\frac{3}{2}(y\mu )2j(ij)+(i+kj+1)(j+ki)$$
$$\lambda _{\mu jk}^{s_1s_2}=\sqrt{\frac{(i+j+k+1)(j+ki)}{4jk}}\chi _{j\frac{1}{2}\mu 1}^{s_1}\kappa _{ky\mu }^{s_2}$$
$$\tau _{\mu jk}^{s_1s_2}=\sqrt{\frac{(i+j+k+2)(j+ki+1)}{(2j+1)(2k+1)}}\kappa _{j+\frac{1}{2}\mu 1}^{s_1}\chi _{ky\mu }^{s_2}$$
$$\xi _{\mu jk}^{s_1s_2}=\sqrt{\frac{(i+jk)(i+kj+1)}{2j(2k+1)}}\chi _{j\frac{1}{2}\mu 1}^{s_1}\chi _{ky\mu }^{s_2}$$
$$\omega _{\mu jk}^{s_1s_2}=\sqrt{\frac{(ij+k)(i+jk+1)}{2k(2j+1)}}\kappa _{j+\frac{1}{2}\mu 1}^{s_1}\kappa _{ky\mu }^{s_2}$$
and $`f_s`$ represents the eigenvalue of the operator $`F`$ corresponding to the state with $`s=(P,Q)`$, given by Eq. (55). This recurrence formula acts within the frame $`(k,\mu ,j)`$ associated to a single line $`(i,y)`$, and relates the isoscalar factors from the points represented in Fig. 9.
X. The canonical base of the space $`V(P,Q)`$
The canonical basis of the space $`V(P,Q)`$ was constructed in sect. VII by direct product, and has the expression of Eq. (128). The unknown isoscalar factors $`\alpha _\mu ^{iy(PQ)}\alpha _{\mu \frac{P}{3}+\frac{\mu }{2}\frac{Q}{3}\frac{y\mu }{2}}^{iy(PQ)(P0)(0Q)}`$ appearing in this formula can be determined by using Eq. (202). In this case $`s_1=(P,0)`$, $`s_2=(0,Q)`$, and the 3D lattice of Fig. 6 reduces to a one-dimensional lattice. Therefore, $`\kappa ^{s_1}=0`$, $`\chi ^{s_2}=0`$, and Eq. (202) takes the simple form
$$\rho _\mu ^{PQ}\alpha _\mu ^{iy(PQ)}=\lambda _\mu ^{PQ}\alpha _{\mu 1}^{iy(PQ)}+\lambda _{\mu +1}^{PQ}\alpha _{\mu +1}^{iy(PQ)},$$
(203)
where
$$\rho _\mu ^{PQ}=(\frac{P+Q}{3}+i\frac{y}{2}+\mu +1)(\frac{P+Q}{3}i\frac{y}{2}+\mu )$$
(204)
$$+(\frac{P}{3}\mu )(\frac{Q}{3}+y\mu ),$$
$$\lambda _\mu ^{PQ}=\sqrt{\frac{P+Q}{3}+i\frac{y}{2}+\mu +1}$$
(205)
$$\times \sqrt{(\frac{P+Q}{3}i\frac{y}{2}+\mu )(\frac{P}{3}\mu +1)(\frac{Q}{3}+y\mu +1)}.$$
By using the notations
$$a_\mu =\frac{P+Q}{3}+i\frac{y}{2}+\mu +1,$$
(206)
$$b_\mu =\frac{P+Q}{3}i\frac{y}{2}+\mu ,$$
and
$$c_\mu =\frac{P}{3}\mu ,d_\mu =\frac{Q}{3}+y\mu ,$$
Eq. (203) takes the form
$$a_\mu b_\mu \alpha _\mu ^{iy(PQ)}+c_\mu d_\mu \alpha _\mu ^{iy(PQ)}=\sqrt{a_\mu b_\mu c_{\mu 1}d_{\mu 1}}\alpha _{\mu 1}^{iy(PQ)}$$
(207)
$$+\sqrt{a_{\mu +1}b_{\mu +1}c_\mu d_\mu }\alpha _{\mu +1}^{iy(PQ)}.$$
This equality is satisfyed if
$$\sqrt{a_\mu b_\mu }\alpha _\mu ^{iy(PQ)}=\sqrt{c_{\mu 1}d_{\mu 1}}\alpha _{\mu 1}^{iy(PQ)},$$
(208)
which is a recurrence relation with the solution
$$\alpha _\mu ^{iy(PQ)}=\alpha _{\stackrel{~}{\mu }}^{iy(PQ)}$$
(209)
$$\times \sqrt{\frac{(\frac{P+Q}{3}+i\frac{y}{2}+1+\stackrel{~}{\mu })!(\frac{P+Q}{3}i\frac{y}{2}+\stackrel{~}{\mu })!(\frac{P}{3}\stackrel{~}{\mu })!(\frac{Q}{3}+y\stackrel{~}{\mu })!}{(\frac{P+Q}{3}+i\frac{y}{2}+1+\mu )!(\frac{P+Q}{3}i\frac{y}{2}+\mu )!(\frac{P}{3}\mu )!(\frac{Q}{3}+y\mu )!}}.$$
The maximum value $`\stackrel{~}{\mu }`$ of $`\mu `$ is determined only by Eq. (189), and is equal to $`(\mu )_{max}`$ given by Eq. (136). The minimum value of $`\mu `$ is
$$(\mu )_{min}=\mathrm{max}(i+\frac{y}{2}\frac{P+Q}{3},y2\frac{Q}{3},2\frac{P}{3}).$$
On the rows through the nodes of the line $`d`$ one obtains $`(\mu )_{max}=(\mu )_{min}=P/3`$, and the sum over $`\mu `$ in Eq. (128) reduces to a single term,
$$|PQ\mathrm{\Delta }i_3\tau >=\alpha _{\frac{P}{3}}^{\mathrm{\Delta }\tau (PQ)}\underset{m}{}C_{mi_3mi_3}^{\frac{P}{2}\mathrm{\Delta }\frac{P}{2}\mathrm{\Delta }}\xi _{Pm}^{\frac{P}{3}}\eta _{Qi_3m}^{\tau \frac{P}{3}}.$$
(210)
The unitarity condition implies $`|\alpha _{\frac{P}{3}}^{\mathrm{\Delta }\tau (PQ)}|=1`$, and the action of $`L_+`$ on the vector $`|PQ\mathrm{\Delta }\mathrm{\Delta }\tau >`$ shows that all these factors have the same sign, and can be considered equal to 1. The factors $`\alpha _{\frac{P}{3}}^{iy(PQ)}`$ can be calculated by using Eq. (201), which takes the form
$$\chi _{i\frac{1}{2}y1}^{PQ}\alpha _\mu ^{i\frac{1}{2}y1(PQ)}=\sqrt{\frac{(i+j+k+2)(i+jk)}{(2i+1)(2j+1)}}\chi _{j\mu }^P\alpha _{\mu +1}^{iy(PQ)}$$
(211)
$$+\sqrt{\frac{(i+jk)(j+k+i1)}{(2i+1)(2k+1)}}\kappa _{k+\frac{1}{2}y\mu 1}^Q\alpha _\mu ^{iy(PQ)}.$$
In the case $`y\frac{PQ}{3}`$, for $`\mu =P/3`$ this equality becomes
$$\sqrt{\frac{P+2Q}{3}+i+\frac{y}{2}+1}\alpha _{\frac{P}{3}}^{i\frac{1}{2}y1(PQ)}=\sqrt{y\frac{PQ}{3}}\alpha _{\frac{P}{3}}^{iy(PQ)}.$$
(212)
Finally, by solving this recurrence relation with the initial condition $`\alpha _{P/3}^{\mathrm{\Delta }\tau (PQ)}=1`$, one obtains
$$\alpha _{\frac{P}{3}}^{iy(PQ)}=\sqrt{\frac{(\frac{P+2Q}{3}+i+\frac{y}{2}+1)!(\frac{P+2Q}{3}i+\frac{y}{2})!}{(P+Q+1)!(\frac{QP}{3}+y)!}}.$$
(213)
In the case $`y\frac{PQ}{3}`$, by using Eq. (208) and (211) one can find similarly
$$\alpha _{y+\frac{Q}{3}}^{iy(PQ)}=\sqrt{\frac{(\frac{2P+Q}{3}i\frac{y}{2})!(\frac{2P+Q}{3}+i\frac{y}{2}+1)!}{(P+Q+1)!(\frac{PQ}{3}y)!}},$$
(214)
such that in both cases, Eq. (209) takes the form
$$\alpha _\mu ^{iy(PQ)}=\sqrt{\frac{(\frac{P+2Q}{3}+i+\frac{y}{2}+1)!}{(P+Q+1)!}}$$
(215)
$$\times \sqrt{\frac{(\frac{P+2Q}{3}i+\frac{y}{2})!(\frac{2P+Q}{3}i\frac{y}{2})!(\frac{2P+Q}{3}+i\frac{y}{2}+1)!}{(\frac{P+Q}{3}+i\frac{y}{2}+1+\mu )!(\frac{P+Q}{3}i\frac{y}{2}+\mu )!(\frac{P}{3}\mu )!(\frac{Q}{3}+y\mu )!}}.$$
One can check that the factors determined by this equation have the important symmetry property
$$\alpha _\mu ^{iy(PQ)}=\alpha _{\mu y}^{iy(QP)}.$$
(216)
This property will be used in the following to prove the equivalence between the representations $`D(P,Q)^{}`$ and $`D(Q,P)`$. By using Eq. (128) one obtains
$$|QPii_3y>=\underset{\mu m}{}\alpha _\mu ^{iy(QP)}C_{mi_3mi_3}^{jki}\xi _{Qm}^\mu \eta _{Pi_3m}^{y\mu }.$$
(217)
If the basis elements $`\xi _{Qm}^\mu `$ and $`\eta _{Pi_3m}^{y\mu }`$ appearing here are expressed in terms of $`(\eta _{Qm}^\mu )^{}`$, respectively $`(\xi _{Pi_3+m}^{y+\mu })^{}`$, by using Eq. (100), and then we introduce the new summation indices $`n=i_3+m`$ and $`\nu =y+\mu `$, then Eq. (217) becomes
$$|QPii_3y>=\underset{\mu n}{}\alpha _{\nu y}^{iy(QP)}$$
(218)
$$\times C_{ni_3ni_3}^{jki}(1)^{\frac{PQ}{3}+i_3+\frac{y}{2}}(\xi _{Pn}^\nu \eta _{Qi_3n}^{y\nu })^{}.$$
By using now Eq. (216), and the symmetry properties of the $`SU(2)`$ Clebsch-Gordan coefficients,
$$C_{ni_3ni_3}^{jki}=C_{ni_3ni_3}^{kji},$$
(219)
one arrives at
$$|QPii_3y>=(1)^{\frac{PQ}{3}+i_3+\frac{y}{2}}|PQii_3y>^{}.$$
(220)
The sign of this expression is given by the number $`t`$, called triality , defined as
$$t=(PQ)mod3,t=1,0,1$$
(221)
such that up to a negligible global phase factor, Eq. (220) becomes
$$|QPii_3y>=(1)^{\frac{t}{3}+i_3+\frac{y}{2}}|PQii_3y>^{}.$$
(222)
The $`SU(3)`$ representations which map the elements of the discrete subgroup $`Z_3`$ onto the unity matrix $`I`$ are faithful representations of the factor group $`SU(3)/Z_3`$. For these representations the integers $`m_1`$, $`m_2`$, and $`m_3`$ of Eq. (42) should satisfy the equality
$$(m_1+m_2+m_3)mod3=(PQ)mod3=0,$$
(223)
which means that both the electric charge $`i_3+y/2`$ which appears in Eq. (222), and the hypercharge $`y`$, are integers.
XI. The symmetric form of the canonical basis in the spaces $`V(P,0)`$ and $`V(0,Q)`$.
The vectors $`\xi _{Pi_3}^y`$ of the canonical basis of the space $`V(P,0)`$ can be expressed as a symmetrized product of the basis elements of $`V(3,0)`$, according to Eq. (98),
$$\xi _{Pi_3}^y=a^P(i_3,y)(x_1)^{i+i_3}(x_2)^{ii_3}(x_3)^{\frac{P}{3}y},i=\frac{P}{3}+\frac{y}{2}.$$
(224)
The unknown factors $`a^P(i_3,y)`$ can be calculated by recurrence. Acting on both sides of Eq. (224) with the operator $`I_+`$ one obtains
$$a^P(i_3,y)=\sqrt{\frac{i+i_3+1}{ii_3}}a^P(i_3+1,y),$$
(225)
which has the solution
$$a^P(i_3,y)=\sqrt{\frac{(2i)!}{(ii_3)!(i+i_3)!}}a_y^P,$$
(226)
where $`a_y^Pa^P(i,y)`$. By choosing in Eq. (224) $`i_3=i`$, and acting on both sides with $`K_+`$, one ontains
$$a_y^P=\sqrt{\frac{\frac{P}{3}y+1}{2\frac{P}{3}+y}}a_{y1}^P.$$
(227)
By recurrence, this gives
$$a_y^P=\sqrt{\frac{P!}{(\frac{2P}{3}+y)!(\frac{P}{3}y)!}}a_{\frac{P}{3}}^P=$$
(228)
$$\sqrt{\frac{P!}{(2i)!(\frac{P}{3}y)!}}a_{\frac{P}{3}}^P.$$
Then, by using Eq. (226) with $`a_{P/3}^P=1`$, one obtains
$$a^P(i_3,y)=\sqrt{\frac{P!}{(ii_3)!(i+i_3)!(\frac{P}{3}y)!}}.$$
(229)
The symmetric basis of the space $`V(0,Q)`$ is expressed by Eq. (103), where all the quantities are now determined.
XII. The isoscalar factors $`\alpha _{\mu jk}^{iy(P_1+P_2,0)(P_1,0)(P_2,0)}`$ and $`\alpha _{\mu jk}^{iy(0,Q_1+Q_2)(0,Q_1)(0,Q_2)}`$
If $`s_1=(P_1,0)`$, $`s_2=(P_2,0)`$, then $`\kappa ^{s_1}=\kappa ^{s_2}=0`$, and Eq. (202) takes the simple form
$$\rho _\mu ^{P_1P_2}\alpha _\mu ^{y(P_1+P_2,0)}=\xi _\mu ^{P_1P_2}\alpha _{\mu 1}^{y(P_1+P_2,0)}+\xi _{\mu +1}^{P_1P_2}\alpha _\mu ^{y(P_1+P_2,0)},$$
(230)
where $`\alpha _\mu ^{y(P_1+P_2,0)}\alpha _{\mu jk}^{iy(P_1+P_2,0)(P_1,0)(P_2,0)}`$,
$$j=\frac{P_1}{3}+\frac{\mu }{2},k=\frac{P_2}{3}+\frac{y\mu }{2},i=\frac{P_1+P_2}{3}+\frac{y}{2},$$
$$\rho _\mu ^{P_1P_2}=(\frac{2P_1}{3}+\mu )(\frac{P_2}{3}y+\mu )+(\frac{P_1}{3}\mu )(\frac{2P_2}{3}+y\mu ),$$
and
$$\xi _\mu ^{P_1P_2}=\sqrt{(\frac{P_1}{3}\mu +1)(\frac{P_2}{3}y+\mu )(\frac{2P_1}{3}+\mu )(\frac{2P_2}{3}+y\mu +1)}.$$
The extreme values of the index $`\mu `$ are given by Eq. (189), as
$$\mathrm{max}(2\frac{P_1}{3},y\frac{P_2}{3})\mu \mathrm{min}(\frac{P_1}{3},y+2\frac{P_2}{3}),$$
(231)
because in this case $`i=j+k`$.
Following the same steps as in sect. XI, it is possible to show that the solutions of Eq. (230) satisfy the relationship
$$\sqrt{(\frac{P_1}{3}\mu )(\frac{2P_2}{3}+y\mu )}\alpha _\mu ^{y(P_1+P_2,0)}=$$
(232)
$$=\sqrt{(\frac{2P_1}{3}+\mu +1)(\frac{P_2}{3}y+\mu +1)}\alpha _{\mu +1}^{y(P_1+P_2,0)},$$
which gives
$$\alpha _\mu ^{y(P_1+P_2,0)}=\alpha _{\stackrel{~}{\mu }}^{y(P_1+P_2,0)}\times $$
(233)
$$\sqrt{\frac{(\frac{P_1}{3}\stackrel{~}{\mu })!(\frac{2P_2}{3}+y\stackrel{~}{\mu })!(\frac{2P_1}{3}+\stackrel{~}{\mu })!(\frac{P_2}{3}y+\stackrel{~}{\mu })!}{(\frac{P_1}{3}\mu )!(\frac{2P_2}{3}+y\mu )!(\frac{2P_1}{3}+\mu )!(\frac{P_2}{3}y+\mu )!}}.$$
Here $`\stackrel{~}{\mu }=(\mu )_{max}=P_1/3`$ if $`y(P_12P_2)/3`$ and $`\stackrel{~}{\mu }=y+2P_2/3`$ if $`y(P_12P_2)/3`$. To calculate $`\alpha _\mu ^{y(P_1+P_2,0)}`$ one can use Eq. (201), which in this case takes the form
$$\sqrt{(\frac{P_1+P_2}{3}y+1)(2\frac{P_1+P_2}{3}+y)}\alpha _\mu ^{y1(P_1+P_2,0)}$$
(234)
$$=\sqrt{(\frac{2P_1}{3}+\mu +1)(\frac{P_1}{3}\mu )}\alpha _{\mu +1}^{y(P_1+P_2,0)}$$
$$+\sqrt{(\frac{2P_2}{3}+y\mu )(\frac{P_2}{3}y+\mu +1)}\alpha _\mu ^{y(P_1+P_2,0)}.$$
By using the initial condition $`\alpha _{\frac{P_1}{3}}^{\frac{P_1+P_2}{3}(P_1+P_2,0)}=1`$, one obtains
$$\alpha _{\frac{P_1}{3}}^{y(P_1+P_2,0)}=\sqrt{\frac{P_2!(2\frac{P_1+P_2}{3}+y)!}{(P_1+P_2)!(\frac{2P_2P_1}{3}+y)!}}$$
(235)
if
$$y\frac{P_12P_2}{3}$$
and
$$\alpha _{\frac{2P_2}{3}+y}^{y(P_1+P_2,0)}=\sqrt{\frac{P_1!(\frac{P_1+P_2}{3}y)!}{(P_1+P_2)!(\frac{P_12P_2}{3}y)!}}$$
(236)
if
$$y\frac{P_12P_2}{3}.$$
In both cases, the isoscalar factors given by Eq. (233) are
$$\alpha _\mu ^{\mu (P_1+P_2,0)}=$$
(237)
$$\sqrt{\frac{P_1!P_2!(\frac{P_1+P_2}{3}y)!(2\frac{P_1+P_2}{3}+y)!}{(P_1+P_2)!(\frac{P_1}{3}\mu )!(\frac{2P_2}{3}+y\mu )!(\frac{2P_1}{3}+\mu )!(\frac{P_2}{3}y+\mu )!}}.$$
By taking the complex conjugate of the equality
$$\xi _{Q_1+Q_2i}^y=\underset{\mu }{}\alpha _\mu ^{y(Q_1+Q_2,0)}\xi _{Q_1j}^\mu \xi _{Q_2ij}^{y\mu },$$
(238)
and making use of Eq. (100), one obtains
$$\eta _{Q_1+Q_2i}^y=\underset{\mu }{}\alpha _\mu ^{y(Q_1+Q_2,0)}\eta _{Q_1j}^\mu \eta _{Q_2(ij)}^{y\mu }.$$
(239)
This result allows to express the isoscalar factors $`\alpha _\mu ^{y(0,Q_1+Q_2)}`$ in terms of the ones given by Eq. (237), by identifying summation terms in Eq. (239) and (186) for the case $`s_1=(0,Q_1)`$, $`s_2=(0,Q_2)`$, $`s=(0,Q_1+Q_2)`$, $`i_3=i=j+k`$, such that
$$\alpha _\mu ^{y(0,Q_1+Q_2)}=\alpha _\mu ^{y(Q_1+Q_2,0)}.$$
(240)
XIII. The isoscalar factors $`\alpha _{\mu jk}^{00(0,0)(P,Q)(Q,P)}`$
These isoscalar factors can be obtained by using the conjugation relationship of Eq. (222) in the expression of the normalized $`SU(3)`$ scalar which can be constructed with the elements of the canonical base of the space $`V(P,Q)`$. This scalar is
$$S=\frac{1}{\sqrt{dimV(P,Q)}}\underset{yi_3jk}{}\delta _{jk}<PQji_3y|PQki_3y>=$$
(241)
$$\frac{1}{\sqrt{dimV(P,Q)}}\underset{yi_3jk}{}\delta _{jk}(1)^{(\frac{t}{3}+m+\frac{\mu }{2})}|PQjm\mu >|QPkm\mu >.$$
The $`SU(2)`$ Clebsch-Gordan coefficient in this case is
$$C_{mm0}^{jk0}=\delta _{jk}\frac{(1)^{jm}}{\sqrt{2j+1}},$$
(242)
and therefore the isoscalar factor is
$$\alpha _{\mu jj}^{00(0,0)(P,Q)(Q,P)}=\sqrt{\frac{2j+1}{dimV(P,Q)}}(1)^{(\frac{t}{3}+\frac{\mu }{2}+j)}.$$
(243)
The maximum value of $`j`$ is $`(P+Q)/2`$, and for this value the sign of the factor determined by Eq. (243) is $`(1)^{Q+t}`$. When $`Q+t`$ is odd, the factor is negative, violating the phase convention used in ref. . However, all these factors are defined up to a constant, which can be chosen such that those with a maximum $`j`$ are positive. The general expression in this case is
$$\alpha _{\mu jj}^{00(0,0)(P,Q)(Q,P)}=\sqrt{\frac{2(2j+1)}{(P+1)(Q+1)(P+Q+2)}}(1)^{(\frac{2t}{3}+\frac{\mu }{2}+j+Q)}.$$
(244)
XIV. The highest weight vectors in the space $`V(1,1)V(1,1)`$.
The Clebsch-Gordan series associated to this product has the form
$$V(1,1)V(1,1)=V(2,2)+V(3,0)+V(0,3)+$$
(245)
$$+V^1(1,1)+V^2(1,1)+V(0,0).$$
For the highest weight $`(\underset{¯}{i},\underset{¯}{y})=(P/2,(P+2Q)/3)`$ of each space $`V(P,Q)`$ appearing in this sum we can construct a diagram $`(k,\mu ,j)`$ similar to the one pictured in Fig. 6. This gives a geometrical representation of the range of the indices resulted from Eq. (196) and (197), making easier to apply Eq. (194), (195) to each term of the sum.
(A). $`V(2,2)`$; $`\underset{¯}{i}=1`$, $`\underset{¯}{y}=2`$, $`\mu =1`$. There is one non-vanishing factor, $`\alpha _{1\frac{1}{2}\frac{1}{2}}^{12(2,2)}`$, and the highest weight vector is
$$|22112>=|11\frac{1}{2}\frac{1}{2}1>|11\frac{1}{2}\frac{1}{2}1>.$$
(246)
(B). $`V(3,0)`$; $`\underset{¯}{i}=3/2`$, $`\underset{¯}{y}=1`$, $`\mu =0,1`$. There are two non-vanishing factors, $`\alpha _{01\frac{1}{2}}^{\frac{3}{2}1(3,0)}`$, $`\alpha _{1\frac{1}{2}1}^{\frac{3}{2}1(3,0)}`$, and in this case Eq. (195) becomes
$$\alpha _{01\frac{1}{2}}^{\frac{3}{2}1(3,0)}+\alpha _{1\frac{1}{2}1}^{\frac{3}{2}1(3,0)}=0.$$
(247)
This equation, combined with the unitarity condition
$$(\alpha _{01\frac{1}{2}}^{\frac{3}{2}1(3,0)})^2+(\alpha _{1\frac{1}{2}1}^{\frac{3}{2}1(3,0)})^2=1,$$
(248)
and the phase convention, determines the values
$$\alpha _{01\frac{1}{2}}^{\frac{3}{2}1(3,0)}=\frac{1}{\sqrt{2}},\alpha _{1\frac{1}{2}1}^{\frac{3}{2}1(3,0)}=\frac{1}{\sqrt{2}}$$
(249)
of the isoscalar factors, and the highest weight vector
$$|30\frac{3}{2}\frac{3}{2}1>=\frac{1}{\sqrt{2}}\underset{m}{}[C_{m\frac{3}{2}m\frac{3}{2}}^{1\frac{1}{2}\frac{3}{2}}|111m0>|11\frac{1}{2}\frac{3}{2}m1>$$
(250)
$$C_{m\frac{3}{2}m\frac{3}{2}}^{\frac{1}{2}1\frac{3}{2}}|11\frac{1}{2}m1>|111\frac{3}{2}m0>].$$
(C). $`V(0,3)`$; $`\underset{¯}{i}=0`$, $`\underset{¯}{y}=2`$, $`\mu =1`$. There is one non-vanishing factor, $`\alpha _{1\frac{1}{2}\frac{1}{2}}^{02(0,3)}`$, of modulus $`1`$, and the highest weight vector is
$$|03002>=\alpha _{1\frac{1}{2}\frac{1}{2}}^{02(0,3)}\underset{m}{}C_{mm0}^{\frac{1}{2}\frac{1}{2}0}|11\frac{1}{2}m1>|11\frac{1}{2}m1>.$$
(251)
The phase of this factor can be determined by noticing that from Eq. (250) we get
$$|30002>=\underset{m}{}C_{mm0}^{\frac{1}{2}\frac{1}{2}0}|11\frac{1}{2}m1>|11\frac{1}{2}m1>,$$
(252)
while from Eq. (222)
$$|03002>^{}=|30002>,$$
(253)
such that $`\alpha _{1\frac{1}{2}\frac{1}{2}}^{02(0,3)}=1`$.
(D). $`V^1(1,1)+V^2(1,1)`$; $`\underset{¯}{i}=1/2`$, $`\underset{¯}{y}=1`$, $`\mu =0,1`$. The general expression of the highest weight vector in this case is
$$|11\frac{1}{2}\frac{1}{2}1>=\underset{\mu jk}{}\alpha _{\mu jk}^{\frac{1}{2}1(11)}\times $$
(254)
$$\underset{m}{}C_{m\frac{1}{2}m\frac{1}{2}}^{jk\frac{1}{2}}|11jm\mu >|11k\frac{1}{2}m1\mu >,$$
where it is necessary to determine the factors
$$\alpha _{01\frac{1}{2}}^{\frac{1}{2}1(11)},\alpha _{00\frac{1}{2}}^{\frac{1}{2}1(11)},\alpha _{1\frac{1}{2}1}^{\frac{1}{2}1(11)},\alpha _{1\frac{1}{2}0}^{\frac{1}{2}1(11)}.$$
The recurrence equations (194), (195) become in this case
$$\alpha _{00\frac{1}{2}}^{\frac{1}{2}1(11)}=\frac{1}{2}(\alpha _{1\frac{1}{2}0}^{\frac{1}{2}1(11)}+\alpha _{1\frac{1}{2}1}^{\frac{1}{2}1(11)}),$$
(255)
$$\alpha _{01\frac{1}{2}}^{\frac{1}{2}1(11)}=\frac{1}{2}(3\alpha _{1\frac{1}{2}0}^{\frac{1}{2}1(11)}\alpha _{1\frac{1}{2}1}^{\frac{1}{2}1(11)}),$$
(256)
which should be solved together with the unitarity condition
$$\underset{\mu jk}{}|\alpha _{\mu jk}^{\frac{1}{2}1(11)}|^2=1.$$
(257)
However, by contrast to the previous situations, this system of equations does not has an unique solution.
A solution of the system is provided by the set used in ref. ,
$$\alpha _{00\frac{1}{2}}^{\frac{1}{2}1(11)}=\alpha _{01\frac{1}{2}}^{\frac{1}{2}1(11)}=\alpha _{1\frac{1}{2}0}^{\frac{1}{2}1(11)}=\alpha _{1\frac{1}{2}1}^{\frac{1}{2}1(11)}=\frac{1}{2}.$$
(258)
An independent set, $`\{\beta _{\mu jk}^{\frac{1}{2}1(11)}\}`$, can be determined by adding to the system of Eq. (255)-(257), the orthogonality condition
$$\underset{\mu jk}{}\alpha _{\mu jk}^{\frac{1}{2}1(11)}\beta _{\mu jk}^{\frac{1}{2}1(11)}=0.$$
(259)
The solution of the system obtained is
$$\beta _{00\frac{1}{2}}^{\frac{1}{2}1(11)}=\alpha _{1\frac{1}{2}0}^{\frac{1}{2}1(11)}=\frac{\sqrt{5}}{10},\beta _{01\frac{1}{2}}^{\frac{1}{2}1(11)}=\beta _{1\frac{1}{2}1}^{\frac{1}{2}1(11)}=3\frac{\sqrt{5}}{10}.$$
(260)
The factors $`\alpha `$ and $`\beta `$ given by Eq. (258) and (260) determine by Eq. (254) the highest weight vectors for the irrep spaces $`V^1(1,1)`$ and $`V^2(1,1)`$, respectively.
(E). $`V(0,0)`$; $`\underset{¯}{i}=0`$, $`\underset{¯}{y}=0`$, $`\mu =1,0,1`$. This space is one-dimensional, and the highest weight vector is the normalized $`SU(3)`$ scalar $`S`$ given by Eq. (241) for $`P=Q=1`$. Therefore, the isoscalar factors can be calculated by using Eq. (244), and have the values
$$\alpha _{011}^{00(00)}=\frac{\sqrt{6}}{4},\alpha _{1\frac{1}{2}\frac{1}{2}}^{00(00)}=\frac{1}{2},\alpha _{000}^{00(00)}=\frac{\sqrt{2}}{4},\alpha _{1\frac{1}{2}\frac{1}{2}}^{00(00)}=\frac{1}{2}.$$
(261)
The same result can be obtained by using the recurrence formulas.
Concluding remarks
The expressions of the isoscalar factors given in Eq. (215) and (237), derived here by using Eq. (202), are the same as in ref. . Also, the isoscalar factors determined in sect. XIV have the values given in the tables .
The symmetry properties of the $`SU(3)`$ Clebsch-Gordan coefficients can be found by using the symmetry properties of the recurrence relationships for the isoscalar factors. However, it is easier to derive them from the integral representation of the Clebsch-Gordan coefficients, given by the projection operators associated to the matrix elements of the irreducible representations , , , .
In general, for the group $`SU(n)`$, $`n>2`$, the canonical factorization $`SU(n)U(1)\times SU(n1)`$ can only solve the problem of the multiplicity for the weights, but not for the representations appearing in the Clebsch-Gordan series. Therefore, whatever method is used for reduction, it is not possible to obtain numerical values for all the Clebsch-Gordan coefficients, without making additional assumptions. Biedenharn and Louck have shown that by using the general property of canonical embedding for $`U(n)`$,
$$U(n^2)U(n)\times U(n)U(n)$$
(262)
it is possible to solve the problem of the multiplicities such that all Clebsch-Gordan coefficients are determined.
The recurrence equations for the $`SU(n)`$ isoscalar factors obtained by using the canonical factorization can be derived by acting on the weight vectors with the representation operators for the $`sl(n,C)`$ generators associated to the $`n1`$ simple roots. The action of the $`sl(n,C)`$ operators associated to the remaining positive roots leads to an additional set of
$$\frac{n(n1)}{2}(n1)=\frac{(n1)(n2)}{2}$$
(263)
equations, which are identically satisfied due to the properties of the $`SU(n1)`$ Clebsch-Gordan coefficients. Therefore, these relationships may provide in general only the dependence of the isoscalar factors on arbitrary parameters, because for $`n>3`$, $`SU(n1)`$ is not simply reducible. In the treatment of a specific quantum system, the use of a complete set of compatible observables might eliminate the ambiguities, leading to well-defined values of these parameters in each physical situation.
Acknowledgements. Thanks are due to Professor Gerry McKeon, for interest and the support provided towards the translation in English of this article.
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# Main Theorem
RATIONAL SOLUTIONS OF THE PAINLEVE’–VI EQUATION
M. Mazzocco
Mathematical Institute
24-29 St Giles, Oxford, UK.
Abstract. In this paper, we classify all values of the parameters $`\alpha `$, $`\beta `$, $`\gamma `$ and $`\delta `$ of the Painlevé VI equation such that there are rational solutions. We give a formula for them up to the birational canonical transformations and the symmetries of the Painlevé VI equation.
1. Introduction. In this paper, we study the general Painlevé sixth equation
$$\begin{array}{cc}\hfill y_{xx}=& \frac{1}{2}\left(\frac{1}{y}+\frac{1}{y1}+\frac{1}{yx}\right)y_x^2\left(\frac{1}{x}+\frac{1}{x1}+\frac{1}{yx}\right)y_x+\hfill \\ & +\frac{y(y1)(yx)}{x^2(x1)^2}\left[\alpha +\beta \frac{x}{y^2}+\gamma \frac{x1}{(y1)^2}+\delta \frac{x(x1)}{(yx)^2}\right],\hfill \end{array}$$
$`PVI`$
where $`x\text{C}\text{l}`$ and $`\alpha ,\beta ,\gamma ,\delta `$ are arbitrary complex parameters. The general solution $`y(x;c_1,c_2)`$ of the PVI equation satisfies the following two important properties (see \[Pain\], \[Gam\]):
1) Painlevé property: the solutions $`y(x;c_1,c_2)`$ may have complicated singularities (i.e. branch points, essential singularities etc.) only at the critical points $`0,1,\mathrm{}`$ of the equation (the so-called fixed singularities). All the other singularities, the position of which depend on the integration constants (the so-called movable singularities), are poles.
2) For generic values of the integration constants $`c_1,c_2`$ and of the parameters $`\alpha ,\beta ,\gamma ,\delta `$, the solution $`y(x;c_1,c_2)`$ can not be expressed via known functions.
The latter property needs to be stated more precisely. Following \[Um1\], \[Um2\], we define known or classical functions to be functions that can be obtained from the field of rational functions $`\text{C}\text{l}(x)`$, by a finite iteration of the following operations:
i) derivation,
ii) quadrature,
iii) arithmetic operations $`+,,\times ,÷`$,
iv) solution of a homogeneous linear ordinary differential equations with classical coefficients,
v) substitution into an Abelian function,
vi) solution of algebraic differential equations of first order with classical coefficients.
According to this definition, Watanabe (see \[Wat\]) proved that for generic values of the integration constants $`c_1,c_2`$ and of the parameters $`\alpha ,\beta ,\gamma ,\delta `$, the solutions $`y(x;c_1,c_2)`$ of PVI are non-classical and classified all the one-parameter families of classical solutions of PVI. Loosely speaking, Watanabe proves that all one-parameter families of classical solutions of PVI are realized for values of the parameters $`\alpha ,\beta ,\gamma ,\delta `$ lying on the walls of the Weyl chamber of the group $`\stackrel{~}{W}`$ of the birational canonical transformations.<sup>1</sup> Recall that $`\stackrel{~}{W}`$ is isomorphic to $`W_a(D_4)`$ the affine extension of the Weyl group of $`D_4`$, (see \[Ok\]). Such one-parameter families of classical solutions, already found by Okamoto, shrink down, by the action of the group $`\stackrel{~}{W}`$, to the following list<sup>2</sup> The group $`\stackrel{~}{W}`$ acts on $`y(x)`$ and on its conjugate momentum $`p(x)`$. In the list i)$`\mathrm{}`$,iv) the conjugate momentum $`p(x)`$ is given by a one-parameter family.
i) $`y(x)\mathrm{}`$, for $`\alpha =0`$,
ii) $`y(x)0`$, for $`\beta =0`$,
iii) $`y(x)1`$, for $`\gamma =0`$,
iv) $`y(x)x`$, for $`\delta =\frac{1}{2}`$,
v) Riccati solutions,
$$\frac{\mathrm{d}}{\mathrm{d}x}y=\frac{y(y1)(yx)}{x(x1)}\left(\frac{\vartheta _1}{y}+\frac{\vartheta _21}{yx}+\frac{\vartheta _3}{y1}\right),$$
$`(1.1)`$
for $`\vartheta _{\mathrm{}}=\vartheta _1+\vartheta _2+\vartheta _3`$, where $`\vartheta _{\mathrm{}},\vartheta _1,\vartheta _2,\vartheta _3`$ are defined by
$$\alpha =\frac{(\vartheta _{\mathrm{}}1)^2}{2},\beta =\frac{\vartheta _1^2}{2},\gamma =\frac{\vartheta _3^2}{2},\delta =\frac{1\vartheta _2^2}{2}.$$
$`(1.2)`$
Solutions (i), (ii), (iii), (iv) are called degenerate.
The theory of the rational and classical solutions of the Painlevé sixth equation has been extensively studied in \[Ai\], \[AMM\], \[Gr\], \[GL\], \[Luk\], \[Um\], \[Wat\]. In this paper we classify all rational solutions of PVI. We prove the following
###### Main Theorem
All rational non-degenerate solutions of PVI belong to the intersection of two or more one-parameter families of classical solutions, i.e. they occur for
$$\vartheta _{\mathrm{}}+\underset{k=1}{\overset{3}{}}\epsilon _k\vartheta _k2\text{ZZ},\epsilon _k=\pm 1,\text{and}\vartheta _k\text{ZZ}\text{ for at least one}k=1,2,3,\mathrm{}.$$
All rational non-degenerate solutions of the PVI equation are equivalent via birational canonical transformations and up to symmetries<sup>2</sup> The symmetries of the Painlevé VI equation are compositions of the following transformations i) $`x1x`$, $`y1y`$, $`\vartheta _1\vartheta _3`$, ii) $`x\frac{1}{x}`$, $`y\frac{1}{y}`$, $`\vartheta _{\mathrm{}}\vartheta _1+1`$, iii) $`x\frac{1}{1x}`$, $`y\frac{qx}{x1}`$, $`\vartheta _1\vartheta _2`$. to the following solutions
$$y(x)=\frac{(x}{(1+\vartheta _2+x+x\vartheta _3)},\text{for}\vartheta _{\mathrm{}}=\underset{k=1}{\overset{3}{}}\vartheta _k\text{and}\vartheta _1=1,$$
$$y(x)=\frac{(\vartheta _2+\vartheta _3x)^2\vartheta _2\vartheta _3x^2}{(\vartheta _2+\vartheta _31)(\vartheta _2+\vartheta _3x)},\text{for}\vartheta _{\mathrm{}}=\underset{k=1}{\overset{3}{}}\vartheta _k\text{and}\vartheta _1=2.$$
Our method to prove this result does not use Umemura’s theory, but the isomonodromy deformation method (see \[Fuchs\], \[Sch\], \[JMU\], \[ItN\], \[FlN\]). The Painlevé VI is represented as the equation of isomonodromy deformation of the two-dimensional auxiliary Fuchsian system
$$\frac{\mathrm{d}}{\mathrm{d}\lambda }\mathrm{\Phi }=\left(\frac{𝒜_1}{\lambda u_1}+\frac{𝒜_2}{\lambda u_2}+\frac{𝒜_3}{\lambda u_3}\right)\mathrm{\Phi },$$
$`(1.3)`$
$`𝒜_j`$ being $`2\times 2`$ matrices independent on $`\lambda `$, and $`u_1,u_2,u_3`$ being pairwise distinct complex numbers. The matrices $`𝒜_j`$ satisfy the following conditions:
$$\mathrm{eigen}\left(𝒜_j\right)=\pm \frac{\vartheta _j}{2}\text{and}𝒜_1𝒜_2𝒜_3=𝒜_{\mathrm{}}:=\frac{1}{2}\left(\begin{array}{cc}\vartheta _{\mathrm{}}& 0\\ 0& \vartheta _{\mathrm{}}\end{array}\right),$$
$`\vartheta _j`$, $`j=1,2,3,\mathrm{}`$ being related to the parameters $`\alpha ,\beta ,\gamma ,\delta `$ of PVI as in (1.2).
The entries of the matrices $`A_i`$ are complicated expressions of $`x,y,y_x`$ and of some quadrature $`R(x,y)dx`$, $`x`$ being the cross ratio of the poles $`x=\frac{u_2u_1}{u_3u_1}`$. The monodromy matrices $`_1`$, $`_2`$ and $`_3`$ of (1.3) remain constant if and only if $`y=y(x)`$ satisfies PVI. To each branch of a solution of PVI corresponds a triple $`_1,_2,_3`$ of monodromy matrices, which is unique up to
$$(_1,_2,_3)L_{\mathrm{}}^1(_1,_2,_3)L_{\mathrm{}},$$
where $`L_{\mathrm{}}`$ is any constant invertible matrix such that $`[L_{\mathrm{}},_{\mathrm{}}]=0`$ (see Section 2).
Following the same strategy as in \[DM\], we describe the procedure of analytic continuation of a branch of a solution of PVI by the action of the pure braid group on its monodromy matrices. Since rational solutions have only one branch, we look for fixed points of this action. We show that a necessary condition to have a rational (non-degenerate) solution is that the corresponding monodromy group is abelian (see Section 3).
Abelian $`2\times 2`$ monodromy groups are reducible. In Section 4, we classify all solutions of PVI having a reducible monodromy group (such solutions where found in \[Hit\] as a reduction of Nahm’s equations for a diffeomorphic group). Then we classify all rational solutions among them.
Acknowledgements. The idea of classifying all rational solutions of PVI came out of a conversation with H. Umemura. I am indebted to B. Dubrovin who introduced me to the theory of Painlevé equations and gave me lots of suggestions. I am grateful to N. Hitchin, who constantly addressed my work and A.C.C. Coolen for kindly hosting me at Kings College London where this paper was started. The author was supported by an EPSRC research assistanship. We thank P. Clarkson and P. Boalch for pointing out some of the typos which have been fixed in this version.
2. The Painlevé VI equation as the isomonodromic deformation equation of a $`2\times 2`$ Fuchsian system. Consider the following Fuchsian system with four pairwise distinct regular singularities at $`u_1,u_2,u_3,\mathrm{}`$:
$$\frac{\mathrm{d}}{\mathrm{d}\lambda }\mathrm{\Phi }=\left(\frac{𝒜_1}{\lambda u_1}+\frac{𝒜_2}{\lambda u_2}+\frac{𝒜_3}{\lambda u_3}\right)\mathrm{\Phi },\lambda \overline{\text{C}\text{l}}\backslash \{u_1,u_2,u_3,\mathrm{}\}$$
$`(2.1)`$
$`𝒜_j`$ being $`2\times 2`$ matrices independent of $`\lambda `$. Assume that the matrices $`𝒜_j`$ satisfy the following conditions:
$$\mathrm{eigen}\left(𝒜_j\right)=\pm \frac{\vartheta _j}{2}\text{and}𝒜_1𝒜_2𝒜_3=𝒜_{\mathrm{}},$$
$`(2.2)`$
for some constants $`\vartheta _j`$, $`j=1,2,3`$ and
$$𝒜_{\mathrm{}}:=\frac{1}{2}\left(\begin{array}{cc}\vartheta _{\mathrm{}}& 0\\ 0& \vartheta _{\mathrm{}}\end{array}\right),\text{ for some constant}\vartheta _{\mathrm{}}0.$$
$`(2.3)`$
The solution $`\mathrm{\Phi }(\lambda )`$ of the system (2.1) is a multi-valued analytic function on $`\text{C}\text{l}\backslash \{u_1,u_2,u_3\}`$, and its multivaluedness is described by the monodromy matrices. To define them, we fix a basis $`\gamma _1,\gamma _2,\gamma _3`$ of loops in $`\pi _1(\overline{\text{C}\text{l}}\backslash \{u_1,u_2,u_3,\mathrm{}\},\mathrm{})`$ as in figure 1, and a fundamental matrix for the system (2.1).
do{
Fig. 1. The branch-cuts $`\pi _j`$ between the ordered singularities $`u_j`$ and the oriented loops $`\gamma _j`$ in the $`\lambda `$-plane.
###### Proposition 2.1
There exists a fundamental matrix of the system (2.1) of the form
$$\mathrm{\Phi }_{\mathrm{}}(\lambda )=\left(\mathrm{𝟏}+𝒪(\frac{1}{\lambda })\right)\lambda ^𝒜_{\mathrm{}}\lambda ^{_{\mathrm{}}},\text{as}\lambda \mathrm{},$$
$`(2.4)`$
where $`\lambda ^{_{\mathrm{}}}:=e^{_{\mathrm{}}\mathrm{log}\lambda }`$, with the choice of the principal branch of the logarithm with the branch-cut along the common direction $`\eta `$ of the cuts $`\pi _1,\pi _2,\pi _3`$ and the matrix entries of $`_{\mathrm{}}`$ given by $`_{\mathrm{}}^{}{}_{11}{}^{}=_{\mathrm{}}^{}{}_{22}{}^{}=0`$ and
$$\begin{array}{cc}& \text{for}\vartheta _{\mathrm{}}\text{ZZ},\text{and for }\vartheta _{\mathrm{}}=0,_{\mathrm{}}^{}{}_{12}{}^{}=_{\mathrm{}}^{}{}_{21}{}^{}=0,\hfill \\ & \text{for}\vartheta _{\mathrm{}}=n\text{ZZ}_+,_{\mathrm{}}^{}{}_{12}{}^{}=_{\mathrm{}}^{}{}_{12}{}^{}(𝒜_{1,2,3},u_{1,2,3}),_{\mathrm{}}^{}{}_{21}{}^{}=0,\hfill \\ & \text{for}\vartheta _{\mathrm{}}=n,n\text{ZZ}_+,_{\mathrm{}}^{}{}_{21}{}^{}=_{\mathrm{}}^{}{}_{21}{}^{}(𝒜_{1,2,3},u_{1,2,3}),_{\mathrm{}}^{}{}_{12}{}^{}=0,\hfill \end{array}$$
$`(2.5)`$
where the functions $`_{\mathrm{}}^{}{}_{12,21}{}^{}(𝒜_{1,2,3},u_{1,2,3})`$ are uniquely determined by $`(𝒜_{1,2,3},u_{1,2,3})`$. Such a fundamental matrix $`\mathrm{\Phi }_{\mathrm{}}(\lambda )`$ is uniquely determined up to
$$\mathrm{\Phi }_{\mathrm{}}(\lambda )\mathrm{\Phi }_{\mathrm{}}(\lambda )L_{\mathrm{}},$$
$`(2.6)`$
where $`L_{\mathrm{}}`$ is any constant invertible matrix such that
$$\lambda ^𝒜_{\mathrm{}}\lambda ^{_{\mathrm{}}}L_{\mathrm{}}\lambda ^𝒜_{\mathrm{}}\lambda ^{_{\mathrm{}}}=\mathrm{𝟏}+\underset{k=1}{\overset{N}{}}\frac{L_{\mathrm{}}^{(k)}}{\lambda ^k},$$
$`(2.7)`$
for some $`L_{\mathrm{}}^{(1)},\mathrm{},L_{\mathrm{}}^{(N)}`$ constant matrices.
The proof can be found in \[Dub\].
The fundamental matrix $`\mathrm{\Phi }_{\mathrm{}}`$ can be analytically continued to an analytic function on the universal covering of $`\overline{\text{C}\text{l}}\backslash \{u_1,u_2,u_3,\mathrm{}\}`$. For any element $`\gamma \pi _1(\overline{\text{C}\text{l}}\backslash \{u_1,u_2,u_3,\mathrm{}\},\mathrm{})`$ denote the result of the analytic continuation of $`\mathrm{\Phi }_{\mathrm{}}(\lambda )`$ along the loop $`\gamma `$ by $`\gamma [\mathrm{\Phi }_{\mathrm{}}(\lambda )]`$. Since $`\gamma [\mathrm{\Phi }_{\mathrm{}}(\lambda )]`$ and $`\mathrm{\Phi }_{\mathrm{}}(\lambda )`$ are two fundamental matrices in the neighbourhood of infinity, they are related by the following relation:
$$\gamma [\mathrm{\Phi }_{\mathrm{}}(\lambda )]=\mathrm{\Phi }_{\mathrm{}}(\lambda )_\gamma $$
for some constant invertible $`2\times 2`$ matrix $`_\gamma `$ depending only on the homotopy class of $`\gamma `$. Particularly, the matrix $`_{\mathrm{}}:=_\gamma _{\mathrm{}}`$, $`\gamma _{\mathrm{}}`$ being a simple loop around infinity in the clock-wise direction, is given by:
$$_{\mathrm{}}=\mathrm{exp}(2\pi i𝒜_{\mathrm{}})\mathrm{exp}(2\pi i_{\mathrm{}}).$$
$`(2.8)`$
The resulting monodromy representation is an anti-homomorphism:
$$\begin{array}{ccc}\pi _1(\overline{\text{C}\text{l}}\backslash \{u_1,u_2,u_3,\mathrm{}\},\mathrm{})& & SL_2(\text{C}\text{l})\\ \gamma & & _\gamma \end{array}$$
$`(2.9)`$
$$_{\gamma \stackrel{~}{\gamma }}=_{\stackrel{~}{\gamma }}_\gamma .$$
$`(2.10)`$
The images $`_j:=_{\gamma _j}`$ of the generators $`\gamma _j`$, $`j=1,2,3`$ of the fundamental group, are called the monodromy matrices of the Fuchsian system (2.1). They generate the monodromy group of the system, i.e. the image of the representation (2.9). Since the loop $`(\gamma _1\gamma _2\gamma _3)^1`$ is homotopic to $`\gamma _{\mathrm{}}`$, the following relation between the generators holds:
$$_{\mathrm{}}_3_2_1=\mathrm{𝟏}.$$
$`(2.11)`$
Observe that if we fix another fundamental matrix $`\mathrm{\Phi }_{\mathrm{}}^{}=\mathrm{\Phi }_{\mathrm{}}L_{\mathrm{}}`$ in the equivalence class defined by (2.6), the monodromy matrices $`_\gamma ^{}`$ with respect to the new fundamental matrix $`\mathrm{\Phi }_{\mathrm{}}^{}`$ are related to the old ones by
$$_j^{}=L_{\mathrm{}}^1_jL_{\mathrm{}},j=1,2,3.$$
Thus the monodromy matrices $`_3`$, $`_2`$, $`_1`$ are uniquely defined up to the ambiguity
$$(_1,_2,_3)(L_{\mathrm{}}^1_1L_{\mathrm{}},L_{\mathrm{}}^1_2L_{\mathrm{}},L_{\mathrm{}}^1_3L_{\mathrm{}}),$$
$`(2.12)`$
where $`L_{\mathrm{}}`$ is given by (2.7). Observe that $`_{\mathrm{}}`$ is invariant w.r.t. (2.12).
I recall the definition of the connection matrices. Near the poles $`u_j`$, the fundamental matrices $`\mathrm{\Phi }_j(\lambda )`$ of the system (2.1), are given by the following
###### Proposition 2.2
There exists a fundamental matrix of the system (2.1) of the form
$$\mathrm{\Phi }_j(\lambda )=G_j\left(\mathrm{𝟏}+𝒪(\lambda u_j)\right)(\lambda u_j)^{J_j}(\lambda u_j)^_j,\text{as}\lambda u_j,$$
$`(2.13)`$
where
$$\begin{array}{cc}& \text{for }\vartheta _j0,J_j=\frac{1}{2}\left(\begin{array}{cc}\vartheta _j& 0\\ 0& \vartheta _j\end{array}\right),\hfill \\ & \text{for }\vartheta _j=0,J_j=J\left(\begin{array}{cc}0& 1\\ 0& 0\end{array}\right).\hfill \end{array}$$
The invertible matrix $`G_j`$ is defined by $`𝒜_j=G_jJ_jG_j^1`$, the diagonal elements of the matrix $`_j`$ are zero and the off-diagonal ones are defined as follows,
$$\begin{array}{cc}& \text{for}\vartheta _j\text{ZZ},\text{and for }\vartheta _j=0_{j}^{}{}_{12}{}^{}=_{j}^{}{}_{21}{}^{}=0,\hfill \\ & \text{for}\vartheta _j=n\text{ZZ}_+,_{j}^{}{}_{12}{}^{}=_{j}^{}{}_{12}{}^{}(𝒜_{1,2,3},u_{1,2,3}),_{j}^{}{}_{21}{}^{}=0,\hfill \\ & \text{for}\vartheta _j=n,n\text{ZZ}_+,_{j}^{}{}_{21}{}^{}=_{j}^{}{}_{21}{}^{}(𝒜_{1,2,3},u_{1,2,3}),_{j}^{}{}_{12}{}^{}=0,\hfill \end{array}$$
$`(2.14)`$
where the functions $`_{j}^{}{}_{21,12}{}^{}(𝒜_{1,2,3},u_{1,2,3})`$ are uniquely determined by $`𝒜_{1,2,3}`$ and $`u_{1,2,3}`$. The choice of the branch of $`\mathrm{log}(zu_j)`$ is along $`\eta `$ as above. The fundamental matrix $`\mathrm{\Phi }_j(\lambda )`$ is uniquely determined up to the ambiguity:
$$\mathrm{\Phi }_j(\lambda )\mathrm{\Phi }_j(\lambda )L_j$$
$`(2.15)`$
where $`L_j`$ is any constant invertible matrix such that
$$(\lambda u_j)^{J_j}(\lambda u_j)^_jL_j(\lambda u_j)^{J_j}(\lambda u_j)^_j=\underset{k=0}{\overset{N}{}}L_j^{(k)}(\lambda u_j)^k,$$
$`(2.16)`$
for $`L_j^{(0)}=G_j`$ and for some $`L_j^{(1)},\mathrm{},L_j^{(N)}`$ constant matrices.
The proof can be found in \[Dub\].
Continuing the solution $`\mathrm{\Phi }_{\mathrm{}}(\lambda )`$ to a neighbourhood of $`u_j`$, along, say, the right-hand-side of the branch-cut $`\pi _j`$, one obtains another fundamental matrix around $`u_j`$, that must be related to $`\mathrm{\Phi }_j(\lambda )`$ by:
$$\mathrm{\Phi }_{\mathrm{}}(\lambda )=\mathrm{\Phi }_j(\lambda )𝒞_j,$$
$`(2.17)`$
for some invertible matrix $`𝒞_j`$. The matrices $`𝒞_1,𝒞_2,𝒞_3`$ are called connection matrices, and they are defined by (2.17) up to the ambiguity $`𝒞_j𝒞_jL_{\mathrm{}}`$ due to (2.6). The connection matrices are related to the monodromy matrices as follows:
$$_j=𝒞_j^1\mathrm{exp}\left(2\pi iJ_j\right)\mathrm{exp}\left(_j\right)𝒞_j,j=1,2,3.$$
$`(2.18)`$
Thanks to the above relation it follows that
$$\mathrm{eigen}(_j)=\mathrm{exp}(\pm \pi i\vartheta _j).$$
$`(2.19)`$
###### Definition 2.3
The Monodromy data of the Fuchsian system (2.1) are
$$\{(_1,_2,_3)/_{},_1,_2,_3\},$$
where $``$ is the equivalence relation defined by (2.12).
Remark 2.4. For non-integer $`\vartheta _j`$, the correspondent $`_j`$ matrix is zero by definition and we drop it from the set of the monodromy data.
The theory of the deformations the poles of the Fuchsian system keeping the monodromy fixed is described by the following result:
###### Theorem 2.5
Let $`\{(_1,_2,_3)/_{},_1,_2,_3\}`$, be monodromy data of the Fuchsian system:
$$\frac{\mathrm{d}}{\mathrm{d}\lambda }\mathrm{\Phi }^0=\left(\frac{𝒜_1^0}{\lambda u_1^0}+\frac{𝒜_2^0}{\lambda u_2^0}+\frac{𝒜_3^0}{\lambda u_3^0}\right)\mathrm{\Phi }^0,$$
$`(2.20)`$
of the above form (2.2), with pairwise distinct poles $`u_j^0`$, and with respect to some basis $`\gamma _1,\gamma _2,\gamma _3`$ of the loops in $`\pi _1(\overline{\text{C}\text{l}}\backslash \{u_1^0,u_2^0,u_3^0,\mathrm{}\},\mathrm{})`$. If $`_k\pm \mathrm{𝟏}`$ for all $`k=1,2,3,\mathrm{}`$, there exists a neighbourhood $`U\text{C}\text{l}^3`$ of the point $`u^0=(u_1^0,u_2^0,u_3^0)`$ such that, for any $`u=(u_1,u_2,u_3)U`$, there exists a unique triple $`𝒜_1(u)`$, $`𝒜_2(u)`$, $`𝒜_3(u)`$ of analytic matrix valued functions such that:
$$𝒜_j(u^0)=𝒜_j^0,i=1,2,3,$$
and the monodromy matrices of the Fuchsian system
$$\frac{\mathrm{d}}{\mathrm{d}\lambda }\mathrm{\Phi }=A(\lambda ;u)\mathrm{\Phi }=\left(\frac{𝒜_1(u)}{\lambda u_1}+\frac{𝒜_2(u)}{\lambda u_2}+\frac{𝒜_3(u)}{\lambda u_3}\right)\mathrm{\Phi },$$
$`(2.21)`$
with respect to the same basis<sup>1</sup> Observe that the basis $`\gamma _1,\gamma _2,\gamma _3`$ of $`\pi _1(\overline{\text{C}\text{l}}\backslash \{u_1,u_2,u_3,\mathrm{}\},\mathrm{})`$ varies continuously with small variations of $`u_1,u_2,u_3`$. This new basis is homotopic to the initial one, so one can identify them. $`\gamma _1,\gamma _2,\gamma _3`$ of the loops, can be chosen to coincide with the given $`_1`$, $`_2`$, $`_3`$. The matrices $`𝒜_j(u)`$ are the solutions of the Cauchy problem with the initial data $`𝒜_j^0`$ for the following Schlesinger equations:
$$\frac{}{u_j}𝒜_i=\frac{[𝒜_i,𝒜_j]}{u_iu_j},\frac{}{u_i}𝒜_i=\underset{ji}{}\frac{[𝒜_i,𝒜_j]}{u_iu_j}.$$
$`(2.22)`$
The solution $`\mathrm{\Phi }_{\mathrm{}}^0(\lambda )`$ of (2.20) of the form (2.4) can be uniquely continued, for $`\lambda u_i`$ $`i=1,2,3`$, to an analytic function $`\mathrm{\Phi }_{\mathrm{}}(\lambda ,u),uU`$, such that
$$\mathrm{\Phi }_{\mathrm{}}(\lambda ,u^0)=\mathrm{\Phi }_{\mathrm{}}^0(\lambda ).$$
This continuation is the local solution of the Cauchy problem with the initial data $`\mathrm{\Phi }_{\mathrm{}}^0`$ for the following system that is compatible with the system (2.21):
$$\frac{}{u_i}\mathrm{\Phi }=\frac{𝒜_i(u)}{\lambda u_i}\mathrm{\Phi }.$$
Moreover the functions $`𝒜_i(u)`$ and $`\mathrm{\Phi }_{\mathrm{}}(\lambda ,u)`$ can be continued analytically to global meromorphic functions on the universal coverings of
$$\text{C}\text{l}^3\backslash \{diags\}:=\left\{(u_1,u_2,u_3)\text{C}\text{l}^3\right|u_iu_j\text{for}ij\},$$
and
$$\left\{(\lambda ,u_1,u_2,u_3)\text{C}\text{l}^4\right|u_iu_j\text{for}ij\text{and}\lambda u_i,i=1,2,3\},$$
respectively.
The proof of this theorem can be found, for example, in \[Mal\], \[Miwa\], \[Sib\].
Remark 2.6. Observe that the isomonodromic deformations equations preserve the connection matrices $`𝒞_i`$ too.
Remark 2.7. When $`_k=\pm \mathrm{𝟏}`$ for some $`k=1,2,3\mathrm{}`$, the existence statements of Theorem 2.5 are still valid, while the uniqueness ones are lost.
Let me now explain, following \[JMU\], how to rewrite the Schlesinger equations (2.22) in terms of the PVI equation. We can assume $`\vartheta _{\mathrm{}}0`$ without loss of generality (see Remark 2.9 below). Schlesinger equations (2.22) with fixed $`𝒜_{\mathrm{}}`$ are invariant with respect to the gauge transformations of the form:
$$𝒜_iD^1𝒜_iD,i=1,2,3,\text{for any }D\text{ diagonal matrix}.$$
$`(2.23)`$
We introduce two coordinates $`(p,q)`$ on the quotient of the space of matrices satisfying (2.22) with respect to the equivalence relation (2.23) and a coordinate $`k`$ that contains the above gauge freedom:
$$[𝒜(q;u_1,u_2,u_3)]_{12}=0,p=\underset{k=1}{\overset{3}{}}\frac{𝒜_{k_{11}}+\frac{\vartheta _k}{2}}{qu_k},k=\frac{2P(\lambda )[𝒜(\lambda ;u_1,u_2,u_3)]_{12}}{\vartheta _{\mathrm{}}(q\lambda )},$$
where $`𝒜(z;u_1,u_2,u_3)`$ is given in (2.21) and $`P(\lambda )=(\lambda u_1)(\lambda u_2)(\lambda u_3)`$. The matrices $`𝒜_i`$ are uniquely determined by the coordinates $`(p,q)`$, and $`k`$ and expressed rationally in terms of them:
$$\begin{array}{cc}\hfill 𝒜_{i}^{}{}_{11}{}^{}& =\frac{1}{\vartheta _{\mathrm{}}P^{}(u_i)}\{P(q)(qu_i)p^2+P(q)(qu_i)p(\frac{\vartheta _{\mathrm{}}}{qu_i}\underset{k=1}{\overset{3}{}}\frac{\vartheta _k}{qu_k})+\hfill \\ & +(qu_i)\left[\frac{\vartheta _{\mathrm{}}^{}{}_{}{}^{2}}{4}\left(q+2u_i\underset{k=1}{\overset{3}{}}u_k\right)+\underset{k=1}{\overset{3}{}}\frac{\vartheta _k^2}{4}\left(q+2u_k\underset{j=1}{\overset{3}{}}u_j\right)\right]+\hfill \\ & +\frac{qu_i}{2}(\vartheta _1\vartheta _2(qu_3)+\vartheta _1\vartheta _3(qu_2)+\vartheta _2\vartheta _3(qu_1))\frac{\vartheta _{\mathrm{}}}{2}\underset{k=1}{\overset{3}{}}\frac{\vartheta _k}{qu_k}\}\hfill \\ \hfill 𝒜_{i}^{}{}_{12}{}^{}& =\vartheta _{\mathrm{}}k\frac{qu_i}{2P^{}(u_i)},\hfill \\ \hfill 𝒜_{i}^{}{}_{21}{}^{}& =\frac{1}{𝒜_{i}^{}{}_{12}{}^{}}\left(\frac{\vartheta _i^2}{4}𝒜_{i}^{}{}_{11}{}^{2}\right),\hfill \\ \hfill 𝒜_{i}^{}{}_{22}{}^{}& =𝒜_{i}^{}{}_{11}{}^{}\hfill \end{array}$$
$`(2.24)`$
for $`i=1,2,3`$, where $`P^{}(z)=\frac{\mathrm{d}P}{\mathrm{d}z}`$. The Schlesinger equations (2.22) in these variables are:
$$\{\begin{array}{cc}\hfill \frac{q}{u_i}& =\frac{P(q)}{P^{}(u_i)}\left[2p+\frac{1}{qu_i}\underset{k=1}{\overset{3}{}}\frac{\vartheta _k}{qu_k}\right]\hfill \\ \hfill \frac{p}{u_i}& =\{P^{}(q)p^2+[2q+u_i\underset{j}{}u_j\underset{k=1}{\overset{3}{}}\vartheta _k(2q+u_k\underset{j}{}u_j)]p+\hfill \\ & +\frac{1}{4}(\underset{k=1}{\overset{3}{}}\vartheta _k\vartheta _{\mathrm{}})(\underset{k=1}{\overset{3}{}}\vartheta _k+\vartheta _{\mathrm{}}2)\}\frac{1}{P^{}(u_i)},\hfill \end{array}$$
$`(2.25)`$
and
$$\frac{\mathrm{log}(k)}{u_i}=(\vartheta _{\mathrm{}}1)\frac{qu_i}{P^{}(u_i)}.$$
$`(2.26)`$
for $`i=1,2,3`$. The system of the reduced Schlesinger equations (2.25) is invariant under the transformations of the form
$$u_iau_i+b,qaq+b,p\frac{p}{a},a,b\text{C}\text{l},a0.$$
Introducing the following new invariant variables:
$$x=\frac{u_2u_1}{u_3u_1},y=\frac{qu_1}{u_3u_1};$$
$`(2.27)`$
the system (2.25), expressed in the these new variables, gives the PVI equation for $`y(x)`$ with parameters
$$\alpha =\frac{(\vartheta _{\mathrm{}}1)^2}{2},\beta =\frac{\vartheta _1^2}{2},\gamma =\frac{\vartheta _3^2}{2},\delta =\frac{1\vartheta _2^2}{2}.$$
$`(2.28)`$
Remark 2.8. Observe that permutations of the poles $`u_i`$ and of the values $`\vartheta _i`$, $`i=1,2,3,\mathrm{}`$ induce transformations of $`(y,x)`$ of the type $`x1x`$ and $`y1y`$, $`x\frac{1}{x}`$ and $`y\frac{1}{y}`$, $`x\frac{1}{1x}`$ and $`y\frac{yx}{1x}`$ and their compositions. These transformations are the symmetries of the Painlevè VI equation.
Remark 2.9. It is clear from (2.28) that changes of the signs of the parameters $`\vartheta _k`$, $`k=1,2,3`$ and transformations on $`\vartheta _{\mathrm{}}`$ of type $`\vartheta _{\mathrm{}}2\vartheta _{\mathrm{}}`$ give rise to the same PVI equation.
We summarise the results of this section in the following:
###### Theorem 2.10
Branches $`y(x)`$ of solutions of the PVI equation with parameters $`\alpha `$, $`\beta `$, $`\gamma `$, $`\delta `$, considered up to symmetries are in one to one correspondence with local solutions of the Schlesinger equations (2.22) with parameters $`\vartheta _1,\vartheta _2,\vartheta _3,\vartheta _{\mathrm{}}`$ given by (2.28) and $`𝒜_{\mathrm{}}`$ given in (2.3), considered up to diagonal conjugation (2.23) and permutation. This one-to-one correspondence is realized by the formulae (2.24), (2.27). For each branch of a solution of the PVI equation, there exist a unique set of monodromy data $`\{(_1,_2,_3)/_{},_1,_2,_3\}`$. Vice-versa, let $`\{(_1,_2,_3)/_{},_1,_2,_3\}`$ be a set of monodromy data such that
$$\mathrm{eigen}\left(_j\right)=\mathrm{exp}(\pm \pi i\vartheta _j),(_3_2_1)^1=_{\mathrm{}},\mathrm{eigen}\left(_{\mathrm{}}\right)=\mathrm{exp}(\pm \pi i\vartheta _{\mathrm{}}),$$
$$_k\pm \mathrm{𝟏}k=\mathrm{\hspace{0.17em}1},2,3,\mathrm{},$$
with some numbers $`\vartheta _1,\vartheta _2,\vartheta _3,\vartheta _{\mathrm{}}`$, and $`_1,_2,_3`$ satisfying (2.14). If there exists a branch of a solution of the PVI equation with parameters (2.28) such that the Fuchsian system of the form (2.1) given by (2.24) has the prescribed monodromy data $`\{(_1,_2,_3)/_{},_1,_2,_3\}`$ then this branch is unique modulo symmetries.
3. Analytic continuation and rational solutions of the Painlevé VI equation.
In Theorem 2.10, we parameterised branches of generic solutions of PVI by triples of monodromy matrices. Following \[DM\], now we show how these parameters change with a change of the branch in the process of analytic continuation of the solutions along a path in $`\overline{\text{C}\text{l}}\backslash \{0,1,\mathrm{}\}`$. Recall that, as it follows from Theorem 2.5, any solution of the Schlesinger equations can be continued analytically from a point $`(u_1^0,u_2^0,u_3^0)`$ to another point $`(u_1^1,u_2^1,u_3^1)`$ along a path
$$(u_1(t),u_2(t),u_3(t))\text{C}\text{l}^3\backslash \{diags\},0t1,$$
where $`\{diags\}=\{u_1,u_2,u_3|u_i=u_j,\text{for some}ij\}`$ and
$$u_i(0)=u_i^0,\text{and}u_i(1)=u_i^1,$$
provided that the end-points are not the poles of the solution. The result of the analytic continuation depends only on the homotopy class of the path in $`\text{C}\text{l}^3\backslash \{diags\}`$. Particularly, to find all the branches of a solution near a given point $`u^0=(u_1^0,u_2^0,u_3^0)`$ one has to compute the results of the analytic continuation along any homotopy class of closed loops in $`\text{C}\text{l}^3\backslash \{diags\}`$ with the beginning and the end at the point $`u^0=(u_1^0,u_2^0,u_3^0)`$. Let
$$\beta \pi _1(\text{C}\text{l}^3\backslash \{diags\};u^0)$$
be an arbitrary loop. Any solution of the Schlesinger equations near the point $`u^0=(u_1^0,u_2^0,u_3^0)`$, is uniquely determined up to (2.23) by the monodromy matrices $`_1`$, $`_2`$ and $`_3`$, computed with respect to the basis of loops $`\gamma _1`$, $`\gamma _2`$, $`\gamma _3`$. Continuing analytically this solution along the loop $`\beta `$, we arrive at another branch of the same solution near $`u^0`$. This new branch is specified, according to Theorem 2.10, by some new monodromy matrices $`_1^\beta `$, $`_2^\beta `$ and $`_3^\beta `$, computed in the same basis $`\gamma _1`$, $`\gamma _2`$, $`\gamma _3`$. We want to compute these new matrices for any loop $`\beta \pi _1(\text{C}\text{l}^3\backslash \{diags\};u^0)`$. The fundamental group $`\pi _1(\text{C}\text{l}^3\backslash \{diags\};u^0)`$ is isomorphic to the pure (or unpermuted) braid group, $`P_3`$ with three strings (see \[Bir\]).
###### Lemma 3.1
For the generators $`\beta _1`$, $`\beta _2`$, $`\beta _3`$ of the pure braid group $`P_3`$, $`_i^\beta `$ have the following form:
$$_1^{\beta _1}=_1_2_1_2^1_1^1,_2^{\beta _1}=_1_2_1^1,_3^{\beta _1}=_3,$$
$`(2.29)`$
$$\begin{array}{cc}\hfill _1^{\beta _2}& =_1_3_1_3^1_1^1,_3^{\beta _2}=_1_3_1^1,\hfill \\ & _2^{\beta _2}=_1_3_1^1_3^1_2_3_1_3^1_1^1.\hfill \end{array}$$
$`(2.30)`$
$$_1^{\beta _3}=_1_2^{\beta _3}=_2_3_2_3^1_2^1,_3^{\beta _3}=_2_3_2^1,$$
$`(2.31)`$
Proof. This lemma is proved in \[DM\].
###### Lemma 3.2
If a solution of the Painlevé VI equation such that none of the corresponding monodromy matrices $`_1,_2,_3`$ are multiples of the identity is rational then $`_1,_2,_3`$ are fixed points under the above action (2.29), (2.30), (2.31) on the space of triples
$$\{(_1,_2,_3)/_{},_3_2_1=_{\mathrm{}}^1\}$$
the equivalence relation $``$ being defined in (2.12).
Proof. The action (2.29), (2.30), (2.31) of the pure braid group on the triples of monodromy matrices not only describes the structure of the analytic continuation of the solutions of the Schlesinger equations (2.22), but also of the reduced ones (2.25) and thus of the PVI equation. A necessary condition for a solution to be rational is that it has only one branch. Since the action (2.29), (2.30), (2.31) of the pure braid group on the triples of monodromy matrices commutes with the conjugation (2.12), the triple of monodromy matrices $`_1,_2,_3`$ corresponding to a rational solution is unique up to (2.12). $``$
###### Lemma 3.3
If a solution of PVI such that none of the corresponding monodromy matrices $`_1,_2,_3`$ are multiples of the identity is rational then the monodromy matrices all commute
$$[_i,_j]=0,i,j=1,2,3.$$
$`(2.32)`$
Proof. By Lemma 3.2, we have to impose
$$\begin{array}{cc}\hfill L_{\mathrm{}}^{}{}_{1}{}^{1}_1L_{\mathrm{}}^{}{}_{1}{}^{}& =_1^{\beta _1}=_1_2_1_2^1_1^1,\hfill \\ \hfill L_{\mathrm{}}^{}{}_{1}{}^{1}_2L_{\mathrm{}}^{}{}_{1}{}^{}& =_2^{\beta _1}=_1_2_1^1,\hfill \\ \hfill L_{\mathrm{}}^{}{}_{1}{}^{1}_3L_{\mathrm{}}^{}{}_{1}{}^{}& =_3^{\beta _1}=_3,\hfill \end{array}$$
$`(2.33)`$
$$\begin{array}{cc}\hfill L_{\mathrm{}}^{}{}_{2}{}^{1}_1L_{\mathrm{}}^{}{}_{2}{}^{}& =_1^{\beta _2}=_1_3_1_3^1_1^1,\hfill \\ \hfill L_{\mathrm{}}^{}{}_{2}{}^{1}_2L_{\mathrm{}}^{}{}_{2}{}^{}& =_2^{\beta _2}=_1_3_1^1_3^1_2_3_1_3^1_1^1,\hfill \\ \hfill L_{\mathrm{}}^{}{}_{2}{}^{1}_3L_{\mathrm{}}^{}{}_{2}{}^{}& =_3^{\beta _2}=_1_3_1^1,\hfill \end{array}$$
$`(2.34)`$
$$\begin{array}{cc}\hfill L_{\mathrm{}}^{}{}_{3}{}^{1}_1L_{\mathrm{}}^{}{}_{3}{}^{}& =_1^{\beta _3}=_1,\hfill \\ \hfill L_{\mathrm{}}^{}{}_{3}{}^{1}_2L_{\mathrm{}}^{}{}_{3}{}^{}& =_2^{\beta _3}=_2_3_2_3^1_2^1,\hfill \\ \hfill L_{\mathrm{}}^{}{}_{3}{}^{1}_3L_{\mathrm{}}^{}{}_{3}{}^{}& =_3^{\beta _3}=_2_3_2^1,\hfill \end{array}$$
$`(2.35)`$
for some suitable matrices $`L_{\mathrm{}}^{}{}_{1}{}^{}`$, $`L_{\mathrm{}}^{}{}_{2}{}^{}`$, $`L_{\mathrm{}}^{}{}_{3}{}^{}`$ that are diagonal for $`\vartheta _{\mathrm{}}\text{ZZ}`$, are in Jordan form for $`\vartheta _{\mathrm{}}\text{ZZ}`$ and $`_{\mathrm{}}0`$. Then we have to distinguish two cases: i) $`_{\mathrm{}}\pm \mathrm{𝟏}`$ is diagonal, ii) $`_{\mathrm{}}\pm \mathrm{𝟏}`$ is in Jordan form.
i) In this case, the matrices $`L_{\mathrm{}}^{}{}_{1}{}^{},L_{\mathrm{}}^{}{}_{2}{}^{},L_{\mathrm{}}^{}{}_{3}{}^{}`$ are diagonal. If none of the matrices $`_{1,2,3}`$ is diagonal in the basis of $`_{\mathrm{}}`$ diagonal, then the above matrices $`L_{\mathrm{}}^{}{}_{1}{}^{},L_{\mathrm{}}^{}{}_{3}{}^{}`$ must be chosen to be multiples of the identity. Thus, from
$$\begin{array}{cc}\hfill _2& =_2^{\beta _1}=_1_2_1^1,\hfill \\ \hfill _3& =_3^{\beta _3}=_2_3_2^1,\hfill \end{array}$$
it follows immediately that
$$[_2,_1]=0,\text{and}[_2,_3]=0,$$
thus
$$L_{\mathrm{}}^{}{}_{2}{}^{1}_2L_{\mathrm{}}^{}{}_{2}{}^{}=_2^{\beta _2}=_2,$$
i.e. $`L_{\mathrm{}}^{}{}_{2}{}^{}`$ is the identity matrix as well and thus we obtain (2.32). If one of the monodromy matrices $`_{1,2,3}`$ is diagonal in the basis of $`_{\mathrm{}}`$ diagonal, for example $`_1`$, then either $`L_{\mathrm{}}^{}{}_{1}{}^{}`$ is the identity matrix or $`_3`$ is diagonal. If $`L_{\mathrm{}}^{}{}_{1}{}^{}`$ is the identity matrix, from (2.33) we have that $`_2`$ is diagonal as well, and being $`_{\mathrm{}}`$ diagonal as well, $`_3`$ must be diagonal and we find again (2.32). If $`_3`$ and $`_1`$ are both diagonal, since $`_{\mathrm{}}`$ is diagonal, $`_2`$ must be diagonal too and we find again (2.32).
ii) Suppose that $`_{\mathrm{}}`$ is in Jordan form. Then the matrices $`L_{\mathrm{}}^{}{}_{1}{}^{},L_{\mathrm{}}^{}{}_{2}{}^{},L_{\mathrm{}}^{}{}_{3}{}^{}`$ have the form $`\left(\begin{array}{cc}1& a_i\\ 0& 1\end{array}\right)`$ for some constants $`a_1,a_2,a_3`$. If none of the matrices $`_1,_2,_3`$ is in Jordan form, all matrices $`L_{\mathrm{}}^{}{}_{1}{}^{},L_{\mathrm{}}^{}{}_{2}{}^{},L_{\mathrm{}}^{}{}_{3}{}^{}`$ are the identity matrix and we obtain (2.32) as above. Suppose that one of the monodromy matrices, for example $`_1`$ is in Jordan form. Then, either $`L_{\mathrm{}}^{}{}_{1}{}^{}`$ is the identity and $`[_2,_1]=0`$ or $`_\mathcal{3}`$ is in Jordan form. Reasoning as above we obtain that all matrices $`_1,_2,_3`$ are in Jordan form and thus commute. $``$
4. Reducible monodromy groups.
###### Theorem 4.1
All solutions of PVI corresponding to reducible monodromy groups are equivalent via birational canonical transformations and symmetries to the following one-parameter family of solutions, realized for $`\vartheta _{\mathrm{}}=(\vartheta _1+\vartheta _2+\vartheta _3)`$
$$y=\frac{(1+\vartheta _1+\vartheta _2x\vartheta _2x)u(x)x(x1)u_x(x)}{(1+\vartheta _1+\vartheta _2+\vartheta _3)u(x)}$$
$`(4.1)`$
where $`u(x)=u_1(x)+\nu u_2(x)`$, $`u_1(x),u_2(x)`$ being two linear independent solutions of the following hypergeometric equation
$$x(1x)u_{xx}(x)+\left[2+\vartheta _1+\vartheta _2(4+\vartheta _1+2\vartheta _2+\vartheta _3)x\right]u_x(x)(2+\vartheta _1+\vartheta _2+\vartheta _3)(\vartheta _2+1)u(x)=0.$$
$`(4.2)`$
Proof. For reducible monodromy groups there exists a basis in which all monodromy matrices are upper triangular. We can always perform a change of basis in order that $`_{\mathrm{}}`$ has the form (2.8) and all monodromy matrices have the form:
$$_k=\left(\begin{array}{cc}\mathrm{exp}(\pi i\vartheta _k)& \\ 0& \mathrm{exp}(\pi i\vartheta _k)\end{array}\right).$$
It then follows, by the relation (2.11), that $`\vartheta _{\mathrm{}}+\epsilon _k_k\vartheta _k=2N`$, $`\epsilon _k=\pm 1`$, $`N\text{ZZ}`$. By means of the birational canonical transformations and symmetries of the PVI equation, we can always assume that $`\epsilon _k=+1`$ and $`N=0`$. Perform the following gauge transformation on the Fuchsian system
$$\mathrm{\Phi }\stackrel{~}{\mathrm{\Phi }}=\mathrm{\Pi }_{k=1}^3(\lambda u_k)^{\frac{\vartheta _k}{2}}\mathrm{\Phi },𝒜_k\stackrel{~}{𝒜}_k=𝒜_k\frac{\vartheta _k}{2}\mathrm{𝟏}.$$
The new residue at infinity is
$$\stackrel{~}{𝒜}_{\mathrm{}}=\left(\begin{array}{cc}0& 0\\ 0& \frac{\vartheta _{\mathrm{}}+\vartheta _k}{2}\end{array}\right),$$
and the new monodromy matrices are
$$\stackrel{~}{}_k=\mathrm{exp}(\pi i\vartheta _k)_k=\left(\begin{array}{cc}1& \\ 0& \mathrm{exp}(2\pi i\vartheta _k)\end{array}\right),k=1,2,3,$$
and
$$\stackrel{~}{}_{\mathrm{}}=\left(\begin{array}{cc}1& 0\\ 0& \mathrm{exp}[\pi i(\vartheta _{\mathrm{}}+_k\vartheta _k)]\end{array}\right)\mathrm{exp}(_{\mathrm{}}).$$
Since all monodromy matrices have the first column given by $`\left(\begin{array}{c}1\\ 0\end{array}\right)`$, the new Fuchsian system admits a non-zero single valued vector solution $`\stackrel{~}{Y}`$. This solution is analytic at $`u_1`$, $`u_2`$ and $`u_3`$ because all residue matrices $`\stackrel{~}{𝒜}_k`$, $`k=1,2,3`$ have a zero eigenvalue. At infinity $`\stackrel{~}{Y}`$ is necessarily constant. Thus, near each $`u_j`$, one has
$$\frac{\stackrel{~}{𝒜}_j}{\lambda u_j}\stackrel{~}{Y}+𝒪(1)=0,$$
that implies that $`\stackrel{~}{Y}`$ has the form $`\left(\begin{array}{c}a\\ 0\end{array}\right)`$, for some $`a0`$ and all residue matrices $`𝒜_k`$ are upper triangular. Correspondingly $`p=_{k=1}^3\frac{\vartheta _k}{(qu_k)}`$ and the solution $`y(x)`$ of PVI is given by (4.1). $``$
Remark 4.2. Observe that the appearance of the hypergeometric equation in Theorem 4.1 is quite natural. In fact, when all the residue matrices $`𝒜_k`$ are upper-triangular, the Schlesinger equations, and thus the Painlevé VI equation, linearise.
###### Lemma 4.3
The one-parameter family of classical solutions (4.1) contains at least one rational solution if and only if one of the values $`\vartheta _1`$, $`\vartheta _2`$, $`\vartheta _3`$, $`\vartheta _{\mathrm{}}`$ is an integer.
Proof. The one-parameter family of classical solutions (4.1) contains at least one algebraic solution if and only if the corresponding Riccati equation
$$y_x(x)=\frac{1+\vartheta _1+\vartheta _2+\vartheta _3}{x(x1)}y^2\frac{1+\vartheta _1+\vartheta _2+\vartheta _1x+\vartheta _3x}{x(x1)}y+\frac{\vartheta _1}{x1}$$
has at least one algebraic solution, i.e., if and only the corresponding hypergeometric equation (4.2) is integrable in the sense of differential Galois theory (see \[Mor\]). This happens if and only if and only if the parameters $`\lambda `$, $`\mu `$,$`\nu `$ of the hypergeometric equation belong to the Schwartz-Kimura table (see \[Mor\]). In particular rational solutions are occur only when at least one of the numbers $`\mu \lambda +\nu `$, $`\mu \lambda \nu `$, $`\mu +\lambda +\nu `$, $`\mu +\lambda \nu `$ is an odd integer. In our case $`\lambda =1\vartheta _1\vartheta _2`$, $`\mu =1+\vartheta _1+\vartheta _3`$, $`\nu =1\vartheta _3\vartheta _2`$ so the one-parameter family of classical solutions (4.1) contains at least one rational solution if and only if one of the values $`\vartheta _1`$, $`\vartheta _2`$, $`\vartheta _3`$, $`\vartheta _{\mathrm{}}`$ is an integer. $``$
To conclude the proof of the main theorem, we use the following lemma
###### Lemma 4.4
All solutions of PVI such that one or more monodromy matrices $`_1,_2,_3,_{\mathrm{}}`$ is a multiple of the identity are either degenerate or are equivalent via birational canonical transformations to the one-parameter family (4.1) of Theorem 4.1.
Proof. We omit the proof of this Lemma that can be found for example in \[Ma1\].
The classification is concluded observing that if one or more of the monodromy matrices is a multiple of the identity, then the corresponding solution is equivalent via birational canonical transformations to a solution with reducible monodromy groups. If none of the monodromy matrices is a multiple of the identity, then in order to have a rational solution the monodromy group must be Abelian. All Abelian $`2\times 2`$ monodromy groups are reducible. Thus all rational non-degenerate solutions of the PVI equation belong to the one parameter family of classical solutions given in Theorem 4.1, up to birational canonical transformations. Then we can use Lemma 4.3 to classify them.
BIBLIOGRAPHY
\[Ai\] H. Airault, Rational solutions of Painlevé equations, Stud. Appl. Math. 61, (1979) 31–53.
\[AMM\] D.A. Albrecht, E.L. Mansfield, A.E. Milne, Algoritms for special integrals of ordinary differential equations, J. Phys. A: Math. Gen. 29 (1996) 973–991.
\[Bir\] J.S. Birman, Braids, Links, and Mapping Class Groups, Ann. Math. Stud. Princeton University (1975).
\[Dub\] B. Dubrovin, Painlevé Transcendents in Two-Dimensional Topological Field Theory, The Painlevé Property One Century Later, Robert Conte editor, CRM Series in Mathematical Physics, (1999).
\[DM\] B. Dubrovin, M.Mazzocco, Monodromy of certain Painlevé VI transcendents and Reflection Groups, Invent. Math., 141 (2000) 55–147.
\[FlN\] H. Flashka and A.C. Newell, Monodromy and Spectrum Preserving Deformations, Comm. Math. Phys. 76 (1980) 67–116.
\[Fuchs\] R. Fuchs, Über Lineare Homogene Differentialgleichungen Zweiter Ordnung mit im drei im Endrichen Gelegene Wesentlich Singulären Stellen, Math. Ann. 63 (1907) 301–321.
\[Gam\] B. Gambier, Sur les Equations Differentielles du Second Ordre et du Primier Degrè dont l’Integrale est a Points Critiques Fixes, Acta Math. 33, (1910) 1–55.
\[Gar\] R. Garnier, Sur des Equations Differentielles du Troisieme Ordre dont l’Integrale Generale est Uniforme et sur una Classe d’Equations nouvelles d’Ordre Superieur dont l’Integrale Generale a ses Points Critiques Fixes, Ann. Sci. Ecole Norm. Sup. 3 No.29 (1912) 1–126.
\[Gr\] V.I. Gromak, Single-parameter families of solutions of Painlevé equations, Diff. Eqns. 14 (1978), 1510–1513.
\[GL\] V.I. Gromak, N.A. Lukashevich, Special classes of solutions of Painlevé equations, Diff. Eqns. 18 (1982) 317–326.
\[Hit\] N. Hitchin, Hypercomplex Manifolds and the Space of Framings, The Geometric Universe, Oxford Univ. Press, Oxford, 1998, 9–30.
\[Ince\] E.L. Ince, Ordinary Differential Equations, Dover Publications, New York (1956).
\[ItN\] A.R. Its and V.Yu. Novokshenov, The Isomonodromic Deformation Method in the Theory of Painlevé Equations, Springer Lect. Notes Math. 1191 (1986).
\[JMU\] M. Jimbo, T. Miwa and K. Ueno, Monodromy Preserving Deformation of the Linear Ordinary Differential Equations with Rational Coefficients I, II, III, Physica D, 2 (1981), no. 2, 306–352, 2 (1981), no. 3, 407–448, 4 (1981/82), no. 1, 26–46.
\[Luk\] N.A. Lukashevich, Elementary solutions of certain Painlevé equations, Differ. Uravn. 1 (1965) 731–735.
\[Mal\] B. Malgrange, Sur les Deformations Isomonodromiques I, Singularités Régulières, Seminaire de l’Ecole Normale Superieure 1979–1982, Progress in Mathematics 37, Birkhäuser, Boston (1983) 401–426.
\[Ma\] M. Mazzocco, Picard and Chazy Solutions to PVI Equation, Math. Ann., 321 (2001), 157–195.
\[Ma1\] M. Mazzocco, The geometry of the classical solutions of the Garnier systems, I.M.R.N., 12 (2002) 613–646.
\[Miwa\] T. Miwa, Painlevé Property of Monodromy Preserving Deformation Equations and the Analyticity of $`\tau `$-function, Publ. RIMS, Kyoto Univ. 17 (1981) 703–721.
\[Mor\] J.J. Morales, Differential Galois theory and non-integrability of Hamiltonian systems Progress in Mathematics, 179, (1999).
\[Ok\] K. Okamoto, Studies on the Painlevé Equations I, Sixth Painlevé Equation, Ann. Mat. Pura Appl. 146 (1987) 337–381.
\[Pain\] P. Painlevé, Sur les Equations Differentielles du Second Ordre et d’Ordre Superieur, dont l’Interable Generale est Uniforme, Acta Math. 25 (1902) 1–86.
\[Sch\] L. Schlesinger, Über eine Klasse von Differentsial System Beliebliger Ordnung mit Festen Kritischer Punkten, J. fur Math. 141 (1912), 96–145.
\[Schw\] H.A. Schwartz, Über Diejenigen Fälle in Welchen die Gaussische Hypergeometrische Reihe einer Algebraische Funktion iheres vierten Elementes Darstellit, Crelle J. 75 (1873) 292–335.
\[Sib\] Y. Sibuya, Linear Differential Equations in the Complex Domain: Problems of Analytic Continuation, AMS TMM 82 (1990).
\[Um\] H. Umemura, Special polynomials associated with the Painlevé equations I, to appear in Painlevé transcendents, CRM Montreal, Canada (1996).
\[Um1\] H. Umemura, Irreducebility of the First Differential Equation of Painlevé, Nagoya Math. J., 117 (1990) 231–252.
\[Um2\] H. Umemura, Second proof of the irreducebility of the First Differential Equation of Painlevé, Nagoya Math. J., 117 (1990) 125–171.
\[Wat\] H. Watanabe, Birational Canonical Transformations and Classical Solutions of the Sixth Painlevé Equation, Ann. Scuola Norm. Sup. Pisa Cl. Sci. 27 (1998), no. 3-4, 379–425 (1999).
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# Quantum Phase Fluctuations Responsible for Pseudogap
\[
## Abstract
The effect of ordering field phase fluctuations on the normal and superconducting properties of a simple 2D model with a local four-fermion attraction is studied. Neglecting the coupling between the spin and charge degrees of freedom an analytical expression has been obtained for the fermion spectral function as a single integral over a simple function. From this we show that, as the temperature increases through the 2D critical temperature and a nontrivial damping for a phase correlator develops, quantum fluctuations fill the gap in the quasiparticle spectrum. Simultaneously the quasiparticle peaks broaden significantly above the critical temperature, resembling the observed pseudogap behavior in high-$`T_c`$ superconductors.
PACS numbers: 74.25.-q, 74.40.+k, 74.72.-h
Key words: pseudogap, phase fluctuations
\]
The pseudogap, or depletion of the single particle spectral weight around the Fermi level, is the most striking demonstration that cuprate superconductors are not described by the BCS scenario of superconductivity. The pseudogap opens in the normal state as the temperature is lowered below a crossover temperature $`T^{}`$ and extends over a wide range of temperatures in the underdoped cuprates. ARPES and the scanning tunneling spectroscopy (STS) (see Refs. in ) provide particularly complete information about the pseudogap behavior and show a smooth crossover from the pseudo- to superconducting (SC) gap. The transition from SC to normal behavior appears to be driven by phase fluctuations and is well described by the Berezinskii-Kosterlitz-Thouless (BKT) theory of vortex-pair unbinding.
There are currently many possible explanations for the unusual behavior of HTSC. One of these is based on the nearly antiferromagnetic Fermi liquid model . Another explanation, proposed by Anderson, relies on the separation of the spin and charge degrees of freedom. The third approach, which we follow in this paper, relates the observed anomalies to precursor SC fluctuations, and in particular fluctuations in the phase of the complex ordering field, as originally suggested by Emery and Kivelson .
The phase diagram for a simple microscopic 2D model (see (1) below) which formalizes the scenario of has been studied in. The Green’s function (GF) for this model was derived in using the correlator $`\mathrm{exp}(i\theta (x)/2)\mathrm{exp}(i\theta (0)/2)`$ for the phase fluctuations in the classical (static) approximation and neglecting the coupling between the spin and charge degrees of freedom.
The associated spectral function (SF) $`A(\omega ,𝐤)=(1/\pi )\mathrm{Im}G(\omega +i0,𝐤)`$ has also been derived analytically in . Being proportional to the intensity of ARPES , the SF encodes information about the quasiparticles and pseudogap. The result of showed that, while the temperature behavior of the quasiparticle peaks is in correspondence with experiment , the gap in the spectrum remains unfilled. In an earlier paper filling of the gap has been achieved as a result of a Doppler shift in the quasiparticle excitation spectrum. This Doppler shift originated from the semiclassical coupling of the mean field $`d`$-wave quasiparticles to the supercurrents induced by classically fluctuating unbound vortex-antivortex pairs. It was shown in that the shift can be identified with the coupling of the spin and charge degrees of freedom.
The purpose of the present paper is to show analytically that filling of the gap can also result from the quantum (dynamical) phase fluctuations, even when the coupling between charge and spin degrees of freedom is entirely neglected. As pointed out in the mechanism of the gap filling has to be understood yet even for the simplest attractive Hubbard model. Thus we contribute into the discussion as to which mechanism or mechanisms lead to pseudogap filling but within the scenario based on phase fluctuations, by studying the mechanism proposed in . The advantage of our calculation is that the SF is obtained as a single integral with an analytical integrand, no numerical analytical continuation was necessary resulting in far greater accuracy.
Let us consider the continuum version of the 2D attractive Hubbard model defined by the Hamiltonian density:
$$=\psi _\sigma ^{}(x)\left(\frac{^2}{2m}+\mu \right)\psi _\sigma (x)V\psi _{}^{}(x)\psi _{}^{}(x)\psi _{}(x)\psi _{}(x),$$
(1)
where $`x=\text{r},\tau `$ denotes the space and imaginary time variables, $`\psi _\sigma (x)`$ is a fermion field with spin $`\sigma =,`$, $`m`$ is the effective fermion mass, $`\mu `$ is the chemical potential, and $`V`$ is an effective local attraction constant; we take $`\mathrm{}=k_B=1`$. Clearly the Hamiltonian (1) is too simple to be adequate for systems as complex as cuprate HTSC. However, it has proved itself as a very convenient model for both numerical, in particular Monte Carlo simulations, and analytical approaches which does exhibit gap-like behavior above $`T_c`$ (see also Refs. in the review ). Moreover one may use the model to obtain a fully analytic treatment of the pseudogap properties, and apply such results to obtain a better understanding of more complex and less tractable models.
The calculation of the GF is performed in Nambu variables and the fermions are treated as composite objects, comprising both spin and charge parts:
$$\mathrm{\Psi }^{}(x)=\left(\begin{array}{cc}\psi _{}^{}(x)\hfill & \hfill \psi _{}(x)\end{array}\right)=\mathrm{{\rm Y}}^{}(x)\mathrm{exp}[i\tau _3\theta (x)/2],$$
(2)
where $`\mathrm{{\rm Y}}^{}(x)`$ is the Nambu spinor of neutral fermions. Substituting (2) into the standard definition of the GF $`G(x)=\mathrm{\Psi }(x)\mathrm{\Psi }^{}(0)`$ one obtains the GF of the charged (observed) fermions
$$G_{\alpha \beta }(x)=\underset{\alpha ^{},\beta ^{}}{}𝒢_{\alpha ^{}\beta ^{}}(x)(e^{i\tau _3\frac{\theta (x)}{2}})_{\alpha \alpha ^{}}(e^{i\tau _3\frac{\theta (0)}{2}})_{\beta ^{}\beta },$$
(3)
as the product of the GF for the neutral fermions $`𝒢_{\alpha \beta }(x)=\mathrm{{\rm Y}}_\alpha (x)\mathrm{{\rm Y}}_\beta ^{}(0)`$ and the phase correlator $`\mathrm{exp}(i\tau _3\theta (x)/2)\mathrm{exp}(i\tau _3\theta (0)/2)`$. For the frequency-momentum representation of Eq. (3) one has
$`G(i\omega _n,\text{k})`$ $`=`$ $`T{\displaystyle \underset{m=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{d^2p}{(2\pi )^2}\underset{\alpha ,\beta =\pm }{}P_\alpha 𝒢(i\omega _m,\text{p})P_\beta }`$ (5)
$`\times D_{\alpha \beta }(i\omega _ni\omega _m,\text{k}\text{p}),`$
where $`P_\pm =\frac{1}{2}(\widehat{I}\pm \tau _3)`$ are the projectors; $`\widehat{I}`$ and $`\tau _3`$ are unit and Pauli matrices; $`D_{\alpha \beta }(i\mathrm{\Omega }_m,\text{q})`$ is the Fourier transformation of the phase correlator $`D_{\alpha \beta }(x)=\mathrm{exp}(i\alpha \theta (x)/2)\mathrm{exp}(i\beta \theta (0)/2)`$; $`\omega _n=(2n+1)\pi T`$ and $`\mathrm{\Omega }_n=2n\pi T`$ are respectively odd and even Matsubara frequencies. The GF of neutral fermions is taken in the mean-field approximation (see )
$$𝒢(i\omega _n,\text{p})=\frac{i\omega _n\widehat{I}+\tau _3\xi (\text{p})\tau _1\rho }{\omega _n^2+\xi ^2(\text{p})+\rho ^2},\xi (\text{p})=\frac{\text{p}^2}{2m}\mu $$
(6)
with p being a 2D vector and $`\rho \rho (x)`$, where $`\rho (x)`$ is the modulus of a complex ordering field $`\mathrm{\Phi }(x)=\rho (x)\mathrm{exp}[i\theta (x)]`$. Note that $`\mathrm{\Phi }(x)=0`$ for $`T0`$ as it should be in 2D in accordance with the Coleman-Mermin-Wagner-Hohenberg theorem, while the value of $`\rho `$ is allowed to be nonzero. Note also that the GF (5) does not contain the symmetry violating term $`\tau _1`$ which could originate from the terms $`P_\pm 𝒢(i\omega _n,\text{k})P_{}`$ since the correlators $`D_+=D_+=0`$ .
The representations (3), (5) for the fermion GF with decoupled spin and charged degrees of freedom are appropriate when one can neglect the fluctuations of the modulus $`\rho (x)`$ and when the energy of the phase distortions is smaller than the energy gain due to nontrivial $`\rho `$. For the present s-wave model this means that the condition $`\rho T`$ should be satisfied, so that the modulus-phase representation (2) should be useful even for $`T`$ close to $`T_{\mathrm{BKT}}`$ and allows one to study the evolution of the SC gap to the pseudogap. The region of temperatures $`T`$ where the condition $`\rho T`$ is satisfied depends crucially on the relative size of the pseudogap region $`(T^{}T_{\mathrm{BKT}})/T^{}`$ which may be reasonably large in 2D for intermediate coupling . The pseudogap in the current work is the result of the low dimensionality of the system and does not need the existence of preformed local pairs. The nonzero value of $`\rho `$ is due to average local density of Cooper pairs which do not have coherence above $`T_{\mathrm{BKT}}`$.
The representation (5) shows (see the explanation after Eq. (6)) that the GF for the charged fermions is defined by the correlator $`D(x)D_{++}(x)=D_{}(x)`$. The asymptotic form of the correlator at large distances, describing also the temporal decay of correlations, is given by the form
$$D^R(t,\text{r})=\mathrm{exp}(\gamma t)(r/r_0)^{T/8\pi J}\mathrm{exp}(r/\xi _+(T)).$$
(7)
Here $`t`$ is the real time, $`\gamma `$ is a decay constant, $`r_0(2/T)(J/K)^{1/2}`$ is the scale for the algebraic decay of correlations in the SC BKT phase ($`T<T_{\mathrm{BKT}}`$, where $`T_{\mathrm{BKT}}`$ is the temperature of the BKT transition), $`\xi _+(T)`$ is the phase coherence length for $`T>T_{\mathrm{BKT}}`$ ($`\xi _+(TT_{\mathrm{BKT}}^+)\mathrm{}`$). The constants $`J`$ and $`K`$ are the bare (mean-field) superfluid stiffness and compressibility, respectively which have been calculated in .
Previously, using the representation (2), only the classical (static $`\mathrm{\Omega }_n=\gamma =0`$) fluctuations have been considered analytically . The Fourier transform of (7) for this case is
$$D(i\mathrm{\Omega }_n,\text{q})\delta _{n,0}C[\text{q}^2+(1/\xi _+)^2]^\alpha /T,$$
(8)
where $`C4\pi (\mathrm{\Gamma }(\alpha )/\mathrm{\Gamma }(1\alpha ))(2/r_0)^{2\alpha 2}`$ and $`\alpha 1T/16\pi J`$. For $`TT_{\mathrm{BKT}}`$ the value of $`\alpha 1T/32T_{\mathrm{BKT}}`$. The presence of $`\alpha 1`$ in (8) is related to the preexponent factor $`(r/r_0)^{T/8\pi J}`$ in (7). This factor was not included in the analysis of , but the treatment of includes this factor.
One can now extend the analysis to the case of quantum (dynamical) phase fluctuations. We propose the following generalization of (8):
$$D(i\mathrm{\Omega }_n,\text{q})=\frac{C(v^2)^\alpha }{T[v^2\text{q}^2+(v/\xi _+)^2+\mathrm{\Omega }_n^2+2\gamma |\mathrm{\Omega }_n|]^\alpha },$$
(9)
where $`v`$ is the velocity of the Bogolyubov excitations. Recall that in 2D $`v=v_F/\sqrt{2}`$, where $`v_F`$ is the Fermi velocity. The asymptotic form of the retarded GF corresponding to the GF (9) is
$`D^R(t,\text{q})`$ (10)
$`\{\begin{array}{cc}t^{\alpha 1}e^{\gamma t},& v^2q^2>\gamma ^2v^2\xi _+^2\\ t^{\alpha 1}e^{t(\gamma \sqrt{\gamma ^2v^2\xi _+^2v^2q^2})},& v^2q^2<\gamma ^2v^2\xi _+^2\end{array}`$ (13)
$`t+\mathrm{}.`$ (14)
Eq. (9) can be regarded as a convenient (for analytical studies) generalization of the phenomenological dependence (7) which for nonzero $`\gamma `$ includes the decay of the phase correlations due to the presence of free vortices above $`T_{\mathrm{BKT}}`$. One can see that for $`v^2q^2<\gamma ^2v^2\xi _+^2`$ the decay rate is less than $`\gamma `$ and it is minimal for $`q=0`$. This means that, for large distances, phase fluctuations do not feel the pair vortices which have smaller size.
In the present work both $`\gamma `$ and $`\xi _+`$ are phenomenological parameters which can be derived from the theory of the BKT transition. Since in the SC BKT phase the vortices are confined, one can state that there is a critical slowing down of the phase fluctuations when the temperature approaches $`T_{\mathrm{BKT}}`$, i.e. $`\gamma (TT_{\mathrm{BKT}}^+)=0`$. In fact the detailed theory of BKT transition predicts that $`\gamma (T)\xi _+^z(T)`$, where $`z`$ is the dynamical exponent and $`\xi _+^1(T)\xi _0^1\mathrm{exp}(b/\sqrt{TT_{\mathrm{BKT}}})`$, where $`b`$ is a positive constant. In what follows we just restrict ourselves by comparing two cases: $`\gamma 0`$ above $`T_{\mathrm{BKT}}`$ and $`\gamma =0`$ for $`T<T_{\mathrm{BKT}}`$.
We note that there is now experimental evidence for the vortex-pair unbinding (BKT) nature of the SC transition in the Bi-cuprates. We stress also that, despite the rather simple form of the GF, it takes into account the presence of vortices while the self-consistent $`T`$-matrix approximation (see e.g. ) cannot describe them.
The SF associated with GF (5) can readily be expressed in terms of the corresponding SF for the GF (6) and (9) as
$`A(\omega ,\text{k})={\displaystyle \frac{1}{\pi }}\mathrm{Im}G_{11}(\omega +i0,\text{k})={\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}d\omega ^{}[{\displaystyle \frac{1}{1+e^{\omega ^{}/T}}}`$ (15)
$`{\displaystyle \frac{1}{1e^{(\omega ^{}\omega )/T}}}]{\displaystyle }{\displaystyle \frac{d^2p}{(2\pi )^2}}a_F(\omega ^{},\text{p})a_B(\omega \omega ^{},\text{k}\text{p}).`$ (16)
where these SF are
$`a_F(\omega ,\text{p})=(\omega +\xi (\text{p}))\delta (\omega ^2\xi ^2(\text{p})\rho ^2)\text{sgn}\omega ,`$ (17)
$`a_B(\mathrm{\Omega },\text{q})={\displaystyle \frac{1}{\pi }}\text{Im}D^R(\mathrm{\Omega },\text{q}).`$ (18)
The $`\delta `$-function in $`a_F`$ allows one to perform the angular integration in (16). Finally integrating over momentum the SF can be expressed as a single integral:
$`A(\omega ,𝐤)={\displaystyle \frac{1}{2\pi ^2T}}{\displaystyle \frac{\mathrm{\Gamma }(\alpha )}{\mathrm{\Gamma }(1\alpha )}}\left({\displaystyle \frac{2}{mr_0^2}}\right)^{\alpha 1}{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}𝑑\omega ^{}\mathrm{sgn}(\omega ^{}+\omega )`$ (19)
$`\times \left[{\displaystyle \frac{1}{1e^{\omega ^{}/T}}}{\displaystyle \frac{1}{1+e^{(\omega ^{}+\omega )/T}}}\right]{\displaystyle \frac{\theta [(\omega ^{}+\omega )^2\rho ^2]}{\sqrt{(\omega ^{}+\omega )^2\rho ^2}}}`$ (20)
$`\times \{[I(\omega ,𝐤,\omega ^{})(\omega ^{}+\omega +\sqrt{(\omega ^{}+\omega )^2\rho ^2})`$ (21)
$`\times \theta (\mu +\sqrt{(\omega ^{}+\omega )^2\rho ^2})]`$ (22)
$`+[\sqrt{(\omega ^{}+\omega )^2\rho ^2}\sqrt{(\omega ^{}+\omega )^2\rho ^2}]\},`$ (23)
where we have introduced a function $`I(\omega ,𝐤,\omega ^{})`$:
$`I(\omega ,𝐤,\omega ^{})=`$ (24)
$`\pi \mathrm{Im}\left[(x_{}aib)^\alpha F({\displaystyle \frac{1}{2}},\alpha ;1;{\displaystyle \frac{x_+x_{}}{a+ibx_{}}})\right],`$ (25)
$`a={\displaystyle \frac{1}{2m}}\left({\displaystyle \frac{\omega ^{\mathrm{\hspace{0.25em}2}}}{v^2}}\xi _+^2\right),b={\displaystyle \frac{\gamma \omega ^{}}{mv^2}},`$ (26)
$`x_\pm =\left(\sqrt{{\displaystyle \frac{k^2}{2m}}}\pm \sqrt{\mu +\sqrt{(\omega +\omega ^{})^2\rho ^2}}\right)^2.`$ (27)
For $`\gamma =0`$ the expression (23) can be rewritten in the following form
$`A(\omega ,𝐤)={\displaystyle \frac{1}{2\pi ^2T}}{\displaystyle \frac{\mathrm{\Gamma }(\alpha )}{\mathrm{\Gamma }(1\alpha )}}\left({\displaystyle \frac{2}{mr_0^2}}\right)^{\alpha 1}{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}𝑑\omega ^{}\mathrm{sgn}\omega ^{}`$ (28)
$`\times \mathrm{sgn}(\omega ^{}+\omega )\left[{\displaystyle \frac{1}{1e^{\omega ^{}/T}}}{\displaystyle \frac{1}{1+e^{(\omega ^{}+\omega )/T}}}\right]`$ (29)
$`\times {\displaystyle \frac{\theta [(\omega ^{}+\omega )^2\rho ^2]}{\sqrt{(\omega ^{}+\omega )^2\rho ^2}}}\theta (\omega ^2v^2\xi _+^2)`$ (30)
$`\times \{[\pi (ax_{})^\alpha F({\displaystyle \frac{1}{2}},\alpha ;1;{\displaystyle \frac{x_+x_{}}{ax_{}}})\theta (ax_+)`$ (31)
$`+(x_+x_{})^{1/2}(ax_{})^{1/2\alpha }B({\displaystyle \frac{1}{2}},1\alpha )`$ (32)
$`\times F({\displaystyle \frac{1}{2}},{\displaystyle \frac{1}{2}};{\displaystyle \frac{3}{2}}\alpha ;{\displaystyle \frac{ax_{}}{x_+x_{}}})\theta (x_+a)\theta (ax_{})]`$ (33)
$`\times \left(\omega ^{}+\omega +\sqrt{(\omega ^{}+\omega )^2\rho ^2}\right)\theta (\mu +\sqrt{(\omega ^{}+\omega )^2\rho ^2})`$ (34)
$`+(\sqrt{(\omega ^{}+\omega )^2\rho ^2}\sqrt{(\omega ^{}+\omega )^2\rho ^2})\}`$ (35)
which describes the case of the non-damped dynamical phase fluctuations. While Eq. (35) appears to be more complicated than the more general Eq. (23), it proves useful for making a general statement about the gap filling.
Expressions (23) and (35) present the main result of the paper for the SF in the case of quantum phase fluctuations. Before proceeding to the discussion of their numerical integration, we recap briefly the case of classical phase fluctuations. The static SF, obtained in , is reproduced from (35) after changing the variable $`\omega ^{}v\omega ^{}`$ and taking formally the limit $`v0`$. It was discussed in detail in and we would like to stress here only two main points. The first is that the SF is identically zero inside the gap ($`A_{cl}(\omega ,\text{k})=0`$ for $`|\omega |<\rho `$) . The second is that in addition to the usual quasiparticle peaks it reveals extra peaks at $`\omega =\pm \rho `$.
Let us now return to the discussion of the results of numerical integration based on Eqs. (23) and (35) which are shown in Fig. 1.
a) For $`T<T_{\mathrm{BKT}}`$ there are two highly pronounced quasiparticle peaks at $`\omega =\pm E(\text{k})\pm \sqrt{\xi ^2(\text{k})+\rho ^2}`$. Since below $`T_{\mathrm{BKT}}`$ both $`\xi _+^1`$ and $`\gamma `$ are zero, the width of the peaks is almost entirely controlled by $`\alpha `$ whose deviation from 1 gives the non-Fermi liquid behavior of the GF . Since the value of $`\alpha `$ is very close to 1 this non-Fermi liquid behavior is hardly distinguishable from ordinary widening of the quasiparticle peaks due to damping. Furthermore, because $`\alpha 1`$ as $`T0`$ the width of the peaks is temperature dependent so that the peaks become sharper as $`T`$ decreases resembling the data of ARPES for the anti-node direction.
We note also that $`kk_F`$ in Fig. 1 so that the quasiparticle peaks are not equal to each other in contrast to the case of symmetrized ARPES data . The reason for choosing $`kk_F`$ in Fig. 1 was to prove that the extra peaks present in the SF calculated for the classical phase fluctuations are now absent.
The gap in the SF remains almost unfilled and has “U”-like shape. We stress, however, that in contrast to the static case the SF (23) is nonzero even for $`|\omega |<\rho `$ and this is evidently related to the presence of the quantum (dynamical) phase fluctuations, i.e. the terms with $`\mathrm{\Omega }_n0`$. In particular, using Eq. (35) one can check analytically that even for $`\gamma =0`$ that due to these terms $`A(\omega =0,𝐤=𝐤_F)0`$. We estimated the temperature $`T_{cl}`$ of the quantum to classical crossover in the way similar to that of . For the low carrier densities this gives $`T_{cl}T_{\mathrm{BKT}}`$, while for the high carrier densities, where the pseudogap phase shrinks we obtain $`T_{cl}\rho (T=0)T^{}`$. Thus, the crossover temperature is always too large (in the absence of dissipation) for classical fluctuations to play a significant role at low $`T`$.
b) For $`T>T_{\mathrm{BKT}}`$ the quasiparticle peaks at $`\omega \pm E(\text{k})`$ are far less prominent. This is caused by the fact that $`\xi _+`$ is now nonzero due to the presence of vortices. Increasing the temperature further decreases the value of $`\xi _+`$ so that the quasiparticle peaks become even less pronounced. This behavior of the quasiparticle peaks reproduces qualitatively the ARPES studies for the anti-node direction which show clearly that the quasiparticle SF broadens dramatically when passing from the SC to normal state. The width of the quasiparticle peaks is controlled primarily by $`\xi _+`$ above $`T_{\mathrm{BKT}}`$ since this width remains practically constant as $`\gamma `$ changes. Since $`\xi _+`$ is the phase coherence length, the current model supports the premise that the quasiparticle peaks are related to the phase coherence, not to the energy gap. The deviation of $`\alpha `$ from 1 becomes, however, sizable at $`TT_{BKT}`$.
c) As discussed above, for $`T>T_{\mathrm{BKT}}`$ the value of $`\gamma `$ is nonzero and increases with increasing temperature. This increase of $`\gamma `$ (typically up to values of the order of $`0.5\mu `$), together with the decrease of $`\xi _+`$, causes filling of the gap and changes its shape from “U”-like to “V”-like. In other words, the quasiparticle peaks grow “shoulders” which eventually fill the gap. We stress, however, that the increase of $`\gamma `$ in fact diminishes the effect of the quantum phase fluctuations pushing $`T_{cl}`$ down which is seen from the reducing of $`A(\omega =0,𝐤=𝐤_F)`$ as $`\gamma `$ grows (see Fig. 1, where the bottom parts of the curves cross each other), but the abovementioned “shoulders” simultaneously emerge. Thus the filling of the gap at $`T>T_{\mathrm{BKT}}`$ in the present model is caused by quantum phase fluctuations in the presence of nonzero dissipation because the damping $`\gamma `$ contributes only when $`\mathrm{\Omega }_n0`$.
d) Due to the smooth dependence of $`\xi _+^1`$ and $`\gamma `$ on $`T`$ as the temperature passes through $`T_{\mathrm{BKT}}`$ there is no sharp transition at $`T_{\mathrm{BKT}}`$ in agreement with experiment . A similar filling of the gap was obtained in for $`\gamma =0.5\mu `$ where the correlator (7) was used for the numerical computation of the self-energy of the fermions and the subsequent extraction of the SF from the fermion GF. A recent Monte Carlo simulation also shows similar behavior for the quasiparticle peaks and the filling of the gap. However it is the analytical character of the present work, which relies on the explicit introduction of the charge and spin degrees of freedom (2) for the Nambu spinors, that enables us to unambiguously state the correspondence between the parameters of the model and the observed features of the SF (23).
In spite of some similarities between the results observed and those just obtained, the latter are only illustrative since we have considered a model with non-retarded $`s`$-wave attraction. However, it is likely that for $`d`$-wave pairing, the properties obtained can be used for the description of the systems in the anti-node direction on the Fermi surface. It is however essential to consider fluctuations in the modulus of the order parameter to extend the analysis to the nodal directions on the Fermi surface. The value of $`\gamma =0.5\mu `$ which results in a filled gap appears to be too large to be due to the vortex-vortex interaction. This leads one to the conclusion that the mechanism considered here for gap filling may well not be the only possible mechanism and that other interactions which lead to the filling of the gap above $`T_c`$, see , are also important.
In summary we studied the effect of the fluctuations in the phase of the complex ordering field on the properties of a 2D system with four-fermion attraction. The fermion SF has been given as a single integral over a function known in closed form. Through the use of analytical techniques, we have been able to demonstrate that the quantum phase fluctuations in the presence of dissipation lead to the filling of the pseudogap even if one ignores the spin-charge coupling.
S.G.Sh thanks Prof. H. Beck and Dr. M. Capezzali for stimulating discussion. This work was partly supported by the NRF, South Africa (R.M.Q and S.G.Sh), by SCOPES-project 7UKPJ062150.00/1 (V.P.G., V.M.L. and S.G.Sh) and by the research project 2000-061901.00/1 (S.G.Sh) of the Swiss National Science Foundation.
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# Spin-lattice relaxation in the mixed state of the high-𝑇_𝑐 cuprates: electronic spin-flip scattering versus spin-fluctuations
## A
The general form for $`1/T_1T`$ in the presence of a supercurrent is given by
$`{\displaystyle \frac{1}{T_1T}}`$ $`=`$ $`{\displaystyle \frac{A}{N}}\left({\displaystyle \frac{1}{v_Fv_\mathrm{\Delta }}}\right)^2{\displaystyle \underset{n,m}{}}\{(𝐪_{n,m})G_1(D_m,D_n,T)`$ (A2)
$`+(𝐪_{n,m+2})G_2(D_m,D_n,T)\}`$
where $`D_m=𝐯_F^{(m)}𝐩_s`$, $`𝐯_F^{(m)}`$ is the Fermi velocity at the node in the $`m`$’th quadrant of the Brillouin zone, $`v_\mathrm{\Delta }=|\mathrm{\Delta }_𝐤/𝐤|`$ at the nodes, and $`𝐪_{n,m}`$ is the momentum connecting nodes $`n`$ and $`m`$. For ESF scattering we find
$`A`$ $`=`$ $`k_B(\mathrm{}^2\gamma _n\gamma _e)^2/(2\mathrm{});`$ (A3)
$`(𝐪_{n,m})`$ $`=`$ $`F_{O(2,3)}(𝐪_{n,m}),`$ (A4)
while for the ASF mechanism on has
$`A`$ $`=`$ $`(\alpha g)^2k_B(\mathrm{}^2\gamma _n\gamma _e)^2/(2\mathrm{});`$ (A5)
$`(𝐪_{n,m})`$ $`=`$ $`{\displaystyle \frac{F_{O(2,3)}(𝐪_{n,m})}{(\xi ^2+(𝐪_{n,m}𝐐)^2)^2}}.`$ (A6)
In both cases, $`G_{1,2}`$ is given by
$`G_1(D_m,D_n,T)`$ $`=`$ $`{\displaystyle \frac{D_nD_m+ϵ^2}{T^2}}n_F(ϵ)+{\displaystyle \frac{\pi ^2}{6}}`$ (A9)
$`+{\displaystyle \frac{D_n+D_m}{T}}\left({\displaystyle \frac{ϵ[1n_F(ϵ)]}{T}}+\mathrm{ln}[n_F(ϵ)]\right)`$
$`2{\displaystyle _0^{ϵ/T}}𝑑xxn_F(x)`$
where $`n_F(ϵ)=[\mathrm{exp}(ϵ/T)+1]^1`$ is the Fermi function, $`ϵ=max\{D_n,D_m\}`$, and
$`G_2(D_m,D_n,T)`$ $`=`$ $`{\displaystyle \frac{D_nD_m}{T^2}}\left[n_F(D_m)n_F(D_n)\right]`$ (A14)
$`+{\displaystyle \frac{D_nD_m}{T}}\{{\displaystyle \frac{D_n[n_F(D_n)1]}{T}}+\mathrm{ln}[n_F(D_n)]`$
$`+{\displaystyle \frac{D_m[n_F(D_m)1]}{T}}\mathrm{ln}[n_F(D_m)]\}`$
$`+\left({\displaystyle \frac{D_n}{T}}\right)^2n_F(D_n)\left({\displaystyle \frac{D_m}{T}}\right)^2n_F(D_m)`$
$`2{\displaystyle _{D_m/T}^{D_n/T}}𝑑xxn_F(x)\mathrm{for}D_n+D_m0`$
$`G_2(D_m,D_n,T)`$ $``$ $`0\mathrm{for}D_n+D_m>0.`$ (A15)
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# Bound state equation for 4 or more relativistic particles.
## 1 Introduction.
The Bethe-Salpeter equation is the usual tool for computing relativistic bound states \[1-20\]. The principal difficulty of this equation comes from the presence of N-1 (for N particles) unphysical degrees of freedom: the relative time-energy degrees of freedom. In a recent work , we built a 3D reduction of the two-fermion Bethe-Salpeter equation around an unspecified instantaneous approximation of the Bethe-Salpeter kernel, the difference with the exact kernel being corrected in a series of higher-order contributions to the final 3D potential. This potential being not explicitly hermitian, we performed a second symmetrizing series expansion at the 3D level. The result of the combination of both series gave us an hermitian potential which turned out to be independent of the starting instantaneous approximation, and was in fact a compact expression of the potential that Phillips and Wallace compute order by order. We found that this 3D potential was also obtainable directly by starting with an approximation of the free propagator, based on integrals in the relative energies instead of the more usual $`\delta `$constraint (integrating propagator-based reduction). In this form, the method was easily generalizable to a system of N particles, consisting of any combination of bosons and fermions. Special cases of N=2 or 3-fermion systems were examined. Taking the retarded part of the full propagator at equal times, following Logunov and Tavkhelidze (for two particles) and Kvinikhidze and Stoyanov (for three particles) also leads to the same 3D equations.
For more detailed explanations about our integrating propagator-based reduction, about related problems like Lorentz covariance, cluster separability or continuum dissolution, or for comparisons with other approaches, we refer to our preceding work .
In the present work we test our reduction method in the four-body problem. The first step is of course the writing of the starting Bethe-Salpeter equation itself. It is an homogeneous equation for a Bethe-Salpeter amplitude, which is derived by identifying the residues of the bound state poles in the inhomogeneous Bethe-Salpeter equation for the full propagator $`G.`$ The interactions are introduced via the Bethe-Salpeter kernel $`K,`$ which is such that the inhomogeneous Bethe-Salpeter equation reproduces, by iterations, the full propagator $`G`$ as deduced by the Feynman’s graphs method. The writing of the Bethe-Salpeter kernel $`K`$ is straightforward in the two and three-body problems. In the four-body problem, however, we meet a new difficulty due to the mutual unconnectedness of the three pairs of two-body irreducible kernels (12)(34), (13)(24) and (14)(23): adding simply the six two-body irreducible kernels (plus the four irreducible three-body kernels and the irreducible four-body kernel) would lead to an overcounting of some graphs in the expansion of $`G.`$ We would get for example the sequences of kernels (12)(34) and (34)(12), which are in fact a unique graph or part of graph (unlike sequences like (12)(23) and (23)(12) which are different). It is however not very difficult to correct this overcounting by including three counterterms in $`K.`$ Computing then the first terms of the 3D potential by our 3D reduction method, we find that, at the often used approximation in which we keep only the two-body kernels, the (relative and total) energy-independent part of these kernels and the positive-energy part of the propagators, the final 3D potential is simply the sum of the six two-particle potentials.
It is clear that the difficulty of the overcounting of the sequences of unconnected kernels will not be so easily overcomed for $`N5.`$ This suggested us to modify our 3D reduction method in order to avoid the explicit writing of the Bethe-Salpeter kernel $`K.`$ Our 3D potential is now written directly in a straightforward way for any number of particles, with simply the additional prescription of removing the duplicating diagrams which appear when $`N4.`$
In section 2, we recall our 3D reduction method of for the N-body Bethe-Salpeter equation. In section 3 we apply it to the four-body problem. In section 4 we suggest a general 3D reduction method starting directly from the expansion of the full propagator $`G,`$ as given by Feynman graphs, without the need of writing explicitly the kernel $`K`$ of the Bethe-Salpeter equation. Section 5 is devoted to conclusions.
## 2 Integrating propagator-based reduction for the N-body Bethe-Salpeter equation.
The inhomogeneous and homogeneous Bethe-Salpeter equations can be written
$$G=G^0+G^0KG$$
(1)
$$\mathrm{\Phi }=G^0K\mathrm{\Phi }$$
(2)
where $`G`$ is the full propagator and $`\mathrm{\Phi }`$ the Bethe-Salpeter amplitude. The free propagator $`G^0`$ for a system of $`f=0,1\mathrm{}N`$ fermions and b = N-f bosons is
$$G^0=G_1^0\mathrm{}G_f^0G_{f+1}^0\mathrm{}G_N^0,$$
(3)
the propagators of the fermion i and the boson j being respectively
$$G_i^0=\frac{1}{p_{i0}h_i+iϵh_i}\beta _i,$$
(4)
$$G_j^0=\frac{1}{p_{j0}^2E_j^2+iϵ}=\frac{1}{2E_j}\underset{\sigma _j}{}\frac{\sigma _j}{p_{j0}\sigma _jE_j+iϵ\sigma _j},$$
(5)
where $`p_i`$ is the 4-momentum of particle i , and
$$h_i=\stackrel{}{\alpha }_i.\stackrel{}{p}_i+\beta _im_i,E_i=(\stackrel{}{p}_i^{\mathrm{\hspace{0.17em}2}}+m_i^2)^{\frac{1}{2}},\sigma _j=\pm 1.$$
(6)
We do not specify the reference frame in which we write noncovariant quantities like $`h_i`$ or $`E_i.`$ Our 3D reduction will in fact be frame-dependent. Practically, we shall choose the global rest frame. The kernel $`K`$ will be chosen in such a way that equation (1) gives by iterations the usual expansion of $`G`$ in terms of Feynman graphs. It is related to the kernel $`K^{}`$ defined in the same way with the dressed propagators
$$G=G^0+G^0K^{}G$$
(7)
where $`G^0`$ is the product of the dressed propagators
$$G_i^0=\frac{1}{\gamma _ip_im_i\mathrm{\Sigma }_i+iϵ}\text{for a fermion,}$$
(8)
$$G_i^0=\frac{1}{p_{i0}^2E_i^2\mathrm{\Sigma }_i+iϵ}\text{for a boson,}$$
(9)
and $`\mathrm{\Sigma }_i`$ the renormalized self-energy function. The transfer of the self-energies from the propagator to the kernel is achieved with
$$K=K^{}+\mathrm{\Sigma },\mathrm{\Sigma }=(G^0)^1(G^0)^1.$$
(10)
In the two-body problem $`K^{}`$ is simply the sum of all irreducible Feynman graphs. In the three-body problem, it is
$$K^{}=K_{12}^{}(G_3^0)^1+K_{23}^{}(G_1^0)^1+K_{31}^{}(G_2^0)^1+K_{123}^{},$$
(11)
where $`K_{ij}^{}`$ is the sum of the two-body (ij) irreducible Feynman graphs while $`K_{123}^{}`$ is the sum of the three-body connected irreducible Feynman graphs. In the four-body problem and beyond, there appear mutually unconnected kernels in the expression of $`K^{}`$ (like $`K_{12}^{}`$ and $`K_{34}^{}`$ for example), and we shall have to be careful to avoid overcountings in the expansion of $`G`$ (see next section).
In ref. , we built a 3D reduction of the two-fermion Bethe-Salpeter equation around an instantaneous approximation of the Bethe-Salpeter kernel, the difference with the exact kernel being corrected in a series of higher-order contributions to the final 3D potential. This potential being not explicitly hermitian, we performed a second symmetrizing series expansion at the 3D level. The result of the combination of both series turned out to be independent of the starting instantaneous approximation, and was also obtainable directly by starting with an approximation of the free propagator, based on integrals in the relative energies instead of the more usual $`\delta `$constraints. Furthermore, the method was easily generalizable to a system of N particles, consisting of any combination of bosons and fermions. In the present work, we shall however present our 3D reduction in yet another way. Let us first define the 3D free propagator
$$𝑑p_0G^0(p_0)\delta (P_0\underset{i=1}{\overset{N}{}}p_{i0})𝑑p_{10}\mathrm{}𝑑p_{N0}G^0(p_{10},\mathrm{}p_{N0})$$
$$=\frac{(2i\pi )^{N1}}{\omega }\tau g^0\beta $$
(12)
with
$$g^0=\frac{1}{P_0S+iϵP_0},S=E(\mathrm{\Lambda }^+\mathrm{\Lambda }^{}),\beta =\beta _1\mathrm{}\beta _f,$$
(13)
$$\mathrm{\Lambda }^\pm =\mathrm{\Lambda }_1^\pm \mathrm{}\mathrm{\Lambda }_f^\pm ,\mathrm{\Lambda }_i^\pm =\frac{E_i\pm h_i}{2E_i},\tau =\mathrm{\Lambda }^++()^{f+1}\mathrm{\Lambda }^{},$$
(14)
$$E=\underset{i=1}{\overset{N}{}}E_i,\omega =\mathrm{\hspace{0.17em}2}^bE_{f+1}\mathrm{}E_N$$
(15)
for $`f0,`$ so that
$$\tau g^0=\frac{\mathrm{\Lambda }^+}{P_0E+iϵ}+()^{f+1}\frac{\mathrm{\Lambda }^{}}{P_0+Eiϵ}.$$
(16)
When $`f=0,`$ (bosons only), we have no $`\tau `$ and no $`\beta `$ (we can replace them by 1 in (12)) and
$$g^0=\frac{1}{P_0E+iϵ}\frac{1}{P_0+Eiϵ}=\frac{2E}{P_0^2E^2+iϵ}.$$
(17)
Let us now define a 3D full propagator in the $`\tau ^2`$ subspace, with $`\tau g^0`$ as free limit:
$$g=\frac{1}{(2i\pi )^{N1}}\tau \sqrt{\omega }𝑑p_0^{}𝑑p_0G(p_0^{},p_0)\beta \tau \sqrt{\omega }.$$
(18)
The integrations with respect to the relative energies preserve the positions of the bound state poles of $`G`$ (it can also be shown that the physical particle-particle, particle-bound state or bound state-bound state scattering amplitudes are also preserved ). If we write $`G`$ as
$$G=G^0+G^0TG^0,T=K(\mathrm{\hspace{0.17em}1}G^0K)^1$$
(19)
we get
$$g=\tau g^0+\tau g^0<T>\tau g^0$$
(20)
with the definition
$$<A>=\frac{1}{(2i\pi )^{N1}}\frac{\tau ^2\sqrt{\omega }}{g^0}𝑑p_0^{}𝑑p_0G^0(p_0^{})A(p_0^{},p_0)G^0(p_0)\beta \frac{\tau ^2\sqrt{\omega }}{g^0}.$$
(21)
To the 3D full propagator $`g`$ we can associate the 3D equation
$$\varphi =g^0\tau V\varphi $$
(22)
where $`V`$ is such that
$$g^1=\tau (g^0)^1V.$$
(23)
The comparison with (20) gives
$$<T>=V(1\tau g^0V)^1$$
(24)
or, conversely
$$V=<T>(1+\tau g^0<T>)^1=<K(1G^0K)^1>(\mathrm{\hspace{0.17em}1}+\tau g^0<K(1G^0K)^1>)^1$$
$$=<K(1G^0K)^1(\mathrm{\hspace{0.17em}1}+>\tau g^0<K(1G^0K)^1)^1>$$
$$=<K(\mathrm{\hspace{0.17em}1}G^0K+>\tau g^0<K)^1>=<K(\mathrm{\hspace{0.17em}1}G^RK)^1>=<K^T>$$
(25)
with
$$K^T=K(1G^RK)^1=K+KG^RK+\mathrm{}$$
(26)
$$G^R=G^0G^I,G^I=>\tau g^0<$$
(27)
or, making the dependence on the relative energies explicit:
$$G^R(p_0^{},p_0)=G^0(p_0)\delta (p_0^{}p_0)G^I(p_0^{},p_0)$$
(28)
with
$$G^I(p_0^{},p_0)=G^0(p_0^{})\beta \frac{\tau \omega }{(2i\pi )^{N1}g^0}G^0(p_0).$$
(29)
By substracting $`G^I`$ from $`G^0,`$ we remove the leading terms which come from the residues of the positive-energy poles of $`G^0.`$ This ensures in principle the decreasing of the terms of the series (26).
The choice (18) of $`g`$ as the integral of $`G`$ between two $`\tau `$ operators defines our 3D reduction (in the two-body case it is also Phillips and Wallace’s reduction). A different approach, in the two-body case, the constraint approach, consists in defining $`<T>`$ by fixing the relative-energy arguments of $`T(p_0^{},p_0),`$ for example by putting one initial particle and one final particle on their positive-energy mass shell . There is no problem with this approach for $`N=2.`$ For $`N=3,`$ we should put two initial and two final particles on their positive-energy mass shell. For the two-body terms, this would give a constraint too much. This difficulty can be avoided by applying different constraints to the different terms of $`T`$ . This leads to a set of coupled equations which can not, without approximations, be reduced to a single equation.
In , we derived equation (22) from the inhomogeneous Bethe-Salpeter equation (2), with the definition
$$\varphi =\tau \sqrt{\omega }𝑑p_0\mathrm{\Phi }(p_0)$$
(30)
A simpler reduction method, which we shall adopt from now, can be obtained by keeping only the positive-energy part $`\mathrm{\Lambda }^+`$ of $`\tau .`$ In this case, the equation obtained after the reduction of the Dirac spinors into Pauli spinors will be
$$(P_0E)\phi =\left[\underset{i=1}{\overset{f}{}}\sqrt{\frac{2E_i}{E_i+m}}\right]v\left[\underset{i=1}{\overset{f}{}}\sqrt{\frac{2E_i}{E_i+m}}\right]\phi .$$
(31)
where $`v`$ is the large-large part of $`V.`$ This simpler reduction method can be defined by the choice of $`g`$ as the integral of $`G,`$ now between two positive-energy projectors $`\mathrm{\Lambda }^+`$ (instead of $`\tau `$ in (18)). This 3D full propagator is, in configuration space, the retarded part of $`G`$ taken at equal times. It has been the starting point of a 3D reduction by Logunov and Tavkhelidze (in the two-fermion case), followed by Kvinikhidze and Stoyanov (three fermions) and Khvedelidze and Kvinikhidze (four fermions).
## 3 The four-body problem.
In this section 3, we shall neglect, for simplicity, the radiative corrections (omitting thus the $`\mathrm{\Sigma }^{}s`$ and all the primes in the expression of $`K).`$
### 3.1 Bethe-Salpeter equation for four particles.
For four particles, we have
$$G^0=G_1^0G_2^0G_3^0G_4^0$$
(32)
$$K=K_{12,34}+K_{13,24}+K_{14,23}$$
$$+K_{123}(G_4^0)^1+K_{124}(G_3^0)^1+K_{134}(G_2^0)^1+K_{234}(G_1^0)^1$$
$$+K_{1234},$$
(33)
with
$$K_{12,34}=K_{12}(G_3^0G_4^0)^1+K_{34}(G_1^0G_2^0)^1K_{12}K_{34}$$
(34)
and similarly for $`K_{13,24}`$ and $`K_{14,23}.`$ The three lines of (33) contain the two, three and four-body kernels respectively. The last term of (34) is a counterterm which has no equivalent in the two and three-body problems. If we write the expansion of $`G`$
$$G=G^0+G^0KG^0+G^0KG^0KG^0+\mathrm{}$$
(35)
and collect the terms containing one $`K_{12}`$ with one $`K_{34},`$ we get indeed
$$G=\mathrm{}G^0K_{12}K_{34}G^0+G^0K_{12}(G_3^0G_4^0)^1G^0K_{34}(G_1^0G_2^0)^1G^0$$
$$+G^0K_{34}(G_1^0G_2^0)^1G^0K_{12}(G_3^0G_4^0)^1G^0+\mathrm{}$$
(36)
Since $`K_{12}`$ and $`K_{34}`$ commute mutually, the two last terms of (36) correspond to the same graph which appears thus twice, while the first term is again the same graph with a minus sign. This mechanism works at all orders. The kernel $`K_{12,34}`$ is indeed given by
$$K_{12,34}=(G^0)^1(G_{12,34})^1$$
(37)
where $`G_{12,34}`$ is the full propagator obtained by keeping only the interactions $`K_{12}`$ and $`K_{34}.`$ It factorizes into
$$(G_{12,34})^1=(G_{12})^1(G_{34})^1=\left[(G_1^0G_2^0)^1K_{12}\right]\left[(G_3^0G_4^0)^1K_{34}\right]$$
(38)
so that (37) gives (34). This problem of commutating kernels was not present in the two and three-body problems. It has been easily overcomed here in the four-body problem. It will become much more complicated in the five-body problem and beyond .
### 3.2 3D reduction of the four-body Bethe-Salpeter equation.
For definiteness we shall work on the Bethe-Salpeter equation for four fermions and perform the reduction based on the positive-energy part of $`g^0.`$ Then:
$$\mathrm{\Lambda }^+g^0=\frac{\mathrm{\Lambda }^+}{P_0E+iϵ}$$
(39)
$$<A>=\frac{1}{(2i\pi )^3}\mathrm{\Lambda }^+(P_0E)𝑑p_0^{}𝑑p_0G^0(p_0^{})A(p_0^{},p_0)G^0(p_0)\beta \mathrm{\Lambda }^+(P_0E)$$
(40)
and the 3D equation will be
$$\varphi =g^0V\varphi $$
(41)
with
$$V=<K^T>=<K>+<KG^RK>+\mathrm{}$$
$$=<K>+<K(G^0>g^0<)K>+\mathrm{}$$
$$=<K>+<KG^0K><K>g^0<K>+\mathrm{}$$
(42)
We shall now compute the two and four-vertex terms, keeping only the two-fermion kernels:
$$<K^T>^{(4)}=<K_{12,34}>+\mathrm{}(3terms)$$
$$+\left[<K_{12,34}G^0K_{14,23}><K_{12,34}>g^0<K_{14,23}>\right]+\mathrm{}(6terms)$$
$$+\left[<K_{12,34}G^0K_{12,34}><K_{12,34}>g^0<K_{12,34}>\right]+\mathrm{}(3terms)$$
(43)
where
$$<K_{12,34}>=<K_{12}(G_3^0G_4^0)^1>+<K_{34}(G_1^0G_2^0)^1><K_{12}K_{34}>$$
(44)
$$<K_{12,34}G^0K_{14,23}>^{(4)}=<K_{12}(G_3^0G_4^0)^1G^0K_{23}(G_1^0G_4^0)^1>$$
$$+<K_{12}(G_3^0G_4^0)^1G^0K_{14}(G_2^0G_3^0)^1>+<K_{34}(G_1^0G_2^0)^1G^0K_{23}(G_1^0G_4^0)^1>$$
$$+<K_{34}(G_1^0G_2^0)^1G^0K_{14}(G_2^0G_3^0)^1>$$
(45)
$$<K_{12,34}G^0K_{12,34}>^{(4)}=<K_{12}(G_3^0G_4^0)^1G^0K_{12}(G_3^0G_4^0)^1>$$
$$+<K_{34}(G_1^0G_2^0)^1G^0K_{34}(G_1^0G_2^0)^1>+<K_{12}(G_3^0G_4^0)^1G^0K_{34}(G_1^0G_2^0)^1>$$
$$+<K_{34}(G_1^0G_2^0)^1G^0K_{12}(G_3^0G_4^0)^1>$$
(46)
The various terms of (44)-(46) are represented in figure 1. The two last terms of (46) are equal to $`<K_{12}K_{34}>.`$ In the expansion of $`G`$ this term will thus appear only once, with a plus sign.
Let us now examine the counterterms $`<K>g^0<K>,`$ which could be represented by the 8 last graphs of figure 1, with a vertical line separating them into two parts. In the two and three-particle problems we have a cancellation between the corresponding contributions of $`<KG^0K>`$ and of $`<K>g^0<K>,`$ when the two-body kernels are approximated by (relative and total) energy independent kernels and the propagators by positive-energy propagators. This cancellation of the leading terms leads in principle to a decreasing of the contributions of the higher-order terms of the series $`<K>+<KG^RK>+\mathrm{}.`$ We shall check this cancellation here in (43) by assuming that the two-body kernels contain positive-energy projectors and are independent on the energies ($`K_{ij}=\mathrm{\Lambda }_i^+\mathrm{\Lambda }_j^+K_{ij}\beta _i\beta _j\mathrm{\Lambda }_i^+\mathrm{\Lambda }_j^+`$ and is independent of $`p_{0ij},P_{0ij}`$ – we shall speak of the ”positive-energy instantaneous approximation”, although ”instantaneous” usually refers only to the independence on the two-body relative energy). It is then easy to verify that the cancellation occurs for the 4th to the 9th graphs of figure 1. We remain thus with
$$<K^T>^{(4)}=<K_{12}(G_3^0G_4^0)^1>+<K_{34}(G_1^0G_2^0)^1>+\mathrm{}$$
$$+<K_{12}K_{34}><K_{12}(G_3^0G_4^0)^1>g^0<K_{34}(G_1^0G_2^0)^1>$$
$$<K_{34}(G_1^0G_2^0)^1>g^0<K_{12}(G_3^0G_4^0)^1>+\mathrm{}$$
(47)
The two-vertex contributions become simply the sum of the six two-fermion potentials:
$$V=V_{12}+V_{34}+V_{13}+V_{24}+V_{14}+V_{23},$$
(48)
$$V_{ij}=2i\pi \beta _i\beta _jK_{ij}.$$
(49)
For $`<K_{12}K_{34}>,`$ we have
$$<K_{12}K_{34}>=\frac{1}{2i\pi }𝑑P_{120}𝑑P_{340}\delta (P_0P_{120}P_{340})$$
$$(P_0E^{})\left[\frac{1}{P_{120}E_1^{}E_2^{}+iϵ}V_{12}\frac{1}{P_{120}E_1E_2+iϵ}\right]$$
$$\left[\frac{1}{P_{340}E_3^{}E_4^{}+iϵ}V_{34}\frac{1}{P_{340}E_3E_4+iϵ}\right](P_0E).$$
(50)
Let us perform the integration with respect to $`P_{120}`$ and close the integration path clockwise. We have to consider the poles at $`P_{120}=E_1+E_2`$ and $`P_{120}=E_1^{}+E_2^{}.`$ We obtain
$$<K_{12}K_{34}>=\frac{P_0E^{}}{E_1+E_2E_1^{}E_2^{}}V_{12}\frac{1}{P_0E_1E_2E_3^{}E_4^{}}V_{34}$$
$$+V_{12}\frac{1}{E_1^{}+E_2^{}E_1E_2}V_{34}\frac{P_0E}{P_0E_1^{}E_2^{}E_3E_4}.$$
(51)
Writing then
$$P_0E^{}=(P_0E_1E_2E_3^{}E_4^{})+(E_1+E_2E_1^{}E_2^{})$$
(52)
$$P_0E=(P_0E_1^{}E_2^{}E_3E_4)+(E_1^{}+E_2^{}E_1E_2)$$
(53)
in (51), we see that the contributions of the first terms cancel mutually, so that we remain with
$$<K_{12}K_{34}>=V_{12}V_{34}\left[\frac{1}{P_0E_1E_2E_3^{}E_4^{}}+\frac{1}{P_0E_1^{}E_2^{}E_3E_4}\right].$$
(54)
These two terms are cancelled by the two last terms of (47) respectively.
We have thus seen that, in the case of instantaneous positive-energy kernels, the contribution of the four-vertex terms vanishes, while the contribution of the two-vertex terms is simply the sum of the six two-fermion potentials. Let us now examine a typical set of six-vertex graphs, the ones which contain two $`K_{12}`$ with one $`K_{34}.`$ Their contributions to the potential are represented in figure 2, with a vertical line when $`G^0`$ is to be replaced by $`>g^0<.`$ After the mutual cancellations of some identical contributions, we obtain the corresponding contribution to $`<T>,`$ minus the four graphs containing one $`>g^0<,`$ plus the three graphs containing two $`>g^0<.`$ Let us now consider again the case of instantaneous potentials with positive-energy propagators. We have shown above that, in a sequence $`K_{12}K_{34},`$ we must consider the two possible orders. Here, in a sequence $`K_{12}K_{12}K_{34},`$ it can be shown that we must consider the three possible orders, so that the first graph of figure 2, e.g., becomes equal to the sum of the three last ones. In figure 3, we expand the graphs of figure 2 in the case of instantaneous kernels with positive-energy propagators, and we see that the sum of the resulting graphs is zero.
## 4 Bypassing the Bethe-Salpeter equation.
In section 3, we presented a 3D reduction method for the N-body Bethe-Salpeter equation. The starting homogeneous equation was written in terms of a kernel $`K,`$ which was designed to reproduce the full propagator $`G`$ by iterations of the inhomogeneous equation. We have seen that the writing of this kernel $`K`$ was straightforward for $`N=2,3,`$ less straightforward for $`N=4,`$ and increasingly complicated for $`N5.`$ Our investigations on the $`N=4`$ case in section 3 suggests us a possible way of avoiding the explicit writing of the Bethe-Salpeter equation. We saw in section 2 that the 3D potential is given by
$$V=<T(1+G^IT)^1>=<K(1G^RK)^1>$$
(55)
where $`G^R=G^0G^I,`$ while $`K`$ contain counter-terms in order to avoid the apparition of topologically identical diagrams in the expansion of $`T=K(1G^0K)^1.`$ It is however possible to adopt a simpler algorithm for $`T:`$
$$T=\left[K^{IR}(1G^0K^{IR})^1\right]^{SC}.$$
(56)
By $`K^{IR}`$ we denote the sum of the irreducible interactions (as (33-34) without the three counter-terms) and by $`SC`$ (single counting), we mean that the diagrams which appear two or more times in the expansion of $`T`$ must be kept only once. The 3D potential will then be written
$$V=<K^{IR}(1G^RK^{IR})^1>^{SC}$$
(57)
in which, after the expansion of the series and the splitting of $`G^R`$ into $`G^0G^I,`$ we shall remove the duplicating diagrams. This makes the writing of the 3D potential as straightforward for $`N4`$ as it was for $`N=`$1 or 2. As an example in the four-particle case:
$$(K_{12}G^RK_{34}+K_{34}G^RK_{12})^{SC}$$
$$=(K_{12}G^0K_{34}+K_{34}G^0K_{12}K_{12}G^IK_{34}K_{34}G^IK_{12})^{SC}$$
$$=K_{12}G^0K_{34}K_{12}G^IK_{34}K_{34}G^IK_{12}.$$
(58)
With positive-energy instantaneous interactions, the term in $`G^0`$ will be cancelled by the two terms in $`G^I.`$
The writing of the Bethe-Salpeter kernel could be done in a quite similar way:
$$K=\left\{T(1+G^IT)^1\right\}_{G^IG^0}=\left\{\left[K^{IR}(1[G^0G^I]K^{IR})^1\right]^{SC}\right\}_{G^IG^0}.$$
(59)
Here, $`G^I`$ is simply a temporary renaming of $`G^0,`$ which indicates that terms like $`AG^IB`$ and $`BG^IA`$ must always be kept both. We know that $`K,`$ unlike $`V,`$ contain only a finite number of parts . For $`N=4,`$ e.g., eqs. (33-34) show that $`K`$ contains 11 subkernels (groups of irreducible graphs) in $`K^{IR}`$ plus 3 counterterms, coming from the second order in $`K^{IR}`$ (we get one of these counterterms by replacing $`G^I`$ by $`G^0`$ in (58)).
The operator $`T`$ is directly given by the Feynman graphs if we neglect the radiative corrections. If not, we can use
$$G^0+G^0TG^0=G=G^0+G^0T^{}G^0$$
(60)
with
$$G^0=G^0(1\mathrm{\Sigma }G^0)^1=\underset{i}{}G_i^0(1\mathrm{\Sigma }_iG_i^0)^1.$$
(61)
The operators $`G`$ and $`T^{}`$ are directly given by the Feynman graphs. For $`T,`$ we have
$$T=\mathrm{\Sigma }(1G^0\mathrm{\Sigma })^1+(1\mathrm{\Sigma }G^0)^1T^{}(1G^0\mathrm{\Sigma })^1.$$
(62)
## 5 Conclusions
In a previous work , we were in search of a 3D reduction method for the two-fermion Bethe-Salpeter equation based on an unspecified positive-energy instantaneous approximation of the Bethe-Salpeter kernel. We performed a series expansion around this approximation, followed by an integration on the relative energy and a second series expansion, at the 3D level, in order to render the resulting 3D potential symmetric. After combining both series, we found that we had in fact built a kind of propagator-based reduction, using an integration with respect to the relative energy, in full contrast with the usual constraining propagator-based reductions, which use constraints. Furthermore, this method was easily generalisable to systems consisting in any number of fermions and/or bosons.
In the present work, the increasing difficulty of writing the Bethe-Salpeter kernel for $`N4`$ suggested us a direct way of writing the 3D potential without the need of first writing the Bethe-Salpeter kernel explicitly. This writing of the 3D potential is straightforward and valid for all $`N,`$ with the simple prescription of removing the duplicating graphs which appear when $`N4.,`$
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# Spin glasses without time-reversal symmetry and the absence of a genuine structural glass transition
## I Introduction
Despite over two decades of work, the controversy concerning the nature of the ordered phase of short-range Ising spin glasses continues. Monte-Carlo simulations of three- and four-dimensional systems appear to be providing evidence for replica symmetry breaking (RSB) in these systems (for recent reviews see ). However, recent developments have cast doubt on this interpretation of the Monte Carlo data. In a series of papers on the Ising spin glass within the Migdal-Kadanoff approximation (MKA), we showed that the equilibrium Monte Carlo data in three and four dimensions that had been interpreted in the past as giving evidence for RSB can actually be interpreted quite easily within the droplet picture, with apparent RSB effects being attributed to a crossover between critical behaviour and the asymptotic droplet-like behaviour for small system sizes. We also showed that system sizes well beyond the reach of current simulations would probably be required in order to unambiguously see droplet-like behaviour. Very recently, a third view on the nature of the low-temperature phase if spin glasses has emerged. In this picture, there exist droplet excitations in short scales, but on large scales there are system-wide excitations that cost only a finite energy in the thermodynamic limit, and that have a surface whose fractal dimension is less than the space dimension . It remains to be seen whether these excitations survive at larger system sizes. Within MKA, fractal excitations are not possible, and the signatures of these excitations found in Monte-Carlo simulations are therefore not present in MKA.
There is a close connection between the question of the nature of the spin glass phase and that of the existence of a phase transition in a spin glass in an external field. Mean-field theory predicts a phase transition to a spin glass phase with RSB along the so-called de Almeida-Thouless line. The droplet picture predicts no transition. The reason is that in the presence of a field time reversal symmetry is broken, and there is no symmetry left that could possibly be broken by a phase transition, except for replica symmetry. Monte-Carlo simulation of the spin glass in a field show some evidence of a phase transition, and in particular of non-self-averaging, however, the situation is complicated by the presence of large finite-size effects due to crossover phenomena. For this reason, Parisi et al studied a different system that also has broken time-reversal symmetry, but is expected not to have strong crossover effects. This system is the three-spin model, where the two-spin products in the Ising spin glass without field are replaced by three-spin products. The numerical evidence for a phase transition in the four-dimensional system seems good, if a quantity that measures the degree of non-self-averaging is studied.
It is the purpose of this paper to study the two mentioned systems without time-reversal symmetry in MKA, in order to check whether a similar degree of non-self-averaging could be produced by finite-size effects. The flow of the coupling constants shows that for both systems there exists only one attractive fixed point which corresponds to a paramagnet in a random field, and that there is consequently no phase transition within MKA. Nevertheless, in both systems the coupling constants increase initially for sufficiently low temperatures, indicating that for small system sizes there might be the appearance of a phase transition. Then, we looked at the non-self-averaging parameter in both models within MKA, for various system sizes and parameter values. We found a behaviour similar to the one reported for the Monte-Carlo simulations, and apparent RSB for system sizes similar to theirs. Furthermore, for the spin glass in a field, the maximum of the non-self-averaging parameter as function of the field (for fixed system size) marks a line that can be interpreted as a remnant of the de Almeida-Thouless line.
Some insights into what might be expected in the finite dimensional three-spin model can be obtained from the mean-field solution of the p-spin model (for a review of which see ). It has an analytical solution in the limit where $`p`$ goes to infinity which can be obtained by a one-step replica symmetry breaking scheme. For our purposes the solution is best understood in terms of metastable states which can be identified with the solutions of Thouless-Anderson-Palmer (TAP)-like equations . The partition function is obtained from the integral
$$Z=_{f_{min}}^{f_{max}}𝑑f\mathrm{exp}(N\sigma (f)Nf\beta ),$$
(1)
where $`\beta =1/k_BT`$ and $`N`$ is the number of spins in the system. $`f`$ is the free energy density of a TAP state. Solutions of the TAP equations exist for $`f_{min}<f<f_{max}`$, and the number of solutions at free energy density $`f`$ is exponentially large and given by $`\mathrm{exp}(N\sigma (f))`$, with $`\sigma (f)`$ vanishing at both $`f_{min}`$ and $`f_{max}`$. For temperatures in the interval $`0<T<T_S`$ the integral is dominated as $`N\mathrm{}`$ by the lower limit of the integral, i.e. states whose free energy $`Nf`$ differs from the state of lowest energy by only a finite amount. For temperatures in the interval $`T_S<T<T_D`$ the integral can again be done by steepest descents and is dominated by some value of $`f`$ lying in the interval $`f_{min}<f<f_{max}`$. As the temperature approaches $`T_D`$ this value tends to $`f_{max}`$. In the temperature interval $`T_S<T<T_D`$ an exponentially large number of states contribute to the thermodynamics in contrast to the situation below $`T_S`$ where only a finite number contribute. Above $`T_D`$, only the trivial paramagnetic state contributes. There are thus two phase transitions at mean-field level. The lower temperature transition at $`T_S`$ is accompanied by singularities in the free energy but at the higher temperature transition $`T_D`$ the free energy is smooth and the presence of a transition is best inferred from singularities in the dynamics.
Now for a finite dimensional system metastable states with $`f>f_{min}`$ are unstable. (Imagine in such a state converting a block of spins of linear dimension $`L`$ to have the orientations that they would have in the lowest state; the free energy gain will be or order $`L^d`$; the energy cost of creating such a region will be no more than the cost of breaking all the bonds at the surface of the region, $`L^{d1}`$. Thus the possibility of nucleating lower free energy states prevents the existence of metastable states in finite dimensional systems). As the transition at the higher temperature $`T_D`$ involves the metastable states only, one deduces that it will not exist in a finite dimensional system. The only transition which could possibly survive to finite dimensions is the one associated with $`T_S`$. Our studies however indicate that it too probably does not occur in finite dimensional systems.
Our study of the three-spin model had another, perhaps physically more significant motivation. The three-spin model and its cousins have been extensively studied at mean-field level as models of structural glasses . The higher temperature transition, $`T_D`$, whose signature is purely dynamical (and very similar to that in the mode-coupling approach to real glasses ), while the transition at the lower temperature, $`T_S`$, is associated the Kauzmann temperature, $`T_K`$ , the temperature at which the configurational entropy of the glass goes to zero. It has been a common belief of many workers for several decades that there is no genuine transition at $`T_K`$. Recently this belief has been strongly reinforced by the novel Monte Carlo simulation of Santen and Krauth , who found no evidence of a genuine transition in a simulation in which the particles could be properly equilibriated. Our work strengthens the argument that no genuine transition analogous to $`T_S`$ or $`T_K`$ will exist in finite dimensions. Nevertheless we can see echoes of the mean-field results in our calculations. It is our belief that future work should focus on how the singularities in mean-field results are ”rounded-off” in finite dimensions.
This paper is organized as follows: First, we define the models and quantities studied in this paper. Then, we describe the MKA for the two models studied in this paper. We use two different methods in order to make sure that the results do not depend on the particular implementation of the MKA. Next, we study the three-spin model and give our results for the flow of the coupling constants, and for the non-self-averaging parameter. In section V, we discuss the spin glass in an external field. Again, we give results for the flows and the degree of non-self-averaging. We also give scaling arguments based on the droplet picture that explain most of our findings. Finally, we summarize and discuss our results.
## II Models and Definitions
The Edwards-Anderson spin-glass Hamiltonian $``$ in the presence of a uniform external magnetic field $`h`$ is given by
$$\beta (\sigma )=\underset{i,j}{}J_{ij}\sigma _i\sigma _j+h\underset{i}{}\sigma _i,$$
(2)
where the Ising spins $`\sigma _i`$ can take the values $`\pm 1`$, and the nearest-neighbor couplings $`J_{ij}`$ are independent from each other and Gaussian distributed with a standard deviation $`J`$. For convenience, the Boltzmann-factor $`\beta `$ is absorbed into the couplings and fields. Without a field $`h`$, the model has a low-temperature phase with nonvanishing correlations $`S_iS_j`$ even for spins that are far apart. According to the droplet picture, this phase is unique (up to a global flip of all spins), and it is destroyed as soon as the field is turned on. The reason is that the field induces regions of a sufficiently large radius to flip if the magnetization of this region opposes the field. The radius $`r`$ of these regions is obtained from the condition that the gain in magnetic energy, $`hr^{d/2}`$ becomes comparable to the loss in coupling energy, $`Jr^\theta `$, leading to
$$r(J/h)^{1/(d/2\theta )}.$$
(3)
Here, $`d`$ is the dimension of the system, and $`\theta `$ the scaling dimension of domain walls. Beyond the radius $`r`$, the long-range correlations of the spin-glass phase are destroyed.
In contrast, the RSB picture predicts the existence of infinitely many different phases of comparable free energy in the absence of a field. With increasing field $`h`$, the number of phases decreases, and it becomes one at the de Almeida-Thouless line $`h_c(J)`$. At the critical spin-glass transition $`J_c`$, the critical field $`h_c`$ vanishes, and it diverges to infinity as the coupling strength $`J`$ diverges (i.e., as the temperature goes to zero).
The numerical analysis of the spin glass in a field is hampered by strong crossover effects. Crossover effects are expected to be less strong in the three-spin model, because it has no tunable parameter that can restore time reversal symmetry and lead to strong crossover effects when small . In the most easily tractable version of this model, each site is occupied by two Ising spins, $`\sigma _i^{(1)}`$ and $`\sigma _i^{(2)}`$, and the Hamiltonian is given by
$`\beta (\sigma )`$ $`=`$ $`{\displaystyle \underset{ij}{}}(J_{ij}^{(1)}\sigma _i^{(1)}\sigma _i^{(2)}\sigma _j^{(1)}+J_{ij}^{(2)}\sigma _i^{(1)}\sigma _i^{(2)}\sigma _j^{(2)}`$ (4)
$`+`$ $`J_{ij}^{(3)}\sigma _i^{(1)}\sigma _j^{(1)}\sigma _j^{(2)}+J_{ij}^{(4)}\sigma _i^{(2)}\sigma _j^{(1)}\sigma _j^{(2)}),`$ (5)
where $`ij`$ are nearest-neighbor pairs, and the couplings $`J_{ij}^{(n)}`$ are chosen independently from a Gaussian distribution with zero mean and width $`J`$. When the signs of all spins are reversed, the sign of the Hamiltonian changes also, indicating the violation of time-reversal symmetry.
If finite-dimensional systems have no RSB, this model has no phase transition since there is no symmetry that could be broken. On the other hand, if RSB occurs in finite-dimensional spin glasses, the three-spin model could show a phase transition at some critical coupling strength $`J_c`$.
It has proven useful to consider two identical copies (replicas) of the system, with the spins $`\{\sigma _i\}`$ and $`\{\tau _i\}`$, and to measure overlaps between them. This gives information about the structure of the low-temperature phase, in particular about the number of pure states. The main quantity studied in this paper is the parameter $`A`$ which measures the degree of non-self averaging, and is defined by
$$A=\frac{\left[(qq)^2^2\right]}{\left[(qq)^2\right]^2}1,$$
(6)
where $`\mathrm{}`$ and $`\left[\mathrm{}\right]`$ denote the thermodynamic and disorder average respectively. The overlap $`q`$ between the two replicas is given by
$$q=\frac{1}{N}\underset{i}{}\sigma _i\tau _i$$
for the Ising spin glass in a field, and by
$$q=\frac{1}{2N}\underset{i}{}(\sigma _i^{(1)}\tau _i^{(1)}+\sigma _i^{(2)}\tau _i^{(2)})$$
for the three-spin model. $`N`$ is the number of sites in the system.
$`A`$ is most easily evaluated by introducing a coupling between the two replicas, and by differentiating with respect to it. The Hamiltonian for the coupled system is then
$$\beta _ϵ(\sigma ,\tau )=\beta (\sigma )\beta (\tau )+ϵNq.$$
The mean overlap is given by the expression
$$q=\left[\frac{1}{N}\frac{}{ϵ}\mathrm{ln}Z\right]_{ϵ=0},$$
and its variance by
$$(qq)^2=\left[\frac{1}{N}\frac{^2}{ϵ^2}\mathrm{ln}Z\right]_{ϵ=0}.$$
$`Z`$ is the partition function.
In systems with RSB, the probability distribution $`P(q)`$ of $`q`$ is broad, and $`A`$ has a nonzero limit in the limit of infinite system size. On the other hand, in the absence of RSB, each sample has only one, sample-independent value of $`q`$, and $`A`$ vanishes in the thermodynamic limit. Consequently, an $`A`$ that increases with increasing system size could be taken as an indicator of RSB. However, we will see in this paper that even systems without RSB can show an increasing $`A`$ over a wide range of system sizes and parameters.
## III The Migdal Kadanoff Approximation
The Migdal-Kadanoff approximation (MKA) is a real-space renormalization group that gives approximate recursion relations for the various coupling constants. An exact decimation, which consists in taking the trace over all those spins that do not belong to the coarse-grained lattice, generates higher-order couplings between spins of more that two sites, and is therefore untractable. In order to circumvent this problem, the MKA moves the bonds of a hypercubic lattice before each decimation step in such a way that no higher-order couplings can be generated. If the bond-moving shall be symmetric with respect to the different space directions, one ends up with the scheme represented in Figure 1.
In a $`d`$-dimensional lattice, $`2^{d1}`$ bonds are superimposed as a consequence of bond-moving. In the absence of field terms (i.e. of terms that couple only spins that sit on the same site $`i`$, and are therefore not in a clear way associated with bonds), the $`2^{d1}`$ coupling constants of each of the $`d`$ bundles of bonds per coarse-grained unit cell simply add up, and the “naked” spins that are left behind have no couplings. Taking the trace over the $`d`$ spins that sit on the $`d`$ main bonds leads to coarse-grained coupling constants between neighboring spins on the coarse-grained lattice. Taking the trace over the naked spins, gives only a constant contribution to the partition function, which can be neglected. This decimation procedure is iterated $`n`$ times on a lattice of linear size $`L=2^n`$, until a single unit cell is left over. Assuming periodic boundary conditions, one can then take the trace over the final spin.
The flow of the coupling constants in this scheme results from alternating the addition of $`2^{d1}`$ bonds with linking two of these new bonds together and taking the trace over the middle spin. Essentially the same flow results when a decimation is done on a hierarchical lattice that is constructed iteratively by replacing each bond by $`2^d`$ bonds, as indicated in Fig. 2. The total number of bonds after $`n`$ iterations is $`2^{dn}`$. Spin decimation on such a lattice is done by taking the trace over the spins that are highest on this hierarchy (i.e., that were added last during the construction procedure). At each decimation step, first the trace is taken over the middle spin of two linked bonds, and then $`2^{d1}`$ bonds are added together to form a new bond, until the lowest level is reached and the trace over the remaining two spins is calculated . Apart from the fact that the order of bond-adding and decimation is reversed, the recursion relations for the coupling constants are obtained by the same procedure as for the bond-moving algorithm.
This equivalence between a bond-moving procedure for a hypercubic lattice and a hierarchical lattice hold no longer when field terms are present. In both of our models, these field terms are present either from the beginning (Ising spin glass in a field), or they are initially absent, but are generated during the decimation procedure (three-spin model). For the bond-moving procedure, one has to decide whether the fields shall also be moved (and to where), or whether they shall remain at their sites, or whether part of them shall be moved. This creates a certain freedom in the renormalization scheme, and the most plausible choice is determined by the requirement that the flows of the fields near the zero-temperature fixed point and at the infinite-temperature fixed point shall be those of the hypercubic lattice. However, if the main results of the MKA shall be generic, they should not depend on the precise implementation of the bond-moving algorithm. Otherwise one might doubt that the MKA reflects the features of the real system. For this reason, we have performed the MKA for a variety of different implementations, and we found our main conclusions concering $`A`$ to be independent of the implementation. In the following sections, we will give our results for the hierarchical lattice and for a bond-moving scheme that has the correct flow of the field terms.
The treatment of a hierarchical lattice in the presence of field terms is straightforward. In order to understand that the flows of the fields are different on a hierarchical lattice compared to a hypercubic lattice, let us consider a situation where the flows of the couplings go to zero with increasing iteration number, which is the situation that we will encounter below for both models. As long as the couplings are nonzero, each decimation step generates a contribution to the fields at the sites that are left over. The two corner spins, which are left over until the end, receive consequently $`2^{n(d1)}`$ field contributions from the first decimation, $`2^{(n1)(d1)}`$ from the second iteration, and so on, until the couplings are virtually zero. For twice the system size, i.e., for a lattice with $`n+1`$ levels, the mean of the field contribution to the corner spins due to decimations is larger by a factor $`2^{(d1)}`$, and so is the variance of the field contribution. Even though the couplings go to zero after a certain number of iterations, the fields keep growing.
In contrast, on a hypercubic lattice, the fields must remain constant as soon as the couplings have become zero. Clearly, this can only be achieved if fields terms are not moved to the sites that will not be traced over. On the other hand, field terms must not be left with the “naked” spins. The reason is that near the zero-temperature fixed point where the couplings are very large, all fields must add up under renormalization. Fields must therefore always stay with spins that are coupled to other spins. For this reason, fields should be moved to those $`d`$ sites that sit at the middle of the main bonds. Even with this restriction, there remains some freedom in choosing which field should be moved where. In our simulations, we treated field terms as belonging to bonds. The initial fields were evenly distributed between the ends of all bonds, and the fields generated during decimations naturally end up at the ends of bonds. When a bond was moved, we moved all its field terms to that end that was to be traced over next.
For the Ising spin glass in a field, the recursion of four different parameters must be considered when studying the flow diagram and thermodynamic quantities. These are the two-spin coupling, the two fields on the two ends of a bond, and a constant. If one evaluates quantities related to the overlap between two replicas, each site has two spins, leading to 16 parameters. The same number of parameters occurs for the three-spin model, if only one replica is needed, as, e.g., for the flows and the phase diagram. For the evaluation of $`A`$, we need two replicas, leading 256 couplings. Luckily, the decimation step can treat all 256 parameters with the same formula, which involves a $`256\times 256`$ matrix that is calculated once at the beginning of the program.
## IV The Three-Spin Model
This model was studied using Monte-Carlo simulations in four dimensions in , and evidence for RSB was found. The authors of found in particular that the non-self-averaging parameter $`A`$ is small for large temperatures, and becomes large for smaller temperatures. Curves for different system size $`L=3,4,5,6`$ cross nearly at the same temperature, and below this temperature $`A`$ increases with increasing $`L`$. Thus, the degree of non-self averaging increases with the system size, just as can be expected for a replica-symmetry breaking transition. Monte-Carlo simulations are usually done with couplings $`J=\pm 1`$. The precise distribution of the couplings should however not affect the universality class.
Analytical results were obtained for the $`p`$-spin model in mean-field theory, where one-step RSB was found. This means that the ground state has a nonzero probability of being occupied below a critical temperature $`T_S`$ (see Section I). This mean-field scenario is fundamentally different from the full RSB claimed to be seen in Monte-Carlo simulations of the four-dimensional system. Thus, the argument usually employed for spin glasses that mean-field like behaviour can be found in finite-dimensional short-range systems, fails here.
In the following, we show using MKA that the assumption of the absence of crossover effects in this model is incorrect, and that $`A`$ might at low temperatures and for small system sizes increase with increasing system size even if the system is self-averaging in the thermodynamic limit. We mainly focus on the case of four dimensions, but report also some results in d=2 and 3. Let us first discuss the flow of the coupling constants as the system is renormalized. Because each bond is connected to 4 spins, the flow of 16 coupling constants has to be considered. In order to obtain this flow, we iterated the recursion relation on a set of $`10^6`$ bonds. At each iteration, each of the new set of $`10^6`$ bonds was generated by randomly choosing $`2^d`$ bonds from the old set. For a hierarchical lattice, where the generated fields remain at that end of a bond at which they are generated, we first took the trace over the inner spins of each of the $`2^{d1}`$ pairs of bonds, and than we added the resulting bonds; for the bond-moving procedure described in the previous section, we first generated two bunches of $`2^{d1}`$ bonds each, then moved the fields of all but the “original” bond of each bunch to the inner spin, and took the trace over the inner spin. Figure 3 shows the flow of the width of the three-spin couplings for different initial width in four dimensions, for the two different algorithms.
One can see that for weak coupling (or, equivalently, high temperature) the coupling strength decreases quickly with increasing system size. However, for stronger coupling or lower temperatures, the coupling strength increases during the first few iterations, and decreases afterwards. The maximum is reached between the 3rd and 4th iteration, or between $`L=8`$ and $`L=16`$. For sufficiently strong coupling, the curves reach an asymptotic shape. On the hierarchical lattice, where the fields grow without bounds, the 3-spin couplings decrease to zero faster than with bond-moving. Furthermore, curves for the hierarchical lattice seem to correspond roughly to those of the bond-moving procedure if the three-spin couplings are divided by a number around 3. The reason is that the first step during the bond-moving procedure summarizes 8 bonds to one new bond. The width of the three-spin coupling is therefore increased by a factor of $`\sqrt{8}`$ in four dimensions. In order to compare to the hierarchical lattice or to Monte-Carlo simulations on a hypercubic lattice, one should divide the coupling strength of the bond-moving procedure by $`\sqrt{8}`$.
If one considered only systems of sizes up to 8, one would get the illusion of a phase transition with a $`(1/J)_c`$ around 3 or 4, a value which is not far from the one given for $`T_c`$ in . (Note that these authors have kept the coupling strength fixed at $`\pm 1`$, and varied the temperature. Their $`T`$ corresponds therefore to our $`1/J`$.)
Figure 4 shows the flow of the widths of the different coupling constants for an initial width of the three-spin coupling $`J=2`$.
One can see that the strengths of the field and of the on-site two-spin coupling (which can also be viewed as a “field”) increase rapidly and without limits for the hierarchical lattice, and that they saturate at a finite value in the bond-moving case. The two- three- and four-spin couplings increase during the first few iterations, and then decrease again. Thus, our three-spin model corresponds on large scales to a system only with random fields and random couplings between the $`\sigma `$ and $`\tau `$ spins on the same site. There are no couplings between spins on different sites on large enough scales, but sites are independent from each other. Only on small scales could one get the impression that the system has long-range correlations. However, these system sizes are exactly the ones studied in .
The crossover regime becomes larger with increasing dimension. Figure 5 shows the flow of the three-spin coupling for an initial value $`J=10`$ in $`d=2,3,4`$ dimensions. Clearly, the strength of the increase and the range of system sizes over which this increase occurs increases with increasing dimension. One can therefore expect that in even higher dimensions, the apparent phase transition becomes more pronounced.
Next, let us study the non-self-averaging parameter $`A`$. As explained in the previous section, $`A`$ can be evaluated by introducing a coupling between two identical replicas of the system. Since there are now 8 spins associated with each bond, the number of couplings that have to be evaluated in MKA approximation is $`2^8=256`$. Figure 6 shows $`A`$ as function of the coupling strength for different system sizes up to 16, in 4 dimensions. Larger system sizes could not be studied due to limitations in computer time. One can see that $`A`$ increases with increasing system size whereever it is appreciably different from zero, and reaches large values. This figure gives the impression that the system shows non self-averaging. Of course, for larger system sizes, $`A`$ must eventually decrease again since we know that there is self-averaging in the thermodynamic limit. In contrast to the Monte-Carlo simulation results , our curves for $`A`$ do not intersect at a coupling strength and $`A`$ value of the order 1. We have performed a similar simulation in $`d=2`$ dimensions and found that $`A`$ increases as the system size increases over the range $`L=2,4,8,16`$. However, for $`L=32`$ and $`L=64`$, $`A`$ decreases. If we assume that the system size for which $`A`$ is largest increases with each dimension by a factor 2, as it does for the flow of the couplings, we can expect that in $`d=4`$ the system sizes for which a decrease in $`A`$ can be seen is beyond $`L=64`$.
To summarize this section, we have shown that the three-spin model, even in situations where we know that it self-averages in the thermodynamic limit, can show indications of non-self-averaging at those system sizes typically studied in simulations. Evidence for non-self-averaging found in Monte-Carlo simulations must therefore be taken with caution as it might be misleading.
## V Ising Spin Glass in a Magnetic Field
Monte-Carlo simulations in four dimensions show some indication of RSB . Just as for the three-spin model and for the spin glass without external field, these findings may again be due to finite-size effects and to the closeness to the critical temperature. Indeed, an investigation of the ground-state structure of a spin glass in a magnetic field shows no indication of RSB. (See, however, the discussion in .)
In order to test for finite-size effects, we studied the spin glass in a field using MKA, and determined the non-self-averaging parameter $`A`$ as function of the system parameters. We found that the degree of non-self-averaging can be large for the system sizes typically used in simulations, in particular when the contribution of the field to the free energy is comparable to that of the couplings. While most published Monte-Carlo simulations were done in four dimensions, we chose to study the MKA in three dimensions, in order to be able to go to larger system sizes. Because there are three parameters to be varied (the system size, the field, and the two-spin couplings), many data points had to be collected, and this is done faster in 3 dimensions. Of course, we expect that the results of the MKA are similar in four dimensions, if the exponents for 3 dimensions are replaced with those for 4 dimensions. Just as for the three-spin model, the apparent non-self-averaging should become even stronger in 4 dimensions.
First, let us study the flows of the couplings and fields. The decimation procedure leads to the creation of random fields, while the mean value of the field is not changed. Figure 7 shows the flow of the two-spin coupling $`J`$ for various initial values, and for a fixed field $`h=0.1`$. For initial couplings larger than the critical coupling (1.13 for the hierarchical lattice and 0.55 for bondmoving), the coupling strength decreases immediately. However, if the initial coupling strength is sufficiently deep in the low-temperature phase, it increases first, until the random field has become strong enough to have a reducing effect on the coupling strength. Ultimately, the flow goes to a fixed point where the coupling strength is zero. On the hierarchical lattice, the width of the field keeps growing indefinitely, while it saturates in the bond-moving case, as discussed in section III. Clearly, there is no phase transition in the presence of an external field, but there are strong crossover effects if the field is small.
As mentioned in Section II, the droplet picture predicts that beyond a length scale $`r(J/h)^{1/(d/2\theta )}`$ the contributions of the field and of the couplings to the free energy become comparable, and we expect that the strength of the couplings decreases beyond this scale. In order to test this prediction, we have plotted in Figure 8 the iteration number for which the two-spin coupling is largest versus the logarithm of $`J/h`$. It should follow a law
$$\mathrm{log}_2L=\frac{1}{\mathrm{ln}2((d/2)\theta )}\mathrm{ln}(J/h)+C1.15\mathrm{ln}(J/h)+C,$$
with a suitable constant $`C`$. As the figure shows, the data for bond-moving agree nicely with this prediction. for the hierarchical lattice, the slope is larger and has a value around 1.4. This might be due to the fact that the field increases faster on the hierarchical lattice, leading to an earlier reduction in the coupling strength than predicted by the scaling theory.
Next, let us discuss the quantity $`A`$ which is a measure of the degree of non self-averaging. Figure 9 shows our results in the absence of a magnetic field. In the high-temperature phase as well as in the spin-glass phase $`A`$ descreases with increasing system size and approaches zero, just as one can expect in the absence of RSB. At the critical coupling strength $`J_c`$, $`A`$ remains constant with increasing system size, its value being $`A0.13`$ for the hierarchical lattice, and $`A0.15`$ for bondmoving. The constancy of $`A`$ at the critical point can be explained from the scaling behaviour of the overlap distribution function $`P(q)`$. Critical scaling implies
$$\left[P(q)\right]=L^{\beta /\nu }\left[\stackrel{~}{P}(qL^{\beta /\nu })\right]$$
and
$$\left[P(q)P(q^{})\right]=L^{2\beta /\nu }\left[\stackrel{~}{P}(qL^{\beta /\nu })\stackrel{~}{P}(q^{}L^{\beta /\nu })\right],$$
with $`\beta `$ being the order parameter critical exponent, and $`\nu `$ the correlation length exponent. Introducing the variable $`y=qL^{\beta /\nu }`$, we then obtain
$$A=\frac{y^2y^2\stackrel{~}{P}(y)\stackrel{~}{P}(y^{})𝑑y𝑑y^{}}{\left(y^2\stackrel{~}{P}(y)𝑑y\right)^2}1$$
independently of $`L`$.
For low temperatures $`T=1/J`$, $`A`$ seems to follow the law $`ATL^\theta `$ with $`\theta 0.24`$. This can be derived by the following argument: At low temperatures, most samples have a value of $`q^2`$ close to 1, and only a fraction $`p`$ proportional to $`kTL^\theta `$ of all samples have system-wide excitations and have therefore some other value $`q^2=x<1`$. We therefore obtain
$$\left[q^2\right]1p+p[x]$$
and
$$\left[q^2^2\right]1p+p[x^2],$$
leading to
$$Ap(1+[x^2]2[x])kTL^\theta .$$
In the presence of a magnetic field, we expect $`A`$ to decrease always to zero for large system sizes, because the system is always in the high-temperature phase without long-range correlations. However, as we will show in the following, $`A`$ can become nevertheless very large for certain combinations of the system size, the field, and the two-spin coupling strength. One can therefore easily get the impression that the system is not self-averaging, while in reality $`A`$ increases only over a limited range of system sizes or parameters.
Figure 10 shows our results for $`A`$ with a magnetic field $`h=0.2`$. For $`J>J_c`$, the values are larger than without field, and they increase with increasing system size and decreasing temperature. We expect that as the system size increases further, $`A`$ will reach a maximum and then decrease again. For fields stronger than $`h=0.5`$, we see this reversal in the trend of $`A`$ already for the system sizes studied in the simulations. For $`J<J_c`$, Fig. 10 shows that the curves for different $`L`$ intersect each other, such that for high temperatures self-averaging is better for larger system sizes. Thus, the behaviour of $`A`$ for weak fields seems to be qualitatively similar to that of the three-spin model.
For given system size and coupling strength $`J>J_c`$, there exists always a value of $`h`$ for which $`A`$ has a maximum. The height of this maximum is higher for larger system sizes and for lower temperatures $`1/J`$. Figure 11 shows the field for which $`A`$ is largest as function of $`1/J`$. The data are in good agreement with a dependence $`h_{max}J`$, and $`h_{max}L^{1.26}`$ for bond-moving. This means that $`A`$ is largest when $`hJL^{(d/2\theta )}`$. For the hierarchical lattice, the fit to the data is best for a dependence $`h_{max}J`$, and $`h_{max}L^{0.93}`$. Just as in Figure 8, the effective value of $`d/2\theta `$ appears to be larger on the hierarchical lattice than for bond-moving. We suspect that this is due to the fact that the field grows indefinitely on the hierarchical lattice.
These results can be understood if one considers the effect of the field on the overlap distribution $`P(q)`$. Without field, $`P(q)`$ is a symmetric function, and varies considerably in shape for different samples for the system sizes typically used in simulations. This feature is seen in Monte-Carlo simulations as well as in MKA . A magnetic field changes the shape of $`P(q)`$ and moves the weight more and more towards positive $`q`$. In the limit $`h\mathrm{}`$, all spins are aligned with the field, leading to $`P(q)=\delta (1)`$. Along the boundary line $`L(J/h)^{1/(d/2\theta )}`$, where the field is not yet strong enough to destroy all features of the low-temperature phase, we can expect that at least some samples still have large droplets than can be flipped without much free energy cost. In Monte-Carlo simulations , one finds indeed for certain intermediate parameter values a $`[P(q)]`$ that has a pronounced peak at some large $`q`$ value, and a long and thin tail that extends almost all way down to $`q=1`$. The authors point out that this feature results from most samples having only the main peak, and other samples having an additional second peak for some other value of $`q`$. They go on to argue that this is a non self-averaging feature characteristic of RSB, and that it should not be expected if there was no RSB. However, they also admit that their simulations do not show a second peak in $`[P(q)]`$ at a value $`q_{min}`$, which is expected from Mean-Field Theory. Although we have not determined $`P(q)`$ in the presence of a magnetic field within MKA, we can conclude from the behaviour of $`A`$ that $`P(q)`$ must have in MKA approximation exactly the same features that we just described for the Monte-Carlo simulations. Indeed, it is easy to show that $`A`$ becomes large if most samples have a $`P(q)`$ with one narrow peak at $`q_0`$, and some samples have two peaks. For those samples with one peak, we have a small variance of $`q`$,
$$\chi _sq^2q_0^2,$$
which is essentially sample-independent. For those samples with two peaks, we have a large variance $`\chi _l`$ which is different for each sample. If the fraction of samples with two peaks is $`p`$, we obtain
$`A`$ $`=`$ $`{\displaystyle \frac{(1p)\chi _s^2+p[\chi _l^2]}{((1p)\chi _s+p[\chi _l])^2}}1`$
$``$ $`{\displaystyle \frac{p[\chi _l^2]}{\chi _s^2+2p\chi _s[\chi _l]+p^2[\chi _l]^2}},`$
where we have only kept the leading terms. As long as $`p`$ is not much smaller than $`[\chi _l]/\chi _s`$, $`A`$ is of the order $`1/p`$. Thus, if $`\chi _s`$ is small and $`p`$ is small but not too small, $`A`$ is large. The second condition is satisfied if the field is such that a small fraction of samples have a second peak in $`P(q)`$, the first condition is better satisfied for larger $`L`$ or smaller $`T`$. This explains why we observe the maximum of $`A`$ for those $`h`$ values where the contribution of the field to the free energy is comparable to that of the couplings, and why the maximum of $`A`$ is larger for larger systems and lower temperatures. Of course, for some even larger value of $`L`$, we expect $`P(q)`$ to start having less sample-to-sample fluctuations, because the samples should become self-averaging. Then the argument will break down, and $`A`$ should remain small. However, this range of system sizes seems to be beyond the reach of our simulations.
In conclusion, we have shown that there exists a wide range of parameters over which the degree of non-self-averaging appears large for system sizes typically used in computer simulations. We expect our results to be valid beyond MKA. As we have argued for systems without a field , the apparent non-self-averaging must be attributed to the influence of the zero-field critical point. This influence reaches surprisingly far and creates a line in the $`hJ`$ plane along which non-self-averaging is particularly large, and which is somewhat reminiscent of the de Almeida-Thouless line. We are here in agreement with Huse and Fisher who argued already almost 10 years ago that Monte-Carlo simulation data for a spin glass in a magnetic field are strongly affected by the critical point.
## VI Discussion
We have shown that for the three-spin model as well as for the spin glass in a magnetic field, a large degree of non-self-averaging found in computer simulations does not represent unequivocal evidence for RSB, but can be caused by finite-size effects.
It seems, however, that a study of the non-self-averaging parameter $`A`$ using Monte-Carlo simulations might be able to discriminate between RSB and the droplet picture for three- and four-dimensional spin glasses. As we have shown, $`A`$ has in zero field a maximum at $`T_c`$, and decreases again with decreasing temperature in MKA. If there was a low-temperature phase with RSB, the low-temperature value for $`A`$ should probably be larger than the critical value, and $`A`$ should therefore increase with decreasing temperature. Also, for temperatures below $`T_c`$, we found that $`A`$ has its maximum not at zero field, but at some finite field value. The degree of self-averaging decreases deep in the supposed low-temperature phase. We expect a similar behaviour from the Monte-Carlo simulations. This would be a hint that non-self-averaging is strongest along the boundary between the field-dominated and coupling-dominated regime, and not in the region where one would expect a low-temperature phase with RSB.
###### Acknowledgements.
This work began when HB and BD were at the Department of Physics, University of Manchester, supported by EPSRC Grants GR/K79307 and GR/L38578. BD also acknowledges support from the Minerva foundation, and from the German Science Foundation (DFG, Grant Dr300/2-1), and she acknowledges the hospitality of the ICTP in Trieste during a short visit there. MAM thanks Dr A. Cavagna for useful discussions on the mean-field limit.
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# Untitled Document
Cuprates: experiment vs. quasiparticle theory
Indeed, one can hardly disagree with most of the points raised by P.W. Anderson in his Reference Frame of the February 2000 issue of Physics Today.
However, one can take a different viewpoint than his on the cuprates.
Experimentally, the ’notorious ”linear T” resistivity’ etc. has only been observed at temperatures T of the order of a few degrees Kelvin at most, and never all the way down to (practically) zero temperature. Somehow, many people in the field make the implicit extrapolation of the linear regime down to T=0. I think this is an important issue, which has been all too often sidestepped, especially in efforts to model the (normal state of the) cuprates.
Now, regarding the irrelevance (or not!) of quasiparticle theory thereon. Conceivably, ’ ”proofs” that purport to show otherwise, are not very useful since they ignore the possibility of anomalies …’. But, are there experimental indications pointing to this possibility? It is easy to imagine that a long discussion may ensue on the matter.
Nevertheless, if one is willing to give good old quasiparticle perturbative methods a chance, the ’notorious ”linear T” resistivity’ and optical conductivity can be analytically understood - at least for overdoped and optimally doped materials. (Stripes apparently dominate the physics of the underdoped regime, and a smooth crossover between these regimes can be envisaged, but that is another story.) One obtains a one-particle scattering rate which goes like x<sup>2</sup>, x=max{T,energy}, for x$``$0. This scattering rate becomes linear in energy if the latter exceeds a crossover value x<sub>o</sub>, or linear in T if T exceeds (x<sub>o</sub>/4). (x<sub>o</sub> is the difference between the chemical potential and the energy of the van Hove singularity.) The result holds true everywhere in the Brillouin zone . This prediction was directly supported by the ARPES expts. of Johnson et al. . Experiment also indicates that x<sub>o</sub> is roughly in the range 50-350 <sup>o</sup>K, not in disagreement with the observed resistivity etc.
A couple more points related to traditional Fermi liquid effects in the cuprates. In the model presented above, phonons play no essential role, if any, in the T dependence of the scattering rate. They merely provide momentum dissipation, thus yielding a finite resistivity . On the other hand, the unconventional effects of impurity scattering, e.g. in the low T resistivity , merit further investigation before any conclusions can be drawn.
George Kastrinakis
Dept. of Chemical Engineering
University of Cambridge
Cambridge CB2 3RA, U.K.
References
G. Kastrinakis, Physica C 317-319, 497 (1999); e-print cond-mat/0005485, Physica C (in print).
T. Valla et al., Science 285, 2110 (1999); e-print cond-mat/0003407.
G.S. Boebinger et al., Phys. Rev. Lett. 77, 5417 (1996).
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# Thermal stability and nova cycles in permanent superhump systems
## 1 Introduction
### 1.1 Regular superhumps
Superhumps were initially observed in SU UMa systems (Vogt 1974; Warner 1975). This subclass of dwarf novae, shows brighter longer outbursts (so called superoutbursts) in addition to the normal outbursts observed in regular dwarf novae (U Gem systems). A quasi-periodic variation, systematically a few per cent larger than the orbital period, is usually detected during SU UMa superoutbursts, and was nicknamed ‘$`\mathrm{𝐬𝐮𝐩𝐞𝐫𝐡𝐮𝐦𝐩}`$’ (see la Dous 1993; Warner 1995a for reviews of SU UMa systems and CVs in general). The typical value of the peak-to-trough amplitude of the superhump is about 20-40 per cent. There is a trend of increasing positive period excess (of the superhump period over the orbital one) with the orbital period (Stolz & Schoembs 1984; Patterson 1999).
Vogt (1982) first suggested that the accretion disc around the white dwarf in Z Cha (and other SU UMa systems) develops an elliptical shape during superoutbursts. Osaki (1985) tested the motion of single particles in an eccentric disc, and found a relation between the superhump period excess as a fraction of the binary period and the binary period. Whitehurst (1988) used hydrodynamic simulations to show that the tidal instability explains the formation of the eccentric disc. Osaki (1989) used two instabilities in the accretion disc to create a uniform theory for non-magnetic CVs. According to this model, ‘The Thermal – Tidal Disc Instability Model’, the superhump periodicity is the beat of the orbital period of the binary system with the period of the apsidal precession of the accretion disc. Further hydrodynamical simulations showed that the tidal instability can occur only if the disc radius exceeds a certain value, the 3:1 resonance radius (Whitehurst & King 1991). This requires that superhumps can appear only in binary systems with a small mass ratio q=($`M_2/M_1)`$$``$0.33, although observationally the limit is probably bigger (Retter & Hellier 2000; Retter et al. 2000).
During the last decade it has been found that superhump behaviour is common among other classes of binary systems as well as in the SU UMa stars. According to the theory, there are two requirements for the presence of superhumps: an extreme mass ratio and a large accretion disc radius. These two conditions are naturally met in superoutbursts of SU UMa systems, but are satisfied in other systems as well. Thus superhumps have also been detected in the ultra-short orbital period AM CVn systems (Patterson et al. 1993a; Patterson, Halpern & Shambrook 1993; Provencall et al. 1997; Patterson et al. 1997a; Harvey et al. 1998; Solheim et al. 1998; Patterson 1999; Skillman et al. 1999); in SW Sex stars (Patterson & Skillman 1994; Patterson 1999) and during bright outbursts of a few low-mass X-ray transients (White 1989; Charles et al. 1991; Bailyn 1992; Zhang & Chen 1992; Kato, Mineshige & Hirata 1995; O’Donoghue & Charles 1996).
### 1.2 Permanent superhumps
In many systems, which do not show eruptions, superhumps are observed. An outburst is thus not a necessary condition for this phenomenon. Patterson & Richman (1991) termed this behaviour as ‘$`\mathrm{𝐩𝐞𝐫𝐦𝐚𝐧𝐞𝐧𝐭}`$ $`\mathrm{𝐬𝐮𝐩𝐞𝐫𝐡𝐮𝐦𝐩}`$’. During the last few years permanent superhumps have been found in almost twenty CVs (see Patterson 1999 for an observational review on permanent superhump systems). Typical full amplitudes of permanent superhumps are about 5-15 per cent, but they are highly variable, and sometimes even disappear from the light curve (Patterson, personal communication). Their name is thus somewhat misleading.
In several systems, a quasi-stable periodicity, a few per cent shorter than the orbital period has been observed. These variations are called ‘$`\mathrm{𝐧𝐞𝐠𝐚𝐭𝐢𝐯𝐞}`$ $`\mathrm{𝐬𝐮𝐩𝐞𝐫𝐡𝐮𝐦𝐩𝐬}`$’. In a few systems they appear simultaneously with the positive superhumps (Patterson et al. 1997b; Arenas et al. 2000; Retter & Hellier 2000; Retter et al. 2000), but in other cases they are the only kind of superhump observed. An alternation between the positive and negative superhumps has been observed in a few objects (e.g. Skillman et al. 1998). The negative superhump deficit over the orbital period seems to be correlated with the orbital period in a manner similar to the Stolz & Schoembs (1984) relation, however with a shallower trend. It has been suggested that negative superhumps are formed by the precession of the accretion disc in the azimuthal axis (Patterson et al. 1993c; Patterson 1999), but there are some theoretical difficulties with this idea (Murray & Armitage 1998; Wood, Montgomery & Simpson 2000).
Osaki (1996) further proposed that only the values of the orbital period and the mass transfer rate determine the basic differences among the four major subclasses of non-magnetic CVs (U Gem systems, SU UMa stars, permanent superhumpers and nova-likes). Fig. 1 is taken from his paper. The permanent superhump systems are thought to be in a state of higher accretion rate (implying large accretion disc radii) than in regular SU UMa systems. Osaki’s model describes pretty well most superhump observations, although his suggestion that the period gap separates between systems with tidally unstable accretion discs and tidally stable objects is violated by the presence of many permanent superhump systems above the gap (Patterson 1999; Retter & Hellier 2000; Retter et al. 2000; see also Table 1 and Fig. 2).
### 1.3 Superhumps in classical nova systems
The permanent superhump model has been invoked so far for three classical novae: V603 Aql (Patterson & Richman 1991; Patterson et al. 1993c; Patterson et al. 1997b), CP Pup (White & Honeycutt 1992; White, Honeycutt & Horne 1993; Thomas 1993; Patterson & Warner 1998) and V1974 Cyg (Retter, Leibowitz & Ofek 1997; Skillman et al. 1997). In all cases an alternative magnetic explanation has been proposed as well (Haefner & Metz 1985; White et al. 1993; Balman, Orio & Ogelman 1995; Semeniuk et al. 1995; Olech et al. 1996), however the arguments for the superhump explanation seem much stronger. The periods of the three novae fit well in the Stolz & Schoembs (1984) diagram, and in two cases negative superhumps have been observed in addition to the presence of positive superhumps (Patterson et al. 1997b; Patterson 1999). Nova V4633 Sgr 1998 is another permanent superhump candidate (Lipkin, personal communication; see also Lipkin & Leibowitz 2000).
Retter & Leibowitz (1998) introduced a simple way of testing the thermal stability state of accretion discs in CVs. Employing this method on the three permanent superhump novae they found that these systems are indeed thermally stable, while the progenitor of V1974 Cygni was located below the critical line for stability. This result led Retter & Leibowitz to suggest that if the decline from outburst in V1974 Cyg towards the pre-nova magnitude continues, the post-nova should evolve into a regular SU UMa system with superhumps appearing only during superoutbursts. Alternatively, its disc might stay optically thick, keeping above the thermal stability limit, and continue to show superhumps permanently in its light curve. Retter & Leibowitz thus proposed that non-magnetic classical novae can be progenitors of permanent superhump systems. Retter, Naylor & Leibowitz (1999) and Retter & Naylor (2000) developed this proposal. In this work we elaborate and extend these ideas.
### 1.4 Existing nova cycle models for CVs
The different subclasses of CVs share a similar configuration, namely a binary system with a primary white dwarf and a Roche-lobe filling secondary red dwarf. Spectra of old novae usually show strong continua and emission lines, very similar to nova-like systems (Bode & Evans 1989; Warner 1995a). A few observations further support a possible connection between classical novae and dwarf novae. Two old novae experience regular dwarf nova outbursts a few decades after their eruption – the peculiar nova, GK Per 1901 (Sabbadin & Bianchini 1983), and V446 Her 1960 (Honeycutt, Robertson & Turner 1995; Honeycutt et al. 1998). Dwarf nova outbursts a few decades before the nova eruption of V446 Her have probably been observed as well. Livio (1989) and Warner (1995a) listed a few other cases, however, the evidence for such a transition in these systems is rather poor.
Robinson (1975) compared the magnitudes of 18 old novae with the values of their progenitors. He found no significant difference between these numbers, and concluded that all novae return to their pre-outburst luminosities. The observational properties of old novae are very similar to those of nova-likes, which have thermally stable accretion discs (Warner 1995b). Therefore, it seems that the outburst does not alter the thermal stability at all, at least for time scales of the order of a few decades. An exception in Robinson’s sample is V446 Her, mentioned above. It is, however, the only clear case of a regular nova<sup>1</sup><sup>1</sup>1We exclude GK Per as it is a very atypical nova – its orbital period is very long ($``$2 d), and its secondary star is believed to be a sub-giant (e.g. Dougherty et al. 1996) unlike red dwarf companions in classical novae. that had dwarf nova outbursts a few decades after (and probably even before) the nova eruption. Recently, it was found that V446 Her has an orbital period near 5 h (Thorstensen & Taylor 2000) – typical for classical novae.
Vogt (1990) investigated a sample of 97 old novae. His results seem to confirm the idea that there is no systematic difference in the brightness of pre-novae and post-novae. He also showed that the brightness of old novae tends to fade slowly in the decades following their outbursts. This finding was confirmed by another study (Duerbeck 1992). A possible interpretation of the decrease in the nova light is a future transition to a different phase (dwarf nova).
Therefore it was suggested that nova outbursts link different CV subclasses. The nova cycle of the subgroups of CVs is, however, still debatable, and several scenarios have been offered for the connections among these subclasses.
The ‘hibernation scenario’ (Shara et al. 1986; Prialnik & Shara 1986; Shara 1989) suggests that dwarf novae $``$ nova-likes $``$ novae $``$ nova-likes $``$ dwarf novae $``$ ‘hibernation’ $``$ dwarf novae etc. However, it was later proposed that the ‘hibernation’ phase ($`\dot{M}`$$``$0) might not exist at all (Livio 1989), thus dwarf novae $``$ nova-likes $``$ novae $``$ nova-likes $``$ dwarf novae… The typical time scales for the transitions were estimated as a few centuries – millenia.
An alternative view to the ‘hibernation scenario’ and to the ‘modified / mild / modern hibernation scenario’ was presented by Mukai and Naylor (1995). They suggested that nova-likes and dwarf novae constitute different classes of pre-nova systems. Therefore, there are two possibilities: 1. nova-likes $``$ novae $``$ nova-likes… 2. dwarf novae $``$ novae $``$ dwarf novae… Nova-likes should have more frequent nova outbursts than dwarf novae because their mass transfer rates are larger than those of dwarf novae. The critical mass for the thermonuclear runaway is thus achieved much faster. Transitions between the two phases are allowed on the long term scale.
It seems that the observations of old novae have not been able to judge between the various models (Naylor et al. 1992). Furthermore, these scenarios were suggested when there were essentially no known non-magnetic novae below the period gap, and before the discovery of the permanent superhump class. In this work we extend the models to the short orbital period CVs and test them by the observations accumulated so far.
## 2 The location of the permanent superhump systems in the ($`𝐏_{\mathrm{𝐨𝐫𝐛}}`$, $`\dot{𝐌}`$) plane
To test Osaki’s suggestion (1996), that the accretion discs in permanent superhump systems are thermally stable (unlike the discs of SU UMa systems that are unstable in quiescence, and become quasi-stable only during superoutbursts), we locate these systems in the ($`P_{orb},\dot{M}`$) plane. We basically follow the method developed by Retter & Leibowitz (1998). It essentially uses $`m_V`$ to estimate $`\dot{M}`$, and assumes that the accretion disc is the dominant light source in the V band. It is also assumed that the disc is not kept at a high temperature due to irradiation by the white dwarf, which is relatively hot in post-novae. A modification that we add to these calculations is the effect of the inclination angle, i, on the visual magnitude, expressed by Warner (1987):
$$M_i=2.5log[(1+1.5cos(i))cos(i)]$$
(1)
where $`M_i`$ is the magnitude change as a function of i. The resulting equations equivalent to equations (7) (for the critical instability value) and (8) (for the calculation of accretion rates) of Retter & Leibowitz (1998) are:
$$(m_V)_{crit}=2.164.25logP_{orb}3.33logM_1+5logd+A_VM_i$$
(2)
$$\dot{M_{17}}=(10^{\frac{m_V+M_iA_V0.69}{2.5}})\frac{d^2}{M_1^{4/3}}$$
(3)
where our symbols are identical with those of Retter & Leibowitz.
In Table 1, we present values of the relevant parameters of permanent superhump systems compiled from various sources. The list of objects is primarily based on Patterson (1999). Only ‘conventional’ permanent superhump systems (i.e. only positive superhumpers) were chosen. Systems showing only negative superhumps in their light curves were rejected from our sample. Note that evidence for the presence of positive superhumps in TV Col has been presented by Retter & Hellier (2000) and Retter et al. (2000). The permanent superhump AM CVn systems were also excluded. The mass transfer rates were calculated by equation (3) using the values in the table as input parameters.
The results are depicted in Fig. 2. The consistency of the method for estimating the accretion rate was discussed by Retter & Leibowitz (1998). The inclination effect that is added in this study is not severe for the three systems discussed in that work, and the values that we obtain for the accretion rates of the three permanent superhump novae are rather similar to those deduced by Retter & Leibowitz. We thus conclude that the modified version is consistent with other methods of estimating mass transfer rates.
## 3 Discussion
### 3.1 Are permanent superhump systems thermally stable?
One of our first aims in this paper was to check whether the accretion discs in permanent superhump systems are indeed thermally stable as was proposed by Osaki (1996). The results found in the previous section, and presented in Fig. 2 seem to support this suggestion. Among the eight objects in our sample with enough data, all have relatively high accretion rates. These values exceed in general the typical mass transfer rates in dwarf nova systems ($`10^{15}`$-$`10^{16}`$gr/sec), and are of the same order of the typical numbers in nova-likes and old novae ($`10^{17}`$-$`10^{18}`$gr/sec – Warner 1995a). Four systems (CP Pup, V1974 Cyg, TT Ari and V603 Aql) are located above the critical line for thermal instability. The permitted ranges of accretion rates of three of the other permanent superhumpers (BK Lyn, AH Men and TV Col) straddle the thermal instability line and are thus consistent with their being stable. Note that the distance estimate for BK Lyn given by Dhillon et al. (2000) is a lower limit thus implying higher values of accretion rates (see the corresponding arrow in Fig. 2). V592 Cas is the only permanent superhump system found below the instability line according to our calculations. Its model-dependent distance estimate (Huber et al. 1998) is very small, and is almost certainly incorrect as its distance to interstellar reddening ratio is about a factor smaller than in the other objects in Table 1. The interstellar reddening thus implies a distance of $``$10 times larger, and therefore mass transfer rates $``$100 times larger (see the corresponding arrow in Fig. 2).
There are many observational and theoretical uncertainties in calculating these mass transfer rates. In fact even the location of the critical thermal instability line itself is controversial. Warner’s (1995b) border line, for example, is placed about a factor of two below Osaki’s (1996) critical line. We are looking, however, at effects of the order of $``$10. We thus conclude that the accretion discs of permanent superhump systems are most likely thermally stable, and that they indeed form a unique subgroup of CVs, with different physical parameters (namely short orbital periods and high mass transfer rates) from other CV subclasses. This finding further supports the suggestion mentioned above that the distance estimate to V592 Cas quoted in Table 1 was underestimated, and a better measurement should raise its location above the critical line. Similarly, the upper limits on the distances of BK Lyn, AH Men and TV Col are preferred to the lower values.
### 3.2 Angular momentum conservation in permanent superhump systems
Magnetic braking can account for the mean mass transfer rates of systems above the period gap, but is believed to occur only for systems with orbital periods larger than about 2.7 h. Braking by gravitational radiation can explain the typical accretion rates in dwarf novae below the gap, but there seems to be a serious problem for the short period permanent superhump systems (Fig. 2). The mass accretion rates estimated for these systems in Section 2 are $``$$`10^{18}`$gr/sec, and thus more than two orders-of-magnitude higher than those yielded by gravitational radiation. The presence of permanent superhump systems below the gap is therefore not understood within the current models, if their mass transfer rates represent mean secular values.
One solution to this problem is to invoke an extra source of angular momentum loss below the period gap, the obvious candidate being magnetic braking. It is often stated that the secondary stars in CVs below the gap lack the radiative core required to anchor a magnetic field. In fact, single M stars (which are typical secondary stars in short orbital period CVs – Smith & Dhillon 1998), which are fully convective, do show magnetic activity (Fleming, Schmitt & Giampapa 1995). Whilst invoking such braking would preclude certain explanations of the period gap itself, which rely on the cessation of magnetic activity (e.g. Verbunt 1984), it would allow a natural explanation for the presence of high accretion rates in systems with orbital periods below the period gap. Therefore, we suggest that the mere existence of permanent superhump systems below the gap might be an argument that magnetic braking does not cease at the gap.
A different solution to the problem is provided by the possibility of mass transfer cycles. If the permanent superhump phase is very short lived, and systems below the gap spend most of their time as SU UMa systems, perhaps even hibernating for a while, then the mean mass transfer rate might not exceed the Gravitational Radiation values. An argument against such a possibility is that King et al. (1996) could not produce self sustaining mass transfer cycles in short period systems. However, as we shall show in the next section, there may be emerging observational evidence that nova explosions drive such cycles below the gap.
### 3.3 Nova cycle scenarios for non-magnetic short orbital period CVs
Post-novae with orbital periods above the gap usually return to their pre-outburst brightness within a few decades (Section 1.4). So far only two non-magnetic novae below the period gap, CP Pup and V1974 Cyg, have been found. Both have permanent superhumps in their light curves suggesting that their discs are thermally stable. Table 2 presents the pre-outburst and post-eruption magnitudes of the two systems. The observations are consistent with the two post-novae being brighter than their progenitors. Moreover, the pre-nova magnitude of V1974 Cyg implies that its disc was thermally unstable and therefore the progenitor of the nova should have been an SU UMa system (Retter & Leibowitz 1998). Thus, it is the first example of a nova that has changed the thermal stability state of its accretion disc<sup>2</sup><sup>2</sup>2It is unknown whether V446 Her and GK Per, the two novae that show dwarf nova outbursts in their light curves (Section 1.4), had nova-like stages before or after their nova outbursts. Note that it is unclear whether accretion discs survive nova eruptions.. This transition might be explained by a temporary change in the mass transfer rate due to the hot post-nova white dwarf irradiating the secondary star or the disc itself. Alternatively, V1974 Cyg is the first direct evidence for mass transfer cycles, in this case driven by the nova explosion. CP Pup is a candidate for the same behaviour. The nova cycle scenarios proposed for long orbital period CVs (Section 1.4) could thus be applicable to the short period systems, but possibly with some minor modifications.
## 4 Summary and conclusions
Our results can be summarized as follows:
1. We established the idea that the accretion discs in permanent superhump systems are thermally stable.
2. If the high values of accretion rates found in the short orbital period permanent superhumpers represent their secular mean values, then there is a problem with the mechanism that removes angular momentum from the systems. It might be solved by extending the magnetic braking mechanism to below the gap, or by invoking another way to lose angular momentum in these CVs. Alternatively, systems below the period gap might have mass transfer cycles, and the permanent superhump stage should be short compared with the full CV cycle.
3. We suggest that nova cycle scenarios similar to those proposed for the long orbital period CVs should be applied to the short period systems. However, they would be slightly modified if the superhumps observed in the post-novae below the gap represent a true long-term change in the mass transfer rate rather than a transient increase in the disc luminosity due to irradiation by the hot white dwarf. The observations of the two non-magnetic novae below the gap seem to support mass transfer cycles driven by the nova outburst.
## 5 Acknowledgments
Hans Ritter, the referee, is acknowledged for many valuable comments. We also thank Dina Prialnik, Elia Leibowitz, Mike Shara, Joe Patterson, Rob Jeffries, Michael Friedjung, Coel Hellier and Sandi Catalan for fruitful discussions. TN is a PPARC advanced fellow. AR is supported by PPARC.
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