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# 1 Introduction
## 1 Introduction
The coupling of scalar matter fields with bidimensional dilaton gravity, originally proposed by Callan, Giddings, Harvey and Strominger (CGHS) has attracted attention due to connections with black hole physics , specially within the context of formation and evaporation of black holes. In cosmology the dilaton field appears to be important, and so several mechanisms for cosmological dilaton production has been discussed in the literature . The CGHS model has its origin in the dimensional reduction of the more realistic, four dimensional, Einstein-Hilbert gravitation with a spherically symmetric metric. The dilaton field in the resulting model is a relic of the angular variables hidden in this procedure. Also, it belongs to the class of models first discussed by Jackiw and Teitelboim , which lend themselves to a gauge theoretical formulation of the problem.
The formation of black holes in bidimensional dilaton gravity coupled to scalar field was considered recently in different contexts. For instance, the scalar matter field was already considered in the form of sine-Gordon and quartic potentials , and these investigations have motived us to introduce new systems for the scalar matter.
In this paper we investigate the formation of black hole in two different systems. The first system considers a sixth-order potential for the matter field that couples to the bidimensional dilaton gravity. This is a new possibility, and although we yet have polynomial potential, the nature of the solutions are different from the one that appear with the quartic potential already investigated, since in the sixth-order potential the kink connects the symmetric vacuum to an asymmetric one. The second system we shall deal with is a system of two coupled fields recently introduced . This system presents very interesting properties , and here we show that the second field adds further effects.
This work is organized as follows. In the next Sec. 2 we present the CGHS model with a single field system, given by a sixth-order potential, in Sec. 2.1. The two field system as the source of matter is worked out in Sec. 2.2 and again we have black hole solutions, now controlled by an enlarged number of parameters. Discussions and final remarks are introduced in Sec 3.
## 2 Models of Matter Fields
In this section we deal with two different examples of coulping dilaton gravity to matter fields. The first example represents a single field system that contains up to sixth power self-interaction terms. The second example is different since it contains two coupled scalar fields, and is defined by a potential that contains up to the quartic power in the fields, and presents an enlarged set of parameters.
### 2.1 System with a single field
Here we consider the model described by the action
$`S`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle }d^2x\sqrt{\overline{g}}\{e^{2\varphi }[\overline{R}+4(\overline{}\varphi )^2+4\lambda ^2]`$ (1)
$`{\displaystyle \frac{1}{2}}(\overline{}f)^2+2\mu ^2f^2(f^2a^2)^2e^{2\varphi }\},`$
where $`\overline{g}`$, $`\varphi `$ and $`f`$ are the metric, dilaton and matter fields respectively. $`\overline{R}`$ is the scalar curvature and $`\lambda ^2`$ is a cosmological constant. This action is the usual action, except that the last term contains a specific sixth-order potential for the self-interacting matter field.
We can be write Eq. $`(1)`$ in a different form, using a rescaled metric tensor $`g_{\mu \nu }`$ in such a way that
$$\overline{g}_{\mu \nu }=e^{2\varphi }g_{\mu \nu }.$$
(2)
In this case, the action given by Eq. $`(1)`$ turns into
$`S`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle }d^2x\sqrt{g}[e^{2\varphi }R+4\lambda ^2`$ (3)
$`{\displaystyle \frac{1}{2}}(f)^2+2\mu ^2f^2(f^2a^2)^2].`$
Note that this conformal reparametrization of the field eliminates the kinetic term for the dilaton that appears in the original action.
The equations of motion that follow from Eq. $`(3)`$ are
$$^2(e^{2\varphi })4\lambda ^22\mu ^2f^2(f^2a^2)^2=0,$$
(4)
and
$$R=0.$$
(5)
Equation $`(5)`$ implies that $`g_{\mu \nu }=\eta _{\mu \nu }`$ and using the fact that in two dimensions we can always put the metric in the conformal gauge
$$\overline{g}_{\mu \nu }=e^{2\rho }\eta _{\mu \nu },$$
(6)
we conclude that $`\rho =\varphi `$.
In order to investigate the solutions of this model it is useful to introduce light-cone coordinates $`x^\pm =t\pm x`$. In these coordinates the line element constructed with the metric tensor $`g_{\mu \nu }`$ is given by
$$ds^2=dx^+dx^{}.$$
(7)
In terms of these coordinates, the action given by Eq. $`(3)`$ can be cast to the form
$`S`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}{\displaystyle }d^2x[(2e^{2\varphi }_+_{}\rho +\lambda ^2)`$ (8)
$`{\displaystyle \frac{1}{2}}_+f_{}f{\displaystyle \frac{1}{2}}\mu ^2f^2(f^2a^2)^2].`$
The field equations of motion are given by
$`_+^2\left(e^{2\varphi }\right)+{\displaystyle \frac{1}{2}}\left(_+f\right)^2`$ $`=`$ $`0,`$ (9)
$`_{}^2\left(e^{2\varphi }\right)+{\displaystyle \frac{1}{2}}\left(_{}f\right)^2`$ $`=`$ $`0,`$ (10)
$`_+_{}\left(e^{2\varphi }\right)+\lambda ^2{\displaystyle \frac{1}{2}}\mu ^2f^2(f^2a^2)^2`$ $`=`$ $`0,`$ (11)
$`_+_{}f+\mu ^2f(f^2a^2)(3f^2a^2)`$ $`=`$ $`0,`$ (12)
and these are the equations we have to deal with.These results follow very naturally in the above procedure, and this should be contrasted to the procedure used in Refs. where one has to appropriately choose the gauge to get to such r esult. Here we see that the equation of motion for the matter field presents solutions that can be cast to the form, working with standard coordinates $`(x,t)`$
$$f^2(x,t)=\frac{1}{2}a^2\{1+\mathrm{tanh}[\alpha ((x\overline{x})+v(t\overline{t}))]\},$$
(13)
where $`\alpha `$ is a constant and $`(\overline{x},\overline{t})`$ represents the center of the kink. In the light-cone coordinates we can write the above solutions as
$$f^2(x^+,x^{})=\frac{1}{2}a^2\{1+\mathrm{tanh}[\alpha _+(x^+\overline{x}^+)\alpha _{}(x^{}\overline{x}^{})]\},$$
(14)
where $`\alpha _\pm =\frac{\alpha }{2}\left(v\pm 1\right)`$.
In order to study black hole formation we use this $`f`$ into the other equations of motion for $`\varphi `$. Here we can write, for instance,
$$_+_{}\left(e^{2\varphi }\right)=\frac{\mu ^2a^6}{16}\left[1+\mathrm{tanh}^3\mathrm{}\mathrm{tanh}^2\mathrm{}\mathrm{tanh}\mathrm{}\right]\lambda ^2,$$
(15)
where we have set $`\mathrm{}=\delta \overline{\delta }`$, with $`\delta =\alpha _+x^+\alpha _{}x^{}`$ and $`\overline{\delta }=\alpha _+\overline{x}^+\alpha _{}\overline{x}^{}`$. This equation can be integrated to give
$$e^{2\varphi }=C_1+b(x^+)+d(x^{})\lambda ^2x^+x^{}+\frac{\mu ^2a^6}{16\alpha _+\alpha _{}}\left[\frac{1}{2}\mathrm{tanh}\mathrm{}\mathrm{ln}\mathrm{cosh}\mathrm{}\right].$$
(16)
The functions $`b(x^+)`$ and $`d(x^{})`$ are determined by the constraint equations, the two first equations of motion $`(8)`$ and $`(9)`$. Here we get
$`b(x^+)`$ $`=`$ $`bx^++C_2,`$ (17)
$`d(x^{})`$ $`=`$ $`dx^{}+C_3.`$ (18)
Therefore, the dilaton field can be determined up to constants in the form
$$e^{2\varphi }=C+bx^++dx^{}\lambda ^2x^+x^{}+\frac{\mu ^2a^6}{16\alpha _+\alpha _{}}\left[\frac{1}{2}\mathrm{tanh}\mathrm{}\mathrm{ln}\mathrm{cosh}\mathrm{}\right],$$
(19)
where $`b`$, $`d`$, and $`C`$ are constants, and in the following we choose $`b=d=0`$, for simplicity.
Let us now investigate the geometric nature of this solution generated by original system. Toward this goal, let us divide spacetime into the three regions: $`\mathrm{}=\delta \overline{\delta }<<1`$, $`\mathrm{}0`$, and $`\mathrm{}>>1`$. In the first region the dilaton field becomes
$$e^{2\varphi }C\lambda ^2\left(x^++\frac{\mu ^2a^6}{16\lambda ^2\alpha _+}\right)\left(x^{}\frac{\mu ^2a^6}{16\lambda ^2\alpha _{}}\right),$$
(20)
where the constant $`C`$ is taken as
$$C=\frac{\mu ^2a^6}{16\alpha _+\alpha _{}}\left(\overline{\delta }+\frac{\mu ^2a^6}{16\lambda ^2}\mathrm{ln}2\frac{1}{2}\right).$$
(21)
In the region where $`\mathrm{}=\delta \overline{\delta }0`$ we get
$$e^{2\varphi }\frac{\mu ^2a^6}{16\alpha _+\alpha _{}}\left(\overline{\delta }+\frac{\mu ^2a^6}{16\lambda ^2}\mathrm{ln}2\frac{1}{2}\right)\lambda ^2x^+x^{}.$$
(22)
With regard to the initial metric $`\overline{g}_{\mu \nu }`$, these solutions represent, where this metric is defined , the linear dilaton vacuum.
In the third region we have $`\mathrm{}=\delta \overline{\delta }>>1`$, and now we obtain
$$e^{2\varphi }\frac{\mu ^2a^6}{16\alpha _+\alpha _{}}(2\overline{\delta })\lambda ^2\left(x^+\frac{\mu ^2a^6}{16\lambda ^2\alpha _+}\right)\left(x^{}+\frac{\mu ^2a^6}{16\lambda ^2\alpha _{}}\right),$$
(23)
which represents, with respect to the original background,the geometry of a black hole with mass
$$\frac{\lambda \mu ^2a^6}{8\alpha _+\alpha _{}}\overline{\delta },$$
(24)
after shifting $`x^+`$ by $`\mu ^2a^6/16\lambda ^2\alpha _+`$ and $`x^{}`$ by $`\mu ^2a^6/16\lambda ^2\alpha _{}`$.
### 2.2 System with two fields
Let us now consider another system, defined by
$`S`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle }d^2x\sqrt{g}\{[e^{2\varphi }R++4\lambda ^2]`$ (25)
$`{\displaystyle \frac{1}{2}}(f)^2{\displaystyle \frac{1}{2}}(g)^2+U(f,g),\}.`$
This action can be written in terms of $`\overline{g}_{\mu \nu }`$ by doing the inverse transformation that correspondes to Eq. $`(2)`$.
Using light-cone coordinates the above action can be cast to the form
$`S`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}{\displaystyle }d^2x[(2e^{2\varphi }_+_{}\rho +\lambda ^2)`$ (26)
$`+{\displaystyle \frac{1}{2}}_+f_{}f+{\displaystyle \frac{1}{2}}_+g_{}g{\displaystyle \frac{1}{4}}U(f,g)].`$
This new system is defined via the potential
$$U(f,g)=\frac{1}{2}\mu ^2(f^2a^2)^2+\mu \nu (f^2a^2)g^2+\frac{1}{2}\nu ^2g^4+2\nu ^2f^2g^2.$$
(27)
This potential identifies a system of two real scalar fields that was recently investigated in , where it was shown to present interesting static field configurations. As we can see from the above potential, we are now dealing with a richer system and we want to explore how this enlarged system change the simpler picture of black hole formation in systems of a single field.
We follow as in the former system.In this case we get
$`_+^2(e^{2\varphi })+{\displaystyle \frac{1}{2}}(_+f)^2+{\displaystyle \frac{1}{2}}(_+g)^2`$ $`=`$ $`0,`$ (28)
$`_{}^2(e^{2\varphi })+{\displaystyle \frac{1}{2}}(_{}f)^2+{\displaystyle \frac{1}{2}}(_{}g)^2`$ $`=`$ $`0,`$ (29)
$`_+_{}(e^{2\varphi })+\lambda ^2{\displaystyle \frac{1}{4}}U(f,g)`$ $`=`$ $`0,`$ (30)
$`_+_{}f+{\displaystyle \frac{1}{2}}\mu ^2(f^2a^2)f+\nu (\nu +{\displaystyle \frac{1}{2}}\mu fg^2`$ $`=`$ $`0,`$ (31)
$`_+_{}g+{\displaystyle \frac{1}{2}}\mu \nu (f^2a^2)g+{\displaystyle \frac{1}{2}}\nu ^2g^3+\nu ^2f^2g`$ $`=`$ $`0.`$ (32)
The above Eqs. $`(31)`$ and $`(32)`$ correspond to the matter field equations of motion. They can be solved to give two different pair of solutions :
$$f=a\mathrm{tanh}\{\mu a[\alpha _+(x^+\overline{x}^{})\alpha _{}(x^{}\overline{x}^{})]\},$$
(33)
and $`g=0`$, and also
$$f=a\mathrm{tanh}\{2\nu a[\alpha _+(x^+\overline{x}^+)\alpha _{}(x^{}\overline{x}^{})]\},$$
(34)
and
$$g=\pm a\left(\frac{\mu }{\nu }2\right)^{1/2}\mathrm{sech}\{2\nu a[\alpha _+(x^+\overline{x}^+)\alpha _{}(x^{}\overline{x}^{})]\},$$
(35)
valid for $`\mu /\nu >2`$. Here we note that the limit $`\nu \mu /2`$ transforms the second pair of solutions into the first one, that presents $`g=0`$. This is interesting since the investigation of the second pair of solutions allows getting results valid for the first pair, with presents $`g=0`$ and so leads to the case of just one field. See below for further details.
For the second pair of solutions we can cast the dilaton field in the form
$$e^{2\varphi }=C\lambda ^2x^+x^{}A\mathrm{tanh}^2[2\nu a(\delta \overline{\delta })]B\mathrm{ln}\mathrm{cosh}[2\nu a(\delta \overline{\delta })].$$
Like in the former case, in the above expression we have also chosen $`b=d=0`$ in $`b(x^+)=bx^++C_2`$ and $`d(x^+)=dx^++C_3`$, which follow from $`(28)`$ and $`(29)`$. The expressions for $`A`$ and $`B`$ are given by
$$A=\frac{a^2}{12\nu ^2\alpha _+\alpha _{}}\left[\frac{\mu ^2}{4}+\frac{1}{4}(\mu 2\nu )(7\mu +10\nu )\right],$$
(36)
and
$$B=\frac{a^2}{12\nu ^2\alpha _+\alpha _{}}\left[\mu ^2+(\mu 2\nu )(\mu 4\nu )\right].$$
(37)
Let us now investigate the geometric nature of the black hole generated in this system. Here we can write the dilaton field ,for $`\mathrm{\Delta }>>1`$, as
$$e^{2\varphi }4aB\nu \overline{\delta }\lambda ^2\left(x^+2\frac{\nu }{\lambda ^2}Aa\alpha _{}\right)\left(x^{}+2\frac{\nu }{\lambda ^2}Aa\alpha _+\right),$$
(38)
which represents the geometry of a black hole with regard to the original background spacetime.The mass of the black hole is given by
$$4aB\lambda \nu \left(\overline{\delta }+\frac{\nu }{\lambda ^2}Aa\alpha _+\alpha _{}\right).$$
(39)
after shifting $`x^+`$ by $`2(\nu /\lambda ^2)aB\alpha _{})`$ and $`x^{}`$ by $`2(\nu /\lambda ^2)aB\alpha _+`$.
For $`\mathrm{\Delta }=\overline{\delta }\delta <<1`$,
$$e^{2\varphi }\lambda ^2\left(x^++2\frac{\nu }{\lambda ^2}Aa\alpha _{}\right)\left(x^{}2\frac{\nu }{\lambda ^2}Aa\alpha _+\right),$$
(40)
where the constant $`C`$ was taken as
$$C=A+2aB\nu \overline{\delta }B\mathrm{ln}2\frac{4\nu ^2a^2B^2\alpha _+\alpha _{}}{\lambda ^2}.$$
(41)
In the region of $`\mathrm{\Delta }0`$ we have
$$e^{2\varphi }AB\mathrm{ln}2+2aB\nu \overline{\delta }\frac{4a^2B^2\nu ^2\alpha _+\alpha _{}}{\lambda ^2}\lambda ^2x^+x^{},$$
(42)
As in the first system, these last two solutions give us the linear dilaton vacuum in the region where the original background is defined.
We recall that the limit $`\nu \mu /2`$ changes the second pair of solutions that we have been considering to the first one, simpler, that presents $`g=0`$, as introduced above. For this reason, we can investigate this first pair of solutions by just setting $`\nu \mu /2`$ in the above results.
This procedure is interesting since it leads to the case with just one field, more precisely to the case where the matter system is described by the $`\varphi ^4`$ model, but this was already investigated in Ref. . Despite slightly different notations, it is not hard to see that the limit $`\nu \mu /2`$ correctly change the above results into the results given in for the case of a single field, in the $`\varphi ^4`$ model for the matter contents.
## 3 Comments and Conclusions
In this work we have investigated bidimensional dilaton gravity coupled to two different matter field systems. The first system is a single field system that contains self-interactions up to the sixth power. It is of the same kind of the sine-Gordon and $`\lambda \varphi ^4`$ models already investigated. The results of these papers, together with the present work show that the black hole that appears is qualitatively the same, but with different masses that depend on the parameters associated with the various solutions. The second system is a system of two fields, and contains up to the fourth power in the fields. This system is richer since it is defined in a enlarged space of parameters, which contains the space of parameters of the $`\lambda \varphi ^4`$ system as a particular case. The soliton solutions that appear in this case also contributes to the generation of black holes, but these black holes are quantitatively different from black holes that appear in systems of a single field, the difference being controlled by the enlarged set of parameters that defines the two field matter system.
Evidently, one can have more examples, for instance investigating the case where the scalar matter is described by the two field system that contains up to the sixth power in the fields, as also investigated in . We think that it is interesting to investigate this kind of mechanism for different matter field potentials, in order to understand mechanisms of formation of black holes within the context of dilaton gravity. Furthermore, there is also the newer context of the Einstein-Maxwell-dilaton-axion system with general dilaton coupling. And this naturally leads to another interesting issue, that concerns investigating the coupling of a dilaton-axion gravitational field with two coupled scalar fields, in the form introduced in . The two cases worked out above lead us to different scenario concerning formation of black holes, and so it turns out to be interesting to analyze quantum or at least semiclassical versions since these black holes may present new termodinamics properties, namely the Hawking-Bekenstein radiation. Work on this and in other related issues is now in progress.
D.B. would like to thank Roman Jackiw for comments and for reading the manuscript.
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# Structural Properties of the Sliding Columnar Phase in Layered Liquid Crystalline Systems
## I Introduction
The search for non-viral vectors for transport of DNA across cell and nuclear membranes in gene therapy has led to the study of DNA-cationic-lipid complexes. DNA and mixtures of neutral (zwitterionic) and cationic lipids dispersed in water self-assemble into spheroidal complexes that can attain micron sizes near the isoelectric point where there is compensation between DNA and lipid charge. Fluorescence tagging of DNA and lipids indicate that both species are dispersed uniformly throughout the complexes. X-ray scattering experiments reveal two bulk structures for the complexes with close association between DNA and lipids. When helper lipids that favor spontaneous curvature of lipid membranes or that decrease membrane charge density are added, a hexagonal inverted micellar structure forms with DNA molecules captured in the water holes of the hexagonal lattice. When only neutral and cationic lipids are used, DNA is intercalated in galleries between lamellae of a lamellar lyotropic phase formed by the lipids as depicted schematically in Fig. 1. This article will investigate the possible equilibrium phases of these lamellar DNA lipid complexes. It will focus primarily on the properties of one phase, the sliding columnar phase, characterized by strong orientational but weak positional correlation between DNA strands in neighboring galleries.
DNA molecules are semi-flexible polymers that when confined to a two-dimensional plane tend to form locally aligned structures with a preferred inter-molecular separation that can be modeled as two-dimensional smectic liquid crystals. Thus, lamellar DNA-lipid complexes can be viewed as a $`3D`$ stack of $`2D`$ smectics as depicted in Fig. 1. This figure establishes the coordinate conventions that will be used throughout this paper. DNA strands align on average along the $`x`$ direction, and the normal to lamellae is along the $`y`$ axis. The normal to the $`2D`$ smectic lattices is along the $`z`$ direction. The average spacing between lamellae is $`a`$, and that between DNA strands is $`d`$. The wavenumbers associated with these lengths are, respectively, $`k_0=2\pi /a`$ and $`q_0=2\pi /d`$.
If the lamellae are assumed not to be disrupted by dislocations or rips, then the following possible equilibrium phases are easily identified:
* Columnar (C) Phase: In this phase, DNA strands are aligned on average along the $`x`$ axis, and their centers occupy positions on a $`2D`$ crystal lattice in the $`yz`$ plane. Since the standard Coulomb repulsion between DNA strands favors staggering of smectic lattices in neighboring galleries, the columnar lattice is normally expected to be centered rectangular as observed in experiments on complexes in which membranes are in the $`L_\beta ^{}`$ phase rather than the more disordered $`L_\alpha `$ phase. We will, however, consider both simple rectangular and centered rectangular columnar lattices. A simple rectangular lattice may occur if the effective interaction between DNA strands in different galleries is attractive, as is the case between DNA strands in solutions with polyvalent salts . The Columnar phase is favored at low temperature. It has the same symmetry as a columnar discotic liquid crystal phase. It is characterized by a $`2D`$ elastic energy with a nonvanishing shear modulus for relative displacements of DNA lattices in different galleries and nonvanishing moduli for compression of both lamellae and DNA smectic lattices. It has long-range positional order in the $`2D`$ $`yz`$-plane and associated Bragg peaks in its x-ray scattering profile at reciprocal lattices vectors $`𝐆_{m_1,m_2}=[0,m_1k_0/2,m_2q_0](m_1,m_2)`$ where $`m_1`$ and $`m_2`$ and integers and $`m_1`$ is even for a simple rectangular lattice and $`m_1+m_2`$ is even for a centered rectangular lattice as shown Fig. 2.
* Nematic Lamellar (NL) Phase: In this phase, the periodic positional order of the columnar phase is destroyed by dislocations in the DNA smectic lattices, but the long-range orientational order of DNA strands is maintained. This phase is characterized by a long-wavelength elastic energy with a lamellar compression modulus, an anisotropic lamellar bending modulus, and orientational rigidities (Frank elastic constants) opposing spatially dependent variation of DNA alignment direction. Both the shear modulus for relative displacement of DNA lattices in different galleries and the compression modulus for DNA lattices vanish, and there is exponential decay of DNA positional correlations. The x-ray scattering profile of the NL phase exhibits lamellar power-law $`(2l,0)`$ peaks at $`𝐆_{2l,0}=(0,lk_0,0)`$. If columnar-phase positional correlations are well developed, it will also exhibit Lorentzian peaks at $`𝐆_{m_1,m_2}`$ for $`m_20`$ as depicted in Fig. 2. If thermal fluctuations are sufficiently strong that these correlations are not well developed in centered rectangular systems, then the x-ray scattering profile could exhibit a broad $`(0,1)`$ peak in the vicinity of $`(0,0,q_0)`$ rather than the expected pair of $`(1,1)`$ and $`(1,1)`$ peaks. In this case, the positions of the x-ray scattering peaks would appear to indicate a tendency to form a simple rectangular rather than the ground-state centered rectangular structure.
* Isotropic Lamellar (IL) Phase: In this phase, orientational as well as positional order of DNA lattices is lost. Macroscopically this phase is identical to a multi-component isotropic lamellar phase. It is included for completeness and will not be considered further in this paper.
* Decoupled $`2D`$ Smectic (DS) Phase: This phase occurs only if there are absolutely no interactions between DNA lattices in neighboring galleries. Its elasticity and correlations are thus those of independent $`2D`$ smectic lattices. Since there are always interactions between galleries, this phase will not exist in real systems. It is, however, a useful limit to consider since systems with weak coupling between layers will behave as though they are decoupled at sufficiently short length scales.
In addition to the above phases, which are straightforward to identify, there is the possibility of another phase with unusual properties:
* Sliding Columnar (SC) Phase: This phase has properties intermediate between those of the columnar and lamellar nematic phases. Its elastic energy is distinguished from that of the NL phase by the presence of a nonvanishing modulus for compression of DNA lattices. In-plane smectic correlations die off as $`\mathrm{exp}(\mathrm{const}.\mathrm{ln}^2r)`$ as a function of separation $`r`$ rather than exponentially as in the NL phase. Correlations between smectic lattices in different galleries die off exponentially with layer-number difference. The x-ray structure factor exhibits power-law lamellar $`(0,lk_0,0)`$ peaks and well defined DNA $`(m_1,m_2)`$ peaks with $`m_20`$ that are sharper than the corresponding Lorentzian peaks in the NL phase, as depicted in Fig. 2. As in the NL phase, $`(1,1)`$ and $`(1,1)`$ peaks may merge to produce a single $`(0,1)`$ DNA peak if thermal fluctuations are sufficiently strong.
Rather stringent conditions must be met before the SC phase can be thermodynamically stable. Thermal fluctuations must be strong enough to destroy the interlayer shear coupling present in the columnar phase but not so strong that dislocation proliferation destroys the smectic compressibility to create a nematic lamellar phase. The SC phase is stable only for temperatures $`T`$ lying above a decoupling temperature $`T_d`$ at which the C phase becomes unstable and below a KT-melting temperature $`T_{KT}`$ above which the NL phase becomes stable. Thus a necessary condition for the SC phase to be stable is $`T_d<T_{KT}`$. In Sec. IV, we will show that this condition is violated for the nearest neighbor models we consider here. Elsewhere, we show that appropriately chosen interactions between further neighbor planes can stabilize the SC phase.
The focus of the paper is the rather unusual properties of the SC phase. We will, therefore, assume throughout most of the paper that the SC phase does exist. This approach is justified because it can under appropriate conditions be an equilibrium phase as discussed in Ref. and because even if it is not thermodynamically stable, there is a range of length scales over which correlations functions will exhibit SC behavior.
This paper is composed of six sections, of which this is the first, and seven appendices, which mostly present mathematical details. Section II derives the Hamiltonians for the columnar, sliding columnar, and nematic lamellar phases, based on a model in which the dominant interactions are couplings between each DNA lattice and the two membranes on either side of it. Section III presents the correlation functions for the SC phase, including x-ray scattering intensities. Section IV addresses the stability of the SC phase. It calculates $`T_d`$, and it derives general expressions for the interaction between dislocations before calculating $`T_{KT}`$. Section V presents a discussion of the various important length scales. Finally, Sec. VI provides an overview of results. Appendix A presents details of the derivation of the Hamiltonians in Sec. II. Appendices B through F provide details of the calculations of various SC phase correlations functions, and App. G presents details of calculations of interactions between dislocations.
## II Model for Lamellar Phases of DNA-lipid complexes
### A Definition of Variables
As depicted in Fig. 1, lamellar phases of DNA-lipid complexes consist of a periodic stack of planar lipid bilayer membranes with spacing $`a`$ separated by galleries intercalated with DNA strands that form a local $`2D`$ smectic lattice with preferred spacing $`d`$. Since our primary interest is in the nature of possible ordering of the DNA strands, we will assume that the lipid membranes are free of defects such as dislocations or focal conic structures that destroy their integrity. We can, therefore, specify a layer by its integer layer number $`n`$, and we can specify positions in the $`xz`$ plane by the vector $`𝐫=(x,z)`$. We take the equilibrium height of the midpoint of bilayer membrane $`n`$ be $`na`$. The height of membrane $`n`$ at position $`𝐫`$ is then
$$H_{\mathrm{mem}}^n(𝐫)=na+h^n(𝐫),$$
(2.1)
where $`h^n(𝐫)`$ is the Lagrangian height variable measuring the displacement of layer $`n`$ from its ideal height. The $`n`$th DNA lattice lies between membranes $`n`$ and $`n+1`$. Its equilibrium height is, therefore, $`(n+\frac{1}{2})a`$, and its height at position $`𝐫`$ is
$$H_{\mathrm{DNA}}^n(𝐫)=(n+\frac{1}{2})a+u_y^n(𝐫),$$
(2.2)
where $`u_y^n(𝐫)`$ measure deviations from equilibrium height.
The DNA lattice in layer $`n`$ can be described in terms of a Fourier expansion of its density. It is important to keep track of the relative phases of mass-density waves in adjacent layers. In an ideal columnar lattice, the phase of the first mass-density wave in layer $`n`$ is simply $`q_0(zz^n)`$, where
$$q_0=2\pi /d$$
(2.3)
and $`z^n`$ specifies the preferred position of lattices in different galleries. In a rectangular lattice, $`z^n=0`$, and in a centered rectangular lattice, $`z^n=nd/2`$. Phase changes at position $`𝐫`$ relative to the above ideal phases are produced by local translations described by the Eulerian displacement variable $`u_z^n(𝐫)`$. Thus the DNA density-wave expansion in layer $`n`$ is
$$\rho _{\mathrm{DNA}}^n(𝐫)=\underset{k}{}\rho ^{nk}(𝐫)e^{ikq_0[zz^nu_z^n(𝐫)]}.$$
(2.4)
If each DNA lattice is perfect with identical rods of linear density $`\lambda `$ and separation $`d`$, then $`\rho ^{nk}(𝐫)=\lambda /d`$ for every $`n`$ and $`k`$. Thermal fluctuations will lead to reductions in $`\rho ^{nk}`$ that are larger for larger $`k`$. Thus, $`\rho _{DNA}^n(𝐱)`$ can be approximated by
$$\rho _{DNA}^n(𝐱)\rho _0+\psi e^{iq_0[zz^nu_z^n(𝐫)]}+\mathrm{c}.\mathrm{c}.,$$
(2.5)
where $`\psi `$, assumed to be independent of $`n`$, is the complex amplitude of the first non-trivial density wave.
The total mass density of DNA strands at positions $`𝐱=(x,y,z)(𝐫,z)`$ is
$`\rho _{\mathrm{DNA}}(𝐱)`$ $`=`$ $`{\displaystyle \underset{nk}{}}\rho ^{nk}(𝐫)e^{ikq_0[zz^nu_z^n(𝐫)]}`$ (2.7)
$`\times f_{DNA}^k[y(n+{\displaystyle \frac{1}{2}})au_y^n(𝐫)],`$
where $`f_{DNA}^k(y)`$ is the form factor along the layer normal for a DNA mass-density wave of wavenumber $`kq_0`$. If the DNA stands are lines with no width, the $`f_{DNA}^k(y)=\delta (y)`$. If the strands are modeled as cylinders of radius $`r_{DNA}`$, then the Fourier transform of $`f_{DNA}^k(y)`$ is
$$f_{DNA}^k(q_y)=\frac{2J_1(\sqrt{q_y^2+k^2q_0^2}r_{DNA})}{\sqrt{q_y^2+k^2q_0^2}r_{DNA}},$$
(2.8)
where $`J_1(x)`$ is the Bessel function of order $`1`$. The membrane density is
$$\rho _{\mathrm{mem}}(𝐱)=\underset{n}{}\rho _L^0f_{\mathrm{mem}}[ynah^n(𝐫)],$$
(2.9)
where $`\rho _L^0`$ is the area density of the lipid and $`f_{\mathrm{mem}}(y)`$ is the lipid membrane form factor, which typically is equal to $`1/w`$ for $`w/2<y<w/2`$ where $`w`$ is the membrane width. X-ray scattering experiments probe $`\rho (𝐪)\rho (𝐪)`$, where $`\rho (𝐪)`$ is the Fourier transform of a linear combination of $`\rho _{\mathrm{DNA}}(𝐱)`$ and $`\rho _{\mathrm{mem}}(𝐱)`$ \[See Sec. III D\].
### B Hamiltonian for Coupled DNA-Lipid Layers
We are now in a position to derive the Hamiltonian describing elastic fluctuations of lamellar DNA-lipid complexes. Our goal is to develop the simplest model consistent with symmetry. We begin with individual membranes and DNA layers. They are characterized by layer bending energies, which can be expressed as
$`^{\mathrm{bend}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{n}{}}a{\displaystyle }d^2r[K_{3d}[(_x^2+_z^2)h^n]^2`$ (2.11)
$`+K_{2d}[(_x^2u_y^n)^2+(_x^2u_z^n)^2]],`$
where $`aK_{3d}=\kappa _{\mathrm{mem}}`$ is the bending rigidity of the lipid membranes and $`adK_{2d}=\kappa _{DNA}`$ is the bending rigidity of an individual DNA strand. We adopt a convention here with a factor of the layer spacing $`a`$ multiplying sums over $`n`$ to facilitate the continuum limit: $`ad^2rd^3x`$. Interactions between DNA strands lead to a preferred separation between strands within a layer and to a compressional elastic energy, characterized by a modulus $`B_{2d}`$, for changing this separation. There is also a preferred separation between a given DNA layer and the two lipid membranes above and below it. Harmonic deviations from this separation are characterized by a modulus $`B_{3d}`$. Changing the distance between membranes on either side of a DNA layer will lead to an expansion or a compression of its DNA smectic lattice, which is described by a coupling with strength $`B_{uh}`$ between inplane strain and membrane separation. Combining all of these effects into a single compression energy, we obtain
$`^{\mathrm{com}}`$ $`=`$ $`{\displaystyle \underset{n}{}}a{\displaystyle }d^2r[{\displaystyle \frac{1}{2}}B_{2d}(u_{zz}^n)^2`$ (2.14)
$`+(B_{uh}/a)u_{zz}^n(h^{n+1}h^n)`$
$`+(B_{3d}/a^2)[(h^{n+1}u_y^n)^2+(u_y^nh^n)^2]],`$
where
$$u_{zz}^n=_zu_z^n[(_xu_z^n)^2+(_zu_z^n)^2]/2$$
(2.15)
is the nonlinear $`2D`$ Eulerian strain and where factors of $`a^1`$ were introduced to facilitate the continuum limit, e.g., $`(h^{n+1}h^n)/a_yh(𝐱)`$. There are additional compressional energies involving interactions proportional to $`(h^{n+1}h^n)^2`$, $`(u_y^{n+1}u_y^n)^2`$, $`(u_{zz}^{n+1}u_{zz}^n)^2`$, and other further neighbor terms. These are smaller than those considered in Eq. (2.14), and we will neglect them. They are, however, needed to achieve $`T_d<T_{KT}`$ and a stable SC phase.
Neighboring DNA strands within a gallery prefer to be parallel. This leads to a Frank-like orientational energy
$$^{\mathrm{orien}}=\frac{1}{2}\underset{n}{}ad^2rK_z(_z\theta ^n)^2,$$
(2.16)
where $`\theta ^n(𝐫)_xu_z^n(𝐫)`$ is the angle that the $`n`$th DNA lattice at $`𝐫`$ makes with the $`x`$-axis. In all but the NL phase, this orientational interaction is subdominant to those in Eq. (2.11), and we will ignore it.
Finally, there are interactions between DNA lattices in neighboring galleries that favor parallel alignment and spatial registry of the lattices. These interactions are described by a sum of layer-coupling Hamiltonians
$$^{\mathrm{int}}=\underset{n}{}(_n^\theta +_n^u).$$
(2.17)
The angular coupling is
$$_n^\theta =V^\theta d^2r\mathrm{cos}[2(\theta ^n\theta ^{n+1})].$$
(2.18)
In all but the isotropic lamellar phase, there is long-range angular order, and we can replace $`_n_n^\theta `$ by
$$^{\mathrm{rot}}=\frac{1}{2}K_y\underset{n}{}ad^2r\left(\frac{_xu_z^{n+1}_xu_z^n}{a}\right)^2,$$
(2.19)
where $`K_y4aV^\theta e^{4(\theta ^n)^2}`$. The interaction $`_n^u`$ in Eq. (2.17), favoring spatial registry, arises from interactions between DNA densities in Eq. (2.5). The phases $`q_0z^n`$ \[Eq. (2.4)\] are chosen so that energy is minimized when the remaining phase, $`\beta ^n(𝐫)=q_0[zu_z^n(𝐫)]`$, of mass-density waves in neighboring DNA lattices are equal at the points of closest approach of the lattices \[Fig. 3\]. The point at $`𝐫^{n+1}`$ in layer $`n+1`$ closest to the point at $`𝐫^n`$ in layer $`n`$ lies along the normal $`𝐍^n(𝐫)(_xu_y^n,1,_zu_y^n)`$. Thus $`𝐫^{n+1}=𝐫^n+\delta 𝐫^{n+1}`$, where $`\delta 𝐫^{n+1}=a𝐍_{}(a_xu_y^n,0,a_zu_y^n)`$ where $`𝐍_{}`$ is the projection of $`𝐍`$ onto the $`xz`$ plane. Energy is minimized when $`\beta ^{n+1}(𝐫^{n+1})=q_0(za_zu_y^nu_z^{n+1})`$ is equal to $`\beta ^n(𝐫^n)`$. The resulting energy is
$$_n^u=V^ud^2r\mathrm{cos}[q_0(u_z^{n+1}u_z^n+a_zu_y^n)],$$
(2.20)
where $`V^u`$ is proportional to the square $`|\psi |^2`$ of the mass-density-wave amplitude. As required, this energy is invariant with respect to spatially uniform translations described by an $`n`$\- and $`𝐫`$-independent displacement of $`u_z^n(𝐫)`$. It is also invariant to harmonic order with respect to uniform rotations in which $`u_z^{n+1}u_z^n=a_zu_y^n`$. The predominant interaction between DNA in neighboring strands is electrostatic, and we expect $`V^u(\lambda _e^2/d)e^{2\pi a/d}`$, where $`\lambda _e`$ is the charge per unit length on a DNA strand.
To summarize, our complete model Hamiltonian coupling membrane height displacements $`h^n`$ and DNA lattice displacements $`u_y^n`$ and $`u_z^n`$ is thus
$$^{\mathrm{tot}}=^{\mathrm{bend}}+^{\mathrm{com}}+^{\mathrm{int}},$$
(2.21)
with entries defined by Eqs. (2.11), (2.14), and (2.17) to (2.20).
#### 1 Effective Hamiltonian for DNA Displacements
Though we will discuss membrane height fluctuations, our primary focus will be on properties of the DNA lattices. It is, therefore, useful to integrate out the membrane height variables $`h^n(𝐫)`$ to obtain an effective Hamiltonian for DNA displacements only. This operation is carried out in Appendix A. The long wavelength result is
$$_{DNA}^{\mathrm{tot}}=^{SC}+^u$$
(2.22)
where $`^u=_n_n^u`$ and $`^{SC}`$ is the sliding phase Hamiltonian,
$`^{SC}={\displaystyle \frac{1}{2}}{\displaystyle \underset{n}{}}a{\displaystyle d^2r[B_{2d}(u_{zz}^n)^2+K_{2d}(_x^2u_z^n)^2]}`$ (2.23)
$`+{\displaystyle \frac{1}{2}}{\displaystyle \underset{n}{}}a{\displaystyle d^2r[(B_{3d}/a^2)(u_y^{n+1}u_y^n)^2+K^{\alpha \beta }(_\alpha _\beta u_y^n)^2]}`$ (2.24)
$`+{\displaystyle \frac{1}{2}}{\displaystyle \underset{n}{}}a{\displaystyle }d^2r[B_{zz}[(u_{zz}^{n+1}u_{zz}^n)/a]^2`$ (2.25)
$`+2(B_{uh}/a)u_{zz}^n(u_y^{n+1}u_y^{n1})]`$ (2.26)
$`{\displaystyle \frac{1}{2}}{\displaystyle \underset{n}{}}a{\displaystyle d^2rK_y[(_xu_z^{n+1}_xu_z^n)/a]^2}`$ (2.27)
where
$$B_{zz}=\frac{B_{uh}^2}{4B_{3d}},$$
(2.28)
the summation convention on $`\alpha =x,z`$ and $`\beta =x,z`$ is understood, and the bending-rigidity tensor is
$$K^{\alpha \beta }=\left(\begin{array}{cc}K^{xx}& K^{xz}\\ K^{zx}& K^{zz}\end{array}\right)=\left(\begin{array}{cc}K_{3d}+K_{2d}& K_{3d}\\ K_{3d}& K_{3d}\end{array}\right).$$
(2.29)
Note that integrating out $`h^n`$ creates interactions between nearest-neighbor and next-nearest-neighbor $`u_z^n`$ and $`u_y^n`$ displacements not present in the original model of Eq. (2.21). These interactions, however, have particular ratios determined by the parameters in $`^{\mathrm{tot}}`$. Of particular importance to us in what follows is the fixed relations between the coefficients of $`(u_y^{n+1}u_y^n)^2`$, $`(u_{zz}^{n+1}u_{zz}^n)^2`$, and $`u_{zz}^n(u_y^{n+1}u_y^{n1})`$. In a more general model with further-neighbor interactions, the simple relation of Eq. (2.28) among $`B_{zz}`$, $`B_{uh}`$, and $`B_{3d}`$ would not hold.
Using $`_{DNA}^{\mathrm{tot}}`$, we can construct the long-wavelength Hamiltonian for each of the phases listed in the introduction as detailed in the following.
#### 2 The Columnar Phase
The columnar crystal is characterized by strong coupling of displacements in different layers. Its elastic Hamiltonian can be obtained by expanding the cosine in $`_n^u`$ \[Eq. (2.20)\] about one of its minima. Performing this operation, taking the continuum limit, and retaining only the lowest order terms in a gradient expansion, we obtain the familiar elastic Hamiltonian for a rectangular columnar lattice:
$`^C`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle d^3x[B_{2d}u_{zz}^2+B_{3d}u_{yy}^2+B_{uh}u_{zz}u_{yy}+B_{yz}u_{yz}^2]}`$ (2.31)
$`+{\displaystyle d^3x[K_{2d}(_x^2u_z)^2+(K_{2d}+K_{3d})(_x^2u_y)^2]}`$
where $`𝐱=(𝐫,na)`$, $`u_{zz}(𝐱)=u_{zz}^{n=y/a}(𝐫)`$, and $`B_{yz}V^uq_0^2a^2`$. Note that the shear elastic modulus $`B_{yz}`$ arises from $`_n^u`$ and is zero when $`V^u`$ is zero. The cross compression $`B_{uh}`$-term arises from the $`u_{zz}^n(u_y^{n+1}u_y^n)`$ term in $`_{DNA}^{\mathrm{tot}}`$, which was generated by the interaction of $`u_{zz}^n`$ with the membrane height field $`h^n`$
#### 3 Sliding Columnar Phase
The sliding columnar phase is characterized by a vanishing shear modulus $`B_{yz}`$ for relative displacement of DNA lattices. Its dominant fluctuations are, therefore, described by the Hamiltonian $`^{SC}`$ of Eq. (2.23) obtained by setting $`V^u=0`$ in $`_{DNA}^{\mathrm{tot}}`$. The $`V^u`$ coupling is irrelevant when the sliding phase is stable as detailed in the following \[See Sec. IV\].
#### 4 Nematic Lamellar Phase
Positional couplings between layers and the smectic compression modulus $`B_{2d}`$ both vanish in the nematic lamellar phase, and
$`^{\mathrm{NL}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle d^3x[K_x(_x\theta )^2+K_y(_y\theta )^2+K_z(_z\theta )^2]}`$ (2.33)
$`+{\displaystyle \frac{1}{2}}{\displaystyle }d^3x[B_{3d}u_{yy}^2++K^{\alpha \beta }(_\alpha _\alpha u_y)^2].`$
where $`K_x=K_{2d}`$ and $`K_z`$ is a Frank elastic constant introduced in Eq. (2.16).
## III Correlations in the Sliding Columnar Phase
In the preceding section, we derived a general Hamiltonian capable of describing the phases of lamellar DNA-lipid complexes. In this section, we will investigate correlations in the sliding columnar phase. The most unusual correlations in the sliding columnar phase are those involving displacements of DNA strands within the layers. We will thus begin by considering the effective Hamiltonian for $`u_z^n`$ obtained by integrating out $`u_y^n`$ from $`^{SC}`$ \[Eq. (2.23)\]. From this we will calculate all correlation functions of $`u_z^n`$. We will then consider membrane height correlations described by an effective Hamiltonian for $`h^n`$ in which $`u_y^n`$ and $`u_z^n`$ have been integrated out of $`^{\mathrm{tot}}`$ \[Eq. (2.21)\]. In the long-wavelength limit fluctuations in $`u_y^n`$ are the same as those in $`h^n`$. Finally, we will discuss the relevancy of the shear coupling $`V^u`$, Eq. (2.20). This involves a consideration of both $`u_z^n`$ and $`_zu_y^n`$.
### A Correlations in $`u_z^n`$
#### 1 Effective Hamiltonian for $`u_z^n`$
Since $`u_y^n`$ and $`u_z^n`$ are harmonically coupled in $`_{DNA}^{\mathrm{tot}}`$, we can integrate over $`u_y^n`$ exactly to obtain an effective Hamiltonian for $`u_z^n`$ fluctuations alone. This calculation is carried out in Appendix A. In the long-wavelength limit, the resulting Hamiltonian is
$`_z^{SC}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{n}{}}a{\displaystyle }d^2r[B(u_{zz}^n)^2+K(_x^2u_z^n)^2`$ (3)
$`+{\displaystyle \frac{K_y}{a^2}}[_x(u_z^nu_z^{n+1})]^2]`$
where $`K=K_{2d}`$ and
$$B=B_{2d}(B_{uh}^2/B_{3d}).$$
(3.4)
Note that there is no term proportional to $`(u_{zz}^{n+1}u_{zz}^n)^2`$ in Hamiltonian $`_z^{SC}`$ even though there is one in $`^{SC}`$. This is a result of the special relation \[Eq. (2.28)\] among energy coefficients in $`^{SC}`$.
The Hamiltonian of Eq. (III A 1) is the sum of two parts: a sum of elastic energies for independent $`2D`$ smectics \[Eq. (3)\] and a term coupling angles in neighboring layers \[Eq. (3)\]. The $`2D`$ smectic energy is a function of the full nonlinear strain $`u_{zz}^n`$. When $`K_y=0`$ and there is no coupling between layers, nonlinearities in $`u_{zz}^n`$ lead to important renormalizations of the long-wavelength elastic constant of a $`2D`$ smectic. When the interlayer coupling term is present, nonlinearities in $`u_{zz}^n`$ also lead to long wavelength renormalization of elastic constants in the SC phase. These renormalizations are only logarithmic, however, and we will ignore them in this article.
We, therefore, consider the harmonic limit of $`_z^{SC}`$, which is conveniently expressed in Fourier space. Introducing
$$\frac{d^3q}{(2\pi )^2}=_{\pi /a}^{\pi /a}\frac{dq_y}{2\pi }_{\mathrm{\Lambda }_x}^{\mathrm{\Lambda }_x}\frac{dq_x}{2\pi }_{\mathrm{\Lambda }_z}^{\mathrm{\Lambda }_z}\frac{dq_z}{2\pi }$$
(3.5)
where $`\mathrm{\Lambda }_x`$ and $`\mathrm{\Lambda }_z2\pi /d`$ are wavenumber cutoffs, and
$$u_z^n(𝐫)=\frac{d^3q}{(2\pi )^3}e^{i(q_yna+𝐪_{}𝐫)}u_z(𝐪),$$
(3.6)
where $`𝐪_{}=(q_x,0,q_z)`$, we obtain
$$_z^{SC}=\frac{1}{2}\frac{d^3q}{(2\pi )^3}[Bq_z^2+Kq_x^4+K_y(q_y)q_x^2q_y^2]|u_z(𝐪)|^2.$$
(3.7)
Here
$$K_y(q_y)=K_yp(q_ya)$$
(3.8)
with
$$p(q_ya)=2\frac{(1\mathrm{cos}q_ya)}{(q_ya)^2},$$
(3.9)
which tends to unity as $`q_ya0`$ so that $`K(q_y=0)=K_y`$.
Two important length scales can be obtained from Eq. (III A 1) by comparing the orientational interaction energy with the $`2D`$ smectic compression and bending energies. The lengthscales
$`x^{}={\displaystyle \frac{a}{\mu _y}}\mathrm{and}z^{}={\displaystyle \frac{a^2}{\mu _y^2\lambda }},`$ (3.10)
with $`\mu _y=\sqrt{K_y/K}`$ and $`\lambda =\sqrt{K/B}`$, separate two-dimensional from three-dimensional behavior. At lengthscales within a gallery less than $`x^{}`$ and $`z^{}`$, the $`2D`$ compression and bending energies are large compared to the orientational interaction, and the DNA lattices behave like independent $`2D`$ smectics. On the other hand, at lengthscales greater than $`x^{}`$ and $`z^{}`$, the orientational interaction is significant, and $`3D`$ sliding behavior occurs.
Before proceeding to discuss the sliding phase correlations, we note that the sliding phase elastic free energy in Eq. (III A 1) exhibits a striking local (gauge) translational invariance of the form
$$u_z(x,z,na)u_z(x,z,na)+f(n);$$
(3.11)
here the displacement $`f(n)`$ is an arbitrary function of $`n`$ assuming a different value in each gallery. In other words, the DNA lattices in different galleries can continuously slide relative to each other by arbitrary distances with no elastic energy costs . In a standard columnar phase, this continuous symmetry is spontaneously broken down by the shear coupling in Eq. (2.20). On the other side, in the sliding phase, this shear coupling is irrelevant at long scales, and elastic properties reflect the gauge symmetry in Eq. (III A 1). This symmetry is entirely responsible for the unconventional fluctuation behavior of the sliding phase we discuss in the following. In particular, we find that in the sliding phase, the fluctuations of relative displacements $`u_z^{n+1}(𝐫)u_z^n(𝐫)`$ of DNA lattices in neighboring galleries diverge in thermodynamic limit. As a reflection of the gauge symmetry Eq. (III A 1), a DNA lattice in a gallery is essentially free to flow relative to DNA lattices in neighboring galleries. The sliding phase thus exhibits zero macroscopic shear modulus
#### 2 Definitions of Correlation Functions
Correlations in $`u_z^n`$ follow directly form Eq. (3.7):
$`G_{zz}(𝐫,na)`$ $`=`$ $`u_z^n(𝐫)u_z^0(0)(u_z^n(𝐫))^2g(𝐫,na)`$ (3.12)
$`=`$ $`{\displaystyle \frac{d^3q}{(2\pi )^3}G_{zz}(𝐪)e^{i(𝐪_{}𝐫+q_yna)}},`$ (3.13)
where
$$G_{zz}(𝐪)=\frac{T}{Bq_z^2+Kq_x^4+K_yq_x^2q_y^2p(q_ya)},$$
(3.14)
$$(u_z^n)^2=\frac{d^3q}{(2\pi )^3}G_{zz}(𝐪),$$
(3.15)
and
$`g(𝐫,na)`$ $`=`$ $`{\displaystyle \frac{1}{2}}[u_z^n(𝐫)u_z^0(\mathrm{𝟎})]^2`$ (3.16)
$`=`$ $`{\displaystyle \frac{d^3q}{(2\pi )^3}G_{zz}(𝐪)\left[1e^{i(𝐪_{}𝐫+q_yna)}\right]}.`$ (3.17)
It is useful to decompose $`g(𝐫,na)`$ into three parts:
$$g(𝐫,na)=g^{(1)}(na)+g^{(2)}(𝐫)g^{(3)}(𝐫,na),$$
(3.18)
where
$`g^{(1)}(na)`$ $`=`$ $`{\displaystyle \frac{1}{2}}[u_z^n(\mathrm{𝟎})u_z^0(\mathrm{𝟎})]^2`$ (3.19)
$`=`$ $`{\displaystyle \frac{d^3q}{(2\pi )^3}G_{zz}(𝐪)(1\mathrm{cos}q_yna)}`$ (3.20)
$`g^{(2)}(𝐫)`$ $`=`$ $`{\displaystyle \frac{1}{2}}[u_z^0(𝐫)u_z^0(\mathrm{𝟎})]^2`$ (3.21)
$`=`$ $`{\displaystyle \frac{d^3q}{(2\pi )^3}G_{zz}(𝐪)(1e^{i𝐪_{}𝐫})}`$ (3.22)
$`g^{(3)}(𝐫,na)`$ $`=`$ $`[u_z^n(\mathrm{𝟎})u_z^0(\mathrm{𝟎})][u_z^n(𝐫)u_z^n(\mathrm{𝟎})]`$ (3.23)
$`=`$ $`{\displaystyle \frac{d^3q}{(2\pi )^3}G_{zz}(𝐪)}`$ (3.25)
$`\times (1\mathrm{cos}q_yna)(1e^{i𝐪_{}𝐫}).`$
It is clear that $`g^{(1)}(na)=g(0,na)`$ and $`g^{(2)}(𝐫)=g(𝐫,0)`$. Each of these functions has a characteristic singular behavior as a function of system size and separation, which we will summarize below. Finally, the function
$$g^{(2)}(𝐫,na)=g^{(2)}(𝐫)g^{(3)}(𝐫,na)$$
(3.26)
will appear in our derivation in App. F of interplane density correlations discussed in Sec. III D.
#### 3 Asymptotic Forms for $`K_y0`$
The local fluctuation $`(u_z^n)^2`$ in the SC phase diverges as the square of the log of the system size with a functional form that depends on the order in which the sample dimensions $`L_x`$, $`L_y`$, and $`L_z`$ along the $`x`$, $`y`$ and $`z`$ directions approach infinity:
$$(u_z^n)^2=l_u^2\{\begin{array}{cc}\mathrm{ln}^2\left[8L_x/x^{}\right]\hfill & L_zL_yL_x\hfill \\ \frac{1}{2}\mathrm{ln}^2\left[\alpha _zL_z/z^{}\right]\hfill & L_xL_yL_z\hfill \end{array},$$
(3.27)
where the length $`l_u`$ is defined via
$$l_u^2=\frac{T}{4\pi ^2\sqrt{BK_y}},$$
(3.28)
$`\alpha _z`$ is a number, $`L_zz^{}`$, $`L_xx^{}`$, and terms that do not diverge with system size have been dropped. The calculation of the displacement fluctuations in the limit $`L_z\mathrm{}`$ and $`L_xL_y`$ is detailed in Appendix B. Thus, SC “in-plane” fluctuations are less divergent than $`2D`$ smectic fluctuations that scale as a power-law with system size but more divergent than $`3D`$ Landau-Peierls smectic lamellar fluctuations that grow logarithmically with system size. The mean-square angular fluctuation, $`(\theta ^n)^2=(_xu_z^n)^2`$, is finite, implying that the SC phase has three-dimensional long-range orientational order.
The $`\mathrm{ln}^2L_{x,z}`$ divergence of $`(u_z^n)^2`$ is converted to a $`\mathrm{ln}^2r`$ divergence in the function $`g^{(2)}(𝐫)`$ \[Eq. (3.21)\]. In Appendix D, we outline the calculation of $`g^{(2)}(𝐫)`$ for large $`r`$. The results are
$$g^{(2)}(𝐫)=l_u^2\{\begin{array}{cc}\mathrm{ln}^2\left[8e^\gamma x/x^{}\right]+C_x\hfill & \mathrm{if}z=0\hfill \\ \frac{1}{2}\mathrm{ln}^2\left[32e^\gamma z/z^{}\right]+C_z\hfill & \mathrm{if}x=0,\hfill \end{array}$$
(3.29)
where $`\gamma 0.577`$ is Euler’s constant and $`C_x`$ and $`C_z`$ are constants that depend on $`\mathrm{\Lambda }_x`$ and $`\mathrm{\Lambda }_z`$ but have well-defined $`\mathrm{\Lambda }_{x,z}\mathrm{}`$ limits. Note that the only $`\mathrm{\Lambda }_{x,z}`$ dependence of the correlation function in these two large distance limits is found in the constants $`C_x`$ and $`C_z`$. $`C_x`$ and $`C_z`$ are evaluated in Appendix D in the continuum limit ($`\mathrm{\Lambda }_{x,z}\mathrm{}`$) in which we find $`C_x=0`$ and $`C_z=\pi ^2/8`$.
The function $`g^{(1)}(na)`$ \[Eq. (3.19)\] is by definition zero when $`n=0`$. For all $`n0`$, it diverges logarithmically with system size:
$$g^{(1)}(na)=2l_u^2S_n(0)\mathrm{ln}\left[A_n(\mathrm{\Lambda }_x,\mathrm{\Lambda }_z)\frac{L_x}{x^{}}\right],$$
(3.30)
where
$$S_n(t)=_0^\pi 𝑑u\frac{1\mathrm{cos}(nu)}{\sqrt{t^2+u^2p(u)}},$$
(3.31)
$$S_n(0)=\underset{k=1}{\overset{n}{}}\frac{2}{2k1},$$
(3.32)
and $`A_n(\mathrm{\Lambda }_x,\mathrm{\Lambda }_z)`$ calculated in App. C depends on the layer separation $`n`$ and the ultraviolet cutoffs but has a well-defined $`\mathrm{\Lambda }_{x,z}\mathrm{}`$ limit that is calculated in Appendix C.
The function $`g^{(3)}(𝐫,na)`$ \[Eq. (3.23)\] is also zero at $`n=0`$. It has the same divergences as a function of separation $`r`$ that $`g^{(1)}(na)`$ has with system size:
$$g^{(3)}(𝐫,n)=2l_u^2S_n(0)\{\begin{array}{cc}\mathrm{ln}(D_nx/x^{})\hfill & \text{if }z=0\hfill \\ \mathrm{ln}(E_naz/z^{})\hfill & \text{if }x=0\text{.}\hfill \end{array},$$
(3.33)
where $`D_n`$ and $`E_n`$ are numerical constants that have well defined values in the continuum limit, $`\mathrm{\Lambda }_x,\mathrm{\Lambda }_z\mathrm{}`$, as discussed in App. E.
It is useful to summarize the results of the calculations just presented. When $`n=0`$, i.e., for points in the same layer, $`g(𝐫,0)=g^{(2)}(𝐫)`$ grows with $`r`$ as $`\mathrm{ln}^2r`$. When $`n0`$, i.e., for sites in different layers, $`g(𝐫,na)`$ diverges logarithmically with system size for all $`𝐫`$ because $`g^{(1)}(na)`$ diverges in this way. If $`g^{(1)}(na)`$ is subtracted from $`g(𝐫,n)`$, then the remaining function, $`g^{(2)}(𝐫,na)`$, grows with $`r`$ as $`\mathrm{ln}^2(r/r_n)`$ with $`r_n`$ depending on the coefficient of $`\mathrm{ln}r`$ in $`g^{(3)}(𝐫,na)`$ when $`n0`$.
#### 4 The limit $`K_y=0`$: $`2D`$ Smectic Correlations
When $`K_y=0`$, there are no interactions between planes (when $`V^u=0`$), and the system reduces to a stack of independent $`2D`$ smectic planes. Since experiments are consistent with nearly independent $`2D`$ smectic layers, we will review here well established results for such systems.
Decoupled smectics are described by the Hamiltonian of Eq. (3) obtained by setting $`K_y=0`$ in $`^{SC}`$. This Hamiltonian is a function of the full nonlinear $`2D`$ strain $`u_{zz}^n`$. Nonlinearities lead to important deviations from harmonic behavior. Since there is no coupling between layers, $`G_{zz}(𝐪)`$ is independent of $`q_y`$, and we can define a $`2D`$ correlation function $`G_2(𝐪_{})=G_{zz}(𝐪)/a`$.
At lengthscales less than the nonlinear lengths
$`l_x={\displaystyle \frac{8\pi K_2^{3/2}}{T\sqrt{B_2}}}\mathrm{and}l_z={\displaystyle \frac{l_x^2}{\lambda }},`$ (3.34)
where $`K_2=Ka`$ and $`B_2=Ba`$ are $`2D`$ bend and compression moduli and $`\lambda =\sqrt{K_2/B_2}`$, nonlinearities are unimportant, and $`2D`$ smectic fluctuations are described by the linearized elastic Hamiltonian $`_z^{SC}`$ of Eq. (3.7) with $`K_y=0`$. At lengthscales longer than $`l_x`$ and $`l_z`$, the nonlinear terms in the rotationally invariant strain in Eq. (2.15) lead to renormalized bending and compression moduli $`K_2(𝐪_{})`$ and $`B_2(𝐪_{})`$ that, respectively, diverge and vanish at small wavenumber $`𝐪_{}`$. Note that the nonlinear lengths $`l_x`$ and $`l_z`$ decrease with increasing temperature, and thus nonlinearities become important at high temperatures. In both the harmonic and nonlinear regimes, the Fourier transformed displacement correlation function $`|u_z^n(𝐪)|^2`$ in each gallery were expressed as
$$G_2(𝐪_{})=\frac{T}{B_2(𝐪_{})q_z^2+K_2(𝐪_{})q_x^4}=\frac{T}{B_2}l_z^2\stackrel{~}{q}_x^\gamma Q(\stackrel{~}{q}_z/\stackrel{~}{q}_x^\nu ),$$
(3.35)
where $`\stackrel{~}{q}_{x,z}=q_{x,z}l_{x,z}`$ and
$$\stackrel{~}{q}_x^\gamma Q(\stackrel{~}{q}_z/\stackrel{~}{q}_x^\nu )\{\begin{array}{cc}\stackrel{~}{q}_x^\gamma ,\hfill & \stackrel{~}{q}_z=0\hfill \\ \stackrel{~}{q}_z^{\gamma /\nu }\hfill & \stackrel{~}{q}_x=0\text{.}\hfill \end{array}$$
(3.36)
The scaling form of the correlation function implies that $`K_2(q_x,q_z=0)q_x^{4+\gamma }`$ and $`B_2(q_x=0,q_z)q_z^{2+\gamma /\nu }`$. In the harmonic regime where $`q_{x,z}l_{x,z}>1`$, $`K_2(𝐪_{})=K_2`$ and $`B_2(𝐪_{})=B_2`$ are constants, and the scaling exponents are $`\gamma =4`$ and $`\nu =2`$. In the anharmonic regime $`q_{x,z}l_{x,z}<1`$, the scaling exponents $`\gamma `$ and $`\nu `$ were calculated exactly by mapping the $`2D`$ smectic model with thermal fluctuations onto the KPZ model in $`1+1`$ dimensions . The exponents in the anharmonic regime are $`\gamma =7/2`$ and $`\nu =3/2`$.
The mean-square displacement fluctuations diverge in both regimes with lengths $`L_x`$ and $`L_z`$ of the sample in the $`xz`$ plane:
$$u_z^2(𝐫)=\frac{d^2q_{}}{(2\pi )^2}G_2(𝐪_{})=\lambda ^2\stackrel{~}{L}_x^{2\alpha }f_u^{(1)}(\stackrel{~}{L}_z/\stackrel{~}{L}_x^\nu ),$$
(3.37)
where $`u_z(𝐫)u_z^n(𝐫)`$, $`2\alpha =\gamma 1\nu =1`$ in both regimes, $`\stackrel{~}{L}_{x,z}=L_{x,z}/l_{x,z}`$, $`f_u^{(1)}(0)=\mathrm{const}.`$, and $`f_u^{(1)}(w)w^{2\alpha /\nu }`$ as $`w\mathrm{}`$. This implies that the Debye-Waller factor $`\mathrm{exp}[iq_0u_z]^2=\mathrm{exp}[q_0^2u_z^2]=0`$ in the limit of infinite system size, and there is no long-range positional order at any finite temperature in a $`2D`$ smectic, even when there are no dislocations.
Since the mean-square displacement fluctuations diverge as a power-law with system size, the displacement correlation function
$`g^{2D}(𝐫)`$ $`=`$ $`{\displaystyle \frac{1}{2}}[u_z^n(𝐫)u_z^n(0)]^2`$ (3.38)
$`=`$ $`\lambda ^2|\stackrel{~}{x}|^{2\alpha }f_u^{(2)}(|\stackrel{~}{z}|/|\stackrel{~}{x}|^\nu )`$ (3.39)
diverges algebraically with in-plane separation $`r`$. In Eq. (3.38), $`\stackrel{~}{x}=x/l_x`$, $`\stackrel{~}{z}=z/l_z`$, and the scaling behavior of $`f_u^{(2)}(w)`$ is similar to that of $`f_u^{(1)}(w)`$. In the harmonic limit,
$$q_0^2g^{2D}(𝐫)=\{\begin{array}{cc}(z/\xi _z)^{1/2},\hfill & \text{if }x=0\text{;}\hfill \\ (x/\xi _x),\hfill & \text{is }z=0\text{ ,}\hfill \end{array}$$
(3.40)
where
$$\xi _z=\frac{\lambda B_2^2d^4}{4\pi ^3T^2}$$
(3.41)
is correlation length along $`z`$ and $`\xi _x=\lambda (d/\pi )^2(B_2/T)=2\sqrt{\lambda \xi _z/\pi }`$ is a correlation length along $`x`$. These correlation lengths are measured in x-ray scattering experiments at sufficiently short length scales when layers are decoupled.
Fluctuations in the angle $`\theta =_xu_z`$ are nondivergent because of an additional factor of $`q_x^2`$ in the numerator of (3.37). Finite angular fluctuations imply that $`\mathrm{cos}\theta =\mathrm{exp}[\theta ^2/2]`$ is nonzero, and thus there is long-range orientational order in $`2D`$ smectics when there are no dislocations.
### B Correlations in $`h^n(𝐫)`$ and $`u_y^n(𝐫)`$.
In the absence of order in the DNA strands, lamellar DNA lipid complexes are simply multicomponent isotropic lamellar systems with height fluctuations identical to those of any isotropic lamellar smectic. The presence of orientational order in the DNA lattices introduces anisotropy into the effective bending moduli of the lipid membranes. By integrating out $`u_y^n(𝐫)`$ and $`u_z^n(𝐫)`$ from the the complete Hamiltonian $`^{\mathrm{tot}}`$ of Eq. (2.21), we obtain in App. A an effective Hamiltonian for height fluctuations. In the long-wavelength limit, this effective Hamiltonian is
$`_h^{SC}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{n}{}}a{\displaystyle }d^2r[(B_{\mathrm{lam}}/a^2)(h^{n+1}h^n)^2`$ (3.43)
$`+K^{\alpha \beta }(_\alpha _\beta h^n)^2],`$
where
$$B_{\mathrm{lam}}=B_{3d}(B_{uh}^2/B_{2d})$$
(3.44)
and $`K^{\alpha \beta }`$ is the bending-rigidity tensor defined in Eq. (2.29). The model in Eq. (3.43) is an anisotropic version of the discrete smectic Hamiltonian often used to describe lamellar systems. An effective Hamiltonian for DNA-lattice height displacements $`u_y^n(𝐫)`$ can be obtained by integrating over $`h^n(𝐫)`$ and $`u_z^n(𝐫)`$ in $`^{\mathrm{tot}}`$ of Eq. (2.21). The resulting Hamiltonian is identical in the long-wavelength limit to $`_h`$ in Eq. (3.43). Thus correlations in $`u_y(𝐫)`$ are identical in this model in the long-wavelength limit to those in $`h^n(𝐫)`$.
Fluctuations in $`h^n`$ are determined by the correlation function
$$G_{hh}(𝐪)=\frac{T}{B_{\mathrm{lam}}q_y^2p(q_ya)+K^{\alpha \beta }q_\alpha ^2q_\beta ^2},$$
(3.45)
where $`K^{\alpha \beta }q_\alpha ^2q_\beta ^2=K_{3d}(q_x^2+q_z^2)^2+K_{2d}q_x^4`$. This equation implies that $`[h^n(𝐫)]^2`$ diverges logarithmically with system size as it does in ordinary isotropic lamellar systems. The coefficient of the system-size logarithm depends on the anisotropy in $`K^{\alpha \beta }`$. Similarly, correlations in $`h^n(𝐫)`$ defined via
$`g_h(𝐫,na)`$ $`=`$ $`{\displaystyle \frac{1}{2}}[h^n(𝐫)h^0(\mathrm{𝟎})]^2`$ (3.46)
$`=`$ $`{\displaystyle \frac{1}{2}}[u_y^n(𝐫)u_y^0(\mathrm{𝟎})]^2`$ (3.47)
$`=`$ $`{\displaystyle \frac{d^3q}{(2\pi )^3}G_{hh}(𝐪)[1e^{i(𝐪_{}𝐫+q_yna)}]}`$ (3.48)
diverge logarithmically with $`na`$ and $`𝐫`$:
$$k_0^2g_h(𝐫,na)=\{\begin{array}{cc}\eta _h\mathrm{ln}(na\lambda _x\mathrm{\Lambda }^2)\hfill & \text{if }𝐫=0\hfill \\ 2\eta _h\mathrm{ln}(\stackrel{~}{\mathrm{\Lambda }}(\theta )r)\hfill & \text{if }n=0\text{ ,}\hfill \end{array}$$
(3.49)
where $`\lambda _x=\sqrt{B/K^{xx}}`$, $`k_0=2\pi /a`$, and
$$\eta _h=\frac{k_0^2T}{8\pi ^2\sqrt{B_{\mathrm{lam}}K^{xx}}}I(K^{xz}/K^{xx},K^{zz}/K^{xx}),$$
(3.50)
where
$$I(\mu ,\nu )=\frac{1}{2\pi }_\pi ^\pi \frac{d\theta }{\sqrt{\mathrm{cos}^4\theta +2\mu \mathrm{sin}^2\theta \mathrm{cos}^2\theta +\nu \mathrm{sin}^4\theta }}.$$
(3.51)
In Eq. (3.49), $`\mathrm{\Lambda }`$ and $`\stackrel{~}{\mathrm{\Lambda }}(\theta )`$ are cutoffs that depend on $`\mathrm{\Lambda }_x`$, $`\mathrm{\Lambda }_z`$ and ratios of bending moduli. In addition, $`\stackrel{~}{\mathrm{\Lambda }}(\theta )`$ depends on the angle $`\theta `$ that $`𝐫`$ makes with the $`x`$ axis.
### C Density Correlations
The DNA density correlation function arising from displacements parallel to lipid membranes is
$$S(𝐫,na)=\mathrm{exp}(iq_0[u_z^n(𝐫)u_z^0(\mathrm{𝟎})].$$
(3.52)
This function is easily evaluated in the SC phase when $`V^u=0`$:
$$S(𝐫,na)=e^{q_0^2g(𝐫,na)},$$
(3.53)
where $`g(𝐫,na)`$ is the displacement correlation function defined in Eq. (3.16). Since $`g(𝐫,na)`$ diverges with the system size for all $`n0`$, $`S(𝐫,na)`$ vanishes in the thermodymic limit for all $`n0`$. Thus, DNA densities in different layers are completely uncorrelated in the SC phase when the coupling $`V^u`$ in Eq. (2.20) is set to zero.
#### 1 In-plane Correlations
When $`n=0`$, $`g(𝐫,na)`$ does not diverge in the thermodynamic limit, and we have
$$S_2(𝐫)S(𝐫,0)=e^{iq_0[u_z^n(𝐫)u_z^n(0)]}=e^{q_0^2g^{(2)}(𝐫)},$$
(3.54)
with $`g^{(2)}(𝐫)`$ given by Eq. (3.21). Thus, from Eqs. (3.29) and (3.38),
$$S_2(x,0)=\{\begin{array}{cc}S_xe^{q_0^2l_u^2\mathrm{ln}^2[8e^\gamma x/x^{}]},\hfill & xx^{}\text{;}\hfill \\ e^{q_0^2g^{2D}(x,0)},\hfill & xx^{}\text{,}\hfill \end{array}$$
(3.55)
where $`S_x=e^{q_0^2l_u^2C_x}`$ is a constant, and $`g^{2D}(𝐫)`$ is the displacement correlation function of a $`2D`$ smectic. In the other direction,
$$S_2(0,z)=\{\begin{array}{cc}S_ze^{(q_0^2l_u^2/2)\mathrm{ln}^2[32e^\gamma z/z^{}]},\hfill & zz^{}\text{;}\hfill \\ e^{q_0^2g^{2D}(0,z)},\hfill & zz^{}\text{,}\hfill \end{array}$$
(3.56)
where $`S_z=e^{q_0^2l_u^2C_z}`$ is a constant.
#### 2 Inter-plane Correlations
As we have just seen $`S(𝐫,na)`$ is zero when $`n0`$ if $`V^u=0`$. When $`V^u`$ is not zero, $`S(𝐫,na)`$ is not zero, and it can be calculated perturbatively in a expansion in $`V^u/T`$. Using the decomposition of Eq. (2.22) of $`_{DNA}^{\mathrm{tot}}`$ in to $`^{SC}`$ and $`^u`$, we can express $`S(𝐫,na)`$ as
$$S(𝐫,na)=\frac{e^{^u/T}e^{iq_0[u_z^n(𝐫)u_z^0(0)]}}{e^{^u/T}},$$
(3.57)
where $``$ signifies an average with respect to $`^{SC}`$. This expression is easily expanded in a power series in $`V^u/T`$. The evaluation of even the lowest order term in this expansion cannot be carried out analytically. In App. F, we evaluate the lowest order term in an approximation in which $`K_y=0`$, $`B_{uh}=0`$, and correlations in $`_zu_y^n(𝐫)`$ in different layers are ignored. The result for the Fourier transform of $`S(𝐫,na)`$ is
$$S(𝐪_{},na)=\left(\frac{\stackrel{~}{V}^u}{2T}\right)^n[S_2(𝐪_{})]^{n+1},$$
(3.58)
where $`\stackrel{~}{V}^u=V^ue^{W_y}`$ with
$$W_y=q_0^2a^2(_zu_y^n)^2/2$$
(3.59)
and
$$S_2(𝐪_{})=\frac{d^2q_{}}{(2\pi )^2}e^{i𝐪_{}𝐫}S_2(𝐫)$$
(3.60)
is the Fourier transform of the two-dimensional intra-plane density correlations function of Eq. (3.54). Since the expression, Eq. (3.58), for $`S(𝐪_{},na)`$ was derived under the assumption that $`K_y=0`$, $`S_2(𝐫)`$ is strictly speaking the density correlation of an isolated $`2D`$ smectic. However, a more sophisticated variational approximation, to be presented in a separate publication, yields a result very similar to Eq. (3.58) for $`S(𝐪_{},na)`$ with $`S_2(𝐪_{})`$ replaced by the true intra-plane correlation function with $`K_y0`$ obtained by Fourier transforming $`S(𝐫,0)`$ of Eq. (3.53).
Transforming $`S(𝐪_{},na)`$ back to real space, we obtain
$$S(𝐫,na)=\left(\frac{\stackrel{~}{V}^u}{2T}\right)^n\frac{d^2q_{}}{(2\pi )^2}e^{i𝐪_{}𝐫}e^{(n+1)\mathrm{ln}S_2(𝐪_{})}.$$
(3.61)
The large $`n`$ limit of this expression can be calculated using steepest descents. For fixed $`𝐫`$ and $`n\mathrm{}`$, the result is
$$S(𝐫,na)=\frac{2\pi T}{\stackrel{~}{V}^un\xi _x\xi _z}e^{|n|a/\xi _y}e^{[(x^2/\xi _x^2)+(z^2/\xi _z^2)]},$$
(3.62)
where
$$\xi _\alpha ^2=\frac{1}{S_2(𝐫=0)}d^2rr_\alpha ^2S_2(𝐫)$$
(3.63)
for $`\alpha =x,z`$, and $`r_\alpha =x,z`$ and
$$\xi _y=a\mathrm{ln}\left(\frac{2T}{\stackrel{~}{V}^uS_2(𝐪_{}=0)}\right)^1.$$
(3.64)
Since $`S_2(𝐫)`$ dies off at least as fast as $`e^r`$ when $`K_y=0`$ and as $`\mathrm{exp}(\mathrm{const}.\mathrm{ln}^2r)`$ when $`K_y0`$, the integrals in Eq. (3.63) converge, and $`\xi _x`$ and $`\xi _z`$ are well defined.
### D X-ray Scattering
X-ray scattering experiments probe $`I(𝐪)=\rho (𝐪)\rho (𝐪)/V`$, where $`V`$ is the volume of the system, and $`\rho (𝐪)`$ is the weighted sum of the DNA and membrane densities $`\rho _{DNA}(𝐱)`$ and $`\rho _{\mathrm{mem}}(𝐱)`$. The total scattering intensity can thus be broken up into a DNA contribution $`I_{DNA}(𝐪)`$, a membrane contribution $`I_{\mathrm{mem}}(𝐪)`$, and a DNA-membrane cross term. Scattering from the ordinary columnar phase will exhibit true Bragg peaks at $`𝐆_{lm}=(0,(l+\frac{1}{2}\sigma m)k_0,mq_0)`$. Here $`\sigma =0`$ for a simple rectangular columnar lattice, and $`\sigma =1`$ for a centered rectangular columnar lattice. Since there is no long-range positional order in the SC phase, there will be no Bragg peaks in its x-ray scattering profile. Rather there will be membrane-dominated lamellar peaks at $`𝐆_{l0}=(0,lk_0,0)`$ and DNA-dominated sub-power-law peaks at $`𝐆_{lm}`$ for $`m0`$.
The contribution of lipid membranes to the x-ray scattering intensity is
$`I_{\mathrm{mem}}(𝐪)`$ $`=`$ $`{\displaystyle \frac{1}{a}}|f_{\mathrm{mem}}(q_y)|^2`$ (3.66)
$`\times {\displaystyle \underset{n}{}}e^{iq_yna}{\displaystyle }d^2re^{i𝐪_{}𝐫}S_h(𝐫,na,q_y),`$
where
$`S_h(r,na,q_y)`$ $`=`$ $`e^{iq_y[h^n(𝐫)h^0(0)]}`$ (3.67)
$`=`$ $`e^{q_y^2g_h(𝐫,na)}.`$ (3.68)
and $`g_h(𝐫,na)`$ is evaluated in Eq. (3.46). Thus, for example,
$$S_h(0,na,k_0)(na)^{\eta _h}$$
(3.69)
Except for some anisotropies, this yields power-law Bragg peaks that are essentially identical to those of a normal lamellar smectic.
The contribution of the DNA lattices to the x-ray scattering intensity is
$`I_{DNA}(𝐪)`$ $`=`$ $`\rho _{DNA}(𝐪)\rho _{DNA}(𝐪)`$ (3.70)
$`=`$ $`{\displaystyle \frac{1}{a}}{\displaystyle \underset{nk}{}}|f_y^k(q_y)|^2|\rho ^{nk}|^2e^{iq_yna}e^{ikq_0z^n}`$ (3.72)
$`\times {\displaystyle }d^2re^{i\mathrm{\Delta }𝐪_{}^k𝐫}S_{yz}(𝐫,na,q_y,k),`$
where $`\mathrm{\Delta }𝐪_{}^k=𝐪_{}kq_0𝐞_z`$ and
$$S_{yz}(𝐫,na,q_y,k)=e^{iq_y[u_y^n(𝐫)u_y^n(0)]}e^{iq_0k[u_z^n(𝐫)u_z^0(0)]}.$$
(3.73)
This correlation function is even more complicated to calculate than $`S(𝐫,na)`$ \[Eq. (3.57)\]. When $`B_{uh}=0`$, $`u_y^n`$ and $`u_z^n`$ decouple, and when $`k=\pm 1`$ we have
$$S_{yz}(𝐫,na,q_y,k=\pm 1)=S_h(𝐫,na,q_y)S(𝐫,na),$$
(3.74)
where we used the fact that correlations in $`u_y^n`$ are the same as those in $`h^n`$. Thus, at $`𝐫=0`$, this function dies off as $`n^{\eta _h1}e^{na/\xi _y}`$ when $`q_y=k_0`$.
If we ignore the $`S_h(𝐫,na,q_y)`$ contribution of $`S_{yz}`$ (which would be justified for $`B_{\mathrm{lam}}\mathrm{}`$) and we use the approximation of Eq. (3.61) for $`S(𝐫,na)`$, then the sum over $`n`$ in Eq. (3.72) can be carried out exactly. The result for the dominant contribution with $`k=\pm 1`$ is
$$I_{DNA}(𝐪)=\frac{1}{a}S_2(\mathrm{\Delta }𝐪_{}^1)|f_{DNA}^1(q_y)|^2F(𝐪),$$
(3.75)
where
$$F(𝐪)=\frac{1\left(\frac{V^uS_2(\mathrm{\Delta }𝐪_{}^1)}{2T}\right)^2}{12\frac{V^uS_2(\mathrm{\Delta }𝐪_{}^1)}{T}\mathrm{cos}(q_ya\sigma \pi )+\left(\frac{V^uS_2(\mathrm{\Delta }𝐪_{}^1)}{2T}\right)^2}.$$
(3.76)
The structure function $`F(𝐪)`$ reaches a maximum at $`q_y=mk_0`$ if a rectangular ($`\sigma =0`$) lattice is preferred and at $`q_y=(m+\frac{1}{2})k_0`$ when a centered rectangular lattices ($`\sigma =1`$) is preferred. At these peaks, $`F(𝐪)`$ does not exhibit singular behavior: the function $`S_2(\mathrm{\Delta }𝐪_{})`$ \[Eq. (3.60)\] entering Eq. (3.76) can be expanded in powers of $`\mathrm{\Delta }𝐪_{}`$ to all orders.
Out-of-plane lamellar fluctuations, ignored in Eqs. (3.75) and (3.76) (by assuming $`B_{\mathrm{lam}}\mathrm{}`$), become significant when the inter-layer positional coupling $`V^u`$ is weak. For example, for $`V^u=0`$ and finite $`B_{\mathrm{lam}}`$, we find using Eqs. (3.72) to (3.74),
$`I_{DNA}`$ $`=`$ $`{\displaystyle \frac{1}{a}}|f_{DNA}^1(q_y)|^2`$ (3.78)
$`\times {\displaystyle }d^2re^{\mathrm{\Delta }𝐪_{}^1𝐫}S_h(𝐫,na=0,q_y)S_2(𝐫),`$
Using Eq. (3.68), it is easy to see that $`I_{DNA}(𝐪)`$ has a maximum at $`q_y=0,\mathrm{\Delta }𝐪_{}^1=0`$ (even for infinitely thin DNA, with $`f_D^1NA(q_y)=1`$). This peak remains at $`q_y=0`$ for sufficiently weak nonzero $`V^u`$. Thus, for sufficiently weak interlayer coupling, the scattering structure factor has a form resembling that of a simple rectangular lattice, with a $`(0,1)`$-like peak at $`q_y=0`$ even though the system has short-range centered rectangular order in real space, With increasing strength of inter-layer positional coupling, this peak will bifurcate, at a critical value of $`V^u`$, into two $`(\pm 1,1)`$ peaks at nonzero $`q_y`$ \[as in Fig. 2\]. Such a change of the form factor is not accompanied by a thermodynamic phase transition. It is similar in character to the transition across so-called disordering line in random microemulsions at which the wave-vector maximizing $`S(𝐪)`$ vanishes.
## IV Stability of the Sliding Columnar Phase
Until now, we have treated the sliding columnar phase as though it were thermodynamically stable. We will now examine conditions for this stability. In order for the SC phase to exist, it must be stable against forces that bring about registry between neighboring DNA lattices to produce the columnar phase; and it must be stable with respect to the disordering effect of dislocations that leads to a nematic lamellar phase. In this section, we will investigate these two effects. We find that the SC phase orders into the columnar phase via a roughening-like transition at a decoupling temperature $`T_d`$ and that it melts to the nematic phase via a dislocation unbinding transition at a temperature $`T_{KT}`$. Thus, the SC phase can only exist in a temperature range $`T_d<T<T_{KT}`$, and it cannot exist at all if $`T_{KT}<T_d`$. We calculate $`T_d`$ and $`T_{KT}`$ for our nearest-neighbor model, and we find that
$$\beta =\frac{T_{KT}}{T_d}=\frac{1}{\pi ^2}<1$$
(4.1)
for systems with temperature independent coupling constants. Thus $`T_{KT}<T_d`$, and, for our model, the SC phase is never thermodynamically stable. The introduction of competing nearest-neighbor and next-nearest-neighbor strain couplings can, however, reverse the inequality on $`T_{KT}`$ and $`T_d`$ and stabilize the SC phase. Even when the SC phase is not thermodynamically stable, there is a range of length scales in the lamellar nematic phase in which correlation functions will exhibit SC behavior.
### A Relevance of Translational Coupling
In the sliding columnar phase, fluctuations of relative displacements $`u_z^{n+1}(𝐫)u_z^n(𝐫)`$ of DNA lattices in neighboring galleries grow with increasing systems size. We found in Sec. III A 2 that $`[u_z^{n+1}(𝐫)u_z^n(𝐫)]^2\mathrm{ln}L`$ \[See Eqs. (3.19) and (3.30)\]. Due to this divergence, the translational coupling \[Eq. (2.20)\] may become irrelevant above a critical decoupling temperature. This result is obtained by calculating the expectation value of the translational coupling
$$_n^u=V^ud^2r\mathrm{cos}[q_0(u_z^{n+1}u_z^n+a_zu_y^n)]$$
(4.2)
with respect to the sliding columnar Hamiltonian in (III A 1). Because the cross correlation $`_zu_y^n(u_z^{n+1}u_z^n)`$ is zero, we can evaluate this quantity exactly:
$$_n^u=V^ue^{W_y}d^2r\mathrm{exp}[q_0^2g^{(1)}(a)],$$
(4.3)
where $`W_y=q_0^2a^2(_zu_y^n)^2/2`$ , Using the fact that $`g^{(1)}(na)`$ diverges logarithmically with system size \[Eq. (3.30)\], we find that the translational coupling scales as
$$_n^uV^ue^{W_y}L^{2\eta _d},$$
(4.4)
where
$$\eta _d=2(q_0l_u)^2S_1(0)=\frac{4T}{d^2\sqrt{BK_y}}.$$
(4.5)
The condition $`2\eta _d=0`$ defines the critical decoupling temperature
$$T_d=\frac{d^2\sqrt{BK_y}}{2}.$$
(4.6)
When $`T<T_d`$, the translational coupling scales as system size to a positive power, and $`V^u`$ is relevant. In this case, the system becomes a columnar phase at the longest lengthscales with a nonzero shear modulus for shifting neighboring lattices relative to each other and long-range positional order in the $`yz`$ plane. When $`T>T_d`$, the translational coupling scales as the system size to a negative power, and $`V^u`$ is irrelevant. The system flows to the sliding columnar phase at the longest lengthscales when $`T>T_d`$. Thus, $`T_d`$ marks the transition from the columnar phase to the SC phase as temperature is increased. This transition is of the roughening type. There is no energy cost for shifting neighboring lattices relative to each other, and thus the sliding columnar phase is positionally disordered in the $`yz`$ plane. These conclusions are supported by an RG analysis to be presented in Ref. .
### B Dislocation unbinding
We have just seen how thermal fluctuations weaken interlayer couplings to produce the SC phase from the columnar phase. We will now consider the disordering effects of dislocations. An individual edge dislocation in a $`2D`$ smectic has a finite rather than a logarithmically divergent energy. As a result, there are thermally excited unbound vortex pairs at all temperatures that convert the $`2D`$ smectic into a $`2D`$ nematic at the longest lengthscales. In a sliding columnar phase, each DNA smectic layer experiences an orientational ordering field from its neighbors. As a result, the energy of an individual edge dislocation in a given layer diverges logarithmically with system size, and there can be a Kosterlitz-Thouless dislocation unbinding transition. The low-temperature phase with bound dislocations is the sliding columnar phase, and the high-temperature phase with unbound dislocations is the lamellar nematic phase.
#### 1 Dislocation Energy
The DNA smectic lattice in each layer can have edge dislocation defects in which the displacement field $`u_z^n`$ undergoes a change of $`kd`$, where $`k`$ is an integer, in one circuit around the defect core. If there are dislocations with integer strengths $`k_{n,l}`$ at positions $`𝐫_{n,l}`$ in layer $`n`$, then
$$\mathbf{}_{}\times 𝐯^n(𝐫)=b^n(𝐫)\widehat{y},$$
(4.7)
where $`𝐯^n(𝐫)=\mathbf{}_{}u_z^n(𝐫)`$ and
$$b^n(𝐫)=d\underset{l}{}k_{n,l}\delta ^2(𝐫𝐫_{n,l})$$
(4.8)
is the dislocation density in layer $`n`$.
We can now calculate the energy cost of an arbitrary assembly of edge dislocations in the sliding columnar phase using the SC Hamiltonian
$`^{\mathrm{sc}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{n}{}}a{\displaystyle }d^2r[B(v_z^n)^2+K(_xv_x^n)^2`$ (4.10)
$`+{\displaystyle \frac{K_y}{a^2}}(v_x^nv_x^{n+1})^2]`$
written in terms of $`𝐯^n(𝐫)=(v_x^n,v_z^n)`$. Our strategy is to minimize this Hamiltonian subject to a nonzero dislocation density $`b^n(𝐫)`$. As usual, the calculation is simpler in Fourier space. To transform from real space to Fourier space, we use
$`𝐯(𝐪)`$ $`=`$ $`{\displaystyle \underset{n}{}}a{\displaystyle d^2re^{i(𝐪_{}𝐫+q_yna)}𝐯^n(𝐫)},`$ (4.11)
$`b(𝐪)`$ $`=`$ $`{\displaystyle \underset{n}{}}a{\displaystyle d^2re^{i(𝐪_{}𝐫+q_yna)}b^n(𝐫)}`$ (4.12)
$`=`$ $`ad{\displaystyle \underset{n,l}{}}k_{n,l}e^{i(𝐪_{}𝐫_{n,l}+q_yna)}.`$ (4.13)
We minimize Eq. (4.10) using the Euler-Lagrange equations and find
$`v_x(𝐪)`$ $`=`$ $`{\displaystyle \frac{q_z}{q_x}}{\displaystyle \frac{B}{K(𝐪)q^2}}v_z(𝐪),`$ (4.14)
where
$$K(𝐪)q^2=Kq_x^2+K_yq_y^2p(q_ya).$$
(4.15)
We then employ the constraint, Eq. (4.7), to relate $`v_z(𝐪)`$ to the specified dislocation density $`b(𝐪)`$. In the final step, we insert the expressions for $`𝐯(𝐪)`$ into the Fourier transformed version of Eq. (4.10) and find
$$E_D=\frac{B}{2}_{\pi /a}^{\pi /a}\frac{dq_y}{2\pi }\frac{d^2q_{}}{(2\pi )^2}\frac{K(𝐪)q^2|b(𝐪)|^2}{Bq_z^2+K(𝐪)q^2q_x^2}$$
(4.16)
for the energy of edge dislocations in the sliding columnar phase. Also note that if we set $`K_y=0`$, Eq. (4.16) reduces to the expression for the energy cost for edge dislocations in a 2D smectic.
We now decompose the dislocation energy into a part that diverges with system size and a part that diverges only with the separation between defects. After we insert Eq. (4.12) into Eq. (4.16), we find
$`E_D`$ $`=`$ $`{\displaystyle \frac{\sqrt{BK_y}d^2}{4\pi ^2}}[{\displaystyle \underset{n,n^{}}{}}\sigma _n\sigma _n^{}J_{nn^{}}\mathrm{ln}[B_{nn^{}}(\mathrm{\Lambda }_x)L_x/x^{}]`$ (4.17)
$`+`$ $`{\displaystyle \underset{n,n^{},l,l^{}}{}}k_{n,l}k_{n^{},l^{}}E_{nn^{}}(𝐫_{n,l}𝐫_{n^{},l^{}})],`$ (4.18)
where $`\sigma _n=_lk_{n,l}`$ is the total dislocation charge in the $`n`$th layer,
$$J_n=_0^\pi 𝑑u\sqrt{u^2p(u)}\mathrm{cos}nu=\frac{4}{14n^2},$$
(4.19)
and $`B_n(\mathrm{\Lambda }_x)`$ is given in Appendix G. The first term diverges logarithmically with system size. The interaction energy $`E_n(𝐫)`$ \[Eq. (4.17)\], on the other hand, is written in terms of an integral that vanishes when $`𝐫=0`$, and thus it diverges only with the separation between defects. In the limit of infinite system size, $`J_{nn^{}}`$ is a positive definite matrix. The dislocation energy diverges logarithmically with system size unless the total dislocation charge in each layer is zero, i.e., $`\sigma _n=0`$ for every $`n`$.
The dislocation interaction energy of Eq. (4.17),
$`E_n(𝐫)`$ $`=`$ $`{\displaystyle _0^\pi }𝑑u{\displaystyle _0^{\mathrm{\Lambda }_xx^{}}}{\displaystyle \frac{dt}{t}}\sqrt{t^2+u^2p(u)}\mathrm{cos}nu`$ (4.21)
$`\times \left(1\mathrm{cos}[tx/x^{}]\mathrm{exp}\left[{\displaystyle \frac{z}{z^{}}}t\sqrt{t^2+u^2p(u)}\right]\right),`$
is difficult to calculate for arbitrary separations $`(𝐫,na)`$, however, it can be calculated in the limits $`xx^{}`$, $`z=0`$ and $`zz^{}`$, $`x=0`$. We find that $`E_n(𝐫)`$ scales as
$`E_n(𝐫)=J_n\{\begin{array}{cc}\mathrm{ln}[C_n^x(\mathrm{\Lambda }_x)|x|/x^{}]\hfill & \mathrm{if}xx^{}\mathrm{and}z=0\hfill \\ \mathrm{ln}[C_n^z(\mathrm{\Lambda }_x)|z|/z^{}]\hfill & \mathrm{if}zz^{}\mathrm{and}x=0,\hfill \end{array}`$ (4.24)
where $`C_n^{x,z}(\mathrm{\Lambda }_x)`$ are numbers that depend on the layer index $`n`$ and the cutoff $`\mathrm{\Lambda }_x`$ and are given in Appendix G. Note that the $`\mathrm{\Lambda }_x\mathrm{}`$ limit of $`E_n(𝐫)`$ is not well-defined. Since $`E_n(r)\mathrm{ln}r/(4n^21)`$ for large $`r`$, $`E_n(𝐫)`$ is positive for all $`n>0`$ and negative for $`n=0`$. As a result, like-signed dislocations in different layers attract each other, whereas like-signed dislocations in the same layer repel each other.
The attraction of like-sign dislocations in different galleries makes physical sense. The dislocation excitations of the SC phase can be viewed as closed loops carrying a single value of charge with portions of the loop passing normal to layers and portions passing parallel to layers. Those parts of the loops passing through layers give rise to layer dislocations, and those parts aligned parallel to layers cost no energy because the shear modulus is zero. A direction can be assigned to a loop. A loop section penetrating a layer in the upward direction gives rise to a dislocation of one sign while one penetrating in the downward direction gives rise to a dislocation of the opposite sign. As is apparent from Fig. 4, the expectation is that the lowest energy configuration of a continuous loop will penetrate nearby points in nearby layers in the same direction, i.e., that nearby dislocations in nearby layers will have the same sign. This argument is, of course, not rigorous because it is possible for dislocation lines running parallel to layers to cross to produce nearby dislocations in nearby layers of opposite sign as shown in Fig. 4(e). In fact, it is possible to construct interlayer interactions for which like sign dislocations in neighboring layers repel rather than attract.
#### 2 Dislocation Unbinding Temperature
In the previous section, we showed that the dislocation energy is logarithmically divergent unless there is dislocation charge neutrality in each layer. Consider now a configuration in which layer $`n`$ has a single dislocation of charge $`\sigma _n`$ (which may be zero). The free energy of such a composite dislocation configuration is, by Eq. (4.17),
$`F_D`$ $`=`$ $`E_DTS_D`$ (4.25)
$`=`$ $`\left({\displaystyle \frac{\sqrt{BK_y}d^2}{4\pi ^2T}}{\displaystyle \underset{n,n^{}}{}}J_{nn^{}}\sigma _n\sigma _n^{}2T\right)\mathrm{ln}L+\mathrm{const}..`$ (4.26)
Clearly Eq. (4.25) indicates that composite dislocations pairs unbind for temperatures above the Kosterlitz-Thouless type temperature,
$$T_{KT}[\sigma _n]=\frac{\sqrt{BK_y}d^2}{8\pi ^2}\underset{n,n^{}}{}J_{nn^{}}\sigma _n\sigma _n^{}.$$
(4.27)
If $`\sigma _n`$ is nonzero only in layer $`m`$ ($`\sigma _n=\delta _{nm}`$), the transition corresponds to unbinding of simple dislocations pairs within a single layer. If $`\sigma _n`$ is nonzero in more than one layer, the transition corresponds to the unbinding of composite multilayer dislocations pairs as depicted in Fig. 4b.
We see that each configuration of dislocations $`\{\sigma _n\}`$ will have a different unbinding temperature. The naive expectation is that $`T_{KT}`$ is lowest when there is a single dislocation in a single layer. It is, however, possible that composite dislocations might melt at a lower temperature. Hence, we define $`T_{KT}=\mathrm{min}_{\sigma _n}T_{KT}[\sigma _n]`$. Above this temperature, one of the dislocation configurations (either composite or individual) unbinds, renormalizes the compression modulus $`B`$ to zero, and destroys the in-plane $`2D`$ smectic order of the sliding columnar phase. The lowest unbinding temperature corresponds to the configuration with the lowest $`\mathrm{ln}L_x`$ divergent energy \[See Eqs. (4.27) and (4.17)\]. Since $`J_n<0`$ for all $`n>0`$, $`|J_n|`$ decays with increasing layer separation, and the dislocation energy scales quadratically with the dislocation charge, the lowest energy configuration with $`N`$ defects is a string of $`N`$ defects of strength $`+1`$ with one in each layer, i.e. $`\sigma _n=_{p=1}^N\delta _{n,p}`$. It is then straightforward to show that the lowest of these configurations is an individual defect with $`N=1`$. Therefore, the dislocation unbinding temperature for the sliding columnar phase is
$$T_{KT}=\frac{d^2\sqrt{BK_y}}{2\pi ^2}.$$
(4.28)
The ratio of this temperature to the unbinding temperature $`T_d`$ \[Eq. (4.6)\] yields Eq. (4.1) for $`\beta =T_{KT}/T_d=1/\pi ^2`$.
### C Comments on Lyotropic Systems
An interesting aspect of the present study is the transitions from the sliding to the columnar and the nematic phase we discussed above in terms of the corresponding critical temperatures: the decoupling temperature $`T_d`$ \[Eq. (4.6\], for the columnar to sliding phase transition, and $`T_{KT}`$ \[Eq. (4.28)\] for the sliding to nematic phase transition. In lyotropic systems of interest here \[Refs. -, temperature changes can hardly produce significant effects. Nonetheless, the above transitions can be still triggered by varying the period of DNA lattices (as lattices described in Refs. -). Indeed, by previous discussions, the decoupling condition $`T>T_d`$, is equivalent to $`\eta _d>2`$, with $`\eta _d=4T/d^2(BK_y)^{1/2}`$. Thus, the decoupling transition from the columnar to sliding phase may be triggered by changing the strength of the inter-layer orientational coupling constant $`K_y`$. As evidenced clearly in experiments , the strength of inter-layer couplings significantly increases by increasing the period of DNA lattices. Thus, by the swelling these lattices, one may reach the transition in which the sliding phase changes into a columnar phase, as discussed in Sec. IV A, and, in more detail, in Ref. . Finally, we emphasize that, within the simple model discussed here, there is no true sliding phase, as suggested by Eq. (4.1) that applies to systems with temperature independent coupling constants. It should be stressed however that the restriction of having temperature-independent constants is not essential for this important conclusion about the sliding phase stability. Neither it is essential whether the system is thermotropic or lyotropic in its character. Indeed, by Sec. IV A, the condition to have dislocations bound \[i.e., $`T<T_{KT}`$\] is equivalent to $`\eta _{KT}>2`$, with $`\eta _{KT}=d^2(BK_y)^{1/2}/\pi ^2T`$, for the model discussed here \[Sec. II\]. Note that
$$\eta _d\eta _{KT}=4/\pi ^2<4.$$
(4.29)
This relation implies that the condition for dislocations confinement ($`\eta _{KT}>2`$) and the condition ($`\eta _d>2`$) for sliding phase decoupling cannot be simultaneously realized. Whenever the decoupling condition is realized, $`\eta _d>2`$, $`\eta _{KT}=4/\pi ^2\eta _d<2`$, and dislocations will unbind. This feature turns the sliding phase into a nematic phase at long length scales. Still, at low dislocation densities (what may well be the realistic case, see Sec. V), in such a nematic phase, there is a broad range of length scales exhibiting sliding-phase correlations and other structural properties discussed in this work.
## V Characteristic Lengthscales
Whether or not sliding columnar behavior is seen in recent X-ray scattering experiments on CL-DNA complexes is still an open question. Even though the sliding columnar phase is converted into a nematic lamellar phase at the longest lengthscales, the SC phase may exist at shorter lengthscales determined by the density of edge dislocations. Thus, CL-DNA complexes studied in recent experiments with domain sizes $`L0.1\mu `$m may exhibit sliding columnar behavior. However, in the next round of experiments it will be important to prepare aligned CL-DNA samples because powder-averaging complicates the functional form of the scattering intensity and makes it difficult to identify sliding columnar behavior.
We showed in Sec. III that the density-density correlation function $`S_2(𝐫)`$ \[Eq. (3.54)\] displays different functional forms depending on the magnitude of the in-plane separation $`𝐫`$. The crossover lengthscales for the correlation function are $`l_x`$, $`l_z`$, $`\xi _d`$, $`x^{}`$, and $`z^{}`$. The harmonic $`2D`$ smectic regime is defined by $`x,z<l_{x,z}`$ and the nonlinear $`2D`$ smectic regime is defined by $`x,z>l_{x,z}`$, where the nonlinear lengths $`l_{x,z}`$ are given in Eq. (3.34). To estimate the nonlinear lengths, we must determine $`l_x`$ and $`l_z`$ in terms the experimentally measured quantities $`d`$ and the harmonic correlation length $`\xi _z`$ of Eq. (3.41). $`l_x`$ and $`l_z`$ can be written in terms of $`d`$ and $`\xi _z`$ by solving Eq. (3.41) for $`B_2`$ and using
$$K_2=T\xi _p/2d,$$
(5.1)
where $`\xi _p=500`$Å is the persistence length of DNA. Table I evaluates the nonlinear lengths for two DNA spacings of the experiments of references . It shows that $`L<l_{x,z}`$ and indicates that significant departure from harmonic 2D smectic behavior is not expected in agreement with these experiments.
The finite length of the DNA molecules $`l_{\mathrm{DNA}}16\mu `$m introduces another crossover lengthscale. The density of DNA molecules within a given layer is $`\rho =1/dl_{\mathrm{DNA}}`$. If we assume that each free end of a DNA strand corresponds to a dislocation, then we can estimate the characteristic distance between dislocations to be
$$\xi _d=\sqrt{dl_{\mathrm{DNA}}}.$$
(5.2)
This length is clearly smaller than the actual distance between dislocations because not every end of a DNA strand has to produce a dislocation: a series of DNA strands can align in a row to produce a layer of a $`2D`$ smectic. At lengths scales greater than $`\xi _d`$, the $`2D`$ smectic behaves like a nematic. Using the estimate of Eq. (5.2), we find $`\xi _d0.21\mu `$m and $`0.30\mu `$m when $`d=28`$Å and $`55`$Å, respectively. Note that $`L<\xi _d`$, and thus the subdomains are small enough to possess 2D smectic ordering.
We can also estimate the energy cost for creating hairpin edge dislocations within the 2D smectic lattices. Hairpins cause the DNA director to change by $`\pi `$ over a lattice spacing $`d`$. The energy cost for a hairpin can be estimated from the 2D bending energy. If we ignore relaxation of the $`2D`$ smectic layers, we estimate the bending energy to be
$$\frac{E_{hp}}{2T}=\frac{\pi }{4}\frac{\xi _p}{d},$$
(5.3)
which implies that hairpins are favored on lengthscales greater than
$$\xi _{hp}=d\mathrm{exp}\left(\frac{E_{hp}}{2T}\right).$$
(5.4)
Since $`\xi _{hp}\xi _d>L`$ throughout the experimental range in $`d`$, hairpins are not important for the current set of experiments. We do not yet have accurate estimates of $`x^{}`$ and $`z^{}`$ since the value of the orientational rigidity $`K_y`$ is unknown. Scattering experiments will see sliding columnar behavior on lengthscales less than $`\xi _d`$ if $`x^{},z^{}<\xi _d`$.
## VI Discussion and Conclusion
In this paper, we have introduced the new sliding columnar (SC) phase of matter which may exist in layered systems composed of weakly-coupled 2D smectic lattices. The sliding columnar phase is characterized by weak positional but strong orientational correlations between neighboring 2D smectic lattices. The SC harmonic free energy contains an orientational rigidity that aligns neighboring 2D smectic lattices in addition to in-plane compression and bending moduli. The SC phase is characterized by a vanishing shear modulus for relative displacements of $`2D`$ smectic lattices. The presence of the orientational rigidity fundamentally alters the energy spectrum. In light of this, we have calculated the structural properties of the sliding columnar phase such as the SC displacement correlation function, scattering intensity, and dislocation energy.
Experimental research on layered liquid crystalline systems studied in this work is in progress-. Correlation peaks theoretically discussed here have been observed in a number of scattering experiments on DNA-cationic lipid complexes. Reference reports DNA scattering patterns that clearly reflect the short-range centered-rectangular order depicted in Fig. 2. Reference , on the other hand, reports peaks in a different system at the simple rectangular lattice positions $`(0,0,\pm q_0)`$ \[i.e., at $`(0,1)`$ peaks\]. In spite of this difference in scattering data, it is likely that both systems have short-range centered rectangular order in real space. Indeed, in a strongly fluctuating disordered systems, positions of correlation peaks in $`q`$-space may not always reflect length scales and short-range order in real space. For example, the correlation peak in $`q`$-space in random microemulsions may occur at $`q=0`$ even in situations in which real space data still show a finite structural length scale. As discussed at the end of Sec. III C, for sufficiently small coupling inter-layer positional coupling $`V^u/T`$, the scattering form factor $`I_{DNA}(𝐪)`$ exhibits a maximum at the simple rectangular-lattice points $`(0,0,\pm q_0)`$ even though the system has short-range centered rectangular order in real space. The differences between the positions of the scattering peaks in Refs. and may be explained by the differences in thermal disorder and/or inter-layer couplings strengths of the two systems. The interlayer correlation length reported in Ref. is several times longer than that reported in Ref. . The former system is stiffer and more ordered than the latter. Consequently the $`(1,1)`$ and $`(1,1)`$ correlation peaks in Ref. reflect the centered-rectangular order in real space depicted in Fig. 1, whereas the the $`(0,1)`$ peak reported in Ref. reflects the merging of $`(\pm 1,1)`$ peaks brought about by strong thermal fluctuations and/or sufficiently weak inter-layer couplings.
After completing this work, we learned from Joachim Rädler of experiments showing a continuous variation of the DNA form factor with changing strength of inter-layer coupling $`V^u`$ consistent with the change we describe above and in Sec. III D \[crossover from centered rectangular scattering pattern to one resembling that of a simple rectangular lattice\]. In these experiments, the variation of inter-layer positional coupling is produced by varying the period of DNA lattices in galleries. This coupling generally increases with increasing period of DNA lattices, most likely because the Coulomb interaction increases and the elastic membrane deformation increases, as suggested by electron density maps of the system.
###### Acknowledgements.
LG acknowledges support by Mylar Pharmaceuticals, Inc., and CSO and TCL acknowledge support from the National Science Foundation under grant DMR97–30405. We are grateful for helpful discussions with Ilya Koltover, Joachim Rädler, Cyrus Safinya, Tim Salditt, and John Toner.
## A Derivation of Effective Hamiltonians
In this appendix, we will derive the effective Hamiltonians, $`^{SC}`$ \[Eq. (2.23)\], $`_z^{SC}`$ \[Eq. (III A 1)\], $`_h`$ \[Eq. (3.43)\] obtained, respectively, by integrating out $`h^n(𝐫)`$, $`h^n(𝐫)`$ and $`u_y^n(𝐫)`$, and $`u_z^n(𝐫)`$ and $`u_y^n(𝐫)`$ from $`\stackrel{~}{}=^{\mathrm{bend}}+^{\mathrm{com}}+^{\mathrm{rot}}`$, which we can express in linearized form in Fourier space as
$`\stackrel{~}{}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle }{\displaystyle \frac{d^3q}{(2\pi )^3}}\{A_h(𝐪)|h(𝐪)|^2+{\displaystyle \underset{\sigma }{}}A_\sigma (𝐪)|u_\sigma (𝐪)|^2`$ (1.2)
$`{\displaystyle \underset{\sigma }{}}[\lambda _\sigma (𝐪)u_\sigma (𝐪)h(𝐪)+\lambda _\sigma (𝐪)u_\sigma (𝐪)h(𝐪)]\},`$
where $`\sigma =y,z`$, and
$`A_h(𝐪)`$ $`=`$ $`K_{3d}q_{}^4+4(B_{3d}/a^2)`$ (1.3)
$`A_y(𝐪)`$ $`=`$ $`4(B_{3d}/a^2)+K_{2d}q_x^4`$ (1.4)
$`A_z(𝐪)`$ $`=`$ $`B_{2d}q_z^2+K_{2d}q_x^4+K_yq_x^2q_y^2p(q_ya)`$ (1.5)
$`\lambda _y(𝐪)`$ $`=`$ $`(2B_{3d}/a^2)\left(1+e^{iq_ya}\right)`$ (1.6)
$`\lambda _z(𝐪)`$ $`=`$ $`(B_{uh}/a)iq_z\left(e^{iq_ya}1\right).`$ (1.7)
Integration over $`h(𝐪)`$ yields
$`\stackrel{~}{}^{SC}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle }{\displaystyle \frac{d^3q}{(2\pi )^3}}[A_y^{}(𝐪)|u_y(𝐪)|^2+A_z^{}|u_z(𝐪)|^2`$ (1.9)
$`\mu (𝐪)u_y(𝐪)u_z(𝐪)\mu (𝐪)u_y(𝐪)u_z(𝐪)],`$
where
$`A_\sigma ^{}(𝐪)`$ $`=`$ $`A_\sigma (𝐪){\displaystyle \frac{|\lambda _\sigma (𝐪)|^2}{A_h(𝐪)}}`$ (1.10)
$`\mu (𝐪)`$ $`=`$ $`{\displaystyle \frac{\lambda _y(𝐪)\lambda _z(𝐪)}{A_h(𝐪)}}.`$ (1.11)
In the small $`𝐪_{}`$ limit, these expressions reduce to
$`A_y^{}(𝐪)`$ $`=`$ $`(B_{3d}/a^2)2(1\mathrm{cos}q_ya)+K^{\alpha \beta }q_\alpha ^2q_\beta ^2`$ (1.12)
$`A_z^{}(𝐪)`$ $`=`$ $`B_{2d}q_z^2{\displaystyle \frac{B_{uh}^2}{4B_{3d}}}q_z^22(1\mathrm{cos}q_ya)+K_{2d}q_x^4`$ (1.14)
$`+K_yq_x^2q_y^2p(q_ya)`$
$`\mu (𝐪)`$ $`=`$ $`(B_{uh}/a)q_z\mathrm{sin}q_ya,`$ (1.15)
where $`K^{\alpha \beta }`$ is the bending rigidity tensor of Eq. (2.29). Fourier transformation of Eq. (1.9) back to real space using the low-$`𝐪_{}`$ form of $`A_y(𝐪)`$, $`A_z(𝐪)`$ and $`\mu (𝐪)`$ produces the sliding phase Hamiltonian $`^{SC}`$ of Eq. (2.23).
The effective Hamiltonian for $`u_z`$ is obtained by integrating $`\stackrel{~}{}^{SC}`$ over $`u_y`$. The result is
$$\stackrel{~}{}_z^{SC}=\frac{1}{2}\frac{d^3q}{(2\pi )^3}A_z^{\prime \prime }(𝐪)|u_z(𝐪)|^2,$$
(1.16)
where
$$A_z^{\prime \prime }(𝐪)=A_z^{}(𝐪)\frac{|\mu (𝐪)|^2}{A_y^{}(𝐪)}.$$
(1.17)
In the $`𝐪_{}0`$ limit, $`A_z^{\prime \prime }(𝐪)`$ becomes
$`A_z^{\prime \prime }(𝐪)`$ $``$ $`q_z^2[B_{2d}{\displaystyle \frac{B_{uh}^2}{4B_{3d}}}2(1\mathrm{cos}q_ya)`$ (1.20)
$`{\displaystyle \frac{B_{uh}^2}{B_{3d}}}{\displaystyle \frac{\mathrm{sin}^2q_ya}{2(1\mathrm{cos}qya)}}]`$
$`+K_yq_x^2q_y^2p(q_ya)+K_{2d}q_x^4`$
$`=`$ $`q_z^2\left(B_{2d}{\displaystyle \frac{B_{uh}^2}{B_{3d}}}\right)+K_{2d}q_x^4,`$ (1.22)
$`+K_yq_x^2q_y^2p(q_ya)`$
where we used the identity
$$1\mathrm{cos}\varphi +\frac{\mathrm{sin}^2\varphi }{1\mathrm{cos}\varphi }=2$$
(1.23)
to obtain Eq. (1.22) Transformation of $`\stackrel{~}{}_z^{SC}`$ to real space using the low $`𝐪_{}`$ form of $`A_z^{\prime \prime }`$ produces $`_z^{SC}`$ of Eq. (III A 1). The absence of any $`q_y`$ dependence in the coefficient of $`q_z^2`$ in $`A_z^{\prime \prime }(𝐪)`$ can be traced the equality $`B_{uh}^2/B_{3d}=4B_{zz}`$ \[Eq. (2.28)\] that results from integrating out $`h^n`$ from the original model. Effective Hamiltonians for $`h^n(𝐫)`$ and $`u_y^n(𝐫)`$ can be obtained by respectively integrating out $`u_y^n`$ and $`u_z^n`$ from $`^{\mathrm{tot}}`$ and $`u_z^n`$ from $`^{SC}`$. The results are
$`_h^{SC}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{d^3q}{(2\pi )^3}A_h^{\prime \prime }(𝐪)|h(𝐪)|^2}`$ (1.24)
$`_y^{SC}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{d^3q}{(2\pi )^3}A_y^{\prime \prime }(𝐪)|u_y(𝐪)|^2}.`$ (1.25)
where
$`A_h^{\prime \prime }(𝐪)`$ $`=`$ $`A_h(𝐪){\displaystyle \frac{|\lambda _z(𝐪)|^2}{A_z(𝐪)}}{\displaystyle \frac{|\lambda _y(𝐪)|^2}{A_y(𝐪)}}`$ (1.26)
$`A_y^{\prime \prime }(𝐪)`$ $`=`$ $`A_y(𝐪){\displaystyle \frac{|\lambda _y(𝐪)|^2}{A_h(𝐪)|\lambda _z(𝐪)|^2/A_z(𝐪)}}.`$ (1.27)
In the limit $`𝐪_{}0`$ with $`q_x^2(B_{2d}a^2/K_y)q_z^2`$, both of these Hamiltonian reduce to $`_h`$ in Eq. (3.43).
## B Calculation of $`(u_z^n)^2`$
In this Appendix, the expression for the sliding columnar displacement fluctuations given in (3.27) is derived. To do this, we evaluate the integral of the the Fourier transformed SC correlator $`G_{zz}(𝐪)`$ over all $`q`$-space:
$`(u_z^n)^2`$ $`=`$ $`{\displaystyle \frac{d^3q}{(2\pi )^3}\frac{T}{Bq_z^2+Kq_x^4+K_yq_x^2q_y^2p(q_ya)}},`$ (2.1)
where $`p(u)=2(1\mathrm{cos}u)/u^2`$. The fluctuations diverge at small wavenumbers $`𝐪1/L`$, where $`L`$ is the system size. To calculate how the fluctuations scale with $`L_x`$, we set $`L_z\mathrm{}`$ and $`L_yL_x`$. Note that the SC form for $`G_{zz}(𝐪)`$ is valid only when $`L_xx^{}`$, where $`x^{}=a/\mu _y`$ and $`\mu _y=\sqrt{K_y/K}`$. The first step in the calculation is to perform the integration over $`q_z`$ with $`\mathrm{\Lambda }_z\mathrm{}`$ and then use the fact that $`G_{zz}(𝐪)`$ is an even function of $`𝐪`$ so that the remaining integrals run over only positive $`q_x`$ and $`q_y`$. Note that taking the $`\mathrm{\Lambda }_z\mathrm{}`$ limit does not alter $`q_{x,y}0`$ divergences. The resulting expression
$`(u_z^n)^2`$ $`=`$ $`{\displaystyle \frac{T}{2\pi ^2\sqrt{BK}}}{\displaystyle _{L_x^1}^{\mathrm{\Lambda }_x}}{\displaystyle \frac{dq_x}{q_x}}{\displaystyle _{L_y^1}^{\mathrm{\Lambda }_y}}𝑑q_y{\displaystyle \frac{1}{\sqrt{q_x^2+\mu _y^2q_y^2p(q_ya)}}},`$ (2.2)
where $`\mathrm{\Lambda }_y=\pi /a`$, is made dimensionless by changing variables to $`v=q_xx^{}`$ and $`w=q_ya`$. The $`\mathrm{ln}^2L_x`$ divergence of the displacement fluctuations can be seen immediately by looking at the $`q_{x,y}0`$ limit of this expression.
The SC displacement fluctuations can be written as the sum of a continuum term $`I_c`$ that does not depend on $`p(w)`$ and discrete term $`I_d`$ that does depend on $`p(w)`$.
$$(u_z^n)^2\frac{T}{2\pi ^2\sqrt{BK_y}}\left[I_c+I_d\right],$$
(2.3)
where $`I_c`$ and $`I_d`$ are defined by
$`I_c`$ $`=`$ $`{\displaystyle _{x^{}L_x^1}^{\mathrm{\Lambda }_xx^{}}}{\displaystyle \frac{dv}{v}}{\displaystyle _{aL_y^1}^{a\mathrm{\Lambda }_y}}{\displaystyle \frac{dw}{\sqrt{v^2+w^2}}}`$ (2.4)
$`I_d`$ $`=`$ $`{\displaystyle _{x^{}L_x^1}^{\mathrm{\Lambda }_xx^{}}}{\displaystyle \frac{dv}{v}}{\displaystyle _{aL_y^1}^{a\mathrm{\Lambda }_y}}𝑑w\left[{\displaystyle \frac{1}{\sqrt{v^2+w^2p(w)}}}{\displaystyle \frac{1}{\sqrt{v^2+w^2}}}\right].`$ (2.5)
We first focus on the continuum contribution $`I_c`$. The integral over $`w`$ is straightforward;
$`I_c`$ $`=`$ $`{\displaystyle _{x^{}L_x^1}^{\mathrm{\Lambda }_xx^{}}}{\displaystyle \frac{dv}{v}}\left[f(\mathrm{\Lambda }_ya,v)f(L_y^1a,v)\right]`$ (2.6)
$``$ $`I_c^{(1)}I_c^{(2)},`$ (2.7)
where
$$f(x,v)=\mathrm{ln}\left[x+\sqrt{x^2+v^2}\right].$$
(2.8)
The integral over $`v`$ in $`I_c^{(1)}`$ can be evaluated by separating the function
$`f(\mathrm{\Lambda }_ya,v)`$ $`=`$ $`\mathrm{ln}[2\mathrm{\Lambda }_ya]+\mathrm{ln}\left[{\displaystyle \frac{1}{2}}+{\displaystyle \frac{1}{2}}\sqrt{1+(v/\mathrm{\Lambda }_ya)^2}\right],`$ (2.9)
into a constant term and a term that is well-behaved at small $`v`$. We then insert this expression into (2.6) and find that
$`I_c^{(1)}`$ $`=`$ $`\mathrm{ln}[2\mathrm{\Lambda }_ya]\mathrm{ln}[\mathrm{\Lambda }_xL_x]`$ (2.10)
$`+`$ $`{\displaystyle _{x^{}L_x^1}^{\mathrm{\Lambda }_xx^{}}}{\displaystyle \frac{dv}{v}}\mathrm{ln}\left[{\displaystyle \frac{1}{2}}+{\displaystyle \frac{1}{2}}\sqrt{1+(v/\mathrm{\Lambda }_ya)^2}\right].`$ (2.11)
The first term diverges with system size $`L_x`$, and the second term is nondivergent. $`f(aL_y^1,v)`$ can also be separated into a constant term and a term that depends on $`v`$. We then integrate $`f(aL_y^1,v)`$ over $`v`$ to find
$`I_c^{(2)}`$ $`=`$ $`\mathrm{ln}[2aL_y^1]\mathrm{ln}[\mathrm{\Lambda }_xL_x]`$ (2.12)
$`+`$ $`{\displaystyle _{x^{}L_x^1}^{\mathrm{\Lambda }_xx^{}}}{\displaystyle \frac{dv}{v}}\mathrm{ln}\left[{\displaystyle \frac{1}{2}}+{\displaystyle \frac{1}{2}}\sqrt{1+\left({\displaystyle \frac{L_yv}{a}}\right)^2}\right].`$ (2.13)
The $`L_y`$ dependence in the integrand of the second term can be moved to limits of the integral by changing variables to $`s=L_yv/a`$. In contrast to the previous expression for $`I_c^{(1)}`$ in (2.10), the large $`s`$ part of the integral in (2.12) diverges with system size. The divergence can be isolated by adding and subtracting $`\mathrm{ln}[s/2]/s`$. The resulting expression,
$`I_c^{(2)}`$ $`=`$ $`\mathrm{ln}[2aL_y^1]\mathrm{ln}[\mathrm{\Lambda }_xL_x]+{\displaystyle \frac{1}{2}}\mathrm{ln}^2\left[{\displaystyle \frac{\mathrm{\Lambda }_xL_y}{2\mu _y}}\right]{\displaystyle \frac{1}{2}}\mathrm{ln}^2\left[{\displaystyle \frac{L_y}{2\mu _yL_x}}\right]`$ (2.14)
$`+`$ $`{\displaystyle _{L_y/L_x\mu _y}^{\mathrm{\Lambda }_xL_y/\mu _y}}{\displaystyle \frac{ds}{s}}\left(\mathrm{ln}\left[{\displaystyle \frac{1}{2}}+{\displaystyle \frac{1}{2}}\sqrt{1+s^2}\right]\mathrm{ln}[s/2]\right),`$ (2.15)
has two terms that diverge and two terms that do not diverge with system size. Note that $`L_y/L_x\mu _y`$ is $`𝒪(1)`$ since $`L_xL_y`$. We then subtract $`I_c^{(2)}`$ from $`I_c^{(1)}`$, drop the nondivergent terms, and set $`L_y=L_x`$ to obtain
$`I_c=\mathrm{ln}[\mathrm{\Lambda }_xL_x]\mathrm{ln}[\mathrm{\Lambda }_yL_y]{\displaystyle \frac{1}{2}}\mathrm{ln}^2\left[{\displaystyle \frac{\mathrm{\Lambda }_xL_y}{2\mu _y}}\right]={\displaystyle \frac{1}{2}}\mathrm{ln}^2[2\mu _y\mathrm{\Lambda }_yL_x]`$ (2.16)
for the continuum contribution to the displacement fluctuations.
The discrete contribution is obtained by evaluating
$$I_d=_{x^{}L_x^1}^{\mathrm{\Lambda }_xx^{}}\frac{dv}{v}F(v),$$
(2.17)
where
$$F(v)=_0^\pi 𝑑w\left[\frac{1}{\sqrt{v^2+w^2p(w)}}\frac{1}{\sqrt{v^2+w^2}}\right].$$
(2.18)
In the definition of $`F(v)`$, we can take the lower limit to zero since the small $`w`$ part of integral is well-behaved. In contrast to the continuum term, the discrete term diverges logarithmically with system size. To see this, we expand $`F(v)`$ around $`v=0`$ and find $`F(v)F(0)+av^2+𝒪(v^4)`$, where $`F(0)=\mathrm{ln}[4/\pi ]`$ and $`a`$ is a constant. Thus, $`I_d=\mathrm{ln}[4/\pi ]\mathrm{ln}[\mathrm{\Lambda }_xL_x]`$ plus terms that do not diverge with system size. To make the argument of the $`\mathrm{ln}`$ term in $`I_d`$ match the $`\mathrm{ln}^2`$ term in $`I_c`$, we add and subtract the constant $`\mathrm{ln}[4/\pi ]\mathrm{ln}[2\mu _y\mathrm{\Lambda }_y/\mathrm{\Lambda }_x]`$ to find
$$I_d=\mathrm{ln}[4/\pi ]\mathrm{ln}[2\mu _y\mathrm{\Lambda }_yL_x].$$
(2.19)
We then add $`I_c`$ and $`I_d`$ and find the following expression for the displacement fluctuations in the limit $`L_z\mathrm{}`$ and $`L_xL_y`$:
$$(u_z^n)^2=l_u^2\mathrm{ln}^2\left[8L_x/x^{}\right].$$
(2.20)
## C Calculation of $`g^{(1)}(na)`$
In this appendix, we calculate the divergent part of the position correlation function $`g^{(1)}(na)`$ in (3.30). We perform the $`q_z`$ integration, switch to dimensionless variables, and find
$`g^{(1)}(na)=2l_u^2\left[{\displaystyle _{x^{}L_x^1}^{\mathrm{\Lambda }_xx^{}}}{\displaystyle \frac{dt}{t}}S_n(t)+\overline{A}_n(\mathrm{\Lambda }_x,\mathrm{\Lambda }_z)\right],`$ (3.1)
where $`S_n(t)`$ was defined previously in (3.31) and
$`\overline{A}_n(\mathrm{\Lambda }_x,\mathrm{\Lambda }_z)`$ $`=`$ $`{\displaystyle _0^{\mathrm{\Lambda }_xx^{}}}{\displaystyle \frac{dt}{t}}{\displaystyle _0^\pi }𝑑u{\displaystyle \frac{1\mathrm{cos}nu}{\sqrt{t^2+u^2p(u)}}}`$ (3.2)
$`\times `$ $`\left(1{\displaystyle \frac{2}{\pi }}\mathrm{arctan}\left[{\displaystyle \frac{\mathrm{\Lambda }_zz^{}}{t\sqrt{t^2+u^2p(u)}}}\right]\right).`$ (3.3)
Note that the integral defining $`\overline{A}_n(\mathrm{\Lambda }_x,\mathrm{\Lambda }_z)`$ does not have an infrared divergence as $`t0`$ ($`L_x\mathrm{}`$), and thus it is a well-defined number. In the final step, we isolate the $`\mathrm{ln}L_x`$ divergence in the first term of (3.1) to find
$$g^{(1)}(na)=2l_u^2S_n(0)\mathrm{ln}\left[A_n(\mathrm{\Lambda }_x,\mathrm{\Lambda }_z)L_x/x^{}\right],$$
(3.4)
where
$$A_n(\mathrm{\Lambda }_x,\mathrm{\Lambda }_z)=\mathrm{exp}[c_n^{(1)}+c_n^{(2)}(\mathrm{\Lambda }_x)+c_n^{(3)}(\mathrm{\Lambda }_x,\mathrm{\Lambda }_z)]$$
(3.5)
depends on the layer index $`n`$ and the ultraviolet cutoffs with
$`c_n^{(1)}`$ $`=`$ $`{\displaystyle _0^1}{\displaystyle \frac{dt}{t}}\left[{\displaystyle \frac{S_n(t)}{S_n(0)}}1\right],`$ (3.6)
$`c_n^{(2)}(\mathrm{\Lambda }_x)`$ $`=`$ $`{\displaystyle _1^{\mathrm{\Lambda }_xx^{}}}{\displaystyle \frac{dt}{t}}{\displaystyle \frac{S_n(t)}{S_n(0)}}`$ (3.7)
$`c_n^{(3)}(\mathrm{\Lambda }_x,\mathrm{\Lambda }_z)`$ $`=`$ $`{\displaystyle \frac{\overline{A}_n(\mathrm{\Lambda }_x,\mathrm{\Lambda }_z)}{S_n(0)}}.`$ (3.8)
Note that $`A_n(\mathrm{\Lambda }_x,\mathrm{\Lambda }_z)`$ is well-defined in the limit $`\mathrm{\Lambda }_{x,z}\mathrm{}`$.
## D Calculation $`g^{(2)}(𝐫)`$
In this Appendix, we evaluate the SC position correlation function
$$g^{(2)}(𝐫)=\frac{1}{2}[u_z^0(𝐫)u_z^0(0)]^2$$
(4.1)
between two DNA strands located in layer $`n=0`$ and separated by $`𝐫`$ in the $`xz`$ plane. For general separations, $`g^{(2)}(𝐫)`$ cannot be expressed in closed form. The aim of this Appendix is to calculate $`g^{(2)}(𝐫)`$ along the special directions $`z=0`$, $`xx^{}`$ and $`x=0`$, $`zz^{}`$.
### 1 Large $`x`$, Small $`z`$ Limit
The following expression for $`g^{(2)}(x,0)`$ is obtained by setting $`z`$ to zero in (3.21):
$$g^{(2)}(x,0)=T\frac{d^3q}{(2\pi )^3}\frac{1\mathrm{cos}(q_xx)}{Bq_z^2+Kq_x^4+K_yq_x^2q_y^2p(q_ya)}.$$
(4.2)
The first step in the derivation of $`g^{(2)}(x,0)`$ is to perform the integration over $`q_z`$ with $`\mathrm{\Lambda }_z\mathrm{}`$. The $`q_z`$ integration yields
$$g^{(2)}(x,0)=\frac{T}{2\pi ^2\sqrt{BK_y}}I(x,\mathrm{\Lambda }_x),$$
(4.3)
where
$`I(x,\mathrm{\Lambda }_x)`$ $`=`$ $`{\displaystyle _0^{\mathrm{\Lambda }_x}}𝑑q_x{\displaystyle \frac{1\mathrm{cos}(q_xx)}{q_x}}{\displaystyle _0^{\pi /x^{}}}{\displaystyle \frac{dq_y}{\sqrt{q_x^2+q_y^2p(q_yx^{})}}}.`$ (4.4)
We then decompose $`I(x,\mathrm{\Lambda }_x)I_c(x,\mathrm{\Lambda }_x)+I_d(x,\mathrm{\Lambda }_x)`$ into continuum and discrete contributions as we did previously in Appendix B, where
$`I_c(x,\mathrm{\Lambda }_x)`$ $`=`$ $`{\displaystyle _0^{\mathrm{\Lambda }_x}}𝑑q_x{\displaystyle \frac{1\mathrm{cos}(q_xx)}{q_x}}{\displaystyle _0^{\pi /x^{}}}{\displaystyle \frac{dq_y}{\sqrt{q_x^2+q_y^2}}}`$ (4.5)
$`I_d(x,\mathrm{\Lambda }_x)`$ $`=`$ $`{\displaystyle _0^{\mathrm{\Lambda }_x}}𝑑q_x{\displaystyle \frac{1\mathrm{cos}(q_xx)}{q_x}}{\displaystyle _0^{\pi /x^{}}}𝑑q_y`$ (4.7)
$`\left[{\displaystyle \frac{1}{\sqrt{q_x^2+q_y^2p(q_yx^{})}}}{\displaystyle \frac{1}{\sqrt{q_x^2+q_y^2}}}\right].`$
Since the $`\mathrm{\Lambda }_x\mathrm{}`$ limit is well-defined, we calculate $`I_c(x)I_c(x,\mathrm{})`$ and $`I_d(x)I_d(x,\mathrm{})`$ and drop terms that depend on the finite ultraviolet cutoff.
To calculate the continuum contribution, we first set $`q_x=uq_y`$ and then $`v=q_yx`$. These changes of variables yield
$$I_c(x)=_0^{\pi x/x^{}}\frac{dv}{v}K(v),$$
(4.8)
where
$$K(v)=_0^{\mathrm{}}𝑑u\frac{1\mathrm{cos}(uv)}{u\sqrt{1+u^2}}.$$
(4.9)
The strategy for calculating the $`\mathrm{ln}^2x`$ term in $`I_c(x)`$ is to isolate the part of $`K(v)`$ that scales as $`\mathrm{ln}v`$ for large $`v`$. To this end, we write
$`K(v)`$ $`=`$ $`{\displaystyle _0^1}{\displaystyle \frac{du}{u}}[1\mathrm{cos}(uv)]+{\displaystyle _1^{\mathrm{}}}{\displaystyle \frac{du}{u}}{\displaystyle \frac{1\mathrm{cos}(uv)}{\sqrt{1+u^2}}}`$ (4.10)
$`+`$ $`{\displaystyle _0^1}{\displaystyle \frac{du}{u}}\left[1\mathrm{cos}(uv)\right]\left[{\displaystyle \frac{1}{\sqrt{1+u^2}}}1\right].`$ (4.11)
It is obvious that only the first term has the correct scaling; the remaining terms in $`K(v)`$ are then separated into constants and functions of $`v`$ that are well-behaved either as $`v0`$ or $`v\mathrm{}`$. This partitioning leads to
$$K(v)=\mathrm{ln}(Bv)+\stackrel{~}{K}(v),$$
(4.12)
where $`B=2e^\gamma `$, $`\gamma `$ is Euler’s constant, and
$`\stackrel{~}{K}(v)`$ $`=`$ $`{\displaystyle _v^{\mathrm{}}}𝑑u{\displaystyle \frac{\mathrm{cos}u}{u}}{\displaystyle _1^{\mathrm{}}}{\displaystyle \frac{du}{u}}{\displaystyle \frac{\mathrm{cos}(uv)}{\sqrt{1+u^2}}}`$ (4.13)
$``$ $`{\displaystyle _0^1}𝑑u{\displaystyle \frac{\mathrm{cos}(uv)}{u}}\left[{\displaystyle \frac{1}{\sqrt{1+u^2}}}1\right]`$ (4.14)
scales as $`1/v^2`$ for large $`v`$.
We then insert $`K(v)`$ into (4.8) and break the integral over $`v`$ into small- and large-$`v`$ parts to obtain
$`I_c(x)`$ $`=`$ $`{\displaystyle _0^1}{\displaystyle \frac{dv}{v}}K(v)+{\displaystyle _1^{\pi x/x^{}}}{\displaystyle \frac{dv}{v}}\mathrm{ln}[Bv]`$ (4.15)
$`+`$ $`{\displaystyle _1^{\pi x/x^{}}}{\displaystyle \frac{dv}{v}}\stackrel{~}{K}(v).`$ (4.16)
Next, we evaluate the integral over $`v`$ in the second term, collect constants, and find
$$I_c(x)=\frac{1}{2}\mathrm{ln}^2[2e^\gamma \pi x/x^{}]+A_x,$$
(4.17)
where
$$A_x=\frac{1}{2}\mathrm{ln}^2[2e^\gamma ]+_0^1\frac{dv}{v}K(v)+_1^{\mathrm{}}\frac{dv}{v}\stackrel{~}{K}(v).$$
(4.18)
The second and third terms in $`A_x`$ are finite since $`K(v)`$ scales as $`v^2`$ for small $`v`$ in the former and there is phase cancellation from the $`\mathrm{cos}(uv)`$ factor at large $`v`$ in the later.
We will now calculate the discrete contribution to $`g^{(2)}(x,0)`$. The first step is to rewrite $`I_d(x)`$ in dimensionless form:
$$I_d(x)=_0^{\mathrm{\Lambda }_xx^{}}\frac{dv}{v}[1\mathrm{cos}(vx/x^{})]F(v),$$
(4.19)
where $`F(v)`$ was defined previously in (2.18). We next break the integral over $`v`$ into small- and large-$`v`$ parts and take the $`xx^{}`$ and $`\mathrm{\Lambda }_x\mathrm{}`$ limits to obtain
$`I_d(x)`$ $`=`$ $`F(0){\displaystyle _0^1}{\displaystyle \frac{dv}{v}}[1\mathrm{cos}(vx/x^{})]+{\displaystyle _1^{\mathrm{}}}{\displaystyle \frac{dv}{v}}F(v)`$ (4.20)
$`+`$ $`{\displaystyle _0^1}{\displaystyle \frac{dv}{v}}[F(v)F(0)].`$ (4.21)
Note that taking the $`xx^{}`$ limit removed the $`\mathrm{cos}(vx/x^{})`$ terms from the last two terms in (4.20) due to phase cancellations. It is again obvious that the first term scales logarithmically with $`x/x^{}`$, and thus
$$I_d(x)=\mathrm{ln}[4/\pi ]\mathrm{ln}\left[e^\gamma \frac{x}{x^{}}\right]+B_x,$$
(4.22)
where
$$B_x=_0^1\frac{dv}{v}[F(v)F(0)]+_1^{\mathrm{}}\frac{dv}{v}F(v)$$
(4.23)
is a constant. The last step in the calculation of $`g(x,0)`$ is to add the continuum and discrete terms, $`I_c(x)`$ and $`I_d(x)`$. The final result is
$$g^{(2)}(x,0)=l_u^2\left(\mathrm{ln}^2\left[8e^\gamma \frac{x}{x^{}}\right]+C_x\right),$$
(4.24)
where $`C_x=2(A_x+B_x)\mathrm{ln}^2[4/\pi ]2\mathrm{ln}[4/\pi ]\mathrm{ln}[2\pi ]`$. $`A_x`$ and $`B_x`$ are computed later in this appendix, see Eqs. (4.46) and (4.48). By using these equations, we eventually find that $`C_x`$ in Eq. (4.24) vanishes; $`C_x=0`$.
### 2 Large $`z`$, Small $`x`$ Limit
The calculation of $`g^{(2)}(0,z)`$ is similar to the calculation of $`g^{(2)}(x,0)`$. The expression for $`g^{(2)}(0,z)`$ is obtained by setting $`x`$ to zero in (3.21). The first step in the calculation is to perform the integration over $`q_z`$ with $`\mathrm{\Lambda }_z\mathrm{}`$ which yields
$$g^{(2)}(0,z)=\frac{T}{2\pi ^2\sqrt{BK_y}}I(z,\mathrm{\Lambda }_x),$$
(4.25)
where
$`I(z,\mathrm{\Lambda }_x)`$ $`=`$ $`{\displaystyle _0^{\mathrm{\Lambda }_x}}{\displaystyle \frac{dq_x}{q_x}}{\displaystyle _0^{\pi /x^{}}}𝑑q_y{\displaystyle \frac{1e^{z\lambda q_x\sqrt{q_x^2+q_y^2p(q_yx^{})}}}{\sqrt{q_x^2+q_y^2p(q_yx^{})}}}.`$ (4.26)
In what follows, we set $`\mathrm{\Lambda }_x\mathrm{}`$, drop terms that depend on the finite ultraviolet cutoff, and define $`I(z)I(z,\mathrm{})`$. The second step is to change variables to $`u=q_x/q_y`$ and $`v=\lambda zq_y^2`$ and decompose $`I(z)I_c(z)+I_d(z)`$ into continuum and discrete terms, where
$`I_c(z)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _0^\tau }{\displaystyle \frac{dv}{v}}{\displaystyle _0^{\mathrm{}}}𝑑u{\displaystyle \frac{1e^{vu\sqrt{1+u^2}}}{u\sqrt{1+u^2}}}`$ (4.27)
$`I_d(z)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dv}{v}}\left[F(v,0)F(v,vz/z^{})\right],`$ (4.28)
with $`\tau =\pi ^2z/z^{}`$, $`z^{}=a^2/\mu _y^2\lambda `$,
$$F(v,\tau )=_0^\pi 𝑑u\left[\frac{e^{\tau \sqrt{v^2+u^2p(u)}}}{\sqrt{v^2+u^2p(u)}}\frac{e^{\tau \sqrt{v^2+u^2}}}{\sqrt{v^2+u^2}}\right],$$
(4.29)
and $`F(v,0)`$ is equivalent to $`F(v)`$ defined in Eq. (2.18).
We first focus on the continuum contribution to $`g^{(2)}(z,0)`$. The integral over $`v`$ in $`I_c(z)`$ can be broken into small- and large-$`v`$ parts,
$$I_c(z)=\frac{1}{2}_0^1\frac{dv}{v}J(v)+\frac{1}{2}_1^{\pi ^2z/z^{}}\frac{dv}{v}J(v),$$
(4.30)
where
$$J(v)=_0^{\mathrm{}}𝑑u\frac{1e^{vu\sqrt{1+u^2}}}{u\sqrt{1+u^2}}.$$
(4.31)
The strategy is to extract the part of $`J(v)`$ that scales as $`\mathrm{ln}v`$ for large $`v`$. If $`J(v)\mathrm{ln}v`$ for large $`v`$, $`I_c(z)`$ will scale as $`\mathrm{ln}^2[z/z^{}]`$ as expected. Note that $`J(v)`$ scales as $`v^2`$ for small $`v`$, and thus the first term in (4.30) is a finite constant. After some algebra, we find
$$J(v)=\mathrm{ln}[Dv]+\stackrel{~}{J}(v),$$
(4.32)
where $`D=2e^\gamma `$,
$`\stackrel{~}{J}(v)`$ $`=`$ $`{\displaystyle _{\sqrt{2}v}^{\mathrm{}}}𝑑u{\displaystyle \frac{e^u}{u}}{\displaystyle _1^{\mathrm{}}}𝑑u{\displaystyle \frac{e^{vu\sqrt{1+u^2}}}{u\sqrt{1+u^2}}}`$ (4.33)
$``$ $`{\displaystyle _0^1}𝑑u{\displaystyle \frac{1\mathrm{\Phi }(u)}{u\sqrt{1+u^2}}}e^{vu\sqrt{1+u^2}},`$ (4.34)
and
$$\mathrm{\Phi }(u)=\frac{1+2u^2}{\sqrt{1+u^2}}.$$
(4.35)
We can now insert the expression for $`J(v)`$ in (4.32) into (4.30) and obtain the continuum contribution
$$I_c(z)=\frac{1}{4}\mathrm{ln}^2\left[2e^\gamma \pi ^2\frac{z}{z^{}}\right]+A_z,$$
(4.36)
where
$`A_z`$ $`=`$ $`{\displaystyle \frac{1}{4}}\mathrm{ln}^2\left[2e^\gamma \right]+{\displaystyle \frac{1}{2}}{\displaystyle _0^1}{\displaystyle \frac{dv}{v}}J(v)`$ (4.37)
$`+`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _1^{\mathrm{}}}{\displaystyle \frac{dv}{v}}\stackrel{~}{J}(v)`$ (4.38)
is a constant. Note that $`\stackrel{~}{J}(v)`$ decays exponentially for large $`v`$, and thus the third term in $`A_z`$ is finite.
We now concentrate on the discrete contribution to $`g^{(2)}(z,0)`$. The integral over $`v`$ in $`I_d(z)I_d^{(1)}+I_d^{(2)}`$ can also be broken into small- and large-$`v`$ parts, where
$$I_d^{(1)}=_0^1\frac{dv}{v}\left[F(v,0)F(v,vz/z^{})\right]$$
(4.39)
and $`I_d^{(2)}`$ is an identical expression except the limits on the integral over $`v`$ run from one to infinity. To isolate the $`\mathrm{ln}z`$ term in $`I_d^{(1)}`$, we change variables to $`t=vz/z^{}`$ and take the $`zz^{}`$ limit. In the large $`z`$ limit, (4.39) becomes
$$I_d^{(1)}=F(0,0)\mathrm{ln}\left[\frac{z}{z^{}}\right]+\overline{B}_z,+_0^1\frac{dv}{v}[F(v,0)F(0,0)]$$
(4.40)
where $`F(0,0)=\mathrm{ln}[4/\pi ]`$ and
$$\overline{B}_z=_0^1\frac{dt}{t}\left[F(0,0)F(0,t)\right]_1^{\mathrm{}}\frac{dt}{t}F(0,t)$$
(4.41)
is a constant. The large-$`v`$ contribution to $`I_d`$,
$$I_d^{(2)}=_1^{\mathrm{}}\frac{dv}{v}F(v,0),$$
(4.42)
is simply a constant when $`zz^{}`$. We then collect the discrete contributions and find
$`I_d(z)`$ $`=`$ $`I_d^{(1)}+I_d^{(2)},`$ (4.43)
$`=`$ $`\mathrm{ln}\left[{\displaystyle \frac{4}{\pi }}\right]\mathrm{ln}\left[{\displaystyle \frac{z}{z^{}}}\right]+B_z,`$ (4.44)
where $`B_z=\overline{B}_z+B_x`$ with $`B_x`$ given by Eq. (4.23).
The last step in the calculation of $`g^{(2)}(0,z)`$ is to add $`I_c(z)`$ and $`I_d(z)`$. The final result is
$$g^{(2)}(0,z)=l_u^2\left(\frac{1}{2}\mathrm{ln}^2\left[32e^\gamma \frac{z}{z^{}}\right]+C_z\right),$$
(4.45)
where $`C_z=2(A_z+B_z)2\mathrm{ln}^2[4/\pi ]2\mathrm{ln}[4/\pi ]\mathrm{ln}[2e^\gamma \pi ^2]`$. $`A_z`$ and $`B_z`$ are calculated later in this appendix \[See Eqs. (4.49) and (4.51)\]. By using these equations, we find $`C_z=\pi ^2/8`$.
### 3 Calculation of $`A_x`$, $`B_x`$, $`A_z`$, and $`B_z`$
In Secs. 1 and 2 of this Appendix, we anticipated that the numerical constants $`A_x`$,$`B_x`$,$`A_z`$, and $`B_z`$ have the values:
$`A_x`$ $`=`$ $`{\displaystyle \frac{\pi ^2}{24}},`$ (4.46)
$`B_x`$ $`=`$ $`{\displaystyle \frac{\pi ^2}{24}}+4\mathrm{ln}^2(2)`$ (4.48)
$`[\mathrm{ln}(2)][\mathrm{ln}(\pi )]{\displaystyle \frac{\mathrm{ln}^2(\pi )}{2}},`$
$`A_z`$ $`=`$ $`{\displaystyle \frac{7\pi ^2}{48}},`$ (4.49)
$`B_z`$ $`=`$ $`{\displaystyle \frac{\pi ^2}{12}}+6\mathrm{ln}^2(2)+2\gamma \mathrm{ln}(2)`$ (4.51)
$`[\mathrm{ln}(2)][\mathrm{ln}(\pi )]\gamma \mathrm{ln}(\pi )\mathrm{ln}^2(\pi ).`$
In this section we outline calculations yielding Eqs. (4.46) to (4.51). We begin by deriving more explicit formulas for these constants. Thus, for $`A_x`$, we obtain, by Eqs. (4.10), (4.13), and (4.18),
$$A_x=\frac{1}{2}\mathrm{ln}^2(2e^\gamma )+A_x^{(1)}+A_x^{(2)}+\gamma A_x^{(3)},$$
(4.52)
with
$`A_x^{(1)}`$ $`=`$ $`{\displaystyle _0^1}𝑑v\mathrm{ln}(v){\displaystyle \frac{1\mathrm{cos}(v)}{v}}`$ (4.54)
$`+{\displaystyle _1^{\mathrm{}}}𝑑v\mathrm{ln}(v){\displaystyle \frac{\mathrm{cos}(v)}{v}},`$
$`A_x^{(2)}`$ $`=`$ $`{\displaystyle _0^1}{\displaystyle \frac{du}{u}}\left[{\displaystyle \frac{1}{\sqrt{1+u^2}}}1\right]\mathrm{ln}(u)`$ (4.56)
$`+{\displaystyle _1^{\mathrm{}}}{\displaystyle \frac{du}{u}}{\displaystyle \frac{1}{\sqrt{1+u^2}}}\mathrm{ln}(u),`$
and
$`A_x^{(3)}`$ $`=`$ $`{\displaystyle _0^1}{\displaystyle \frac{du}{u}}\left[{\displaystyle \frac{1}{\sqrt{1+u^2}}}1\right]`$ (4.58)
$`+{\displaystyle _1^{\mathrm{}}}{\displaystyle \frac{du}{u}}{\displaystyle \frac{1}{\sqrt{1+u^2}}}.`$
Likewise, for the constant $`A_z`$, we obtain, by Eqs. (4.31), (4.33), (4.35), and (4.37),
$`A_z`$ $`=`$ $`{\displaystyle \frac{1}{4}}\mathrm{ln}^2(2e^\gamma )`$ (4.61)
$`+{\displaystyle \frac{1}{2}}\left[A_z^{(1)}+{\displaystyle \frac{1}{8}}\mathrm{ln}^2(2)+{\displaystyle \frac{\gamma }{2}}\mathrm{ln}(2)\right]`$
$`+{\displaystyle \frac{1}{2}}\left[A_z^{(2)}+\gamma A_z^{(3)}\right],`$
where
$`A_z^{(1)}`$ $`=`$ $`{\displaystyle _0^1}𝑑x\mathrm{ln}(x){\displaystyle \frac{1e^x}{x}}`$ (4.63)
$`+{\displaystyle _1^{\mathrm{}}}𝑑x\mathrm{ln}(x){\displaystyle \frac{e^x}{x}},`$
$`A_z^{(2)}`$ $`=`$ $`{\displaystyle _0^1}𝑑u{\displaystyle \frac{1\mathrm{\Phi }(u)}{u\sqrt{1+u^2}}}\mathrm{ln}(u\sqrt{1+u^2})`$ (4.65)
$`+{\displaystyle _1^{\mathrm{}}}𝑑u{\displaystyle \frac{1}{u\sqrt{1+u^2}}}\mathrm{ln}(u\sqrt{1+u^2}),`$
$`A_z^{(3)}`$ $`=`$ $`{\displaystyle _0^1}𝑑u{\displaystyle \frac{1\mathrm{\Phi }(u)}{u\sqrt{1+u^2}}}`$ (4.67)
$`+{\displaystyle _1^{\mathrm{}}}𝑑u{\displaystyle \frac{1}{u\sqrt{1+u^2}}}`$
with $`\mathrm{\Phi }(u)`$ defined in Eq. (4.35). We proceed to compute the integrals in Eqs. (4.52) to (4.67). Integrals (4.58) and (4.67) can be done by elementary integration methods yielding
$$A_x^{(3)}=\mathrm{ln}(2),A_z^{(3)}=\frac{\mathrm{ln}(2)}{2}.$$
(4.68)
Other integrals are not elementary, and we outline their calculation in the following. Thus, we compute the integrals $`A_x^{(1)}`$ and $`A_z^{(1)}`$ first by showing that
$$A_x^{(1)}=A_z^{(1)}\frac{\pi ^2}{8}.$$
(4.69)
and, next, by showing that
$$A_z^{(1)}=\frac{1}{2}\left(\frac{d^2\mathrm{\Gamma }}{dz^2}\right)_{z=1},$$
(4.70)
where $`\mathrm{\Gamma }(z)`$ signifies the gamma function.
To obtain Eq. (4.69), consider the complex function $`f(z)=e^z\mathrm{ln}(z)/z`$. In the complex $`z=x+iy`$ plane, $`f(z)`$ is analytic inside the contour $`C`$ made of the following four segments: $`C_1:[z=x;ϵ<x<R]`$, $`C_2:[z=Re^{i\theta };0<\theta <\pi /2]`$, $`C_3:[z=iy;R>y>0]`$, and $`C_4:[z=ϵe^{i\theta };\pi /2>\theta >0]`$. By applying the Cauchy residue theorem to the above $`f(z)`$ along the contour $`C=C_1+C_2+C_3+C_4`$ in the limit $`ϵ0`$ and $`R\mathrm{}`$, we directly obtain the relation (4.69) between $`A_x^{(1)}`$ and $`A_z^{(1)}`$. Next, we demonstrate Eq. (4.70) by differentiating the standard integral representation of the gamma function. This yields
$$\left(\frac{d\mathrm{\Gamma }}{dz}\right)_{z=1}=_0^{\mathrm{}}𝑑x\mathrm{ln}(x)e^x=\gamma ,$$
(4.71)
$$\left(\frac{d^2\mathrm{\Gamma }}{dz^2}\right)_{z=1}=_0^{\mathrm{}}𝑑x\mathrm{ln}^2(x)e^x.$$
(4.72)
We then rewrite Eq. (4.72) as,
$`\left({\displaystyle \frac{d^2\mathrm{\Gamma }}{dz^2}}\right)_{z=1}`$ $`=`$ $`{\displaystyle _0^1}𝑑x(e^x1)\mathrm{ln}^2(x)`$ (4.74)
$`+{\displaystyle _0^1}𝑑x\mathrm{ln}^2(x)+{\displaystyle _1^{\mathrm{}}}𝑑xe^x\mathrm{ln}^2(x)`$
and integrate by parts the first integral \[by writing $`(e^x1)dx=d(1e^xx)`$,etc.\] as well as the last integral \[by writing $`e^xdx=d(e^x)`$, etc.\]. After few elementary integrations, this yields our Eq. (4.70) which, combined with Eq. (4.69), yields values of the integrals $`A_x^{(1)}`$ and $`A_z^{(1)}`$. To complete this calculation we need also the value of the second derivative of $`\mathrm{\Gamma }(z)`$ at $`z=1`$ that enters Eq. (4.70). To compute it, we use a relation from the theory of the gamma function:
$$\frac{d\mathrm{\Gamma }(z)}{dz}=\mathrm{\Gamma }(z)\mathrm{\Psi }(z),$$
(4.75)
with
$`\mathrm{\Psi }(z)`$ $`=`$ $`\gamma +\left(1{\displaystyle \frac{1}{z}}\right)+\left({\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{z+1}}\right)`$ (4.77)
$`+\left({\displaystyle \frac{1}{3}}{\displaystyle \frac{1}{z+2}}\right)+\mathrm{};`$
$`\mathrm{\Gamma }(1)=1`$, $`(d\mathrm{\Gamma }/dz)_{z=1}=\mathrm{\Psi }(1)=\gamma `$. By differentiating Eq. (4.75),
$`\left({\displaystyle \frac{d^2\mathrm{\Gamma }}{dz^2}}\right)_{z=1}`$ $`=`$ $`\left({\displaystyle \frac{d\mathrm{\Gamma }}{dz}}\right)_{z=1}\mathrm{\Psi }(1)+\mathrm{\Gamma }(1)\left({\displaystyle \frac{d\mathrm{\Psi }}{dz}}\right)_{z=1}`$ (4.78)
$`=`$ $`\gamma ^2+\left({\displaystyle \frac{d\mathrm{\Psi }}{dz}}\right)_{z=1}.`$ (4.79)
The use of Eq. (4.77) to compute $`d\mathrm{\Psi }(z)/dz`$ at $`z=1`$ yields
$`\left({\displaystyle \frac{d^2\mathrm{\Gamma }}{dz^2}}\right)_{z=1}`$ $`=`$ $`\gamma ^2+1+{\displaystyle \frac{1}{2^2}}+{\displaystyle \frac{1}{3^2}}+{\displaystyle \frac{1}{4^2}}+\mathrm{}`$ (4.80)
$`=`$ $`\gamma ^2+{\displaystyle \frac{\pi ^2}{6}}.`$ (4.81)
By Eqs. (4.81), (4.69), and (4.70), we finally obtain
$$A_x^{(1)}=\frac{\gamma ^2}{2}\frac{\pi ^2}{24},A_z^{(1)}=\frac{\gamma ^2}{2}+\frac{\pi ^2}{12}.$$
(4.82)
Next, we proceed to compute the integral $`A_x^{(2)}`$ in Eq. (4.56). For this purpose, consider the integrals
$`A_x^{(2)}(t)`$ $`=`$ $`{\displaystyle _0^t}{\displaystyle \frac{du}{u}}\left[{\displaystyle \frac{1}{\sqrt{1+u^2}}}1\right]\mathrm{ln}(u)`$ (4.84)
$`+{\displaystyle _t^{\mathrm{}}}{\displaystyle \frac{du}{u}}{\displaystyle \frac{1}{\sqrt{1+u^2}}}\mathrm{ln}(u)`$
$`A_x^{(3)}(t)`$ $`=`$ $`{\displaystyle _0^t}{\displaystyle \frac{du}{u}}\left[{\displaystyle \frac{1}{\sqrt{1+u^2}}}1\right]`$ (4.86)
$`+{\displaystyle _t^{\mathrm{}}}{\displaystyle \frac{du}{u}}{\displaystyle \frac{1}{\sqrt{1+u^2}}}`$
For $`t=1`$, $`A_x^{(2)}(t=1)=A_x^{(2)}`$ is the desired integral, whereas $`A_x^{(3)}(t=1)=A_x^{(3)}=\mathrm{ln}(2)`$, by (4.68). Further, by Eqs. (4.84) and (4.86), $`dA_x^{(2)}(t)/dt=\mathrm{ln}(t)/t`$, and $`dA_x^{(3)}(t)/dt=1/t`$. By integrating these relations over $`t`$,
$`A_x^{(2)}`$ $`=`$ $`A_x^{(2)}(t)+{\displaystyle \frac{\mathrm{ln}^2(t)}{2}},`$ (4.87)
$`A_x^{(3)}`$ $`=`$ $`A_x^{(3)}(t)+\mathrm{ln}(t),`$ (4.88)
for any $`t>0`$. By Eqs. (4.84) to (4.88), with $`t=ϵ0`$,
$`A_x^{(2)}`$ $`=`$ $`\underset{ϵ0}{lim}\left[{\displaystyle _ϵ^{\mathrm{}}}{\displaystyle \frac{du}{u}}{\displaystyle \frac{1}{\sqrt{1+u^2}}}\mathrm{ln}(u)+{\displaystyle \frac{\mathrm{ln}^2(ϵ)}{2}}\right]`$ (4.89)
$`A_x^{(3)}`$ $`=`$ $`\underset{ϵ0}{lim}\left[{\displaystyle _ϵ^{\mathrm{}}}{\displaystyle \frac{du}{u}}{\displaystyle \frac{1}{\sqrt{1+u^2}}}+\mathrm{ln}(ϵ)\right]`$ (4.90)
To proceed, it is useful to form the difference
$$\mathrm{\Delta }=A_x^{(2)}\frac{1}{2}[A_x^{(3)}]^2.$$
(4.91)
Using here Eqs. (4.89) and (4.90), we find
$$\mathrm{\Delta }=_0^{\mathrm{}}𝑑u\frac{1}{u\sqrt{1+u^2}}_0^u\frac{dv}{v}\left[1\frac{1}{\sqrt{1+v^2}}\right].$$
(4.92)
The integral over $`v`$ here can be done by changing $`vt`$ with $`t=\mathrm{ln}[(\sqrt{1+v^2}+1)/2]`$. This reduces Eq. (4.92) to
$$\mathrm{\Delta }=_0^{\mathrm{}}𝑑u\frac{1}{u\sqrt{1+u^2}}\mathrm{ln}[\frac{\sqrt{1+v^2}+1}{2}].$$
(4.93)
By changing variables via $`ut`$ with $`t=\mathrm{ln}[(\sqrt{1+u^2}+1)/2]`$, we eventually obtain $`\mathrm{\Delta }`$ in a form involving a well known integral,
$$\mathrm{\Delta }=\frac{1}{2}_0^{\mathrm{}}𝑑t\frac{t}{e^t1}=\frac{\pi ^2}{12}.$$
(4.94)
By Eqs. (4.94), (4.91), and (4.68), we obtain
$$A_x^{(2)}=\frac{\pi ^2}{12}+\frac{\mathrm{ln}^2(2)}{2}.$$
(4.95)
Equations (4.52), (4.68), (4.82), and (4.95) yield our final result for the constant $`A_x`$ anticipated in Eq. (4.46).
Next, we compute the integral $`A_z^{(2)}`$ in Eq. (4.65). For this purpose, consider the integral
$`A_z^{(2)}(t)`$ $`=`$ $`{\displaystyle _0^t}𝑑u{\displaystyle \frac{1\mathrm{\Phi }(u)}{u\sqrt{1+u^2}}}\mathrm{ln}(u\sqrt{1+u^2})`$ (4.97)
$`+{\displaystyle _t^{\mathrm{}}}𝑑u{\displaystyle \frac{1}{u\sqrt{1+u^2}}}\mathrm{ln}(u\sqrt{1+u^2}),`$
For $`t=1`$, $`A_z^{(2)}(t=1)=A_z^{(2)}`$ is the desired integral. It is easy to show, along the lines we used to derive Eqs. (4.87) and (4.88), that
$$A_z^{(2)}=A_z^{(2)}(t)+\frac{[\mathrm{ln}(t\sqrt{1+t^2})]^2}{2}\frac{[\mathrm{ln}\sqrt{2}]^2}{2},$$
(4.98)
for any $`t>0`$. Thus, by Eqs. (4.97) and (4.98), with $`t=ϵ0`$,
$`A_z^{(2)}`$ $`=`$ $`\underset{ϵ0}{lim}\left[{\displaystyle _ϵ^{\mathrm{}}}𝑑u{\displaystyle \frac{1}{u\sqrt{1+u^2}}}\mathrm{ln}(u\sqrt{1+u^2})+{\displaystyle \frac{\mathrm{ln}^2(ϵ)}{2}}\right]`$ (4.100)
$`{\displaystyle \frac{\mathrm{ln}^2(2)}{8}}.`$
Next, by Eqs. (4.100) and (4.89),
$`A_z^{(2)}A_x^{(2)}=`$ (4.101)
$`{\displaystyle _0^{\mathrm{}}}𝑑u{\displaystyle \frac{1}{u\sqrt{1+u^2}}}\mathrm{ln}(\sqrt{1+u^2}){\displaystyle \frac{\mathrm{ln}^2(2)}{8}}.`$ (4.102)
By the change of variables $`ut`$, with $`u=\sqrt{e^{2t}1}`$, we find $`A_z^{(2)}A_x^{(2)}=S_1+S_2\mathrm{ln}^2(2)/8`$, where $`S_1`$ and $`S_2`$ are two well known integrals, $`S_1=_0^{\mathrm{}}𝑑tt/(e^t+1)=\pi ^2/12`$, and $`S_2=_0^{\mathrm{}}𝑑tt/(e^{2t}1)=\frac{1}{4}_0^{\mathrm{}}𝑑xx/(e^x1)=\frac{1}{4}(\pi ^2/6)`$. Thus,
$$A_z^{(2)}A_x^{(2)}=\frac{\pi ^2}{8}\frac{\mathrm{ln}^2(2)}{8}.$$
(4.103)
By Eqs. (4.103) and (4.95),
$$A_z^{(2)}=\frac{5\pi ^2}{24}+\frac{3\mathrm{ln}^2(2)}{8}.$$
(4.104)
Equations (4.61), (4.68), (4.82), and (4.104) yield our final result for the constant $`A_z`$ anticipated in Eq. (4.49).
We now proceed to discuss calculations yielding the values of the numerical constants $`B_x`$ and $`B_z`$ quoted in Eqs. (4.48) and (4.51). For $`B_x`$, by (4.23) and (2.18), we obtain
$$B_x=_0^\pi 𝑑w\left[\frac{\left(A_x^{(3)}(t)\right)_{t=1/w\sqrt{p(w)}}}{w\sqrt{p(w)}}\frac{\left(A_x^{(3)}(t)\right)_{t=1/w}}{w}\right].$$
(4.105)
By using (4.88), i.e., $`A_x^{(3)}(t)=A_x^{(3)}\mathrm{ln}(t)=\mathrm{ln}(2)\mathrm{ln}(t)`$, we find
$`B_x`$ $`=`$ $`[{\displaystyle _ϵ^\pi }dw{\displaystyle \frac{\mathrm{ln}(2)+\mathrm{ln}(w\sqrt{p(w)})}{w\sqrt{p(w)}}}`$ (4.107)
$`{\displaystyle _ϵ^\pi }dw{\displaystyle \frac{\mathrm{ln}(2)+\mathrm{ln}(w)}{w}}]_{ϵ0}.`$
Next, in the first integral above, we change from $`w`$ to $`z=w\sqrt{p(w)}=2\mathrm{sin}(w/2)`$, whereas in the second integral we change the variable $`w`$ into $`z`$. After rearranging the expression thus obtained, we find
$`B_x`$ $`=`$ $`{\displaystyle _0^2}𝑑z\left[{\displaystyle \frac{1}{\sqrt{1(z/2)^2}}}1\right]{\displaystyle \frac{\mathrm{ln}(2)+\mathrm{ln}(z)}{z}}`$ (4.109)
$`{\displaystyle _2^\pi }𝑑z{\displaystyle \frac{\mathrm{ln}(2)+\mathrm{ln}(z)}{z}}.`$
For convenience, we change variables via $`zx=z/2`$ and thus obtain
$`B_x`$ $`=`$ $`B^{(1)}+2\mathrm{ln}(2)B^{(2)}{\displaystyle \frac{\mathrm{ln}^2(\pi )}{2}}`$ (4.111)
$`[\mathrm{ln}(2)][\mathrm{ln}(\pi )]+{\displaystyle \frac{3\mathrm{ln}^2(2)}{2}},`$
with
$$B^{(1)}=_0^1𝑑x\left[\frac{1}{\sqrt{1x^2}}1\right]\mathrm{ln}(x),$$
(4.112)
and
$$B^{(2)}=_0^1𝑑x\left[\frac{1}{\sqrt{1x^2}}1\right].$$
(4.113)
A similar formula can be derived for the constant $`B_z=\overline{B}_z+B_x`$ introduced in Sec. D 2 of this Appendix. By Eqs. (4.29) and (4.41)), we find for $`\overline{B}_z`$
$$\overline{B}_z=\left[_ϵ^\pi 𝑑u\frac{\gamma +\mathrm{ln}(u\sqrt{p(u)})}{u\sqrt{p(u)}}_ϵ^\pi 𝑑u\frac{\gamma +\mathrm{ln}(u)}{u}\right]_{ϵ0}.$$
(4.114)
By treating Eq. (4.114) in exactly the same way we treated above Eq. (4.107), we eventually find
$`\overline{B}_z`$ $`=`$ $`B^{(1)}+[\mathrm{ln}(2)+\gamma ]B^{(2)}+{\displaystyle \frac{\mathrm{ln}^2(2)}{2}}`$ (4.116)
$`+\gamma \mathrm{ln}(2){\displaystyle \frac{\mathrm{ln}^2(\pi )}{2}}\gamma \mathrm{ln}(\pi ).`$
To complete our calculation, we need the integrals $`B^{(1)}`$ and $`B^{(2)}`$ in Eqs. (4.112) and (4.113). $`B^{(2)}`$, Eq. (4.113), can be calculated by elementary integration methods yielding
$$B^{(2)}=\mathrm{ln}(2).$$
(4.117)
On the other hand, the integral $`B^{(1)}`$, Eq. (4.112), is not elementary. As detailed below, we find
$$B^{(1)}=\frac{\pi ^2}{24}+\frac{\mathrm{ln}^2(2)}{2}.$$
(4.118)
The values of $`B_x`$ and $`B_z=\overline{B}_z+B_x`$ quoted in Eqs. (4.48) and (4.51), directly follow from Eqs. (4.111), (4.116), (4.117), and (4.118). It remains to outline the calculation yielding the value of $`B^{(1)}`$ in Eq. (4.118). It is obtained by relating the integral $`B^{(1)}`$, Eq. (4.112), to the integral $`A_x^{(2)}`$, Eq. (4.56), computed in Eq. (4.95). We find that
$$B^{(1)}=\frac{\pi ^2}{8}+A_x^{(2)},$$
(4.119)
To derive relation (4.119), we apply the Cauchy residue theorem to the complex function $`f(z)=\mathrm{ln}(z)/z\sqrt{1z^2}`$ along the same contour that has been used before to derive Eq. (4.69) \[see the text following Eq. (4.70)\]. This yields the relation
$`{\displaystyle _ϵ^1}{\displaystyle \frac{dx}{x}}{\displaystyle \frac{1}{\sqrt{1x^2}}}\mathrm{ln}(x)=`$ (4.120)
$`{\displaystyle \frac{\pi ^2}{8}}+{\displaystyle _ϵ^{\mathrm{}}}{\displaystyle \frac{du}{u}}{\displaystyle \frac{1}{\sqrt{1+u^2}}}\mathrm{ln}(u)+\varphi (ϵ),`$ (4.121)
where $`\varphi (ϵ)0`$ as $`ϵ0`$. Equation (4.121) is identical to
$`{\displaystyle _ϵ^1}{\displaystyle \frac{dx}{x}}\left[{\displaystyle \frac{1}{\sqrt{1x^2}}}1\right]\mathrm{ln}(x)=`$ (4.122)
$`{\displaystyle \frac{\pi ^2}{8}}+{\displaystyle _ϵ^{\mathrm{}}}{\displaystyle \frac{du}{u}}{\displaystyle \frac{1}{\sqrt{1+u^2}}}\mathrm{ln}(u)+{\displaystyle \frac{\mathrm{ln}^2(ϵ)}{2}}+\varphi (ϵ).`$ (4.123)
By recalling here Eq. (4.89) and taking the limit $`ϵ0`$ in Eq. (4.123), we eventually obtain the relation (4.119) used to compute $`B^{(1)}`$.
## E Calculation of $`g^{(3)}(𝐫,na)`$
In this appendix, we will outline the evaluation of Eq. (3.33) for $`g^{(3)}(𝐫,na)g^{(3)}(x,z,na)`$. We begin with the expression, Eq. (3.23) for $`g^{(3)}(𝐫,na)`$. We assume the continuum limit $`\mathrm{\Lambda }_z\mathrm{}`$ and integrate over $`q_z`$ to obtain
$`g^{(3)}(𝐫,na)`$ $`=`$ $`2l_u^2{\displaystyle _0^{\mathrm{\Lambda }_xx^{}}}{\displaystyle \frac{dt}{t}}{\displaystyle _0^\pi }𝑑u{\displaystyle \frac{1\mathrm{cos}(nu)}{t\sqrt{t^2+u^2p(u)}}}`$ (5.2)
$`\times \left[1\mathrm{cos}(tx/x^{})e^{(tz/z^{})\sqrt{t^2u^2+p(u)}}\right].`$
This expression can now be used to evaluate various limits.
### 1 $`xx^{}`$ and $`z=0`$
From Eq. (5.2), we have
$$g^{(3)}(x,0,na)=2l_u^2_0^{\mathrm{\Lambda }_xx^{}}\frac{dt}{t}S_n(t)[1\mathrm{cos}(tx/x^{})],$$
(5.3)
where $`S_n(t)`$, defined in Eq. (3.31), is proportional to $`1/t`$ for $`t1`$. Thus the integral in Eq. (5.3) is convergent at large $`t`$, and we can take the continuum limit $`\mathrm{\Lambda }_x\mathrm{}`$. Thus, we have
$`g^{(3)}(x,0,na)`$ $`=`$ $`2l_u^2{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dt}{t}}S_n(t)[1\mathrm{cos}(tx/x^{})]`$ (5.4)
$`=`$ $`2l_u^2{\displaystyle _0^1}{\displaystyle \frac{dt}{t}}[1\mathrm{cos}(tx/x^{})]S_n(t)`$ (5.6)
$`+2ł_u^2{\displaystyle _1^{\mathrm{}}}{\displaystyle \frac{dt}{t}}[1\mathrm{cos}(tx/x^{})]S_n(t)`$
$`=`$ $`2l_u^2[\mathrm{ln}(x/x^{})+\gamma ]S_n(0)`$ (5.9)
$`+2l_u^2{\displaystyle _0^1}{\displaystyle \frac{dt}{t}}[1\mathrm{cos}(tx/x^{})][S_n(t)S_n(0)]`$
$`+2l_u^2{\displaystyle _1^{\mathrm{}}}{\displaystyle \frac{dt}{t}}[1\mathrm{cos}(tx/x^{})]S_n(t)`$
The contributions from the $`\mathrm{cos}(tx/x^{})`$ terms in the integrals in Eq. (5.9) vanish when $`x/x^{}1`$, leaving
$$g^{(3)}(x,0,na)=2l_u^2S_n(0)\mathrm{ln}(D_nx/x^{}),$$
(5.10)
where
$$D_n=e^{\gamma +c_n^{(1)}+c_n^{(2)}(\mathrm{\Lambda }_x=\mathrm{})}$$
(5.11)
with $`c_n^{(1)}`$ and $`c_n^{(2)}(\mathrm{\Lambda }_x)`$ given, respectively, by Eqs. (3.6) and (3.7).
### 2 $`x=0`$, $`zz^{}`$
Again, we will simplify our discussion by taking the limit $`\mathrm{\Lambda }_x\mathrm{}`$ and integrating over $`q_x`$ in Eq. (3.33) to obtain
$`g^{(3)}(0,z,na)=2l_u^2{\displaystyle _0^{\mathrm{}}}𝑑t{\displaystyle _0^\pi }𝑑u{\displaystyle \frac{1\mathrm{cos}(nu)}{t\sqrt{t^2+u^2p(u)}}}`$ (5.12)
$`\times \left[1e^{(tz/z^{})\sqrt{t^2+u^2p(u)}}\right].`$ (5.13)
Changing variables via $`t=s\sqrt{u^2p(u)}`$, we obtain
$$g^{(3)}(0,z,na)=2l_u^2_0^\pi \frac{1\mathrm{cos}(nu)}{\sqrt{u^2p(u)}}J[(z/z^{})u^2p(u)],$$
(5.14)
where $`J(v)`$ is defined in Eq. (4.31). Next, using $`J(v)=\mathrm{ln}(2e^\gamma )+\stackrel{~}{J}(v)`$ \[Eq. (4.32)\], we obtain
$`g^{(3)}(0,z,na)`$ $`=`$ $`2l_u^2[\mathrm{ln}(z/z^{})+\gamma +\mathrm{ln}2]{\displaystyle _0^\pi }𝑑u{\displaystyle \frac{1\mathrm{cos}(nu)}{\sqrt{u^2p(u)}}}`$ (5.17)
$`+2l_u^2{\displaystyle _0^\pi }𝑑u{\displaystyle \frac{1\mathrm{cos}(nu)}{\sqrt{u^2p(u)}}}\mathrm{ln}[u^2p(u)]`$
$`+2l_u^2{\displaystyle _0^\pi }𝑑u{\displaystyle \frac{1\mathrm{cos}(nu)}{\sqrt{u^2p(u)}}}\stackrel{~}{J}[(z/z^{})u^2p(u)]`$
In the limit $`z/z^{}\mathrm{}`$, $`\stackrel{~}{J}[(z/z^{})u^2p(u)]`$ tends to zero for all $`u`$, and we obtain
$$g^{(3)}(0,z,na)=2l_u^2S_n(0)\mathrm{ln}(E_nz/z^{}),$$
(5.18)
where
$$E_u=2e^\gamma \mathrm{exp}\left[\frac{1}{S_n(0)}_0^\pi 𝑑u\frac{1\mathrm{cos}(nu)}{\sqrt{u^2p(u)}}\mathrm{ln}[u^2p(u)]\right]$$
(5.19)
## F Inter-plane Correlations when $`v^u0`$
In this appendix, we will derive Eq. (3.58) from the expression, Eq. (3.53) for $`S(𝐫,na)`$. We begin by looking at the linear term in $`V^u`$ in $`S(𝐫,a)`$:
$`S(𝐫,a)=V^u/2T{\displaystyle }d^2r_1e^{iq_0u_z^0(0)}`$ (6.1)
$`\times (e^{iq_0[u_z^0(𝐫_1)u_z^1(𝐫_1)]}e^{iq_0a_zu_y^1(𝐫_1)}+\mathrm{c}.\mathrm{c}.)e^{iq_0u_z^1(𝐫)}_{SC}`$ (6.2)
$`={\displaystyle \frac{V^u}{2T}}{\displaystyle }d^2r_1e^{iq_0[u_z^1(𝐫)u_z^1(𝐫_1)]}`$ (6.3)
$`\times e^{iq_0[u_z^0(𝐫_1)u_z^0(0)]}e^{iq_0a_zu_y^1(𝐫_1)}_{SC}`$ (6.4)
$`+{\displaystyle \frac{V^u}{2T}}{\displaystyle }d^2r_1e^{iq_0[u_z^1(𝐫)+u_z^1(𝐫_1)]}`$ (6.5)
$`\times e^{iq_0[u_z^0(𝐫_1)+u_z^0(0)]}\times e^{iq_0a_zu_y^1(𝐫_1)}_{SC},`$ (6.6)
where the subscript $`SC`$ indicates that the averages are to be evaluated with respect to the sliding columnar Hamiltonian, $`^{SC}`$ of Eq. (2.23). Since $`^{SC}`$ is quadratic in $`u_z^n(𝐫)`$ and $`u_y^n(𝐫)`$, the averages in this equation can be performed exactly, and $`S_2(𝐫,a)`$ can be expressed as an exponential of correlations functions of $`u_z^n(𝐫)`$ and $`_zu_y^n(𝐫)`$. The second term in Eq. (6.6) has terms in the exponential proportional to $`q_0^2[u_z^n(𝐫)]^2`$ for $`n=0,1`$, which diverge in thermodynamic limit and cause the exponential to vanish. Thus only the first term of Eq. (6.6) survives, and we have
$$S(𝐫,a)=\frac{V^u}{2T}d^2r_1e^{\mathrm{\Phi }(𝐫_1,𝐫)},$$
(6.7)
where
$`\mathrm{\Phi }(𝐫_1,𝐫)={\displaystyle \frac{1}{2}}q_0^2(\stackrel{~}{u}_z^1(𝐫,𝐫_1)+\stackrel{~}{u}_z^0(𝐫_1,0)+a_zu_y^1(𝐫_1)^2`$ (6.8)
$`=g^{(2)}(𝐫𝐫_1,0)+g^{(2)}(𝐫_1,0)`$ (6.9)
$`+g^{(2)}(𝐫,a)+g^{(2)}(0,a)g^{(2)}(𝐫𝐫_1,a)g^{(2)}(𝐫_1,a)`$ (6.10)
$`+{\displaystyle \frac{1}{2}}a^2[_zu_y^1(𝐫_1)]^2`$ (6.11)
$`+a_zu_y^1(𝐫_1)[\stackrel{~}{u}_z^1(𝐫,𝐫_1)\stackrel{~}{u}_z^0(𝐫_1,0)];`$ (6.12)
here, $`g^{(2)}(𝐫,na)`$ is defined in Eq. (3.26), and
$$\stackrel{~}{u}_z^n(𝐫,𝐫^{})=u_z^n(𝐫)u_z^n(𝐫^{}).$$
(6.13)
This expression is quite complex. It, however, simplifies considerably if we set $`K_y=0`$ and $`B_{uh}=0`$. Then $`g^{(2)}(𝐫,a)=0`$, and all cross terms in $`u_z^n(𝐫)`$ and $`u_y^m(𝐫^{})`$ vanish. In this case,
$$\mathrm{\Phi }(𝐫_1,𝐫)=g^{(2)}(𝐫𝐫_1)+g^{(2)}(𝐫_1)+W_y,$$
(6.14)
where $`g^{(2)}(𝐫)=g^{(2)}(𝐫,0)`$ \[See Eq. (3.26)\] and $`W_y=q_0^2a^2(_zu_y^n)^2/2`$.
The generalization of Eqs. (6.7) and (6.8) to $`n>1`$ is straightforward. The leading contribution to $`S(𝐫,na)`$ is
$$S(𝐫,na)=\left(\frac{V^u}{2T}\right)^nd^2r_1\mathrm{}d^2r_ne^{\mathrm{\Phi }(𝐫_1,\mathrm{},𝐫_n,𝐫)},$$
(6.15)
where
$$\mathrm{\Phi }=\frac{1}{2}\left(\underset{m=0}{\overset{n}{}}\stackrel{~}{u}_z^m(𝐫_m,𝐫_{m+1})+a\underset{m=1}{\overset{n}{}}_zu_y^m(𝐫_m)\right)^2_{SC},$$
(6.16)
where $`𝐫_0=0`$ and $`𝐫_{n+1}=𝐫`$. This function can be expressed in terms of the reduced correlation functions $`g^{(2)}(𝐫,n)`$ for $`u_z^n(𝐫)`$ and correlation functions involving $`_zu_y^n(𝐫)`$. In general, it will have terms connecting all pairs of layers, and the evaluation of $`S(𝐫,na)`$ is highly nontrivial.
We can obtain a useful approximation by setting $`K_y=0`$ to eliminate all $`g^{(2)}(𝐫,na)`$ for $`n0`$ and by setting $`B_{uh}=0`$ to eliminate couplings between $`u_z^n`$ and $`u_y^m`$. Even in this approximation, couplings between distant layers arise from $`_zu_y^n(𝐫)_zu_y^m(𝐫^{})`$. This function, however, dies off rapidly with $`nm`$, and we will set it equal to zero when $`nm`$. \[In another publication, we will use a variational procedure to calculate $`S(𝐫,na)`$. The form of $`S(𝐫,na)`$ evaluated with this procedure is very similar to that obtained using the above approximations.\] With these approximations,
$$\mathrm{\Phi }(𝐫_1,\mathrm{},𝐫_n,𝐫)=q_0^2\underset{m=0}{\overset{n}{}}g^{(2)}(𝐫_m𝐫_{m+1})+nW_y$$
(6.17)
and
$`S(𝐫,na)`$ $`=`$ $`{\displaystyle \frac{\stackrel{~}{V}^u}{2T}}^n{\displaystyle d^2r_1\mathrm{}d^2r_nS_2(𝐫_1)}`$ (6.19)
$`\times S_2(𝐫_1𝐫_2)\mathrm{}S_2(𝐫_n𝐫),`$
where
$$S_2(𝐫)=e^{iq_0[u_z^n(𝐫)u_z^n(0)]}$$
(6.20)
is the inplane density correlation function and where
$$\stackrel{~}{V}^u=V^ue^{W_y}.$$
(6.21)
Fourier transforming this equation, we obtain Eq. (3.58) in the text.
## G Interaction Energy between Edge Dislocations
In this appendix, we evaluate the interaction energy $`E(𝐫,na)`$ between two dislocations with separation $`𝐱=(𝐫,na)`$ in the limits $`xx^{}`$, $`z=0`$ and $`zz^{}`$, $`x=0`$. To obtain the $`xx^{}`$, $`z=0`$ limit, we must evaluate
$`E_n(x)`$ $`=`$ $`{\displaystyle _0^{\mathrm{\Lambda }_xx^{}}}{\displaystyle \frac{dt}{t}}{\displaystyle _0^\pi }𝑑u\sqrt{t^2+u^2p(u)}`$ (7.1)
$`\times `$ $`\mathrm{cos}[nu]\left(1\mathrm{cos}[tx/x^{}]\right).`$ (7.2)
We isolate the $`\mathrm{ln}x`$ divergence by adding and subtracting $`\sqrt{u^2p(u)}`$ under the $`u`$-integral. This procedure yields
$`E_n(x)=J_n\mathrm{ln}\left[C_n^x(\mathrm{\Lambda }_x)|x|/x^{}\right]`$ (7.3)
for $`xx^{}`$. In the above expression, $`C_n^x(\mathrm{\Lambda }_x)=e^\gamma B_n(\mathrm{\Lambda }_x)`$ with $`B_n(\mathrm{\Lambda }_x)=\mathrm{\Lambda }_xx^{}e^{\overline{B}_n(\mathrm{\Lambda }_x)}`$,
$$\overline{B}_n(\mathrm{\Lambda }_x)=_0^{\mathrm{\Lambda }_xx^{}}\frac{dt}{t}\left[\frac{J_n(t)}{J_n}1\right],$$
(7.4)
and
$$J_n(t)=_0^\pi 𝑑u\sqrt{t^2+u^2p(u)}\mathrm{cos}nu.$$
(7.5)
The $`zz^{}`$, $`x=0`$ limit is obtained in a similar way from
$`E_n(z)`$ $`=`$ $`{\displaystyle _0^{\mathrm{\Lambda }_xx^{}}}{\displaystyle \frac{dt}{t}}{\displaystyle _0^\pi }𝑑u\sqrt{t^2+u^2p(u)}`$ (7.6)
$`\times `$ $`\mathrm{cos}[nu]\left(1\mathrm{exp}\left[{\displaystyle \frac{z}{z^{}}}t\sqrt{t^2+u^2p(u)}\right]\right).`$ (7.7)
We find that
$`E_n(z)=J_n\mathrm{ln}\left[C_n^z(\mathrm{\Lambda }_x)|z|/z^{}\right]`$ (7.8)
scales logarithmically for $`zz^{}`$, where $`C_n^z(\mathrm{\Lambda }_x)=B_n(\mathrm{\Lambda }_x)e^{\overline{C}_n^z(\mathrm{\Lambda }_x)}`$,
$`\overline{C}_n^z={\displaystyle _0^1}𝑑y\left[{\displaystyle \frac{1F_n(y)}{y}}\right]{\displaystyle _1^{\mathrm{}}}𝑑y{\displaystyle \frac{F_n(y)}{y}},`$ (7.9)
and
$`F_n(y)={\displaystyle \frac{1}{J_n}}{\displaystyle _0^\pi }𝑑u\mathrm{cos}(nu)\sqrt{u^2p(u)}\mathrm{exp}\left[y\sqrt{u^2p(u)}\right].`$ (7.10)
Note that $`C_n^{x,z}(\mathrm{\Lambda }_x)`$ diverge with $`\mathrm{\Lambda }_x`$, and thus $`E_n(𝐫)`$ does not have a well-defined $`\mathrm{\Lambda }_x\mathrm{}`$ limit.
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# Untitled Document
HIERARCHIC THEORY OF CONDENSED MATTER
AND ITS INTERACTION WITH LIGHT:
New Theories of Light Refraction, Brillouin Scattering
and Mössbauer effect
Alex Kaivarainen
JBL, University of Turku, Finland, FIN-20520
URL: http://www.karelia.ru/~alexk
H2o@karelia.ru
Materials, presented in this article are the part of new quantum theory of condensed matter, described in:
. Book by A. Kaivarainen: Hierarchic Concept of Matter and Field. Water, biosystems and elementary particles. New York, 1995; ISBN 0-9642557-0-7 and two articles:
. ”New Hierarchic Theory of Matter General for Liquids and Solids: dynamics, thermodynamics and mesoscopic structure of water and ice” (see URL: http://www.karelia.ru/~alexk \[New articles\]) and:
. Hierarchic Concept of Matter, General for Liquids and Solids: Water and ice (see Proceedings of the Second Annual Advanced Water Sciences Symposium, October 4-6, 1966, Dallas, Texas.
Computerized verification of new described here theories are presented on examples of WATER and ICE, using special computer program: ”Comprehensive Analyzer of Matter Properties (CAMP)” (copyright, 1997, A. Kaivarainen).
CONTENTS
Summary of ”New Hierarchic Theory of Condensed Matter.”
1: New approach to theory of light refraction
1.1. Refraction in gas
1.2. Light refraction in liquids and solids
2: Mesoscopic theory of Brillouin light scattering
2.1. Traditional approach
2.2. Fine structure of scattering
2.3. Mesoscopic approach
2.4. Quantitative verification of hierarchic theory of Brillouin scattering
3: Mesoscopic theory of Mössbauer effect
3.1. General background
3.2. Probability of elastic effects
3.3. Doppler broadening in spectra nuclear gamma-resonance (NGR)
3.4. Acceleration and forces, related to thermal dynamics of molecules and ions. Vibro-gravitational interaction
=========================================================================
Summary of:
New Hierarchic Theory of Condensed Matter
by: A. Kaivarainen
A basically new hierarchic quantitative theory, general for solids and liquids, has been developed.
It is assumed, that unharmonic oscillations of particles in any condensed matter lead to emergence of three-dimensional (3D) superposition of standing de Broglie waves of molecules, electromagnetic and acoustic waves. Consequently, any condensed matter could be considered as a gas of 3D standing waves of corresponding nature. Our approach unifies and develops strongly the Einstein’s and Debye’s models.
Collective excitations, like 3D standing de Broglie waves of molecules, representing at certain conditions the molecular mesoscopic Bose condensate, were analyzed, as a background of hierarchic model of condensed matter.
The most probable de Broglie wave (wave B) length is determined by the ratio of Plank constant to the most probable impulse of molecules, or by ratio of its most probable phase velocity to frequency. The waves B are related to molecular translations (tr) and librations (lb).
As the quantum dynamics of condensed matter does not follow in general case the classical Maxwell-Boltzmann distribution, the real most probable de Broglie wave length can exceed the classical thermal de Broglie wave length and the distance between centers of molecules many times.
This makes possible the atomic and molecular Bose condensation in solids and liquids at temperatures, below boiling point. It is one of the most important results of new theory, which we have confirmed by computer simulations on examples of water and ice.
Four strongly interrelated new types of quasiparticles (collective excitations) were introduced in our hierarchic model:
1. Effectons (tr and lb), existing in ”acoustic” (a) and ”optic” (b) states represent the coherent clusters in general case;
2. Convertons, corresponding to interconversions between tr and lb types of the effectons (flickering clusters);
3. Transitons are the intermediate $`\left[ab\right]`$ transition states of the tr and lb effectons;
4. Deformons are the 3D superposition of IR electromagnetic or acoustic waves, activated by transitons and convertons.
Primary effectons (tr and lb) are formed by 3D superposition of the most probable standing de Broglie waves of the oscillating ions, atoms or molecules. The volume of effectons (tr and lb) may contain from less than one, to tens and even thousands of molecules. The first condition means validity of classical approximation in description of the subsystems of the effectons. The second one points to quantum properties of coherent clusters due to molecular Bose condensation.
The liquids are semiclassical systems because their primary (tr) effectons contain less than one molecule and primary (lb) effectons - more than one molecule. The solids are quantum systems totally because both kind of their primary effectons (tr and lb) are molecular Bose condensates. These consequences of our theory are confirmed by computer calculations.
The 1st order $`\left[gasliquid\right]`$ transition is accompanied by strong decreasing of rotational (librational) degrees of freedom due to emergence of primary (lb) effectons and $`\left[liquidsolid\right]`$ transition - by decreasing of translational degrees of freedom due to Bose-condensation of primary (tr) effectons.
In the general case the effecton can be approximated by parallelepiped with edges corresponding to de Broglie waves length in three selected directions (1, 2, 3), related to the symmetry of the molecular dynamics. In the case of isotropic molecular motion the effectons’ shape may be approximated by cube.
The edge-length of primary effectons (tr and lb) can be considered as the ”parameter of order”.
The in-phase oscillations of molecules in the effectons correspond to the effecton’s (a) - acoustic state and the counterphase oscillations correspond to their (b) - optic state. States (a) and (b) of the effectons differ in potential energy only, however, their kinetic energies, impulses and spatial dimensions - are the same. The b-state of the effectons has a common feature with Frölich’s polar mode.
The $`(ab)`$ or $`(ba)`$ transition states of the primary effectons (tr and lb), defined as primary transitons, are accompanied by a change in molecule polarizability and dipole moment without density fluctuations. At this case they lead to absorption or radiation of IR photons, respectively.
Superposition (interception) of three internal standing IR photons of different directions (1,2,3) - forms primary electromagnetic deformons (tr and lb).
On the other hand, the \[lb$``$tr\] convertons and secondary transitons are accompanied by the density fluctuations, leading to absorption or radiation of phonons.
Superposition resulting from interception of standing phonons in three directions (1,2,3), forms secondary acoustic deformons (tr and lb).
Correlated collective excitations of primary and secondary effectons and deformons (tr and lb), localized in the volume of primary tr and lb electromagnetic deformons, lead to origination of macroeffectons, macrotransitons and macrodeformons (tr and lb respectively).
Correlated simultaneous excitations of tr and lb macroeffectons in the volume of superimposed tr and lb electromagnetic deformons lead to origination of supereffectons.
In turn, the coherent excitation of both: tr and lb macrodeformons and macroconvertons in the same volume means creation of superdeformons. Superdeformons are the biggest (cavitational) fluctuations, leading to microbubbles in liquids and to local defects in solids.
Total number of quasiparticles of condensed matter equal to 4!=24, reflects all of possible combinations of the four basic ones \[1-4\], introduced above. This set of collective excitations in the form of ”gas” of 3D standing waves of three types: de Broglie, acoustic and electromagnetic - is shown to be able to explain virtually all the properties of all condensed matter.
The important positive feature of our hierarchic model of matter is that it does not need the semi-empiric intermolecular potentials for calculations, which are unavoidable in existing theories of many body systems. The potential energy of intermolecular interaction is involved indirectly in dimensions and stability of quasiparticles, introduced in our model.
The main formulae of theory are the same for liquids and solids and include following experimental parameters, which take into account their different properties:
$`\left[1\right]`$\- Positions of (tr) and (lb) bands in oscillatory spectra;
$`\left[2\right]`$\- Sound velocity;
$`\left[3\right]`$\- Density;
$`\left[4\right]`$\- Refraction index (extrapolated to the infinitive wave length of photon$`)`$.
The knowledge of these four basic parameters at the same temperature and pressure makes it possible using our computer program, to evaluate more than 300 important characteristics of any condensed matter. Among them are such as: total internal energy, kinetic and potential energies, heat-capacity and thermal conductivity, surface tension, vapor pressure, viscosity, coefficient of self-diffusion, osmotic pressure, solvent activity, etc. Most of calculated parameters are hidden, i.e. inaccessible to direct experimental measurement.
The new interpretation and evaluation of Brillouin light scattering and Mössbauer effect parameters may also be done on the basis of hierarchic theory. Mesoscopic scenarios of turbulence, superconductivity and superfluity are elaborated.
Some original aspects of water in organization and large-scale dynamics of biosystems - such as proteins, DNA, microtubules, membranes and regulative role of water in cytoplasm, cancer development, quantum neurodynamics, etc. have been analyzed in the framework of Hierarchic theory.
Computerized verification of our Hierarchic concept of matter on examples of water and ice is performed, using special computer program: Comprehensive Analyzer of Matter Properties (CAMP, copyright, 1997, Kaivarainen). The new opto-acoustical device (CAMP), based on this program, with possibilities much wider, than that of IR, Raman and Brillouin spectrometers, has been proposed
(see URL: http://www.karelia.ru/~alexk).
This is the first theory able to predict all known experimental temperature anomalies for water and ice. The conformity between theory and experiment is very good even without any adjustable parameters.
The hierarchic concept creates a bridge between micro- and macro- phenomena, dynamics and thermodynamics, liquids and solids in terms of quantum physics.
1: New approach to theory of light refraction
1.1. Refraction in gas
If the action of photons onto electrons of molecules is considered as a force, activating a harmonic oscillator with decay, it leads to the known classical equations for a complex refraction index (Vuks, 1984).
The Lorentz-Lorenz formula obtained in such a way is convenient for practical needs. However, it does not describe the dependence of refraction index on the incident light frequency and did not take into account the intermolecular interactions. In the new theory proposed below we have tried to clear up the relationship between these parameters.
Our basic idea is that the dielectric penetrability of matter $`ϵ`$, (equal in the optical interval of frequencies to the refraction index squared $`n^2)`$, is determined by the ratio of partial volume energies of photon in vacuum to similar volume energy of photon in matter:
$$ϵ=n^2=\frac{[E_p^0]}{[E_p^m]}=\frac{m_pc^2}{m_pc_m^2}=\frac{c^2}{c_m^2}$$
(1.1)
where $`m_p=h\nu _p/c^2`$ is the effective photon mass, $`c`$ is the light velocity in vacuum, $`c_m`$ is the effective light velocity in matter.
We introduce the notion of partial volume energy of a photon in vacuum $`[E_p^0]`$ and in matter $`[E_p^m]`$ as a product of photon energy $`(E_p=h\nu _p)`$ and the volume $`(V_p)`$ occupied by 3D standing wave of photon in vacuum and in matter, correspondingly:
$$[E_p^0]=E_pV_p^0[E_p^m]=E_pV_p^m$$
(1.2)
The 3D standing photon volume as an interception volume of 3 different standing photons normal to each other was termed in our mesoscopic model as a primary electromagnetic deformon (see Introduction of ).
In vacuum, where the effect of an excluded volume due to the spatial incompatibility of electron shells of molecules and photon is absent, the volume of $`\mathrm{\hspace{0.17em}3}D`$ photon standing wave (primary deformon) is:
$$V_p^0=\frac{1}{n_p}=\frac{3\lambda _p^2}{8\pi }$$
(1.3)
We will consider the interaction of light with matter in this mesoscopic volume, containing a thousands of molecules of condensed matter. It is the reason why we titled this theory of light refraction as mesoscopic one.
Putting (1.3) into (1.2), we obtain the formula for the partial volume energy of a photon in vacuum:
$$[E_p^0]=E_pV_p^0=h\nu _p\frac{9\lambda _p^2}{8\pi }=\frac{9}{4}\mathrm{}c\lambda _p^2$$
(1.4)
Then we proceed from the assumption that waves B of photons can not exist with waves B of electrons, forming the shells of atoms and molecules in the same space elements. Hence, the effect of excluded volume appears during the propagation of an external electromagnetic wave through the matter. It leads to the fact that in matter the volume occupied by a photon, is equal to
$$V_p^m=V_p^0V_p^{\text{ex}}=V_p^0n_M^pV_e^M$$
(1.5)
where $`V_p^{\text{ex}}=n_M^pV_e^M`$ is the excluded volume which is equal to the product of the number of molecules in the volume of one photon standing wave $`(n_M^p)`$ and the volume occupied by the electron shell of one molecule $`(V_e^M)`$.
$`n_M^p`$ is determined by the product of the volume of the photons 3D standing wave in the vacuum (1.3) and the concentration of molecules $`(n_M=N_0/V_0)`$:
$$n_M^p=\frac{9\lambda _p^3}{8\pi }\left(\frac{N_0}{V_0}\right)$$
(1.6)
In the absence of the polarization by the external field and intermolecular interaction, the volume occupied by electrons of the molecule:
$$V_e^M=\frac{4}{3}\pi L_e^3$$
(1.7)
where $`L_e`$ is the radius of the most probable wave $`B(L_e=\lambda _e/2\pi )`$ of the outer electron of a molecule. As it has been shown in (7.5) that the mean molecule polarizability is:
$$\alpha =L_e^3$$
(1.8)
Then taking (1.7) and (1.6) into account, the excluded volume of primary electromagnetic deformon in the matter is:
$$V_p^{\text{ex}}=\frac{9\lambda _p^3}{8\pi }n_M\frac{4}{3}\pi \alpha =\frac{3}{2}\lambda _p^3n_M\alpha $$
(1.9)
Therefore, the partial volume energy of a photon in the vacuum is determined by eq.(1.4), while that in matter, according to (1.5):
$$[E_p^m]=E_pV_p^m=E_p[V_p^0V_p^{\text{ex}}]$$
(1.10)
Putting (1.4) and (1.10) into (1.1) we obtain:
$$ϵ=n^2=\frac{E_pV_p^0}{E_p(V_p^0V_p^{\text{ex}})}$$
(1.11)
or
> $$\frac{1}{n^2}=1\frac{V_p^{\text{ex}}}{V^0}$$
> (1.12)
> Then, putting eq.(1.9) and (1.3) into (1.12) we derive new equation for refraction index, leading from our mesoscopic theory:
>
> $$\frac{1}{n^2}=1\frac{4}{3}\pi n_M\alpha $$
> (1.13)
> or in another form:
$$\frac{n^21}{n^2}=\frac{4}{3}\pi n_M\alpha =\frac{4}{3}\pi \frac{N_0}{V_0}\alpha $$
(1.14)
where: $`n_M=N_0/V_0`$ is a concentration of molecules;
In this equation $`\alpha =L_e^3`$ is the average static polarizability of molecules for the case when the external electromagnetic fields as well as intermolecular interactions inducing the additional polarization are absent. This situation is realized at $`E_p=h\nu _p0`$ and $`\lambda _p\mathrm{}`$ in the gas phase. As will be shown below the value of resulting $`\alpha ^{}`$in condensed matter is bigger.
1.2. Light refraction in liquids and solids
According to the Lorentz classical theory, the electric component of the outer electromagnetic field is amplified by the additional inner field $`(E_{\text{ad}})`$, related to the interaction of induced dipole moments in composition of condensed matter with each other:
$$E_{\text{ad}}=\frac{n^21}{3}E$$
(1.15)
The mean Lorentz acting field $`\overline{F}`$ can be expressed as:
$$\overline{F}=E+E_{\text{ad}}=\frac{n^2+2}{3}E(\text{at }n1,\overline{F}E)$$
(1.16)
$`\overline{F}`$\- has a dimensions of electric field tension and tends to E in the gas phase when $`n1`$.
In accordance with our mesoscopic model, except the Lorentz acting field, the total internal acting field, includes also two another contributions, increasing the molecules polarizability ($`\alpha `$) in condensed matter:
1. Potential intermolecular field, including all the types of Van- der-Waals interactions in composition of coherent collective excitations, even without external electromagnetic field. Like total potential energy of matter, this contribution must be dependent on temperature and pressure;
2. Primary internal field, related with primary electromagnetic deformons (tr and lib). This component of the total acting field also exist without external fields. Its frequencies corresponds to IR range and its action is much weaker than the action of the external visible light.
Let us try to estimate the energy of the total acting field and its effective frequency ($`\nu _f`$) and wavelength $`(\lambda _f)`$, that we introduce as:
$$A_f=h\nu _f=\frac{hc}{\lambda _f}=A_L+A_V+A_D$$
(1.17)
where: $`A_L,A_V`$ and $`A_D`$ are contributions, related with Lorentz field, potential field and primary deformons field correspondingly.
When the interaction energy of the molecule with a photon $`(E_p=h\nu _p)`$ is less than the energy of the resonance absorption, then it leads to elastic polarization of the electron shell and origination of secondary photons, i.e. light scattering. We assume in our consideration that the increment of polarization of a molecule $`(\alpha )`$ under the action of the external photon $`(h\nu _p)`$ and the total active field $`(A_f=h\nu _f)`$ can be expressed through the increase of the most probable radius of the electron’s shell $`(L_e=\alpha ^{1/3})`$, using our (eq. 7.6 from ):
$$\mathrm{\Delta }L_e=\frac{\omega _pm_e}{2\mathrm{}}\alpha $$
(1.18)
where the resulting increment:
$$\mathrm{\Delta }L^{}=\mathrm{\Delta }L_e+\mathrm{\Delta }L_f=\frac{(h\nu _p+A_f)m_e}{2\mathrm{}^2}\alpha $$
(1.18a)
where: $`\alpha =L_e^3`$ is the average polarizability of molecule in gas phase at $`\nu _f`$ 0.
For water molecule in the gas phase:
$$L_e=\alpha ^{1/3}=1.1310^{10}m$$
is a known constant, determined experimentally .
The total increment of polarizability radius $`(\mathrm{\Delta }L^{})`$ and resulting polarizability of molecules ($`\alpha ^{})`$in composition of condensed matter affected by the acting field
$$\alpha ^{}=(L^{})^3$$
can be find from the experimental refraction index (n) using our formula (1.14):
$$L^{}=\left[\frac{3}{4\pi }\frac{V_0}{N_0}\frac{n^21}{n^2}\right]^{1/3}$$
(1.19)
$$\mathrm{\Delta }L^{}=L^{}L_e$$
(1.20)
$$\text{where: }L^{}=(\alpha ^{})^{1/3}$$
From (1.18) we get a formula for the increment of radius of polarizability $`(\mathrm{\Delta }L_f)`$, induced by the total internal acting field:
$$\mathrm{\Delta }L_f=\mathrm{\Delta }L^{}\mathrm{\Delta }L_e=\frac{A_fm_e}{2\mathrm{}^2}\alpha $$
(1.21)
Like total internal acting field energy (1.17), this total acting increment can be presented as a sum of contributions, related to Lorentz field $`(\mathrm{\Delta }L_F)`$, potential field $`(\mathrm{\Delta }L_V)`$ and primary deformons field $`(\mathrm{\Delta }L_D)`$:
$$\mathrm{\Delta }L_f=\mathrm{\Delta }L_L+\mathrm{\Delta }L_V+\mathrm{\Delta }L_D$$
(1.22)
Increment $`\mathrm{\Delta }L_e`$, induced by external photon only, can be calculated from the known frequency ($`\nu _p`$) of the incident light (see 1.18a):
$$\mathrm{\Delta }L_e=\frac{h\nu _pm_e}{2\mathrm{}^2}\alpha $$
(1.23)
It means that $`\mathrm{\Delta }L_f`$ can be found from (1.21) and (1.17), using (1.23). Then from (1.21) we can calculate the energy $`(A_f)`$, effective frequency $`(\nu _f)`$ and wave length $`(\lambda _f)`$ of the total acting field like:
$$A_f=h\nu _f=hc/\lambda _f=2\frac{\mathrm{\Delta }L_f\mathrm{}^2}{m_e\alpha }$$
(1.24)
The computer calculations of $`\alpha ^{};L^{}=L_e+\mathrm{\Delta }L^{}=(\alpha ^{})^{1/3}`$and$`A_f`$ in the temperature range $`(095^0)`$ are presented on Fig.1.1.
One must keep in mind that in general case $`\alpha `$ and $`L`$ are tensors. It means that all the increments, calculated on the base of eq.(1.18a) must be considered as the effective ones. Nevertheless, it is obvious that our approach to analysis of the acting field parameters can give useful additional information about the properties of transparent condensed matter.
> Fig. 1.1. (a)- Temperature dependencies of the most probable outer electron shell radius of $`H_2O(L^{})`$ and the effective polarizability $`\alpha ^{}=(L^{})^3`$ in the total acting field;
>
> (b)- Temperature dependence of the total acting field $`(A_f)`$ energy in water at the wavelength of the incident light $`\lambda _p=5.46110^5cm^1`$. The experimental data for refraction index n(t) were used in calculations. The initial electron shell radius is: $`L_e=\alpha _{H2O}^{1/3}=1.1310^8cm`$. In graphical calculations in Fig.1.1a, the used experimental temperature dependence of the water refraction index were obtained by Frontas’ev and Schreiber .
The temperature dependencies of these parameters were computed using the known experimental data on refraction index $`n(t)`$ for water and presented in Fig.1.1a. The radius $`L^{}`$ in the range $`095^0C`$ increases less than by 1% at constant incident light wavelength $`(\lambda =546.1nm)`$. The change of $`\mathrm{\Delta }L_f`$ with temperature is determined by its potential field component change$`\mathrm{\Delta }L_V`$.
The relative change of this component: $`\mathrm{\Delta }\mathrm{\Delta }L_V/\mathrm{\Delta }L_f(t=0^0C)`$ is about 9%. Corresponding to this change the increasing of the acting field energy $`A_f`$(eq.1.23) increases approximately by $`8kJ/M(`$Fig 1.1 b) due to its potential field contribution.
It is important that the total potential energy of water in the same temperature range, according to our calculations, increase by the same magnitude (Fig.5b in or Fig.3b in \[ 3\]). This fact points to the strong correlation between potential intermolecular interaction in matter and the value of the acting field energy.
It was calculated that, at constant temperature $`(20^0)`$ the energy of the acting field $`(A_f),(eq\mathrm{.1.23})`$ in water practically does not depend on the wavelength of incident light $`(\lambda _p)`$. At more than three time alterations of $`\lambda _p`$: from $`12.5610^5cm`$ to $`3.0310^5cm`$ and the water refraction index $`\left(n\right)`$ from 1.320999 to 1.358100 the value of $`A_f`$ changes less than by 1%.
At the same conditions the electron shell radius L and the acting polarizability $`\alpha ^{}`$ thereby increase from (1.45 to 1.5)$``$10$`{}_{}{}^{10}m`$ and from (3.05 to $`3.274)10^{30}m^3`$ respectively (Fig.1.2). These changes are due to the incident photons action only. For water molecules in the gas phase and $`\lambda _p\mathrm{}`$ the initial polarizability $`(\alpha =L_e^3)`$ is equal to $`1.4410^{24}cm^3`$, i.e. significantly less than in condensed matter under the action of external and internal fields.
Obviously, the temperature change of energy $`A_f`$ (Fig.1.1b) is determined by the internal pressure increasing (section 11.2 of ), related to intermolecular interaction change, depending on mean distances between molecules and, hence, on the concentration $`(N_0/V_0)`$ of molecules in condensed matter.
> Fig. 1.2. Dependencies of the acting polarizability $`\alpha ^{}=(L^{})^3`$ and electron shell radius of water in the acting field $`(L^{})`$ on incident light wavelength $`(\lambda _p)`$, calculated from eq. (1.14) and experimental data $`n(\lambda _p)[6]`$. The initial polarizability of $`H_2O`$ in the gas phase at $`\lambda _p\mathrm{}`$ is equal to $`\alpha =L_e^3=1.4410^{24}cm^3`$. The corresponding initial radius of the $`H_2O`$ electron shell is$`L_e=1.1310^8cm`$.
On the basis of our data, changes of $`A_f,`$ calculated from (1.24) are caused mainly by the heat expansion of the matter. The photon induced increment of the polarizability $`(\alpha \alpha ^{})`$ practically do not change $`A_f.`$
The ability to obtain new valuable information about changes of molecule polarizability under the action of incident light and about temperature dependent molecular interaction in condensed medium markedly reinforce such a widely used method as refractometry.
The above defined relationship between the molecule polarizability and the wave length of the incident light allows to make a new endeavor to solve the light scattering problems.
2. Mesoscopic theory of Brillouin light scattering in condensed matter
2.1. Traditional approach
According to traditional concept, light scattering in liquids and crystals as well as in gases takes place due to random heat fluctuations. In condensed media the fluctuations of density, temperature and molecule orientation are possible.
Density ($`\rho `$) fluctuations leading to dielectric penetrability $`(ϵ)`$ fluctuations are of major importance. This contribution is estimated by means of Einstein formula for scattering coefficient of liquids :
$$R=\frac{Ir^2}{I_0V}=\frac{\pi }{2\lambda ^4}kT\beta _T\left(\rho \frac{ϵ}{\rho }\right)_T$$
(2.1)
where $`\beta _T`$ is isothermal compressibility.
Many authors made attempts to find a correct expression for the variable $`(\rho \frac{ϵ}{\rho }).`$
The formula derived by Vuks is most consistent with experimental data:
$$\rho \frac{ϵ}{\rho }=(n^21)\frac{3n^2}{2n^21}$$
(2.2)
2.2. Fine structure of scattering
The fine structure - spectrum of the scattering in liquids is represented by two Brillouin components with frequencies shifted relatively from the incident light frequency: $`\nu _\pm =\nu _0\pm \mathrm{\Delta }\nu `$ and one unshifted band like in gases $`(\nu _0)`$.
The shift of the Brillouin components is caused by the Doppler effect resulting from a fraction of photons scattering on phonons moving at sound speed in two opposite directions .
This shift can be explained in different way as well. If in the antinodes of the standing wave the density oscillation occurs at frequency ($`\mathrm{\Omega }`$):
$$\rho =\rho _0\mathrm{cos}\mathrm{\Omega }t,$$
(2.3)
then the scattered wave amplitude will change at the same frequency. Such a wave can be represented as a superposition of two monochromatic waves having the frequencies:$`(\omega +\mathrm{\Omega })`$ and $`(\omega \mathrm{\Omega })`$, where
$$\mathrm{\Omega }=2\pi f$$
(2.4)
is the elastic wave frequency at which scattering occurs when the Wolf-Bragg condition is satisfied:
$$2\mathrm{\Lambda }\mathrm{sin}\phi =2\mathrm{\Lambda }\mathrm{sin}\frac{\theta }{2}=\lambda ^{^{}}$$
(2.5)
or
$$\mathrm{\Lambda }=\lambda ^{}/(2\mathrm{sin}\frac{\theta }{2})=\frac{c}{n\nu }(2\mathrm{sin}\frac{\theta }{2})=v_{ph}/f$$
(2.6)
where $`\mathrm{\Lambda }`$ is the elastic wave length corresponding to the frequency $`f;\lambda ^{^{}}=\lambda /n=c/n\nu (\lambda ^{}`$ and $`\lambda `$ are the incident light wavelength in matter and vacuum, respectively$`);\phi `$ is the angle of sliding; $`\theta `$ is the angle of scattering; n is the refraction index of matter; $`c`$is the light speed.
The value of Brillouin splitting is represented as:
$$\pm \mathrm{\Delta }\nu _{MB}=f=\frac{V_{ph}}{\mathrm{\Lambda }}=2\nu \frac{V_{ph}}{c}n\mathrm{sin}\frac{\theta }{2}$$
(2.7)
where: $`\nu n/c=1/\lambda ;n`$ is the refraction index of matter; $`\nu `$ is incident light frequency;
$$v_{ph}=v_S$$
(2.8)
is the phase velocity of a scattering wave equal to hypersonic velocity.
The formula (2.7) is identical to that obtained from the analysis of the Doppler effect:
$$\frac{\mathrm{\Delta }\nu }{\nu }=\pm 2\frac{V_S}{c}n\mathrm{sin}\frac{\theta }{2}$$
(2.9)
According to the classical theory, the central line, which is analogous to that observed in gases, is caused by entropy fluctuations in liquids, without any changes of pressure . On the basis of Frenkel theory of liquid state, the central line can be explained by fluctuations of ”hole” number - cavitational fluctuations .
The thermodynamic approach of Landau and Plachek leads to the formula, which relates the intensities of the central (I) and two lateral $`(I_{MB})`$ lines of the scattering spectrum with compressibility and heat capacities:
$$\frac{I}{2I_{MB}}=\frac{I_p}{I_{\text{ad}}}=\frac{\beta _T\beta _S}{\beta _S}=\frac{C_pC_v}{C_v}$$
(2.10)
where: $`\beta _T`$ and $`\beta _S`$ are isothermal and adiabatic compressibilities; $`C_p`$ and $`C_v`$ are isobaric and isohoric heat capacities.
In crystals, quartz for example, the central line in the fine structure of light scattering is usually absent or very small. However, instead of one pair of shifted components, observed in liquids, there appear three Brillouin components in crystals. One of them used to be explained by scattering on the longitudinal phonons, and two - by scattering on the transversal phonons.
2.3. New mesoscopic approach to problem
In our hierarchic theory the thermal ”random” fluctuations are ”organized” by different types of superimposed quantum excitations.
According to our Hierarchic model, including microscopic, mesoscopic and macroscopic scales of matter (see Introduction of ), the most probable (primary) and mean (secondary) effectons, translational and librational are capable of quantum transitions between two discreet states: $`(ab)_{tr,lb}`$ and $`(\overline{a}\overline{b})_{tr,lb}`$ respectively. These transitions lead to origination/annihilation of photons and phonons, forming primary and secondary deformons.
The mean heat energy of molecules is determined by the value of 3kT, which as our calculations show, has the intermediate value between the discreet energies of a and b quantum states of primary effectons (Fig.19 of ), making, consequently, the non equilibrium conditions in condensed matter. Such kind of instability is a result of ”competition” between classical and quantum distributions of energy .
The maximum deviations from thermal equilibrium and that of the dielectric properties of matter occur when the same states of primary and secondary quasiparticles, e.g. a,ā and b,occur simultaneously. Such a situation corresponds to the A and B states of macroeffectons. The $`(AB)_{tr,lb}`$ transitions represent thermal fluctuations. The big density fluctuations are related to ”flickering clusters” (macroconvertions between librational and translational primary effectons) and the maximum fluctuations correspond to Superdeformons.
Only in the case of spatially independent fluctuations the interference of secondary scattered photons does not lead to their total compensation.
The probability of the event that two spatially uncorrelated events coincide in time is equal to the product of their independent probabilities .
Thus, the probabilities of the coherent (a,ā) and (b,) states of primary and secondary effectons, corresponding to A and B states of the macroeffectons (tr and lib), independent on each other, are equal to:
$$\left(\begin{array}{c}P_M^A\end{array}\right)_{tr,lb}^{\text{ind}}=\left(\begin{array}{c}P_{ef}^a\overline{P}_{ef}^a\end{array}\right)_{tr,lb}^S\left(\frac{1}{Z^2}\right)=\left(\frac{P_M^A}{Z^2}\right)_{tr,lb}$$
(2.11)
$$\left(P_M^B\right)_{tr,lb}^{\text{ind}}=\left(\begin{array}{c}P_{ef}^b\overline{P}_{ef}^b\end{array}\right)_{tr,lb}^S\left(\frac{1}{Z^2}\right)=\left(\frac{P_M^B}{Z^2}\right)_{tr,lb}$$
(2.12)
where
$$\frac{1}{Z}\left(P_{ef}^a\right)_{tr,lb}\text{ and }\frac{1}{Z}\left(\overline{P}_{ef}^a\right)_{tr,lb}$$
(2.13)
are the independent probabilities of a and ā states determined according to formulae (4.10 and 4.18 of ), while probabilities $`\left(P_{ef}^b/Z\right)_{tr,lb}`$ and $`\left(\overline{P}_{ef}^b/Z\right)_{tr,lb}`$ are determined according to formulae (4.11 and 4.19 of );
$`Z`$ is the sum of probabilities of all types of quasiparticles states - eq.(4.2 of ).
The probabilities of molecules being involved in the spatially independent translational and librational macrodeformons are expressed as the products (2.11) and (2.12):
$$\left(\begin{array}{c}P_D^M\end{array}\right)_{tr,lb}^{\text{ind}}=\left[\begin{array}{c}\left(\begin{array}{c}P_M^A\end{array}\right)^{\text{ind}}\left(\begin{array}{c}P_M^B\end{array}\right)^{\text{ind}}\end{array}\right]_{tr,lb}=\frac{P_D^M}{Z^4}$$
(2.14)
Formulae (2.11) and (2.12) may be considered as the probabilities of space-independent but coherent macroeffectons in A and B states, respectively.
For probabilities of space-independent supereffectons in $`A^{}`$ and $`B^{}`$ states we have:
$$\left(\begin{array}{c}P_S^A^{}\end{array}\right)^{\text{ind}}=\left(\begin{array}{c}P_M^A\end{array}\right)_{tr}^{\text{ind}}\left(\begin{array}{c}P_M^A\end{array}\right)_{lb}^{\text{ind}}=\frac{P_S^A^{}}{Z^4}$$
(2.15)
$$\left(\begin{array}{c}P_S^B^{}\end{array}\right)^{\text{ind}}=\left(\begin{array}{c}P_M^B\end{array}\right)_{tr}^{\text{ind}}\left(\begin{array}{c}P_M^B\end{array}\right)_{tr}^{\text{ind}}=\frac{P_S^b^{}}{Z^4}$$
(2.15a)
In a similar way we get from (2.14) the probabilities of spatially independent superdeformons:
$$\left(\begin{array}{c}P_S^D^{}\end{array}\right)^{\text{ind}}=\left(\begin{array}{c}P_M^D\end{array}\right)_{tr}\left(\begin{array}{c}P_M^D\end{array}\right)_{lb}=\frac{P_S^D^{}}{Z^4}$$
(2.16)
The concentrations of molecules, the states of which markedly differ from the equilibrium one and which cause light scattering in composition of spatially independent macroeffectons and macrodeformons, are equal to:
$$\left[N_M^A=\frac{N_0}{Z^2V_0}\left(\begin{array}{c}P_M^A\end{array}\right)\right]_{tr,lb};\left[N_M^B=\frac{N_0}{Z^2V_0}\left(\begin{array}{c}P_M^B\end{array}\right)\right]_{tr,lb}$$
(2.17)
$$\left[N_M^D=\frac{N_0}{Z^4V_0}\left(\begin{array}{c}P_M^D\end{array}\right)\right]_{tr,lb}$$
The concentrations of molecules, involved in a-convertons, b- convertons and Macroconvertons or c-Macrotransitons (see Introduction) are correspondingly:
$$N_M^{ac}=\frac{N_0}{Z^2V_0}P_{ac};N_M^{bc}=\frac{N_0}{Z^2V_0}P_{bc};N_M^C=\frac{N_0}{Z^4V_0}P_{\text{cMt}}$$
(2.18)
The probabilities of convertons-related excitations are the same as used in Chapter 4 of book .
The concentration of molecules, participating in the independent supereffectons and superdeformons:
$$N_M^A^{}=\frac{N_0}{Z^4V_0}P_s^A^{};N_M^B^{}=\frac{N_0}{Z^4V_0}P_S^B^{}$$
(2.19)
$$N_M^D^{}=\frac{N_0}{Z^8V_0}P_S^D^{}$$
(2.20)
where$`N_0`$ and $`V_0`$ are the Avogadro number and the molar volume of the matter.
Substituting (2.17 - 2.20) into well known Raleigh formula for scattering coefficient, measured at the right angle between incident and scattered beams:
$$R=\frac{I}{I_0}\frac{r^2}{V}=\frac{8\pi ^4}{\lambda ^4}\alpha ^2n_\text{M}(cm^1)$$
(2.20a)
we obtain the values of contributions from different states of quasiparticles to the resulting scattering coefficient:
$$\left(R_A^M\right)_{tr,lb}=\frac{8\pi ^4}{\lambda ^4}\frac{(\alpha ^{})^2}{Z^2}\frac{N_0}{V_0}\left(\begin{array}{c}P_M^A\end{array}\right)_{tr,lb};R_A^s=\frac{8\pi ^4}{\lambda ^4}\frac{(\alpha ^{})^2}{Z^4}\frac{N_0}{V_0}P_s^A^{}$$
(2.21)
$$\left(R_B^M\right)_{tr,lb}=\frac{8\pi ^4}{\lambda ^4}\frac{(\alpha ^{})^2}{Z^2}\frac{N_0}{V_0}\left(\begin{array}{c}P_M^B\end{array}\right)_{tr,lb};R_B^s=\frac{8\pi ^4}{\lambda ^4}\frac{(\alpha ^{})^2}{Z^4}\frac{N_0}{V_0}P_s^B^{}$$
(2.22)
$$\left(R_D^M\right)_{tr,lb}=\frac{8\pi ^4}{\lambda ^4}\frac{(\alpha ^{})^2}{Z^2}\frac{N_0}{V_0}\left(\begin{array}{c}P_M^D\end{array}\right)_{tr,lb};R_D^s=\frac{8\pi ^4}{\lambda ^4}\frac{(\alpha ^{})^2}{Z^4}\frac{N_0}{V_0}P_s^D^{}$$
(2.23)
The contributions of excitations, related to $`[tr/lb]`$ convertons are:
$$R_{ac}=\frac{8\pi ^4}{\lambda ^4}\frac{(\alpha ^{})^2}{Z^2}\frac{N_0}{V_0}R_{bc}=\frac{8\pi ^4}{\lambda ^4}\frac{(\alpha ^{})^2}{Z^2}\frac{N_0}{V_0}P_{bc}$$
$$R_{\text{abc}}=\frac{8\pi ^4}{\lambda ^4}\frac{(\alpha ^{})^2}{Z^4}\frac{N_0}{V_0}P_{\text{cMt}}$$
where: $`\alpha ^{}`$ is the acting polarizability determined by eq.(1.24) and (1.25).
The resulting coefficient of the isotropic scattering $`(R_{\text{iso}})`$ is defined as the sum of contributions (2.21-2.23) and is subdivided into three kinds of scattering: caused by translational quasiparticles, caused by librational quasiparticles and by the mixed type of quasiparticles:
$$R_{\text{iso}}=[R_A^M+R_B^M+R_D^M]+[R_A^M+R_B^M+R_D^M]+[R_{ac}+R_{bc}+R_{\text{abc}}]+[R_A^s+R_B^s+R_D^s]$$
(2.24)
Total contributions, related to convertons and superexcitations are correspondingly:
$$R_C=R_{ac}+R_{bc}+R_{\text{abc }}\text{ and }R_S=R_A^s+R_B^s+R_D^s$$
The polarizability of anisotropic molecules having no cubic symmetry is a tensor. In this case, total scattering (R) consists of scattering at density fluctuations $`(R_{\text{iso}})`$ and scattering at fluctuations of the anisotropy $`\left(R_{\text{an}}=\frac{13\mathrm{\Delta }}{67\mathrm{\Delta }}R_{\text{iso}}\right)`$:
$$R=R_{\text{iso}}+\frac{13\mathrm{\Delta }}{67\mathrm{\Delta }}R_{\text{iso}}=R_{\text{iso}}\frac{6+6\mathrm{\Delta }}{67\mathrm{\Delta }}=R_{\text{iso}}K$$
(2.25)
where R$`_{\text{iso}}`$ corresponds to $`eq.(2.24);\mathrm{\Delta }`$ is the depolarization coefficient.
The factor: $`\left(\frac{6+6\mathrm{\Delta }}{67\mathrm{\Delta }}\right)=K`$ was obtained by Cabanne and is called after him. In the case of isotropic molecules when $`\mathrm{\Delta }=0`$, the Cabanne factor is equal to 1.
The depolarization coefficient ($`\mathrm{\Delta }`$) could be determined experimentally as the ratio:
$$\mathrm{\Delta }=I_x/I_z,$$
(2.26)
where$`I_x`$ and $`I_z`$ are two polarized components of the beam scattered at right angle with respect to each other in which the electric vector is directed parallel and perpendicular to the incident beam, respectively. For example, in water $`\mathrm{\Delta }=0.09(`$Vuks, 1977).
According to the proposed theory of light scattering in liquids the central unshifted (like in gases) component of the Brillouin scattering spectrum, is caused by fluctuations of concentration and self-diffusion of molecules, participating in the convertons, macrodeformons (tr and lib) and superdeformons. The scattering coefficients of the central line $`(R_{\text{centr}})`$ and side lines $`(2R_{\text{side}})`$ in transparent condensed matter, as follows from (2.24) and (2.25), are equal correspondingly to:
$$R_{\text{cent}}=K\left[\left(\begin{array}{c}R_D^M\end{array}\right)_{tr}+\left(\begin{array}{c}R_D^M\end{array}\right)_{lb}\right]+K(R_C+R_S)$$
(2.27)
and
$$2R_{\text{side}}=\left(\begin{array}{c}R_A^M+R_B^M\end{array}\right)_{tr}+\left(\begin{array}{c}R_A^M+R_A^M\end{array}\right)_{lb}$$
(2.27a)
where $`K`$ is the Cabanne factor.
The total coefficient of light scattering is:
$$R_t=R_{\text{cent}}+2R_{\text{side}}$$
(2.28)
In accordance with our model the fluctuations of anisotropy (Cabanne factor) should be taken into account for calculations of the central component only. The orientations of molecules in composition of A and B states of Macroeffectons are correlated and their coherent oscillations are not accompanied by fluctuations of anisotropy of polarizability (see Fig.2.1).
The probabilities of the convertons, macrodeformons and superdeformons excitations (eqs.2.14, 4.16, 4.27 in ) are much lower in crystals than in liquids and hence, the central line in the Brillouin spectra of crystals is not usually observed.
The lateral lines in Brillouin spectra are caused by the scattering on the molecules forming (A) and (B) states of spatially independent macroeffectons, as it was mentioned above.
The polarizabilities of the molecules forming the independent macroeffectons, synchronized in (A)<sub>tr,lb</sub> and $`(B)_{tr,lb}`$ states and dielectric properties of these states, differ from each other and from that of transition states (macrodeformons). Such short-living states should be considered as the non equilibrium ones.
In fact we must keep in mind, that static polarizabilities in the more stable ground A state of the macroeffectons are higher than in B state, because the energy of long-term Van der Waals interaction between molecules of the A state is bigger than that of B-state.
If this difference may be attributed mainly to the difference in the long-therm dispersion interaction, then from (1.33) we obtain:
$$E_BE_A=V_BV_A=\frac{3}{2}\frac{E_0}{r^6}\left(\begin{array}{c}\alpha _B^2\alpha _A^2\end{array}\right)$$
(2.29)
where polarizability of molecules in A-state is higher, than that in B-state:
$$\alpha _A^2>\left[\begin{array}{c}\left(\begin{array}{c}\alpha ^{}\end{array}\right)^2\alpha _D^2\end{array}\right]>\alpha _B^2$$
The kinetic energy and dimensions of ”acoustic” and ”optic” states of macroeffectons are the same: $`T_{\text{kin}}^A=T_{\text{kin}}^B`$.
In our present calculations of light scattering we ignore this difference (2.29) between polarizabilities of molecules in A and B states.
But it can be taken into account if we assume, that polarizabilities in (A) and (a), (B) and (b) states of primary effectons are like:
$$\alpha _A\alpha _a\alpha ^{};\alpha _B\alpha _b$$
and the difference between the potential energy of (a) and (b) states is determined mainly by dispersion interaction (eq.2.28).
Experimental resulting polarizability ($`\alpha ^{}\alpha _a`$) can be expressed as:
$$\alpha _a=f_a\alpha _a+f_b\alpha _b+f_t\alpha $$
(2.29a)
where $`\alpha _t\alpha `$ is polarizability of molecules in the gas state (or transition state);
$`f_a`$ $`={\displaystyle \frac{P_a}{P_a+P_b+P_t}};f_b={\displaystyle \frac{P_b}{P_a+P_b+P_t}}\text{;}`$
$`\text{and }f_t`$ $`=f_d={\displaystyle \frac{P_t}{P_a+P_b+P_t}}`$
are the fractions of (a), (b) and transition (t) states (equal to 2.66) as far $`P_t=P_d=P_aP_b`$.
On the other hand from (1.33) at $`r=const`$ we have:
$$\mathrm{\Delta }V_{\text{dis}}^{ba}=\frac{3}{4}\frac{(2\alpha \mathrm{\Delta }\alpha )}{r^6}I_0(r_a=r_b;I_0^aI_0^b)\text{ and }$$
$$\frac{\mathrm{\Delta }V_{\text{dis}}^{ba}}{V^b}=\frac{h\nu _p}{h\nu _b}=\frac{\mathrm{\Delta }\alpha _a}{\alpha }\text{ or }\mathrm{\Delta }\alpha _a=\alpha _a\frac{\nu _p}{\nu _b}$$
(2.29b)
$$\alpha _b=\alpha _a\mathrm{\Delta }\alpha _a=\alpha _a(1\nu _p/\nu _b)$$
where: $`\mathrm{\Delta }\alpha _a`$ is a change of each molecule polarizability as a result of the primary effecton energy changing: $`E_bE_a+h\nu _p`$ with photon radiation; $`\nu _b`$ is a frequency of primary effecton in (b)- state (eq.2.28).
Combining (2.29) and (2.29b) we derive for $`\alpha _a`$ and $`\alpha _b`$ of the molecules composing primary translational or librational effectons:
$$\alpha _a=\frac{f_t\alpha }{1\left(f_a+f_b+f_b\frac{\nu _p}{\nu _b}\right)}$$
(2.30)
$$\alpha _b=\alpha _a\left(1\frac{\nu _p}{\nu _b}\right)$$
(2.30a)
The calculations by means of (2.30) are approximate in the framework of our assumptions mentioned above. But they correctly reflect the tendencies of $`\alpha _a`$ and $`\alpha _b`$ changes with temperature.
The ratio of intensities or scattering coefficients for the central component to the lateral ones previously was described by Landau- Plachek formula (2.10). According to our mesoscopic theory this ratio can be calculated in another way leading from (2.27) and (2.28):
$$\frac{I_{\text{centr}}}{2I_{MB}}=\frac{R_{\text{cent}}}{2R_{\text{side}}}$$
(2.30b)
Combining (2.30) and Landau- Plachek formula (2.10) it is possible to calculate the ratio $`(\beta _T/\beta _S)`$ and $`(C_P/C_V)`$ using our mesoscopic theory of light scattering.
2.4. Factors that determine the Brillouin line width
The known equation for Brillouin shift is (see 2.7):
$$\mathrm{\Delta }\nu _{MB}=\nu _0=2\frac{v_s}{\lambda }n\mathrm{sin}(\theta /2)$$
(2.31)
where: $`v_s`$ is the hypersonic velocity; $`\lambda `$ is the wavelength of incident light, $`n`$ is the refraction index of matter, and $`\theta `$ \- scattering angle.
The deviation from $`\nu _0`$ that determines the Brillouin side line half width may be expressed as the result of fluctuations of sound velocity $`v_s`$ and $`n`$ related to A and B states of tr and lib macroeffectons:
$$\frac{\mathrm{\Delta }\nu _0}{\nu _0}=\left(\frac{\mathrm{\Delta }v_s}{v_s}+\frac{\mathrm{\Delta }n}{n}\right)$$
(2.32)
$`\mathrm{\Delta }\nu _0`$ is the most probable side line width, i.e. the true half width of Brillouin line. It can be expressed as:
$$\mathrm{\Delta }\nu _0=\mathrm{\Delta }\nu _{\mathrm{exp}}F\mathrm{\Delta }\nu _{\text{inc}}$$
where $`\mathrm{\Delta }\nu _{\mathrm{exp}}`$ is the half width of the experimental line, $`\mathrm{\Delta }\nu _{\text{inc}}`$ \- the half width of the incident line, $`F`$ \- the coefficient that takes into account apparatus effects.
Let us analyze the first and the second terms in the right part of (2.32) separately.
The $`v_s`$ squared is equal to the ratio of the compressibility modulus ($`K`$) and density $`(\rho )`$:
$$v_s^2=K^2/\rho $$
(2.33)
Consequently, from (2.33) we have:
$$\frac{\mathrm{\Delta }v_s}{v_s}=\frac{1}{2}\left(\frac{\mathrm{\Delta }K}{K}\frac{\mathrm{\Delta }\rho }{\rho }\right)$$
(2.34)
In the case of independent fluctuations of K and $`\rho :`$:
$$\frac{\mathrm{\Delta }v_s}{v_s}=\frac{1}{2}\left(\left|\frac{\mathrm{\Delta }K}{K}\right|\left|\frac{\mathrm{\Delta }\rho }{\rho }\right|\right)$$
(2.35)
From our equation (1.14) we obtain for refraction index:
$$n^2=\left(1\frac{4}{3}N\alpha ^{}\right)^1,$$
(2.36)
where $`N=N_0/V_0`$ is the concentration of molecules.
From (2.36) we can derive:
$$\frac{\mathrm{\Delta }n}{n}=\frac{1}{2}\left(n^21\right)\left(\frac{\mathrm{\Delta }\alpha ^{}}{\alpha ^{}}+\frac{\mathrm{\Delta }N}{N}\right)$$
(2.37)
where:
$$(\mathrm{\Delta }N/N)=(\mathrm{\Delta }\rho /\rho )$$
(2.38)
and
$$\left(\frac{\mathrm{\Delta }\alpha ^{}}{\alpha ^{}}\right)\left(\frac{\mathrm{\Delta }K}{K}\right)$$
(2.39)
we can assume eq.(2.39) as far both parameters: polarizability $`(\alpha ^{})`$ and compressibility models (K) are related with the potential energy of intermolecular interaction.
For the other hand one can suppose that the following relation is true:
$$\frac{\mathrm{\Delta }\alpha ^{}}{\alpha ^{}}\frac{\overline{E}_{ef}^a3kT}{3kT}=\frac{\mathrm{\Delta }K}{K}$$
(2.40)
where: $`\begin{array}{c}\overline{E}_{ef}^a\hfill \end{array}`$ is the energy of the secondary effectons in (ā) state; $`E_0=3kT`$ is the energy of an ”ideal” quasiparticle as a superposition of 3D standing waves.
The density fluctuations can be estimated as a result of the free volume $`(v_f)`$ fluctuations (see 2.45):
$$\left(\frac{\mathrm{\Delta }v_f}{v_f}\right)_{tr,lb}=\frac{1}{Z}\left(P_D^M\right)_{tr,lb}\left(\mathrm{\Delta }N/N\right)_{tr,lb}$$
(2.41)
Now, putting (2.40) and (2.41) into (2.37) and (2.34) and then into (2.32), we obtain the semiempirical formulae for the Brillouin line half width calculation:
$$\frac{\mathrm{\Delta }\nu _f}{\nu _f}\frac{n^2}{2}\left[\frac{\overline{E}_{ef}^a3kT}{3kT}+\frac{1}{Z}\left(\begin{array}{c}P_D^M\end{array}\right)\right]_{tr,lb}$$
(9.42)
Brillouin line intensity depends on the half-width $`\mathrm{\Delta }\nu `$ of the line in following ways:
for a Gaussian line shape:
$$I(\nu )=I_0^{\mathrm{max}}\mathrm{exp}\left[0.693\left(\frac{\nu \nu _0}{\frac{1}{2}\mathrm{\Delta }\nu _0}\right)^2\right];$$
(2.43)
for a Lorenzian line shape:
$$I(\nu )=\frac{I_0^{\mathrm{max}}}{1+\left[\begin{array}{c}\left(\nu \nu _0\right)/\frac{1}{2}\mathrm{\Delta }\nu _0\end{array}\right]^2}$$
(2.44)
The traditional theory of Brillouin line shape gives a possibility for calculation of $`\mathrm{\Delta }\nu _0`$ taking into account the elastic (acoustic) wave dissipation.
The fading out of acoustic wave amplitude may be expressed as:
$$A=A_0e^{\alpha x}\text{ or }A=A_0e^{\alpha v_s}$$
(2.45)
where $`\alpha `$ is the extinction coefficient; $`x=v_st`$ \- the distance from the source of waves; $`v_s`$ and $`t`$ \- sound velocity and time, correspondingly.
The hydrodynamic theory of sound propagation in liquids leads to the following expression for the extinction coefficient:
$$\alpha =\alpha _s+\alpha _b=\frac{\mathrm{\Omega }^2}{2\rho v_s^3}\left(\frac{4}{3}\eta _s+\eta _b\right)$$
(2.46)
where: $`\alpha _s`$and $`\alpha _b`$ are contributions to $`\alpha `$, related to share viscosity ($`\eta _s`$) and bulk viscosity $`(\eta _b)`$, respectively; $`\mathrm{\Omega }=2\pi f`$ is the angular frequency of acoustic waves.
When the side lines in Brillouin spectra broaden slightly, the following relation between their intensity (I) and shift $`(\mathrm{\Delta }\omega =\omega \omega _0)`$ from frequency $`\omega _0`$, corresponding to maximum intensity $`(I=I_0)`$ of side line is correct:
$$I=\frac{I_0}{1+\left(\frac{\omega \omega _0}{a}\right)},$$
(2.47)
where:
$$a=\alpha v_s.$$
One can see from (2.46) that at $`I(\omega )=I_0/2`$, the half width:
$$\mathrm{\Delta }\omega _{1/2}=2\pi \mathrm{\Delta }\nu _{1/2}=\alpha v_s\text{ and }\mathrm{\Delta }\nu _{1/2}=\frac{1}{2}\pi \alpha v_s$$
(2.48)
It will be shown in Chapter 12 how one can calculate the values of $`\eta _s`$ and consequently $`\alpha _s`$ on the basis of the mesoscopic theory of viscosity.
2.5. Quantitative verification of mesoscopic theory of
Brillouin scattering
The calculations made according to the formula (2.21 - 2.27) are presented in Fig.2.1-2.7. The proposed theory of scattering in liquids, based on our hierarchic concept, is more adequate than the traditional Einstein, Mandelschtamm-Brillouin, Landau-Plachek theories based on classical thermodynamics. It describes experimental temperature dependencies and the $`I_{\text{centr}}/2I_{MB}`$ ratio for water very well (Fig.2.3).
The calculations are made for the wavelength of incident light: $`\lambda _{ph}=546.1nm=5.46110^5cm`$. The experimental temperature dependence for the refraction index (n) at this wavelength was taken from the Frontas’ev and Schreiber paper (1965). The rest of data for calculating of various light scattering parameters of water (density the location of translational and librational bands in the oscillatory spectra) are identical to those used above in Chapter 6.
> Fig. 2.1. Theoretical temperature dependencies of the total scattering coefficient for water without taking into account the anisotropy of water molecules polarizability fluctuations in the volume of macroeffectons, responsible for side lines: $`\left[R(tot)\right]`$ \- eq.(2.27a; 2.28) and taking them into account: $`\left[KR(tot)\right]`$, where $`K`$ is the Cabanne factor (eq.2.25).
>
> Fig. 2.2. Theoretical temperature dependencies of contributions to the total coefficient of total light scattering (R) caused by translational and librational macroeffectons and macrodeformons (without taking into account fluctuations of anisotropy).
> Fig. 2.3. Theoretical temperature dependencies of central to side bands intensities ratio in Brillouin spectra (eq.2.30).
Mesoscopic theory of light scattering can be used to verify the correctness of our formula for refraction index of condensed matter we got from our theory (eq. 1.14):
$$\frac{n^21}{n^2}=\frac{4}{3}\pi \frac{N_0}{V_0}\alpha ^{}$$
(2.48a)
and to compare the results of its using with that of the Lorentz-Lorenz formula:
$$\frac{n^21}{n^2+1}=\frac{4}{3}\pi \frac{N_0}{V_0}\alpha $$
(2.49)
From formula (2.48a) the resulting or effective molecular polarizability squared $`(\alpha ^{})^2`$ used in eq.(2.21-2.23) is:
$$\left(\alpha ^{}\right)^2=\left[\frac{(n^21)/n^2}{(4/3)\pi (N_0/V_0)}\right]^2$$
(2.50)
On the other hand, from the Lorentz-Lorenz formula (2.49) we have another value of polarizability:
$$\alpha ^2=\left[\frac{(n^21)/(n^2+2)}{(4/3)\pi (N_0/V_0)}\right]^2$$
(2.51)
It is evident that the light scattering coefficients (eq.2.28), calculated using (2.50) and (2.51) taking refraction index: $`n=1.33`$should differ more than four times as far:
$$\frac{R(\alpha ^{})}{R(\alpha )}=\frac{(\alpha ^{})^2}{(\alpha )^2}=\frac{(n^21)/n^2}{(n^21)/(n^2+2)}=\left(\frac{n^2+2}{n^2}\right)^2=4.56$$
(2.52)
At $`\mathrm{\hspace{0.33em}25}^0`$ and $`\lambda _{ph}=546nm`$ the theoretical magnitude of the scattering coefficient for water, calculated from our formulae (2.28) is equal (see Fig.2.1) to:
$$R=11.210^5m^1$$
(2.53)
This result of our theory coincides well with the most reliable experimental value (Vuks, 1977):
$$R_{\mathrm{exp}}=10.810^5m^1$$
Multiplication of the side bands contribution $`(2R_{\text{side}})`$ to Cabanne factor increases the calculated total scattering to about 25% and makes the correspondence with experiment worse. This fact confirms our assumption that fluctuations of anisotropy of polarizability in composition of A and B states of macroeffectons should be ignored in light scattering evaluation due to correlation of molecular dynamics in these states, in contrast to that of macrodeformons.
> Fig. 2.4. Theoretical temperature dependencies of the contributions of A and B states of translational Macroeffectons to the total scattering coefficient of water (see also Fig.2.2);
> Fig. 2.5. Theoretical temperature dependencies of the contributions of the A and B states of librational Macroeffectons to the coefficient of light scattering (R).
It follows from the Fig.2.4 and 2.5 that the light scattering depends on $`(AB)`$ equilibrium of macroeffectons because $`(R_A)>(R_B)`$, i.e. scattering on $`A`$ states is bigger than that on $`B`$ states.
> Fig. 2.6. Theoretical temperature dependencies of the contributions to light scattering (central component), related to translational $`(R_D)_{tr}`$ and librational $`(R_D)_{lb}`$ macrodeformons.
Comparing Figs. 2.1; 2.3, and 2.6 one can see that the main contribution to central component of light scattering is determined by $`[lb/tr]`$ convertons $`R_c(`$see eq.2.27).
> Fig. 2.7. Theoretical temperature dependences for temperature derivative $`(dR/dT)`$ of the total coefficient of light scattering of water.
Nonmonotonic deviations of the dependencies $`dR/dT(`$Fig.2.7) reflect the nonmonotonic changes of the refraction index for water $`n_{H_2O}(T)`$, as indicated by available experimental data (Frontas’ev and Schreiber, 1965). The deviations of dependence $`n_{H_2O}(t)`$ from the monotonic way in accordance with hierarchic theory, are a consequence of the nonmonotonic change in the stability of water structure, i.e. nonlinear change of $`(AB)_{tr,lb}`$ equilibrium. Some possible reasons of such equilibrium change were discussed in Chapter 6.
It is clear from (2.52) that the calculations based on the Lorentz-Lorentz formula (2.51) give scattering coefficient values of about 4.5 times smaller than experimental ones. It means that the true $`\alpha ^{}`$ value can be calculated just on the basis of our mesoscopic theory of light refraction (eq.2.50).
The traditional Smolukhovsky-Einstein theory, valid for the integral light scattering only (eq. 2.1), yield values in the range of $`R=8.8510^5m^1`$ to $`R=10.510^5m^1[4`$, 8$`]`$.
All the results, discussed above, mean that our mesoscopic theory of light scattering works better and is much more informative than the conventional one.
2.6. Light scattering in solutions
If the guest molecules are dissolved in a liquid and their sizes are much less than incident light wavelength, they do not radically alter the solvent properties. For this case the described above mechanism of light scattering of pure liquids does not changed qualitatively.
For such dilute solutions the scattering on the fluctuations of concentration of dissolved molecules $`(R_c)`$ is simply added to the scattering on the density fluctuations of molecules of the host solvent (eq.2.28). Taking into account the fluctuations of molecule polarizability anisotropy (see 2.25) the total scattering coefficient of the solution $`(R_S)`$ is:
$$R_S=R_t+R_c$$
(2.54)
Eqs. (2.21 - 2.28) could be used for calculating $`R_t`$ until critical concentrations $`(C_{cr})`$ of dissolved substance when it start to destroy the solvent structure, so that the latter is no longer able to form primary librational effectons. Perturbations of solvent structure will induce low-frequency shift of librational bands in the oscillatory spectrum of the solution until these bands totally disappear.
If the experiment is made with a two-component solution of liquids, soluble in each other, e.g. water-alcohol, benzol-methanol etc., and the positions of translational and librational bands of solution components are different, then at the concentration of the dissolved substance: $`C>C_{cr}`$, the dissolved substance and the solvent (the guest and host) can switch their roles. Then the translational and librational bands pertinent to the guest subsystem start to dominate. In this case, $`R_t`$ is to be calculated from the positions of the new bands corresponding to the ”new” host-solvent. The total ”melting” of the primary librational ”host effectons” and the appearance of the dissolved substance ”guest effectons” is like the second order phase transition and should be accompanied by a heat capacity jump. The like experimental effects take place indeed .
According to our concept, the coefficient R<sub>c</sub> in eq.(2.54) is caused by the fluctuations of concentration of dissolved molecules in the volume of translational and librational macro- and superdeformons of the solvent. If the destabilization of the solvent is expressed in the low frequency shift of librational bands, then the coefficients $`(R_A`$ and $`R_B)_{lb}`$ increase (eq.2.21 and 2.22) with the probability of macro-excitations.. The probabilities of convertons and macro- and superdeformons and the central component of Brillouin spectra will increase also. Therefore, the intensity of the total light scattering increases correspondingly.
The fluctuations of concentration of the solute molecules, in accordance with our model, occur in the volumes of macrodeformons and superdeformons. Consequently, the contribution of solute molecules in scattering $`(R_c`$ value in eq.2.54) can be expressed by formula, similar to (2.23), but containing the molecule polarizability of the dissolved substance (”guest$`\mathrm{"}),`$ equal to $`(\alpha _g^{})^2`$ instead of the molecule polarizability $`(\alpha ^{})`$ of the solvent (”host”), and the molecular concentration of the ”guest” substance in the solution $`(n_g)`$ instead of the solvent molecule concentration $`(n_M=N_0/V_0)`$. For this case $`R_c`$ could be presented as a sum of the following contributions:
$$(R_c)_{tr,lb}=\frac{8\pi ^4}{\lambda ^4}(\alpha _g^{})^2n_g\left[(P_M^D)_{tr,lb}+P_S^D^{}\right]$$
(2.55)
$$R_c^D^{}=\frac{8\pi ^4}{\lambda ^4}(\alpha _g^{})^2n_g(P_S^D^{})$$
(2.55a)
The resulting scattering coefficient $`(R_e)`$ on fluctuations of concentration in (2.54) is equal to:
$$R_c=(R_c)_{tr}+(R_c)_{lb}+R_c^D^{}$$
(2.56)
If several substances are dissolved with concentrations lower than $`(C_{cr})`$, then their $`R_c`$ are summed up additively.
Formulae (2.55) and (2.56) are valid also for the dilute solutions.
Eqs.(2.21-2.28) and (2.54-2.56) should, therefore, be used for calculating the resulting coefficient of light scattering in solutions $`(R_S)`$.
The traditional theory represents the scattering coefficient at fluctuations of concentration as (Vuks, 1977):
$$R_c=\frac{\pi ^2}{2\lambda ^4}\left(\frac{ϵ}{x}\right)^2\mathrm{\Delta }x^2v$$
(2.57)
where $`(ϵ/x)`$ is the dielectric penetrability derivative with respect to one of the components: $`\mathrm{\Delta }\overline{x}^2`$ is the fluctuations of concentration of guest molecules squared in the volume element $`v`$.
The transformation of (2.57) on the basis of classical thermodynamics leads to the formula:
$$R_c=\frac{\pi ^2}{2\lambda ^4N_0}\left(2n\frac{n}{x}\right)\left(\frac{9n^2}{(2n^2+1)(n^2+2)}\right)^2x_1x_2V_{12}f,$$
(2.58)
where $`N_0`$ is the Avogadro number, $`x_1`$ and $`x_2`$ are the molar fractions of the first and second components in the solution, $`V_{12}`$ is the molar volume of the solution, $`f`$ is the function of fluctuations of concentration determined experimentally from the partial vapor pressures of the first $`(P_1)`$ and second $`(P_2)`$ solution components :
$$\frac{1}{f}=\frac{x_1}{P_1}\frac{P_1}{x_1}=\frac{x_2}{P_2}\frac{P_2}{x_2}$$
(2.59)
In the case of ideal solutions
$$\frac{P_1}{x_1}=\frac{P_1}{x_1};\frac{P_2}{x_2}=\frac{P_2}{x_2}\text{ ; and }f=1.$$
For application the mesoscopic theory of light scattering to study of crystals, liquids and solutions, the following information is needed:
1. Positions of translational and librational band maxima in oscillatory spectra;
2. Concentration of all types of molecules in solutions;
3. Refraction index or polarizability in the acting field of each component of solution at given temperature.
Application of our theory to quantitative analysis of transparent liquids and solids yields much more information about properties of matter, its mesoscopic and hierarchic dynamic structure than the traditional one.
3. Mesoscopic theory of Mössbauer effect
3.1. General background
When the atomic nucleus with mass (M) in the gas phase irradiates $`\gamma `$-quantum with energy of
$$E_0=h\nu _0=m_pc^2$$
(3.1)
where: $`m_p`$ is the effective photon mass, then according to the law of impulse conservation, the nuclear acquires additional velocity in the opposite direction:
$$v=\frac{E_0}{Mc}$$
(3.2)
The corresponding additional kinetic energy
$$E_R=\frac{Mv^2}{2}=\frac{E_0^2}{2Mc^2}$$
(3.3)
is termed recoil energy.
When an atom which irradiates $`\gamma `$-quantum is in composition of the solid body, then three situations are possible:
1. The recoil energy of the atom is higher than the energy of atom - lattice interaction. In this case, the atom irradiating $`\gamma `$-quantum would be knocked out from its position in the lattice. That leads to defects origination;
2. Recoil energy is insufficient for the appreciable displacement of an atom in the structure of the lattice, but is higher than the energy of phonon, equal to energy of secondary transitons and phonons excitation. In this case, recoil energy is spent for heating the lattice;
3. Recoil energy is lower than the energy of primary transitons, related to \[emission/absorption\] of IR translational and librational photons $`(h\nu _p)_{tr,lb}`$ and phonons $`(h\nu _{ph})_{tr,lb}`$. In that case, the probability (f) of $`\gamma `$-quantum irradiation without any the losses of energy appears, termed the probability (fraction) of a recoilless processes.
For example, when $`E_R<<h\nu _{ph}(\nu _{ph}`$ \- the mean frequency of phonons), then the mean energy of recoil:
$$E_R=(1f)h\nu _{ph}$$
(3.4)
Hence, the probability of recoilless effect is
$$f=1\frac{E_R}{h\nu _{ph}}$$
(3.5)
According to eq.(3.3) the decrease of the recoil energy $`E_R`$ of an atom in the structure of the lattice is related to increase of its effective mass $`(M`$). In our model $`M`$ corresponds to the mass of the effecton.
The effect of $`\gamma `$-quantum irradiation without recoil was discovered by Mössbauer in 1957 and named after him.
The value of Mössbauer effect is determined by the value of $`f1`$.
The big recoil energy may be transferred to the lattice by portions that are resonant to the frequency of IR photons (tr and lb) and phonons. The possibility of superradiation of IR quanta stimulation as a result of such recoil process is a consequence of our mesoscopic model.
The scattering of $`\gamma `$-quanta without lattice excitation, when $`E_R<<h\nu _{ph}`$, is termed the elastic one. The general expression for the probability of such phononless elastic $`\gamma `$-quantum radiation acts is equal to:
$$f=\mathrm{exp}\left(\frac{4\pi <x^2>}{\lambda _0^2}\right)$$
(3.6)
where $`\lambda _0=c/\nu _0`$ is the real wavelength of $`\gamma `$-quantum; $`<`$x$`{}_{}{}^{2}>`$ \- the nucleus oscillations mean amplitude squared in the direction of $`\gamma `$-quantum irradiation.
The $`\gamma `$-quanta wavelength parameter may be introduced like:
$$L_0=\lambda _0/2\pi ,$$
(3.7)
where: $`L_0=1.3710^5cm`$ for $`Fe^{57}`$, then eq.(3.6) could be written as follows:
$$f=\mathrm{exp}\left(\frac{<x^2>}{L_0^2}\right)$$
(3.8)
It may be shown , proceeding from the model of crystal as a system of 3N identical quantum oscillators, that when temperature (T) is much lower than the Debye one $`(\theta _D)`$ then:
$$<x^2>=\frac{9\mathrm{}^2}{4Mk\theta _D}\left\{1+\frac{2\mathrm{}^2T^2}{3\theta _D^2}\right\},$$
(3.9)
where $`\theta _D=h\nu _D/k`$ and $`\nu _D`$ is the Debye frequency.
From (3.1), (3.3) and (3.7) we have:
$$\frac{1}{L}=\frac{E_0}{\mathrm{}c}$$
(3.10)
where: $`E_0=h\nu =c(2ME_R`$)$`^{1/2}`$ is the energy of $`\gamma `$-quantum
Substituting eqs.(3.9 and 3.10) into eq.(3.8), we obtain the Debye-Valler formula:
$$f=\mathrm{exp}\left[\frac{E_R}{k\theta _D}\left\{\frac{3}{2}+\frac{\pi ^2T^2}{\theta _D}\right\}\right]$$
(3.11)
when $`T0`$, then
$$f\mathrm{exp}\left(\frac{3E_R}{2k\theta _D}\right)$$
(3.12)
3.2. Probability of elastic effects
Mean square displacements $`<`$x$`{}_{}{}^{2}>`$ of an atoms or molecules in condensed matter (eq. 3.8) is not related to excitation of thermal photons or phonons (i.e. primary or secondary transitons). According to our concept, $`<x^2>`$ is caused by the mobility of the atoms forming effectons and differs for primary and secondary translational and librational effectons in $`(a,\overline{a})_{tr,lb}`$ and $`(b,b)_{tr,lb}`$ states.
We will ignore below the contributions of macro- and supereffectons in Mössbauer effect as very small. Then the resulting probability of elastic effects at $`\gamma `$-quantum radiation is determined by the sum of the following contributions (see $`eqs\mathrm{.4.2}4.4`$of $`)`$:
$$f=\frac{1}{Z}\underset{tr,lb}{}\left[\left(\begin{array}{c}P_{ef}^af_{ef}^a+P_{ef}^bf_{ef}^b\end{array}\right)+\left(\begin{array}{c}\overline{P}_{ef}^a\overline{f}_{ef}^a+\overline{P}_{ef}^b\overline{f}_{ef}^b\end{array}\right)\right]_{tr,lb}$$
(3.13)
where: $`P_{ef}^a,P_{ef}^b,\overline{P}_{ef}^a,\overline{P}_{ef}^b`$ are the relative probabilities of the acoustic and optic states for primary and secondary effectons; Z is the total partition function.
These parameters are calculated as described in Chapter 4 of book and in papers cited in the Summary of this article. Each of contributions to resulting probability of the elastic effect can be calculated separately as:
$$\left(\begin{array}{c}f_{ef}^a\end{array}\right)_{tr,lb}=\mathrm{exp}\left[\frac{<\left(\begin{array}{c}x^a\end{array}\right)_{tr,lb}^2>}{L_0^2}\right]$$
(3.14)
$`\left(\begin{array}{c}f_{ef}^a\end{array}\right)_{tr,lb}`$ is the probability of elastic effect, related to dynamics of primary translational and librational effectons in a-state;
$$\left(\begin{array}{c}f_{ef}^b\end{array}\right)_{tr,lb}=\mathrm{exp}\left[\frac{<\left(\begin{array}{c}x^b\end{array}\right)_{tr,lb}^2>}{L_0^2}\right]$$
(3.15)
$`\left(\begin{array}{c}f_{ef}^b\end{array}\right)_{tr,lb}`$ is the probability of elastic effect in primary translational and librational effectons in b-state;
$$\left(\begin{array}{c}\overline{f}_{ef}^a\end{array}\right)_{tr,lb}=\mathrm{exp}\left[\begin{array}{c}\frac{<\left(\begin{array}{c}\overline{x}^a\end{array}\right)_{tr,lb}^2>}{L_0^2}\end{array}\right]$$
(3.16)
$`\left(\begin{array}{c}\overline{f}_{ef}^a\end{array}\right)_{tr,lb}`$ is the probability for secondary effectons in ā -state;
$$\left(\begin{array}{c}\overline{f}_{ef}^b\end{array}\right)_{tr,lb}=\mathrm{exp}\left[\frac{<\left(\begin{array}{c}\overline{x}^b\end{array}\right)_{tr,lb}^2>}{L_0^2}\right]$$
(3.17)
$`\left(\begin{array}{c}\overline{f}_{ef}^b\end{array}\right)_{tr,lb}`$ is the probability of elastic effect, related to secondary effectons in -state.
Mean square displacements within different types of effectons in eqs.(3.14-3.17) are related to their phase and group velocities. At first we express the displacements using group velocities of the waves $`B(v_{gr})`$ and periods of corresponding oscillations $`(T)`$ as:
$$<\left(\begin{array}{c}x^a\end{array}\right)_{tr,lb}^2>=\frac{<(v_{gr}^a)_{tr,lb}^2>}{<\nu _a^2>_{tr,lb}}=<\left(\begin{array}{c}v_{gr}^aT^a\end{array}\right)_{tr,lb}^2>$$
(3.18)
where $`(T^a)_{tr,lb}=(1/\nu _a)_{tr,lb}`$ is a relation between the period and the frequency of primary translational and librational effectons in a-state;
$`(v_{gr}^a=v_{gr}^b)_{tr,lb}`$ are the group velocities of atoms forming these effectons equal in (a) and (b) states.
In a similar way we can express the displacements of atoms forming (b) state of primary effectons (tr and lib):
$$<\left(\begin{array}{c}x^b\end{array}\right)_{tr,lb}^2>=\frac{<(v_{gr}^b)_{tr,lb}^2>}{<\nu _b^2>_{tr,lb}}$$
(3.19)
where $`\nu _b`$ is the frequency of primary translational and librational effectons in b-state.
The mean square displacements of atoms forming secondary translational and librational effectons in ā and states:
$$<\left(\begin{array}{c}\overline{x}^a\end{array}\right)_{tr,lb}^2>=\frac{<(\overline{v}_{gr}^a)_{tr,lb}^2>}{<\overline{\nu }_a^2>_{tr,lb}}$$
(3.20)
$$<\left(\begin{array}{c}\overline{x}^b\end{array}\right)_{tr,lb}^2>=\frac{<(\overline{v}_{gr}^b)_{tr,lb}^2>}{<\overline{\nu }_b^2>_{tr,lb}}$$
(3.21)
$$\text{where: }(\overline{v}_{gr}^a=\overline{v}_{gr}^b)_{tr,lb}$$
Group velocities of atoms in primary and secondary effectons may be expressed using corresponding phase velocities $`(v_{ph})`$ and formulae for waves B length as follows:
$`\left(\begin{array}{c}\lambda _a\end{array}\right)_{tr,lb}`$ $`={\displaystyle \frac{h}{m<v_{gr}>_{tr,lb}}}=\left({\displaystyle \frac{v_{ph}^a}{\nu _a}}\right)_{tr,lb}=`$ (3.22)
$`=\left(\begin{array}{c}\lambda _b\end{array}\right)_{tr,lb}=\left({\displaystyle \frac{v_{ph}^b}{\nu _b}}\right)_{tr,lb}`$ (3)
hence for the group velocities of the atoms or molecules forming primary effectons (tr and lb) squared we have:
$$\left(\begin{array}{c}v_{gr}^{a,b}\end{array}\right)_{tr,lb}^2=\frac{h^2}{m^2}\left(\frac{\nu _{a,b}}{v_{ph}^{a,b}}\right)_{tr,lb}^2$$
(3.23)
In accordance with mesoscopic theory, the wave B length, impulses and group velocities in a and b states of the effectons are equal. Similarly to (3.23), we obtain the group velocities of particles, composing secondary effectons:
$$\left(\begin{array}{c}\overline{v}_{gr}^{a,b}\end{array}\right)_{tr,lb}^2=\frac{h^2}{m^2}\left(\frac{\overline{\nu }_{a,b}}{\overline{v}_{ph}^{a,b}}\right)_{tr,lb}^2$$
(3.24)
Substituting eqs.(3.23) and (3.24) into (3.18-3.21), we find the important expressions for the average coherent displacements of particles squared as a result of their oscillations in the volume of the effectons (tr, lib) in both discreet states (acoustic and optic):
$$<(x^a)_{tr,lb}^2>=(h/mv_{ph}^a)_{tr,lb}^2$$
(3.25)
$$<(x^b)_{tr,lb}^2>=(h/mv_{ph}^b)_{tr,lb}^2$$
(3.26)
$$<(\overline{x}^a)_{tr,lb}^2>=(h/m\overline{v}_{ph}^a)_{tr,lb}^2$$
(3.27)
$$<(\overline{x}^b)_{tr,lb}^2>=(h/m\overline{v}_{ph}^b)_{tr,lb}^2$$
(3.28)
Then, substituting these values into eqs.(3.14-3.17) we obtain a set of different contributions to the resulting probability of effects without recoil:
$$\begin{array}{c}\hfill \left(\begin{array}{c}f_f^a\end{array}\right)_{tr.lb}=\mathrm{exp}[\left(\frac{h}{mL_0v_{ph}^a}\right)^2]_{tr,lb};\\ \hfill \left(\begin{array}{c}f_f^b\end{array}\right)_{tr.lb}=\mathrm{exp}[\left(\frac{h}{mL_0v_{ph}^b}\right)^2]_{tr,lb};\end{array}\}$$
(3.29)
$$\begin{array}{c}\hfill \left(\begin{array}{c}\overline{f}_f^a\end{array}\right)_{tr.lb}=\mathrm{exp}[\left(\frac{h}{mL_0\overline{v}_{ph}^a}\right)^2]_{tr,lb};\\ \hfill \left(\begin{array}{c}\overline{f}_f^b\end{array}\right)_{tr.lb}=\mathrm{exp}[\left(\frac{h}{mL_0\overline{v}_{ph}^b}\right)^2]_{tr,lb};\end{array}\}$$
(3.30)
where the phase velocities $`(v_{ph}^a,v_{ph}^b,\overline{v}_{ph}^a,\overline{v}_{ph}^b)_{tr,lb}`$ are calculated from the resulting sound velocity and the positions of translational and librational bands in the oscillatory spectra of matter at given temperature using eqs.2.69-2.75. The wavelength parameter:
$$L_0=\frac{c}{2\pi \nu _0}=\frac{hc}{2\pi E_0}=1.37510^{11}m$$
for gamma-quanta, radiated by nuclear of $`Fe^{57}`$, with energy:
$$E_0=14.4125\text{ kev }=2.3016710^8\text{erg }$$
Substituting eqs.(3.29) and (3.30) into (3.13), we find the total probability of recoilless effects $`(f_{\text{tot}})`$ in the given substance. Corresponding computer calculations for ice and water are presented on Figs.3.1 and 3.2.
As far the second order phase transitions in general case are accompanied by the alterations of the sound velocity and the positions of translational and librational bands, they should also be accompanied by alterations of f$`_{\text{tot}}`$ and its components.
> Fig. 3.1. Temperature dependences of total probability (f) for elastic effect without recoil and phonon excitation: (a) in ice; (b) in water; (c)-during phase transition. The calculations were performed using eq.(3.13).
>
> Fig. 3.2. (a) - The contributions to probability of elastic effect (f) (see Fig.3.1) for primary $`(f_{ef}^{a,b})_{tr}`$ and secondary $`(\overline{f})_{tr}`$ translational effectons and (b) and those of librational effectons $`(f_{ef}^{a,b})_{lb}`$ and $`(\overline{f})_{lb}`$ near the temperature of \[ice $``$ water\] phase transition.
The total probability (f) and its components, caused by primary and secondary quasiparticles were calculated according to formula (3.13). The value of $`\left(f\right)`$ determines the magnitude of the Mössbauer effect registered by $`\gamma `$-resonance spectroscopy.
The band width caused by recoilless effects is determined by the uncertainty principle and expressed as follows:
$$\mathrm{\Gamma }=\frac{h}{\tau }\frac{10^{27}}{1.410^7}=7.1410^{21}\text{ erg }=4.410^9eV$$
(3.31)
where $`\tau `$ is the lifetime of nucleus in excited state (for $`Fe^{57}\tau =1.410^7s)`$.
The position of the band depends on the mean square velocity of atoms, i.e. on second order Doppler effect. In the experiment, such an effect is compensated by the velocity of $`\gamma `$-quanta source motion relative to absorbent. In the framework of our model this velocity is interrelated with the mean velocity of the secondary effectons diffusion in condensed matter.
3.3. Doppler broadening in spectra of nuclear gamma-resonance (NGR)
Mössbauer effect is characterized by the nonbroadened component of NGR spectra only, with probability of observation determined by eq.(3.13).
When the energy of absorbed $`\gamma `$-quanta exceeds the energy of thermal IR photons (tr,lib) or phonons excitation, the absorbance band broadens as a result of Doppler effect. Within the framework of our mesoscopic concept the Doppler broadening is caused by thermal displacements of the particles during $`[ab`$ and $`\overline{a}\overline{b}]_{tr,lb}`$ transitions of primary and secondary effectons, leading to origination/annihilation of the corresponding type of deformons (electromagnetic and acoustic).
The flickering clusters: $`[lb/tr]`$ convertons (a and b), can contribute in the NGR line broadening also.
In that case, the value of Doppler broadening ($`\mathrm{\Delta }\mathrm{\Gamma }`$) of the band in the NGR spectrum could be estimated from corresponding kinetic energies of these excitations, related to their group velocities (see eq. 4.31). In our consideration we take into account the reduced to one molecule kinetic energies of primary and secondary translational and librational transitons, a-convertons and b-convertons. The contributions of macroconvertons, macro- and superdeformons are much smaller due to their small probability and concentration:
$`\mathrm{\Delta }\mathrm{\Gamma }`$ $`={\displaystyle \frac{V_0}{N_0Z}}\underset{tr,lb}{{\displaystyle }}\left(\begin{array}{c}n_tP_tT_t+\overline{n}_t\overline{P}_t\overline{T}_t\end{array}\right)_{tr,lb}+`$ (3.32)
$`+{\displaystyle \frac{V_0}{N_0Z}}(n_{ef})_{lb}[P_{ac}T_{ac}+P_{bc}T_{bc}]`$
where: $`N_0`$ and $`V_0`$ are the Avogadro number and molar volume;
Z is the total partition function $`(eq\mathrm{.4.2});n_t`$ and $`\overline{n}_t`$ are the concentrations of primary and secondary transitons (eqs.3.5 and 3.7);
$`(n_{ef})_{lb}=n_{\text{con }}`$ is a concentration of primary librational effectons, equal to that of the convertons; $`P_t`$ and $`\overline{P}_t`$ are the relative probabilities of primary and secondary transitons (eqs. 4.26 and $`4.27);P_{ac}`$ and $`P_{bc}`$ are relative probabilities of (a and b) -convertons (see Chapter 4 of );
$`T_t`$ and $`\overline{T}_t`$ are the kinetic energies of primary and secondary transitons, related to the corresponding total energies of these excitations $`(E_t`$ and $`\overline{E}_t)`$, their masses $`(M_t`$ and $`\overline{M}_t)`$ and the resulting sound velocity $`(v_s`$, see eq.2.40) in the following form:
$$(T_t)_{tr,lb}=\frac{_1^3\left(E_t^{1,2,3}\right)_{tr,lb}}{2M_t(v_s^{\text{res}})^2}$$
(3.33)
$$(T_t)_{tr,lb}=\frac{_1^3\left(\overline{E}_t^{1,2,3}\right)_{tr,lb}}{2\overline{M}_t(v_s^{\text{res}})^2}$$
(3.34)
The kinetic energies of (a and b) convertons are expressed in a similar way:
$$(T_{ac})=\frac{_1^3\left(E_{ac}^{1,2,3}\right)_{tr,lb}}{2M_c(v_s^{\text{res}})^2}$$
$$(T_{bc})=\frac{_1^3\left(E_{bc}^{1,2,3}\right)_{tr,lb}}{2M_c(v_s^{\text{res}})^2}$$
where: $`E_{ac}^{1,2,3}`$ and $`E_{bc}^{1,2,3}`$ are the energies of selected states of corresponding convertons; $`M_c`$ is the mass of convertons, equal to that of primary librational effectons.
The broadening of NGR spectral lines by Doppler effect in liquids is generally expressed using the diffusion coefficient (D) at the assumption that the motion of Mössbauer atom has the character of unlimited diffusion :
$$\mathrm{\Delta }\mathrm{\Gamma }=\frac{2E_0^2}{\mathrm{}c^2}D$$
(3.35)
where: $`E_0=h\nu _0`$ is the energy of gamma quanta; c is light velocity and
$$D=\frac{kT}{6\pi \eta a}$$
(3.36)
where: $`\eta `$ is viscosity, (a) is the effective Stokes radius of the atom $`Fe^{57}`$
The probability of recoilless $`\gamma `$-quantum absorption by the matter containing for example Fe<sup>57</sup>, decreases due to diffusion and corresponding Doppler broadening of band $`(\mathrm{\Delta }\mathrm{\Gamma })`$:
$$f_D=\frac{\mathrm{\Gamma }}{\mathrm{\Gamma }+\mathrm{\Delta }\mathrm{\Gamma }}$$
(3.37)
where $`\mathrm{\Delta }\mathrm{\Gamma }`$ corresponds to eq.(3.32).The formulae obtained here make it possible to experimentally verify a set of consequences of our mesoscopic theory using the gamma- resonance method. A more detailed interpretation of the data obtained by this method also becomes possible.
The magnitude of ($`\mathrm{\Delta }\mathrm{\Gamma }`$) was calculated according to formula (3.32). It corresponds well to experimentally determined Doppler widening in the nuclear gamma resonance (NGR) spectra of ice.
> Fig. 3.3. The temperature dependences of the parameter $`\mathrm{\Delta }\mathrm{\Gamma }`$, characterizing the nonelastic effects and related to the excitation of thermal phonons and IR photons: a) in ice; b) in water; c) near phase transition.
3.4. Acceleration and forces, related to thermal dynamics of molecules and ions.
Hypothesis of Vibro-gravitational interaction
During the period of particles thermal oscillations (tr and lb), their instant velocity, acceleration and corresponding forces alternatively and strongly change.
The change of wave B instant group velocity, averaged during the molecule oscillation period in composition of the (a) and (b) states of the effectons, determines the average acceleration:
$$\left[a_{gr}^{a,b}=\frac{dv_{gr}^{a,b}}{dt}=\frac{v_{gr}^{a,b}}{T}=v_{gr}\nu ^{a,b}\right]_{tr,lb}^{1,2,3}$$
(3.38)
We keep in mind that group velocities, impulses and wave B length in (a) and (b) states of the effectons are equal, in accordance with our model.
Corresponding to (3.38) forces:
$$\left[F^{a,b}=ma_{gr}^{a,b}\right]_{tr,lb}^{1,2,3}$$
(3.39)
The energies of molecules in (a) and (b) states of the effectons also can be expressed via accelerations:
$$\left[E^{a,b}=h\nu ^{a,b}=F^{a,b}\lambda =ma^{a,b}\lambda =ma^{a,b}(v_{ph}^{a,b}/\nu ^{a,b})\right]_{tr,lb}^{1,2,3}$$
(3.40)
From (3.40) one can express the accelerations of particles in the primary effectons of condensed matter, using their phase velocities as a waves B:
$$\left[a_{gr}^{a,b}=\frac{h(\nu ^{a,b})^2}{mv_{ph}^{a,b}}\right]_{tr,lb}^{1,2,3}$$
(3.41)
The accelerations of particles in composition of secondary effectons have a similar form:
$$\left[\overline{a}_{gr}^{a,b}=\frac{h(\overline{\nu }^{a,b})^2}{m\overline{v}_{ph}^{a,b}}\right]_{tr,lb}^{1,2,3}$$
(3.42)
These parameters are important for understanding the dynamic properties of condensed systems. The accelerations of the atoms, forming primary and secondary effectons can be calculated, using eqs.(2.74-2.75 of ) to determine phase velocities and eqs. $`(2.27,\mathrm{\hspace{0.33em}2.28},\mathrm{\hspace{0.33em}2.54},\mathrm{\hspace{0.33em}2.55}[1])`$ to find a frequencies.
Multiplying (3.41) and (3.42) by the atomic mass $`m`$, we derive the most probable and mean forces acting upon the particles in both states of primary and secondary effectons in condensed matter:
$$\left[F_{gr}^{a,b}=\frac{h(\nu ^{a,b})^2}{v_{ph}^{a,b}}\right]_{tr,lb}^{1,2,3}\left[\overline{F}_{gr}^{a,b}=\frac{h(\overline{\nu }^{a,b})^2}{\overline{v}_{ph}^{a,b}}\right]$$
(3.43)
The comparison of calculated accelerations with empirical data of the Mössbauer effect - supports the correctness of our approach.
According to eq.(2.54) in the low temperature range, when $`h\nu _a<<kT`$, the frequency of secondary tr and lb effectons in the (a) state can be estimated as:
$$\nu ^a=\frac{\nu _a}{\mathrm{exp}\left(\frac{h\nu _a}{kT}\right)1}\frac{kT}{h}$$
(3.44)
For example, at $`T=200K`$ we have $`\overline{\nu }^a410^{12}s^1`$.
If the phase speed in eq.(3.42) is taken equal to $`\overline{v}_{ph}^a=2.110^5cm/s(`$see Fig.2a) and the mass of water molecule:
$$m=181.6610^{24}g=2.9810^{23}g,$$
then from (3.42) we get the acceleration of molecules in composition of secondary effectons of ice in (a) state:
$$\overline{a}_{gr}^a=\frac{h(\overline{\nu }^a)^2}{m\overline{v}_{ph}^a}=1.610^{16}cm/s^2$$
This value is about $`10^{13}`$ times more than that of free fall acceleration $`(g=9.810^2cm/s^2)`$, which agrees well with experimental data, obtained for solid bodies .
Accelerations of $`H_2O`$ molecules in composition of primary librational effectons $`(a_{gr}^a)`$ in the ice at 200K and in water at 300K are equal to: $`0.610^{13}cm/s^2`$ and $`210^{15}cm/s^2`$, correspondingly. They also exceed to many orders the free fall acceleration.
It was shown experimentally (Sherwin, 1960), that heating of solid body leads to decreasing of gamma-quanta frequency (red Doppler shift) i.e. increasing of corresponding quantum transitions period. This can be explained as the relativist time-pace decreasing due to elevation of average thermal velocity of atoms .
The thermal vibrations of particles (atoms, molecules) in composition of primary effectons as a partial Bose-condensate are coherent. The increasing of such clusters dimensions, determined by most probable wave B length, as a result of cooling, pressure elevation or under magnetic field action (see section 14.6 of ) leads to enhancement of coherent regions.
Each coherently vibrating cluster of particles with big alternating accelerations, like librational and translational effectons is a source of coherent gravitational waves.
The frequency of these vibro-gravitational waves (VGW) is equal to frequency of particles vibrations (i.e. frequency of the effectons in a or b states). The amplitude of VGW $`(A_G)`$ is proportional to the number of vibrating coherently particles $`(N_G)`$ in composition of primary effectons:
$$A_GN_GV_{ef}/(V_0/N_0)=(1/n_{ef})/(V_0/N_0)$$
The resonant long-distance gravitational interaction between coherent clusters of the same body or between that of different bodies is possible. The formal description of this vibro-gravitational interaction (VGI) could be like that of distant macroscopic Van der Waals interaction.
Different patterns of nonlocal Bose-condensate of standing gravitational waves in vacuum represent the field-informational copy of local Bose- condensate of the effectons of condensed matter.
Very important role of proposed here distant resonant VIBRO-GRAVITATIONAL INTERACTION (VGI) in elementary acts of perception and memory can be contributed by coherent primary librational water effectons in microtubules of the nerve cells (see paper ”Hierarchic Model of Consciousness” in URL: http://www.karelia.ru/~alexk and http://arXiv.org/abs/physics/0003045).
> References
>
> . Kaivarainen A. Hierarchic Concept of Matter and Field. Water, biosystems and elementary particles. New York, 1995, pp. 485.
>
> . Kaivarainen A. New Hierarchic Theory of Matter General for Liquids and Solids: dynamics, thermodynamics and mesoscopic structure of water and ice
>
> (see URL: http://www.karelia.ru/~alexk)
>
> . Kaivarainen A. Hierarchic Concept of Matter, General for Liquids and Solids: Water and ice (see Proceedings of the Second Annual Advanced Water Sciences Symposium, October 4-6, 1966, Dallas, Texas.
>
> . Eisenberg D., Kauzmann W. The structure and properties of water. Oxford University Press, Oxford, 1969.
>
> . Frontas’ev V.P., Schreiber L.S. J. Struct. Chem. (USSR$`)`$6, (1966), 512.
>
> . Kikoin I.K. (Ed.) Tables of physical values. Atomizdat, Moscow, 1976 (in Russian).
>
> . Einstein A. Collection of works. Nauka, Moscow, 1965.
>
> . Vuks M.F. Light scattering in gases, liquids and solutions. Leningrad University Press, Leningrad. 1977.
>
> . Vuks M.F. Electrical and optical properties of molecules and condensed matter. Leningrad University Press, Leningrad, 1984.
>
> . Theiner O., White K.O. J.Opt. Soc.Amer. $`\mathrm{\hspace{0.33em}1969},\mathrm{\hspace{0.17em}59},\mathrm{\hspace{0.17em}181}`$.
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> . Wertheim G.K. Mössbauer effect. Academic Press, N.Y. and London. 1964.
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> . Shpinel V.C. Gamma-ray resonance in crystals. Nauka, Moscow, 1969.
>
> . Singvi K., Sielander A. In book: Mössbauer effect. Ed. Kogan Yu. Moscow 1962.
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# Projective generation and smoothness of congruences of order 1Mathematical Subject Classification: 14M15, 14J
## 1 About Congruences of Order $`1`$
Let $`\text{IP}^N`$ be the $`N`$–dimensional projective complex space and $`G(r,N)`$ be the Grassmannian of $`r`$–dimensional subspaces in $`\text{IP}^N`$. We begin compiling some basic facts about $`G(r,N)`$ (see for more details):
It is well known that the cohomology group $`H^i(G(r,N),)`$ is $`0`$ if $`i[0,2(Nr)(r+1)]`$ and the direct sum
$$H^{}(G(r,N),)=\underset{i}{}H^i(G(r,N),)$$
is a graded ring with the cup–product $``$. We can assign the cohomology class $`[X]`$ to each closed subset $`XG(r,N)`$, being satisfied the next properties:
* Two algebraically equivalent subvarieties have the same cohomology class.
* If $`XG(r,N)`$ is irreducible $`((r+1)(Nr)p)`$–dimensional, then $`[X]H^{2p}(G(r,N),)`$.
* For each two subvarieties $`X,YG(r,N)`$, $`[XY]=[X][Y]`$.
Let $`A_0A_1\mathrm{}A_r`$ be a strictly increasing sequence of subspaces of $`\text{IP}^N`$ , $`a_i:=dimA_i`$, $`i=0,\mathrm{},r`$, and $`\mathrm{\Omega }(A_0,\mathrm{},A_r)`$ denote the next subvariety of the grassmannian:
$$\mathrm{\Omega }(A_0,\mathrm{},A_r)=\{\sigma G(r,N)/dim(\text{IP}^r(\sigma )A_i)i,i=0,\mathrm{},r\}$$
Such a subvariety is called a Schubert cycle in $`G(r,N)`$. Regarding $`G(r,N)`$ as a projective variety through the Plücker’s embedding, Schubert cycles can be obtained cutting it with some linear subspaces; for instance, the Schubert cycle $`\mathrm{\Omega }_rG(1,3)`$ parametrizing the lines in $`\text{IP}^3`$ meeting a given one $`r`$ consists of the intersection $`\mathrm{\Omega }_r=G(1,3)T_{G(1,3),[r]}`$.
Given two Schubert cycles $`\mathrm{\Omega }(A_0,\mathrm{},A_r)`$, $`\mathrm{\Omega }(B_0,\mathrm{},B_r)`$, they are algebraically equivalent when $`dimA_i=dimBi`$ for every $`i`$; thus the cohomology class of $`\mathrm{\Omega }(A_0,\mathrm{},A_r)`$ depends only of the integers $`a_i:=dimA_i`$, $`i=0,\mathrm{},r`$, and we denote it by $`\mathrm{\Omega }(a_0,\mathrm{},a_r)`$. They play an important role in the cohomology of $`G(r,N)`$, in fact:
###### Theorem 1.1
$`H^{}(G(r,N),)`$ is a free abelian group with basis
$$\{\mathrm{\Omega }(a_0,\mathrm{},a_r)/\mathrm{\hspace{0.17em}0}a_0<a_1<\mathrm{}<a_rN\}.$$
Moreover, each $`\mathrm{\Omega }(A_0,\mathrm{},A_r)`$ is irreducible of dimension $`{\displaystyle \underset{i=0}{\overset{r}{}}}(a_ii)`$, so we have the next corollary:
###### Corollary 1.2
For every integer $`p`$, $`H^{2p}(G(r,N),)`$ is a free abelian group generated by the elements $`\mathrm{\Omega }(a_0,\mathrm{},a_r)`$ with $`p=[(r+1)(Nr){\displaystyle \underset{i=0}{\overset{r}{}}}(a_ii)]`$. Each $`H^{2p+1}(G(r,N),)`$ is $`0`$.
###### Example 1.3
By the above, we can calculate the basis of the cohomology groups of codimension $`2`$ and dimension $`2`$:
$$\begin{array}{c}H^{2(Nr)(r+1)4}(G(r,N),)=\mathrm{\Omega }_1:\mathrm{\Omega }(0,1,\mathrm{}r1,r+2),\hfill \\ \mathrm{\Omega }_2:\mathrm{\Omega }(0,1,\mathrm{}r2,r,r+1)\hfill \\ H^4(G(r,N),)=\mathrm{\Omega }^1:\mathrm{\Omega }(Nr2,Nr+1,Nr+2,\mathrm{}N),\hfill \\ \mathrm{\Omega }^2:\mathrm{\Omega }(Nr1,Nr,Nr+2\mathrm{}N)\hfill \end{array}$$
Denoting by $`L^i\text{IP}^n`$ a fixed generic $`i`$–dimensional subspace, the elements above correspond to the classes of the Schubert cycles:
$$\begin{array}{c}\mathrm{\Omega }_1=[\{\sigma G(r,N)/L^{r1}\text{IP}(\sigma )L^{r+2}\}]\hfill \\ \mathrm{\Omega }_2=[\{\sigma G(r,N)/L^{r2}\text{IP}(\sigma )L^{r+1}\}]\hfill \\ \mathrm{\Omega }^1=[\{\sigma G(r,N)/\text{IP}(\sigma )L^{Nr2}\mathrm{}\}]\hfill \\ \mathrm{\Omega }^2=[\{\sigma G(r,N)/dim(\text{IP}(\sigma )L^{Nr})1\}]\hfill \end{array}$$
and we can compute the products $`\mathrm{\Omega }_i\mathrm{\Omega }^j`$ as the classes of the intersections of such Schubert cycles; thus $`\mathrm{\Omega }_i\mathrm{\Omega }^j=0`$ if $`ij`$ and $`\mathrm{\Omega }_i\mathrm{\Omega }^j=1`$ if $`i=j`$. Hence the cohomology class of a $`2`$–dimensional irreducible subvariety $`\mathrm{\Sigma }G(r,N)`$ is determined by the number of points of intersection of $`\mathrm{\Sigma }`$ with a generic Schubert cycle in the class of $`\mathrm{\Omega }^1`$ (the number of elements of $`\mathrm{\Sigma }`$ cutting a generic $`(Nr2)`$–dimensional space), and with another one in the class of $`\mathrm{\Omega }_2`$ (the number of elements of $`\mathrm{\Sigma }`$ meeting a generic $`(Nr)`$–dimensional space in lines); these numbers are called, respectively, order ($`ord(\mathrm{\Sigma })`$) and class ($`cl(\mathrm{\Sigma })`$) of $`\mathrm{\Sigma }`$. The pair ($`ord(\mathrm{\Sigma },cl(\mathrm{\Sigma })`$) is called bidegree of $`\mathrm{\Sigma }`$. Moreover, regarding $`\mathrm{\Sigma }`$ as a projective variety through the Plücker embedding of the grassmannian, the degree of $`\mathrm{\Sigma }`$ is exactly $`ord(\mathrm{\Sigma })+cl(\mathrm{\Sigma })`$.
We wish to investigate $`2`$–dimensional families of $`r`$–dimensional subspaces in $`\text{IP}^N`$, called classically congruences. From now on, $`\mathrm{\Sigma }`$ will denote a $`2`$–dimensional reduced and irreducible closed subvariety of $`G(r,N)`$ and $`V(\mathrm{\Sigma })`$ its projective realization,
(1)
$$V(\mathrm{\Sigma }):=\underset{\sigma \mathrm{\Sigma }}{}\text{IP}^r(\sigma )$$
The objective of this work is the classification of congruences of order $`1`$; by the next lemma we will be reduced to studying the case $`N=r+2`$.
###### Lemma 1.4
Let $`\mathrm{\Sigma }G(r,N)`$ be a congruence. If $`ord(\mathrm{\Sigma })=1`$, then there exist a $`(r+2)`$–dimensional projective subspace $`A\text{IP}^N`$ such that $`\text{IP}^r(\sigma )A`$ for all $`\sigma \mathrm{\Sigma }`$.
Proof: Since $`\mathrm{\Sigma }`$ is irreducible, $`V(\mathrm{\Sigma })`$ is irreducible too. Let $`B\text{IP}^N`$ be a generic $`(Nr2)`$–dimensional subspace. As $`ord(\mathrm{\Sigma })=1`$ we have $`BV(\mathrm{\Sigma })=B\text{IP}^r(\sigma )\mathrm{}`$ for only one $`\sigma \mathrm{\Sigma }`$. It follows that $`dimV(\mathrm{\Sigma })=r+2`$ and $`\mathrm{deg}V(\mathrm{\Sigma })=1`$. We finish taking $`A:=V(\mathrm{\Sigma })\text{IP}^{r+2}`$.
According to this lemma, we will only consider $`\mathrm{\Sigma }G(r,r+2)`$, using the next notation for the incidence variety:
###### Definition 1.5
Let $`\mathrm{\Sigma }G(r,r+2)`$ be a congruence. A point $`P\text{IP}^{r+2}`$ is called fundamental for $`\mathrm{\Sigma }`$ if $`dim(p^1(P))1`$. The congruence $`\mathrm{\Sigma }`$ is called degenerate iff $`dimV(\mathrm{\Sigma })<r+2`$, that is iff every point in $`V(\mathrm{\Sigma })`$ is fundamental. For instance, a congruence of order $`1`$ is nondegenerate.
Given a congruence $`\mathrm{\Sigma }`$ in $`\text{IP}^{r+2}`$ and a projective subspace $`L\text{IP}^{r+2}`$, we can obtain a congruence in $`L`$ cutting it with each element of $`\mathrm{\Sigma }`$. The next results provide some properties of such sections.
###### Lemma 1.6
Let $`\mathrm{\Sigma }G(r,r+2)`$ be a congruence of order $`ord(\mathrm{\Sigma })`$ and $`L\text{IP}^{r+2}`$ be a $`(l+2)`$–dimensional generic subspace, $`l1`$; then $`\mathrm{\Sigma }`$ does not meet $`\mathrm{\Omega }_L:=\{\sigma G(r,r+2)/dimL\text{IP}^r(\sigma )l+1\}`$, which is a codimension $`l+2`$ Schubert cycle; thus the projection
$$\rho _L:\mathrm{\Sigma }G(r,r+2)G(l,L),\pi _L(\sigma ):=[\text{IP}^r(\sigma )L]$$
is regular and has as image a surface $`\mathrm{\Sigma }_L`$ such that $`ord(\mathrm{\Sigma }_L)=ord(\mathrm{\Sigma })`$, $`cl(\mathrm{\Sigma }_L)=cl(\mathrm{\Sigma })`$.
Proof: Consider the incidence variety
$$:=\{(L,\sigma )/dim(L\text{IP}^r(\sigma ))l+1\}G(l+2,r+2)\times \mathrm{\Sigma }$$
with projections $`p_1`$ and $`q_1`$. Given $`\sigma \mathrm{\Sigma }`$, $`q_1^1(\sigma )\{LG(l+2,r+2)/dim(\text{IP}^r(\sigma )L)l+1\}`$ is a codimension $`l+2`$ Schubert cycle (because its cohomology class is $`\mathrm{\Omega }(r1l,\mathrm{},r1,r,r+2)`$), so $`dim=dim\mathrm{\Sigma }+dimq_1^1(\sigma )=2+dimG(l+2,\text{IP}^{r+2})l2=dimG(l+2,\text{IP}^{r+2})1`$. Therefore $`p_1()G(l+2,r+2)`$ is a proper closed subset, and if $`LG(l+2,r+2)p_1()`$, then $`dim(L\text{IP}^r(\sigma ))=l`$ for all $`\sigma \mathrm{\Sigma }`$. It follows that the map $`\rho _L`$ defined above is regular. Furthermore it is birrational: a generic point $`P\text{IP}^{r+2}`$ is contained in a finite number of elements of the congruence, $`\text{IP}^r(\sigma _1),\mathrm{},\text{IP}^r(\sigma _n)`$; a generic $`l`$–dimensional subspace through $`P`$, $`W\text{IP}^r(\sigma _1)`$ will not be contained in another $`\text{IP}^r(\sigma _i)`$; so taking $`L\text{IP}^r(\sigma _1)=:W`$, we will have $`L\text{IP}^r(\sigma _1)=WL\text{IP}^r(\sigma _2),\mathrm{},L\text{IP}^r(\sigma _n)`$; thus $`\rho _L^1(W)=\{\sigma _1\}`$, so $`\mathrm{deg}\rho _L=1`$.
By construction, the equality $`ord(\mathrm{\Sigma }_L)=ord(\mathrm{\Sigma })`$ holds; by the birrationality of $`\rho _L`$ $`\mathrm{deg}(\mathrm{\Sigma }_L)=\mathrm{deg}(\mathrm{\Sigma })`$ and so $`cl(\mathrm{\Sigma }_L)=cl(\mathrm{\Sigma })`$ holds too.
###### Remark 1.7
Furthermore, regarding $`G(r,r+2)`$ and $`G(l,l+2)`$ as projective varieties through the corresponding Plücker embeddings, the map $`\rho _L`$ is easily shown to be the projection from the space generated by the Schubert cycle $`\{\sigma G(r,r+2)/dim(\text{IP}^r(\sigma )L)2\}`$.
###### Lemma 1.8
Let $`\mathrm{\Sigma }G(r,r+2)`$ be a nondegenerate congruence, and $`F(\mathrm{\Sigma })\text{IP}^{r+2}`$ its fundamental locus. Then $`dimF(\mathrm{\Sigma })r`$.
Proof: Suppose the lemma were false. We could find an $`(r+1)`$–dimensional irreducible component $`CF(\mathrm{\Sigma })`$. Considering the restriction $`p^1(C)_\mathrm{\Sigma }C\text{IP}^{r+2}`$, we would have $`dimp^1(C)=r+2`$, and so $`p^1(C)=_\mathrm{\Sigma }`$. Consequently $`\mathrm{\Sigma }`$ would be degenerate.
###### Lemma 1.9
Under the hypotheses of (1.6) and with the above notation, it is verified the equality $`F(\mathrm{\Sigma })L=F(\mathrm{\Sigma }_L)`$.
Proof: By construction $`F(\mathrm{\Sigma })LF(\mathrm{\Sigma }_L)`$. If $`PF(\mathrm{\Sigma })LF(\mathrm{\Sigma }_L)`$, $`P`$ would be contained in an irreducible family of elements of $`\mathrm{\Sigma }`$ with the same trace on $`L`$. An $`l`$–dimensional subspace contained in infinite elements of $`\mathrm{\Sigma }`$ is called fundamental. Suppose $`\mathrm{\Sigma }`$ is nondegenerate (the degenerate case is trivial). It is sufficient to show that, if $`D`$ is an irreducible component of $`F(\mathrm{\Sigma })`$ and if $`L`$ is generic, $`DL`$ does not contain fundamental subspaces. On the contrary, suppose $`DL`$ contains fundamental subspaces for all $`L`$; by (1.8) $`dimDr`$, so $`D`$ would be an $`r`$–dimensional projective space whose $`l`$–dimensional subspaces are all fundamental. Considering the incidence variety:
$$𝒥:=\{(l,\sigma )/l\text{IP}^r(\sigma )D\}G(l,D)\times \mathrm{\Sigma }$$
with projections $`p_2`$ and $`q_2`$, we have $`dim𝒥=2+dimq_2^1(\sigma )2+(l+1)(r1l)`$, because $`q_2^1(\sigma )G(1,\text{IP}^r(\sigma )D)`$. Thus $`dim𝒥2+(l+1)(r1l)(l+1)(rl)=dimG(l,D)`$, so either $`p_2`$ is not surjective, or it has generically finite fibers. In both cases, the generic $`l`$–dimensional subspace is not fundamental, a contradiction.
###### Definition 1.10
A point $`P\text{IP}^r(\sigma )`$ is called fixed point of the congruence $`\mathrm{\Sigma }`$ when $`P\text{IP}^r(\sigma )`$ for all $`\sigma \mathrm{\Sigma }`$.
###### Remark 1.11
The locus of fixed points of a congruence $`\mathrm{\Sigma }`$ is a linear subspace that we will denote by $`T(\mathrm{\Sigma })\text{IP}^{r+2}`$, and $`\mathrm{\Sigma }`$ lies on the Schubert’s cycle:
$$\mathrm{\Omega }_{T(\mathrm{\Sigma })}:=\{\sigma G(r,r+2)/\text{IP}^r(\sigma )T(\mathrm{\Sigma })\}$$
Being $`L\text{IP}^{r+2}`$ a generic complementary subspace of $`T(\mathrm{\Sigma })`$, $`k+2:=dimL=r+2dimT(\mathrm{\Sigma })1`$, the maps
(2)
$$\begin{array}{c}\rho _L:\sigma \mathrm{\Omega }_{T(\mathrm{\Sigma })}[\text{IP}^r(\sigma )L]G(k,k+2)\\ \rho ^{T(\mathrm{\Sigma })}:\tau G(k,k+2)[\text{IP}^k(\tau )+T(\mathrm{\Sigma })]\mathrm{\Omega }_{T(\mathrm{\Sigma })}\end{array}$$
provide an isomorphism $`G(k,k+2)\mathrm{\Omega }_{T(\mathrm{\Sigma })}G(r,r+2)`$. Clearly, the congruence $`\rho _L(\mathrm{\Sigma })\mathrm{\Sigma }`$ has no fixed points; roughly speaking, $`\rho _L(\mathrm{\Sigma })`$ is the same surface as $`\mathrm{\Sigma }`$ living in a grassmannian of lesser dimension. Therefore it suffices to study congruences without fixed points.
## 2 Focal Locus of a Congruence in $`G(r,r+2)`$
###### Definition 2.1
Let $`\mathrm{\Sigma }G(r,r+2)`$ be a congruence. Being $`\sigma \mathrm{\Sigma }`$ smooth, a point $`P\text{IP}^r(\sigma )`$ is called focal if $`(dp)_{(P,\sigma )}`$ is not injective. If $`\mathrm{ker}((dp)_{(P,\sigma )})v`$ for some $`vq^{}T_{\mathrm{\Sigma },\sigma }`$, $`v0`$, we say that $`P`$ is focal for the direction $`v`$.
###### Proposition 2.2
Let $`\mathrm{\Sigma }G(r,r+2)`$ be a congruence. If $`P\text{IP}^{r+2}`$ is a fundamental point of $`\mathrm{\Sigma }`$, then $`P`$ is focal for every $`\sigma \mathrm{\Sigma }`$ smooth with $`P\text{IP}^r(\sigma )`$. Furthermore, if $`P\text{IP}^r(\sigma )`$ is focal and if $`ord(\mathrm{\Sigma })=1`$, then $`P`$ is fundamental.
Proof: Being $`P`$ fundamental, there is an irreducible curve $`C\mathrm{\Sigma }`$ such that $`P\sigma `$ for all $`\sigma C`$. Taking a smooth point $`\sigma C`$ and regarding the curve $`\{P\}\times C_\mathrm{\Sigma }`$, we have $`p(\{P\}\times C)=P`$. Assuming that $`\sigma `$ is smooth in $`\mathrm{\Sigma }`$ we can write $`(dp)_{(P,\sigma )}(T_{\{P\}\times C,(P,s)})=0`$. Therefore $`P\text{IP}^r(\sigma )`$ is focal.
For the second part we use that $`ord(\mathrm{\Sigma })=\mathrm{deg}(p)`$. If $`ord(\mathrm{\Sigma })=1`$, the map $`p`$ is dominant and $`p(_\mathrm{\Sigma })=\text{IP}^{r+2}`$ is normal. Hence $`\mathrm{deg}(p)=_{p(x)=P}m_x(p)`$ for every nonfundamental point $`P`$, being $`m_x(p):=dim_{\text{C}\text{ }\text{ }}(𝒪_{_\mathrm{\Sigma },x}/p^{}\text{m}_P)`$ for $`\text{m}_P`$ the maximal ideal of the local ring $`𝒪_{\text{IP}^{r+2},P}`$. Thus, if $`ord(\mathrm{\Sigma })=1`$, $`m_x(p)>1`$ forces $`P`$ to be fundamental. The proof is completed with the next lemma.
###### Lemma 2.3
Let $`\mathrm{\Sigma }G(r,r+2)`$ be a congruence and denote $`x:=(P,\sigma )_\mathrm{\Sigma }`$. $`x`$ is focal if and only if $`m_x(p)2`$.
Proof: Regarding the diagram
it is easily seen that $`dim_{\text{C}\text{ }\text{ }}(𝒪_{_\mathrm{\Sigma },x}/p^{}\text{m}_P)=dim_{\text{C}\text{ }\text{ }}(\text{m}_x/p^{}\text{m}_P)+1`$. Consider the next diagram:
By Nakayama’s lemma $`p^{}`$ is surjective iff $`(dp)_x^{}`$ is surjective, that is iff $`(dp)_x`$ is injective. Hence, $`x`$ is focal iff $`dim_{\text{C}\text{ }\text{ }}(\text{m}_x/p^{}\text{m}_P)>0m_x(p)>1`$.
We continue with some basic facts about the focal locus of a congruence in $`G(r,r+2)`$; we refer the reader to for more details in the case $`r=2`$, but the proofs also work for all $`r`$. Let $`\mathrm{\Sigma }G(r,r+2)`$ be a congruence and $`\sigma \mathrm{\Sigma }`$ an smooth point. For each $`P\text{IP}^r(\sigma )`$ there is an isomorphism $`T_{_\mathrm{\Sigma },(P,\sigma )}T_{\mathrm{\Sigma },\sigma }T_{\text{IP}^r(\sigma ),P}`$.
The focal locus at $`\text{IP}^r(\sigma )`$ can be calculated through the Characteristic Map of Kodaira–Segre–Spencer:
(3)
$$\chi :T_{\mathrm{\Sigma },\sigma }H^0(𝒩_{\text{IP}^r(\sigma ),\text{IP}^{r+2}})H^0(𝒪_{\text{IP}^r(\sigma )}(1))H^0(𝒪_{\text{IP}^r(\sigma )}(1))$$
being $`\chi (v)=0`$ the equations of the focal locus of $`\mathrm{\Sigma }`$ at $`\text{IP}^r(\sigma )`$ in the direction $`vT_{\mathrm{\Sigma },\sigma }`$. An easy computation shows that:
* The focal locus at $`\text{IP}^r(\sigma )`$ in a direction $`v`$ is a projective subspace of dimension $`r2`$ or $`r1`$ and we denote it by $`Z(\chi (v))`$. Taking a base $`\{v_1,v_2\}`$ in $`T_{\mathrm{\Sigma },\sigma }`$, whose images through $`\chi `$ are $`(f_{11},f_{12})`$ and $`(f_{21},f_{22})`$ respectively, $`Z(\chi (\lambda v_1+\mu v_2))`$ will be the set of points $`P\text{IP}^r(\sigma )`$ satisfying the equations:
$$\begin{array}{c}\lambda f_{11}(P)+\mu f_{21}(P)=0\\ \lambda f_{12}(P)+\mu f_{22}(P)=0\end{array}\}$$
* The focal locus at $`\text{IP}^r(\sigma )`$ is a quadric $`Q(\sigma )=_{vT_{\mathrm{\Sigma },\sigma }}Z(\chi (\sigma ))`$, and it is given by the equation
(4)
$$det\left(\begin{array}{cc}f_{11}(P)& f_{21}(P)\\ f_{12}(P)& f_{22}(P)\end{array}\right)=0$$
It is reducible if and only if there exist a direction $`vT_{\mathrm{\Sigma },\sigma }`$ such that $`dimZ(\chi (v))=r1`$. Such directions are called developable, and the number of them is fixed in an open subset of $`\mathrm{\Sigma }`$; this number can be $`0`$, $`1`$, $`2`$ (distinct or coincident) or $`\mathrm{}`$ (if all direction is developable).
* If $`Q(\sigma )`$ is irreducible, it contains a $`1`$–dimensional family of $`(r2)`$–dimensional projective subspaces. Therefore $`rankQ(\sigma )4`$ (see for instance , book V, page 100). $`Q(\sigma )`$ will be a cone with vertex a subspace $`C(\sigma )\text{IP}^r(\sigma )`$ over an smooth conic or an smooth quadric surface. In both cases the focal loci in each direction are subspaces of maximal dimension contained in $`Q(\sigma )`$, and so the intersection of two of them is $`C(\sigma )`$. Hence $`C(\sigma )`$ is the set of points focal for all direction.
* Let $`\mathrm{\Delta }G(r,r+2)`$ be a reduced and irreducible curve. We say that $`\mathrm{\Delta }`$ is a developable system if the focal locus in each $`\text{IP}^r(\sigma )`$, $`\sigma \mathrm{\Delta }`$ has dimension $`r1`$. Following , $`\mathrm{\Delta }`$ parametrizes a family of the next kind: Cone with vertex a subspace $`\text{IP}^k`$ over the family of $`(rk1)`$–osculating spaces to an irreducible curve $`𝒞\text{IP}^{rk+1}`$, $`k=1,0,\mathrm{},r1`$. For instance, a $`1`$–dimensional family of $`\text{IP}^r`$’s lying in an $`(r+1)`$–dimensional space is always developable.
* A nondegenerate congruence $`\mathrm{\Sigma }G(r,r+2)`$ has reducible generic focal quadric iff by the generic $`\sigma \mathrm{\Sigma }`$:
+ is passing $`1`$ developable system, or
+ are passing $`2`$ distinct developable systems, or
+ are passing $`2`$ coincident developable systems, or
+ are passing infinite developable systems.
Being $`\mathrm{\Sigma }G(r,r+2)`$, denote $`F_1(\mathrm{\Sigma })\text{IP}^{r+2}`$ its focal variety, that is, the projective realization of the family of focal quadrics, which is defined in an open set $`U\mathrm{\Sigma }`$, $`_1(U)U`$. In general, the fundamental locus $`F(\mathrm{\Sigma })`$ is a subvariety of $`F_1(\mathrm{\Sigma })`$ (see ). In the case $`ord(\mathrm{\Sigma })=1`$, (2.2) clearly forces the equality $`F_1(\mathrm{\Sigma })=F(\mathrm{\Sigma })`$. Moreover, by (1.8) $`dimF(\mathrm{\Sigma })r`$, $`dimF(\mathrm{\Sigma })<dimF_1(U)=r+1`$. The generic focal quadric can be irreducible or not, so either $`F(\mathrm{\Sigma })`$ is irreducible, or it is the projective realization of two $`2`$–dimensional families of $`\text{IP}^{r1}`$’s, having at least $`2`$ irreducible components.
###### Example 2.4
Let $`C_1`$, $`C_2`$ be two irreducible curves in $`\text{IP}^3`$. We write $`\mathrm{\Sigma }(C_1,C_2)`$ for the set of lines joining points of both curves, and $`\mathrm{\Sigma }(C_1)`$ the family of secant lines to $`C_1`$. Clearly the focal locus of these congruences contains the base curves (furthermore, they will be exactly the focal locus in some cases). If $`C_1`$ and $`C_2`$ are not coplanar, then $`\mathrm{\Sigma }(C_1,C_2)`$ is nondegenerate, and if $`C_1`$ is not a plane curve, then $`\mathrm{\Sigma }(C_1)`$ is nondegenerate. We will show later when these congruences have order $`1`$.
In the next two results we examine the focal locus of a congruence of order $`1`$, concretely its irreducible components and degree.
###### Proposition 2.5
If the focal locus of a congruence $`\mathrm{\Sigma }`$ of order $`1`$ is reducible, $`F(\mathrm{\Sigma })=C_1C_2`$, then both components have dimension $`r`$ and one of them is linear.
Proof: We have shown that the number of irreducible components of $`F(\mathrm{\Sigma })`$ is at most two, with dimensions lesser than or equal to $`r`$ by (1.8). Let $`C_1`$, $`C_2`$ be the irreducible components of $`F(\mathrm{\Sigma })`$. It is sufficient to show that if $`F(\mathrm{\Sigma })`$ is reducible, its section with the generic $`3`$–dimensional space consists of a line and an irreducible curve. Thus using (1.6) and (1.9), we can suppose $`r=1`$. If $`dimC_i=0`$ for some $`i`$, then $`\mathrm{\Sigma }`$ would be the family of lines passing by $`C_i`$, whose focal locus is the point $`C_i`$, so $`C_1`$ and $`C_2`$ are irreducible curves. Every line of $`\mathrm{\Sigma }`$ contains a point in each curve, so $`\mathrm{\Sigma }=\mathrm{\Sigma }(C_1,C_2)`$, and $`\mathrm{\Sigma }\mathrm{\Sigma }(C_1)`$ and $`\mathrm{\Sigma }\mathrm{\Sigma }(C_2)`$ are closed subsets in $`\mathrm{\Sigma }`$ (if $`\mathrm{\Sigma }=\mathrm{\Sigma }(C_1)`$, then every line of $`\mathrm{\Sigma }`$ would contain at least three focal points, two of $`C_1`$ and one of $`C_2`$, so the focal quadric at each line must be the whole line; equivalently, the congruence would be degenerate, so $`ord(\mathrm{\Sigma })=0`$, false.). Being $`K\text{IP}^3`$ a generic plane, it contains $`cl(\mathrm{\Sigma })`$ generators of $`\mathrm{\Sigma }`$, none belonging to $`\mathrm{\Sigma }(C_1)`$ nor $`\mathrm{\Sigma }(C_2)`$. Hence every generator in $`K`$ meets $`C_1`$ and $`C_2`$ exactly one time. If $`\mathrm{deg}C_1,\mathrm{deg}C_22`$, we could take $`l=P_1P_2`$, $`l^{}=Q_1Q_2`$, $`P_i,Q_iC_i`$, $`P_iQ_i`$. The point $`P=ll^{}`$ could not lie on $`C_1`$, $`C_2`$, but as $`ord(\mathrm{\Sigma })=1`$, $`P`$ should be fundamental, a contradiction.
###### Proposition 2.6
If the focal locus $`F(\mathrm{\Sigma })`$ of a congruence of order $`1`$ is irreducible, one of these possibilities is true:
1. $`F(\mathrm{\Sigma })`$ is an $`(r1)`$–dimensional projective space and $`\mathrm{\Sigma }`$ is the family of $`\text{IP}^r`$’s containing $`F(\mathrm{\Sigma })`$.
2. $`dimF(\mathrm{\Sigma })=r`$, and then:
1. either $`F(\mathrm{\Sigma })`$ is a projective space,
2. or $`F(\mathrm{\Sigma })`$ is an irreducible variety of degree $`3`$.
Proof: By hypothesis $`F(\mathrm{\Sigma })`$ is irreducible, containing the focal quadrics. If $`dimF(\mathrm{\Sigma })r1`$, then $`F(\mathrm{\Sigma })=Q(\sigma )\text{IP}^r(\sigma )`$ for all $`\sigma \mathrm{\Sigma }`$. Therefore $`F(\mathrm{\Sigma })`$ is an $`(r1)`$–dimensional projective space.
Suppose $`dimF(\mathrm{\Sigma })=r`$. Using (1.6) and (1.9) we are reduced to the case $`r=1`$. If $`\mathrm{deg}F(\mathrm{\Sigma })=n>1`$, then $`\mathrm{\Sigma }=\mathrm{\Sigma }(F(\mathrm{\Sigma }))`$: let $`\sigma \mathrm{\Sigma }`$ be generic and $`PF(\mathrm{\Sigma })\text{IP}^1(\sigma )\mathrm{}`$; since $`n>1`$, a generic plane $`\pi \text{IP}^1(\sigma )`$ contains another point $`QF(\mathrm{\Sigma })`$, $`QP`$. A $`1`$–dimensional family of lines of $`\mathrm{\Sigma }`$ is passing by $`Q`$, forming a cone of degree $`1`$; one of its generators $`\text{IP}^1(\tau )`$ lies on $`\pi `$, cutting $`\text{IP}^1(\sigma )`$ in a fundamental point $`RF(\mathrm{\Sigma })`$, and $`\text{IP}^1(\sigma )=PR`$.
Since $`\mathrm{\Sigma }=\mathrm{\Sigma }(F(\mathrm{\Sigma }))`$, necessarily $`\mathrm{deg}F(\mathrm{\Sigma })3`$ (otherwise, $`F(\mathrm{\Sigma })`$ would be planar, and $`\mathrm{\Sigma }`$ would be degenerate). Suppose, contrary to our claim, that $`\mathrm{deg}F(\mathrm{\Sigma })4`$. The generic plane $`\pi \text{IP}^3`$ cuts $`F(\mathrm{\Sigma })`$ in $`n4`$ points such that any $`3`$ of them are not collinear (by the Trisecant Lemma). Taking four of them, the diagonal points of the square they form must be fundamental; that is, they belong to $`F(\mathrm{\Sigma })`$; hence there are $`3`$ collinear points of $`F(\mathrm{\Sigma })\pi `$, a contradiction.
We study now the case where $`F(\mathrm{\Sigma })`$ (irreducible or not) contains a projective space.
###### Lemma 2.7
Let $`\pi `$ be an $`r`$–dimensional projective space and consider the Schubert’s cycle:
(5)
$$\mathrm{\Omega }_\pi :=\{\sigma G(r,r+2)/dim\text{IP}^r(\sigma )\pi r1\}G(r,r+2)$$
$`\mathrm{\Omega }_\pi `$ is a cone with vertex $`[\pi ]`$ over a Segre Variety of the form $`\text{IP}^r\times \text{IP}^1\text{IP}^{2r+1}`$
Proof: Regard $`H_\pi =\{[H]\text{IP}_{}^{r+2}{}_{}{}^{}/H\pi \}\text{IP}^1`$ and the projection
(6)
$$\varphi :\mathrm{\Omega }_\pi \{[\pi ]\}\pi ^{}\times H_\pi ,\varphi (\sigma ):=(\text{IP}^r(\sigma )\pi ,\text{IP}^r(\sigma )+\pi )$$
whose fibers are $`\varphi ^1(h,H)=\mathrm{\Omega }_{h,H}=\{\sigma G(r,r+2)/h\text{IP}^r(\sigma )H\}\text{IP}^1`$, and $`[\pi ]\mathrm{\Omega }_{h,H}`$. In this way $`\varphi `$ defines $`\mathrm{\Omega }_\pi `$ as a cone.
###### Remark 2.8
Segre variety $`\pi ^{}\times H_\pi `$ can be geometrically thought in the next way: a line $`l\text{IP}^{r+2}`$ disjoint of $`\pi `$ parametrizes the elements of $`H_\pi `$, so the subvariety of the grassmannian $`\mathrm{\Omega }_\pi \{\sigma /\text{IP}^r(\sigma )l\mathrm{}\}`$ is a Segre variety of the form $`\pi ^{}\times H_\pi `$. Moreover, it is the intersection of two Schubert cycles, one of them not containing $`[\pi ]`$, so it is a hyperplane section of $`\mathrm{\Omega }_\pi `$.
###### Proposition 2.9
Let $`\mathrm{\Sigma }G(r,r+2)`$ be a congruence and $`\pi \text{IP}^{r+2}`$ be an $`r`$–dimensional projective subspace such that $`\mathrm{\Sigma }\mathrm{\Omega }_\pi `$. $`\mathrm{\Sigma }`$ has order $`1`$ if and only if $`\varphi ^1(\pi ^{}\times \{H\})\mathrm{\Sigma }`$ is a line for all $`HH_\pi `$. Thus the map
(7)
$$\alpha :\mathrm{\Sigma }\{[\pi ]\}H_\pi \text{IP}^1$$
given by $`\alpha (\sigma )=\text{IP}^r(\sigma )+\pi `$ defines $`\mathrm{\Sigma }`$ as a rational ruled surface.
Proof: Taking $`HH_\pi `$, $`\alpha ^1(H)=\{\sigma \mathrm{\Sigma }/\text{IP}^r(\sigma )H\}`$ is $`1`$–dimensional and each one of its irreducible components is a developable family of $`\text{IP}^r`$’s (because they are contained in $`H`$), whose focal locus has dimension $`r`$. If $`ord(\mathrm{\Sigma })=1`$, then $`dimF(\mathrm{\Sigma })=r`$, and so the focal locus of those families must be $`(r1)`$–dimensional. Equivalently, each one is a family of $`\text{IP}^r(\sigma )`$’s containing an $`(r1)`$–dimensional space in $`H`$; since $`ord(\mathrm{\Sigma })=1`$, each $`H\pi `$ can only contain one of such families (otherwise, every point in $`H`$ would be focal). So $`\alpha ^1(H)\text{IP}^1`$. Conversely, a congruence constructed in such way has clearly order $`1`$.
###### Proposition 2.10
Let $`\mathrm{\Sigma }\mathrm{\Omega }_\pi G(r,r+2)`$ be a congruence of order $`1`$. If $`\mathrm{\Sigma }`$ has no fixed points, then $`cl(\mathrm{\Sigma })r`$.
Proof: Consider
(8)
$$\beta :\mathrm{\Sigma }\{[\pi ]\}\stackrel{\varphi }{}\pi ^{}\times H_\pi \stackrel{p}{}\pi ^{}$$
given by $`\beta (\sigma )=\text{IP}^r(\sigma )\pi `$. In each hyperplane $`H\pi `$, $`\mathrm{\Sigma }`$ consists of a pencil of $`\text{IP}^r`$’s contained in $`H`$ (with base a certain $`(r1)`$–dimensional projective space). If $`P\text{IP}^{r+2}`$ is a fixed point of $`\mathrm{\Sigma }`$, necessarily $`P\pi `$, and so $`\beta (\mathrm{\Sigma })\pi ^{}(P)=\{h\pi ^{}/Ph\}\pi ^{}`$ will be degenerate. Thus $`\mathrm{\Sigma }`$ has fixed points iff $`\beta (\mathrm{\Sigma })\pi ^{}`$ is degenerate. Now, we have two cases:
1. If $`\varphi (\mathrm{\Sigma })`$ is a curve: $`cl(\mathrm{\Sigma })+1=\mathrm{deg}(\mathrm{\Sigma })=\mathrm{deg}(\varphi (\mathrm{\Sigma }))=\mathrm{\#}(\varphi (\mathrm{\Sigma })(h\times H_\pi +\pi ^{}\times H))`$ where $`h\pi ^{}`$ and $`HH_\pi `$. Hence $`cl(\mathrm{\Sigma })+1=\mathrm{\#}(\varphi (\mathrm{\Sigma })(h\times H_\pi ))+\mathrm{\#}(\varphi (\mathrm{\Sigma })(\pi ^{}\times H))\mathrm{deg}(\beta (\mathrm{\Sigma }))+1dim(\beta (\mathrm{\Sigma }))+1=r+1`$.
2. If $`\varphi (\mathrm{\Sigma })`$ is a surface: $`cl(\mathrm{\Sigma })+1=\mathrm{deg}(\mathrm{\Sigma })\mathrm{deg}(\varphi (\mathrm{\Sigma }))`$. The map $`p:\varphi (\mathrm{\Sigma })\beta (\mathrm{\Sigma })`$ consists of projecting from a certain $`\pi ^{}\times H`$, that contains exactly one generator of $`\varphi (\mathrm{\Sigma })`$. Hence $`cl(\mathrm{\Sigma })+1\mathrm{deg}(\beta (\mathrm{\Sigma }))+2dim\beta (\mathrm{\Sigma })1+2=r+1`$.
## 3 Case I: $`F(\mathrm{\Sigma })`$ is $`r`$–dimensional irreducible of degree $`3`$
First examine the case $`r=1`$: by (2.6) $`F(\mathrm{\Sigma })`$ is a nondegenerate curve of degree $`3`$ in $`\text{IP}^3`$, that is a rational normal cubic, and $`\mathrm{\Sigma }=\mathrm{\Sigma }(F(\mathrm{\Sigma }))`$. Conversely, the congruence of secant lines to a rational normal cubic has clearly order $`1`$. Since a generic plane $`\pi \text{IP}^3`$ cuts $`F(\mathrm{\Sigma })`$ in three noncollinear points, $`cl(\mathrm{\Sigma })=3`$, thus $`\mathrm{deg}(\mathrm{\Sigma })=4`$. Moreover, the secant lines through a fixed point $`PF(\mathrm{\Sigma })`$ form a quadric cone, so $`\mathrm{\Sigma }`$ is not a ruled surface (the lines in $`G(1,3)`$ parametrize linear pencils of lines through a point). Such a quadric cone is parametrized by an irreducible conic $`\mathrm{\Sigma }(P)\mathrm{\Omega }_P`$ where $`\mathrm{\Omega }_P`$ is the family of lines through $`P`$. Given two points $`P,QF(\mathrm{\Sigma })`$, $`\mathrm{\Omega }_P\mathrm{\Omega }_Q=[PQ]`$ so $`\mathrm{\Sigma }(P),\mathrm{\Sigma }(Q)=\mathrm{\Omega }_P+\mathrm{\Omega }_Q=T_{G(1,3),[PQ]}`$ that is $`4`$–dimensional and $`T_{G(1,3),[PQ]}\mathrm{\Sigma }=\mathrm{\Sigma }(P)\mathrm{\Sigma }(Q)`$. Hence $`\mathrm{\Sigma }`$ is nondegenerate in $`\text{IP}^5=G(1,3)`$. By Del Pezzo’s Theorem (see for instance page 525), $`\mathrm{\Sigma }`$ must be the Veronese surface.
By (1.6), if $`r>1`$, a generic $`3`$–dimensional subspace $`L\text{IP}^{r+2}`$ provide us a projection $`\rho _L:\mathrm{\Sigma }\mathrm{\Sigma }_L\text{IP}^5`$, where $`\mathrm{\Sigma }_L`$ is a Veronese surface by the above. Since $`\mathrm{\Sigma }_L`$ is normal and $`\rho _L`$ is regular and birrational, $`\mathrm{\Sigma }`$ will be a Veronese surface too.
$`F(\mathrm{\Sigma })`$ is an $`r`$–dimensional irreducible variety of degree $`3`$ in $`\text{IP}^{r+2}`$. A description of such varieties can be found in : Such a variety is, in general, a cone with vertex a subspace $`C`$ over a variety of the same degree and codimension in a projective subspace complementary to $`C`$. If it is smooth, then it is a rational normal scroll. Thus it contains a $`1`$–dimensional family of disjoint $`\text{IP}^{r1}`$’s, which is only possible if $`r3`$.
Therefore $`F(\mathrm{\Sigma })`$ is one of the next varieties:
* the rational normal cubic in $`\text{IP}^3`$ ($`r=1`$),
* the rational normal cubic surface in $`\text{IP}^4`$, projective realization of the Blowing-Up of the plane in a point ($`r=2`$),
* the rational normal cubic $`3`$–fold in $`\text{IP}^5`$, projection of the Blowing-Up of $`\text{IP}^3`$ in a line from one of its points ($`r=3`$),
* a cone with vertex a subspace $`C`$ over one of the varieties described above.
Such varieties are intersection of $`3`$ quadrics, and they contain exactly a $`2`$–dimensional family of $`(r1)`$–dimensional quadrics. These quadrics are smooth iff $`F(\mathrm{\Sigma })`$ is smooth, othercase they are cones with vertex $`C`$. $`\mathrm{\Sigma }`$ must therefore be the congruence of $`\text{IP}^r`$’s containing such quadrics. Conversely, a congruence constructed in this way has clearly order $`1`$, because cutting it with a $`3`$–dimensional space we obtain the family of secant lines to a irreducible nondegenerate curve of degree $`3`$. Summarizing, we have:
###### Theorem 3.1
Every irreducible subvariety of $`\text{IP}^{r+2}`$ of degree $`3`$ and codimension $`2`$ contains a $`2`$–dimensional family of $`(r1)`$–dimensional quadrics, and the family of $`\text{IP}^r`$’s containing those quadrics is a congruence of order $`1`$ and class $`3`$, parametrized by a Veronese surface $`\mathrm{\Sigma }`$. $`\mathrm{\Sigma }`$ lie on the Schubert cycle $`\mathrm{\Omega }_V=\{\sigma G(r,r+2)/\text{IP}^r(\sigma )V\}`$, being $`V`$ the singular locus of all the quadrics, and $`r2dimVr4`$. Conversely, every congruence of order $`1`$ with nonlinear irreducible focal locus is constructed in this way.
###### Remark 3.2
Thus there are only three cases of congruences in the case II without fixed points, and the rest are cones over one of those; that is, using the isomorphism
$$\mathrm{\Omega }_CG(rdimC1,rdimC+1)$$
there are only three different embeddings of the Veronese surface in a grassmannian as a congruence of bidegree $`(1,3)`$: in $`G(1,3)`$, in $`G(2,4)`$ and in $`G(3,5)`$.
## 4 Case II: $`F(\mathrm{\Sigma })`$ is reducible
In this case, for each $`\sigma \mathrm{\Sigma }`$, the focal quadric in $`\text{IP}^r(\sigma )`$ is reducible, $`Q(\sigma )=\text{IP}_1^{r1}(\sigma )\text{IP}_2^{r1}(\sigma )`$, and the components of $`F(\mathrm{\Sigma })`$ are the projective realizations, $`X`$ and $`\pi `$, of the $`\text{IP}_1^{r1}(\sigma )`$’s and of the $`\text{IP}_2^{r1}(\sigma )`$’s, respectively. By (2.5), $`dimX=dim\pi =r`$ and $`\pi `$ is linear.
###### Proposition 4.1
Under the above assumptions, $`X`$ is an $`r`$–dimensional scroll (resp. a curve, for $`r=1`$) such that each one of its generators (resp. points, for $`r=1`$) is contained in exactly one linear pencil of $`\text{IP}^r`$’s of the congruence.
Proof: For $`r=1`$, we have already seen that $`\mathrm{\Sigma }=\mathrm{\Sigma }(X,\pi )`$ in the proof of (2.5). Being $`PX`$ generic, the lines of the congruence containing $`P`$ lie on the plane $`\pi +P`$, hence they are exactly one linear pencil.
If $`r>1`$, the family of $`\text{IP}_1^{r1}(\sigma )`$’s contained in $`X`$ cannot be $`2`$–dimensional; otherwise $`X`$ would be a projective space and, as $`\mathrm{\Sigma }`$ is nondegenerate, $`dimX\pi =(r2)`$, so $`X\pi =\text{IP}_1^{r1}(\tau )\text{IP}_2^{r1}(\tau )`$ for all $`\tau \mathrm{\Sigma }`$; but the family of $`\text{IP}^{r1}`$’s in $`X`$ containing $`X\pi `$ is not $`2`$–dimensional, a contradiction. Therefore $`X`$ is a scroll and each one of its generators $`\text{IP}_1^{r1}(\sigma )`$ is contained in an infinity of elements of $`\mathrm{\Sigma }`$, all of them lying in $`\text{IP}_1^{r1}(\sigma )+\pi `$; by (2.6) and (2.9) such elements are exactly the linear pencil $`\mathrm{\Omega }_{\text{IP}_1^{r1}(\sigma ),\text{IP}_1^{r1}(\sigma )+\pi }=\alpha ^1(\alpha (\sigma ))`$.
###### Example 4.2
The case $`r=1`$: let $`\pi \text{IP}^3`$ be a line, $`X\text{IP}^3`$ an irreducible curve of degree $`n`$ and suppose $`\mathrm{\Sigma }=\mathrm{\Sigma }(\pi ,X)`$ has order $`1`$. By (2.6), the generic plane $`H\pi `$ contains exactly one point of $`X`$ out of $`\pi `$ (the base point of the pencil of lines of $`\mathrm{\Sigma }`$ contained in $`H`$). Thus $`\pi `$ cuts $`X`$ in $`n1`$ points counted with multiplicity, and the projection
$$\gamma :X\{H\text{IP}_{}^{3}{}_{}{}^{}/H\pi \}\text{IP}^1$$
given by $`\gamma (x)=x+\pi `$ shows that $`X`$ is rational and smooth out of $`\pi `$. Moreover, if $`H\pi `$ is generic, it cuts $`X`$ in $`n`$ distinct points, no two collinear with $`X\pi `$, so $`H`$ contains exactly $`n`$ lines of $`\mathrm{\Sigma }`$; therefore $`cl(\mathrm{\Sigma })=\mathrm{deg}(X)=n`$.
$`X`$ is a projection of a rational normal curve $`\mathrm{\Gamma }_n\text{IP}^n`$ from a space $`Vx_1,\mathrm{},x_{n1}`$ , $`x_1,\mathrm{},x_{n1}\mathrm{\Gamma }_n`$, disjoint of $`\mathrm{\Gamma }_n`$. There are two possibilities: either $`X\text{IP}^3`$ is nondegenerate, equivalently $`dimV=n4`$, or $`X`$ is a plane curve, equivalently $`dimV=n3`$. Conversely, if $`X`$ is such a projection of a rational normal curve, and $`\pi `$ is a line containing the projection of $`x_1,\mathrm{},x_{n1}`$, the congruence $`\mathrm{\Sigma }(\pi ,X)`$ has clearly order $`1`$.
###### Remark 4.3
The curve $`X`$ is smooth out of $`\pi `$, but it is not smooth in general. For instance, consider the congruence $`\mathrm{\Sigma }(X,\pi )`$ where $`X`$ is a nodal cubic in a plane $`L`$ and $`\pi `$ is a line meeting $`L`$ in the double point of $`X`$.
###### Example 4.4
The case $`r>1`$: let $`\mathrm{\Sigma }G(r,r+2)`$ be a congruence of order $`1`$ in the case II. In (4.1) we have shown how $`\mathrm{\Sigma }`$ is constructed from its focal locus $`F(\mathrm{\Sigma })=X\pi `$. According to (1.11), assume that $`\mathrm{\Sigma }`$ has no fixed points. By each generator $`\text{IP}_1^{r1}(\sigma )X`$ passes the linear pencil $`\mathrm{\Omega }_{\text{IP}_1^{r1}(\sigma ),\text{IP}_1^{r1}(\sigma )+\pi }`$. Let us denote $`\mathrm{\Sigma }(\pi ,X)=\mathrm{\Sigma }`$. By (2.6), the projection $`\gamma :X\{H\text{IP}_{}^{r+2}{}_{}{}^{}/H\pi \}\text{IP}^1`$ given by $`\gamma (x)=x+\pi `$ forces $`X`$ to be a rational scroll and smooth out of $`\pi `$. In order to compute $`cl(\mathrm{\Sigma })`$, take a generic $`3`$–dimensional subspace $`L\text{IP}^{r+2}`$ and consider the section $`\mathrm{\Sigma }_L`$, whose focal locus is $`F(\mathrm{\Sigma }_L)=(XL)(\pi L)`$. By (1.6) and the above example, we get: $`cl(\mathrm{\Sigma })=cl(\mathrm{\Sigma }_L)=\mathrm{deg}(XL)=\mathrm{deg}(X)`$.
Being $`\mathrm{deg}X=n`$, $`X`$ is a projection of a rational normal scroll $`R\text{IP}^{n+(r1)}`$ of the same degree. The center $`V`$ of such projection ($`p_V`$) is disjoint of $`R`$, and is contained in the space generated by $`p_V^{}(\pi X)`$. Moreover $`dimV=(n+(r1)dimX1)`$, being $`dimX=r+2`$ or $`r+1`$. Conversely, every $`r`$–dimensional rational normal scroll can be projected to $`\text{IP}^{r+2}`$ in this way, obtaining the focal scroll of a congruence of order $`1`$ in the case II without fixed points.
The next proposition provides a criterion for the smoothness of a congruence in the case II.
###### Theorem 4.5
Let $`\mathrm{\Sigma }=\mathrm{\Sigma }(\pi ,X)G(r,r+2)`$ be a congruence in the case II. $`\mathrm{\Sigma }`$ is a rational ruled surface of degree $`\mathrm{deg}(\mathrm{\Sigma })=\mathrm{deg}(X)+1`$ contained in $`\mathrm{\Omega }_\pi =\{\sigma G(r,r+2)/dim(\text{IP}^r(\sigma )\pi )r1\}`$. If $`\mathrm{\Sigma }`$ has no fixed points, by $`[\pi ]`$ are passing at most $`nr`$ generators of $`\mathrm{\Sigma }`$ (exactly $`n1`$ if $`r=1`$) counted with multiplicity. The singular locus of $`\mathrm{\Sigma }`$ consists of, at most, $`[\pi ]`$ and the multiple generators of $`\mathrm{\Sigma }`$ passing by it. Hence $`\mathrm{\Sigma }`$ is smooth if and only if $`\pi X`$ contains at most one generator of $`X`$. In particular, if $`nr+1`$, $`\mathrm{\Sigma }`$ is always smooth, and if $`r=1`$, then $`\mathrm{\Sigma }`$ is smooth iff $`n2`$.
Proof: If $`r=1`$, consider $`\gamma :X\{H\text{IP}_{}^{3}{}_{}{}^{}/H\pi \}\text{IP}^1`$ defined in (4.2). It is an isomorphism in $`X\{[\pi ]\}`$ and puts in correspondence the $`n1`$ points (counted with multiplicity) of $`\pi X`$ with the planes by $`\pi `$ containing pencils of lines of $`\mathrm{\Sigma }`$ with base point in $`\pi X`$; equivalently $`[\pi ]`$ belongs to those pencils. Thus $`[\pi ]`$ will be contained in $`n1`$ generators of $`\mathrm{\Sigma }`$ counted with multiplicity.
If $`r>1`$ and $`\mathrm{\Sigma }`$ has no fixed points, consider the projection $`p_V:RX`$ given in (4.4), whose center $`V`$ is contained in $`p_V^{}(\pi X)`$. Clearly $`p_V^{}(\pi X)HF`$ where $``$ denotes linear equivalence, $`H`$ the hyperplane section of $`R`$ and $`F`$ one of its generators. The generators of $`R`$ contained in $`p_V^{}(\pi X)`$ are in correspondence with the generators of $`\mathrm{\Sigma }`$ passing by $`[\pi ]`$. Thus, it is enough to show that $`p_V^{}(\pi X)`$ contains at most $`nr`$ generators counted with multiplicity: a hyperplane section $`HR`$ is a divisor of the form $`H=C+F_1+\mathrm{}+F_k`$, where $`F_i`$ are generators and $`C=\overline{_{FF_1,\mathrm{},F_k}(FH)}`$ is an irreducible $`(r1)`$–dimensional scroll with disjoint generators (since $`R`$ is normal); then $`dimC2(r2)+1`$ and $`\mathrm{deg}(C)codim(CC)+1r1`$. Since $`\mathrm{deg}(H)=n`$, $`kn(r1)=nr+1`$, which is our claim.
Being $`n_1,\mathrm{},n_r1`$ integers, $`n=n_1+\mathrm{}+n_r`$, $`r>1`$, let $`R(n_1,\mathrm{},n_r)\text{IP}^{n+(r1)}`$ denote the rational normal scroll generated by the rational normal curves $`\mathrm{\Gamma }_{n_1},\mathrm{},\mathrm{\Gamma }_{n_r}`$ in disjoint spaces (given $`r`$ isomorphisms $`\nu _i:\text{IP}^1\mathrm{\Gamma }_{n_i}`$, $`i=1,\mathrm{},r`$, $`R(n_1,\mathrm{},n_r)`$ is the scroll generated by the spaces $`\nu _1(t),\mathrm{},\nu _r(t)`$, with $`t\text{IP}^1`$); $`H`$ denotes its hyperplane section and $`F`$ one of its generators.
###### Proposition 4.6
The rational normal scroll $`R(n_1,\mathrm{},n_r)`$ can be projected to $`\text{IP}^{r+2}`$ providing the focal scroll $`X`$ of a nonsingular congruence $`\mathrm{\Sigma }(\pi ,X)`$. Hence, given $`nr`$, there exist smooth congruences $`\mathrm{\Sigma }G(r,r+2)`$ of order $`1`$ and class $`n`$ without fixed points. Moreover, if $`nr+2`$ and if $`\mathrm{max}(n_1,\mathrm{},n_r)3`$, $`R(n_1,\mathrm{},n_r)`$ can be projected to $`\text{IP}^{r+2}`$ providing the focal scroll of a singular congruence $`\mathrm{\Sigma }(\pi ,X)`$. Hence, if $`nr+2`$, there exist singular congruences $`\mathrm{\Sigma }G(r,r+2)`$ of order $`1`$ and order $`n`$.
Proof: As we have seen in the proof of (4.5), the singularity of $`\mathrm{\Sigma }(\pi ,X)`$ depends of the number of generators of $`R:=R(n_1,\mathrm{},n_r)`$ contained in $`p_V^{}(\pi X)HF`$; we will thus have to choose properly a divisor $`CHF`$ in each case.
For the first part, it is sufficient to show that $`|HF|`$ contains irreducible divisors: the trace of $`|H|`$ over $`F`$ is the complete linear series $`|𝒪_F(HF)|`$, so we have an exact sequence
$$0H^0(𝒪_R(HF))H^0(𝒪_R(H))H^0(𝒪_F(1))0$$
and $`h^0(𝒪_R(HF))=n`$; moreover the trace of $`|HF|`$ over another generator $`F^{}`$ is $`|𝒪_F^{}(HF^{})|`$, so the exact sequence
$$0H^0(𝒪_R(HFF^{}))H^0(𝒪_R(HF))H^0(𝒪_F^{}(1))0$$
provides $`h^0(𝒪_R(HFF^{}))=nr`$. Making $`F^{}`$ vary in $`R`$, the set of reducible elements of $`|HF|`$ has dimension $`nr`$, hence it is a proper closed subset.
For the second part, we will show that $`k=\mathrm{max}(n_1,\mathrm{},n_r)`$ is the greatest number of generators of $`R`$ contained in a hyperplane (thus, if $`k3`$, a divisor in $`|HF|`$ can contain $`k12`$ generators of $`R`$): let $`F_1,\mathrm{},F_l`$ be generators of $`R`$, $`F_i=P_i^1,\mathrm{},P_i^r`$, $`P_i^j\mathrm{\Gamma }_{n_j}`$, $`i=1,\mathrm{},l`$. They are contained in a hyperplane iff $`F_1+\mathrm{}+F_l\text{IP}^{n+r1}`$, iff $`P_1^1,\mathrm{},P_l^1+\mathrm{}+P_1^r,\mathrm{},P_l^r\text{IP}^{n+r1}`$, iff $`ln_j`$ for some $`j=1,\mathrm{},r`$, iff $`l\mathrm{max}(n_1,\mathrm{},n_r)`$, which is our claim.
The next proposition deals with the existence of congruences with assigned linearly normal model. We will use the notation of chapter V §2.
###### Proposition 4.7
Given $`n1`$ and $`1rn`$, there exist congruences without fixed points $`\mathrm{\Sigma }G(r,r+2)`$ in the case II of order $`1`$ and class $`n`$, with given invariant $`e`$, $`0en1`$, $`ne1(mod2)`$.
Proof: Let $`X_e`$ be the geometrically ruled surface of invariant $`e`$, $`C_0`$ its minimal directrix ($`C_0^2=e`$) and $`F`$ one of its generators. Suppose $`0en1`$ and $`ne1(mod2)`$ and consider the morphisms
$$\begin{array}{cc}\psi _1:X_e\text{IP}^1,& \psi _2:X_e\text{IP}^n\end{array}$$
given by the linear systems $`|F|`$ and $`|C_0+(n+e1)/2F|`$, respectively. Since $`(ne1)/2e`$ the second one is base point free, so both maps are regular. Regard the composition
$$\psi :X_e\stackrel{\psi _1\times \psi _2}{}\text{IP}^1\times \text{IP}^n\stackrel{i}{}\mathrm{\Omega }_\pi \text{IP}^{2n+2}$$
where $`\pi \text{IP}^{n+2}`$ is an $`n`$–dimensional subspace, and we use the identification of $`\mathrm{\Omega }_\pi `$ with the cone with vertex $`[\pi ]`$ over the Segre variety given in (2.7). If $`H\text{IP}^{2n+2}`$ is a hyperplane, then $`\psi ^1(H)=C_0+(n+e+1)/2F`$ and $`\psi ^{}(|H|)=|C_0+(n+e+1)/2F|`$. Hence $`\psi (X_e)`$ is a rational normal ruled surface of degree $`n+1`$ and invariant $`e`$. Since $`\psi _1^1(x)F`$, the condition of (2.9) holds, so $`\mathrm{\Sigma }`$ has order $`1`$. Each generator $`F\psi (X_e)`$ parametrizes a pencil of $`\text{IP}^n`$’s not containing $`\pi `$, equivalently, their focal locus does not lie on $`\pi `$, so $`\psi (X_e)`$ is in the case II, and it is smooth since $`[\pi ]\psi (X_e)`$.
Moreover, for $`r<n`$ it is sufficient to take the generic projections of $`\psi (X_e)`$ according to (1.6).
###### Remark 4.8
Moreover, every congruence in the case II is expected to be a projection of one of those constructed above, because, as rational surfaces, those are their linearly normal models.
## 5 Case III: $`F(\mathrm{\Sigma })`$ is an $`r`$-dimensional projective space
$`F(\mathrm{\Sigma })`$ is the only $`r`$-dimensional projective subspace $`\pi \text{IP}^{r+2}`$ such that $`\mathrm{\Sigma }\mathrm{\Omega }_\pi =\{\sigma G(r,r+2)/dim\text{IP}^r(\sigma )\pi r1\}`$. Moreover, in this case, $`2(\text{IP}^{r1}(\sigma )):=2(\text{IP}^r(\sigma )\pi )`$ is the focal quadric of $`\mathrm{\Sigma }`$ at $`\text{IP}^r(\sigma )`$. The fibers of the map
$$\alpha :\mathrm{\Sigma }\{[\pi ]\}H_\pi \text{IP}^1$$
defined in (2.9) parametrize pencils whose base loci lie in $`F(\mathrm{\Sigma })=\pi `$, so they all contain $`[\pi ]`$. Therefore $`\mathrm{\Sigma }`$ is a rational cone with vertex $`[\pi ]`$. Consider the projection $`\beta :\mathrm{\Sigma }\{[\pi ]\}\pi ^{}`$ defined in (2.10), $`\beta (\sigma )=\text{IP}^r(\sigma )\pi `$: since the fibers of this map are $`\beta ^1(\text{IP}^{r1}(\sigma ))=\alpha ^1(\text{IP}^r(\sigma )+H)`$, $`\beta (\mathrm{\Sigma })`$ is a curve. Thus we have a birrational map
(9)
such that $`\mathrm{\Lambda }\alpha =\beta `$ and so $`\mathrm{\Sigma }`$ is characterized by $`\mathrm{\Lambda }`$ in the next way:
(10)
$$\mathrm{\Sigma }=\underset{H\pi }{}\mathrm{\Omega }_{\mathrm{\Lambda }(H),H}$$
being $`\mathrm{\Omega }_{\mathrm{\Lambda }(H),H}=\{\sigma G(r,r+2)/\mathrm{\Lambda }(H)\text{IP}^r(\sigma )H\}`$.
$`\beta (\mathrm{\Sigma })`$ parametrize a family of hyperplanes in $`\pi `$, hence it is developable, and we can say how it is constructed: suppose $`\mathrm{\Sigma }`$ has no fixed points (otherwise see (1.11)); $`\beta (\sigma )`$ is the family of $`(r1)`$–osculating hyperplanes to a nondegenerate curve $`C\pi `$, which is birrational to $`\beta (\mathrm{\Sigma })`$. Thus we have another birrational map:
that characterizes:
(11)
$$\mathrm{\Sigma }=\underset{H\pi }{}\mathrm{\Omega }_{T_{r1,C,\mathrm{\Lambda }^{}(H)},H}$$
where $`T_{r1,C,\mathrm{\Lambda }^{}(H)}`$ denotes the $`(r1)`$-osculating hyperplane to $`C`$ at $`\mathrm{\Lambda }^{}(H)`$.
Let $`\varphi `$ ($`=\alpha \times \beta `$) be the map defined in (2.7); $`\mathrm{\Sigma }`$ is a cone with vertex $`[\pi ]`$ over the curve $`\varphi (\mathrm{\Sigma }\{[\pi ]\})\pi ^{}\times H_\pi `$, which is a hyperplane section of $`\mathrm{\Sigma }`$. Since this curve is the graph of the map $`\mathrm{\Lambda }:H_\pi \pi ^{}`$, $`\varphi (\mathrm{\Sigma }\{[\pi ]\})`$ is an smooth curve of degree $`n+1`$ (where $`n:=\mathrm{deg}(\beta (\mathrm{\Sigma }))=\mathrm{deg}(\mathrm{\Lambda }(H_\pi ))`$). Thus $`\mathrm{deg}(\mathrm{\Sigma })=\mathrm{deg}(\varphi (\mathrm{\Sigma }\{[\pi ]\}))=n+1=\mathrm{deg}(\mathrm{\Lambda }(H_\pi ))+1`$ and we conclude $`cl(\mathrm{\Sigma })=n`$. Moreover the singular locus of $`\mathrm{\Sigma }`$ is only the vertex $`[\pi ]`$. Summarizing, we have:
###### Theorem 5.1
For every congruence $`\mathrm{\Sigma }G(r,r+2)`$ in the case III without fixed points there exist a regular map
$$\mathrm{\Lambda }:H_\pi \pi ^{}$$
(being $`\pi =F(\mathrm{\Sigma })`$ and $`H_\pi `$ the pencil of hyperplanes containing $`\pi `$) with nondegenerate image $`\mathrm{\Lambda }(H_\pi )`$ such that:
$$\mathrm{\Sigma }=\underset{H\pi }{}\mathrm{\Omega }_{\mathrm{\Lambda }(H),H}$$
being $`\mathrm{\Omega }_{\mathrm{\Lambda }(H),H}=\{\sigma G(r,r+2)/\mathrm{\Lambda }(H)\text{IP}^r(\sigma )H\}`$. Moreover, $`\mathrm{\Sigma }`$ is a cone of degree $`\mathrm{deg}(\mathrm{\Lambda }(H_\pi ))+1`$ ($`cl(\mathrm{\Sigma })=n`$), with vertex $`[\pi ]`$ over an smooth curve. Conversely, a congruence constructed in this way is in the case III.
The curve $`\beta (\mathrm{\Sigma })\pi ^{}`$ is the projection of its normal model $`\mathrm{\Gamma }_n(\text{IP}^n)^{}`$ from a subspace $`V(\text{IP}^n)^{}`$, hence the map $`\mathrm{\Lambda }`$ factorizes through $`\mathrm{\Gamma }_n`$. Regard $`\pi \text{IP}^n`$, and everything contained in a projective space $`\text{IP}^{n+2}`$ such that $`\pi `$ is obtained cutting $`\text{IP}^n`$ with the $`(r+2)`$–dimensional space containing the congruence $`\mathrm{\Sigma }`$ (allowing us to identify the hyperplanes in $`\text{IP}^{r+2}`$ containing $`\pi `$ with the hyperplanes in $`\text{IP}^{n+2}`$ containing $`\text{IP}^n`$), we will have a map $`\mathrm{\Lambda }_n:H_{P^n}(\text{IP}^n)^{}`$ providing a congruence $`\mathrm{\Sigma }_nG(n,\text{IP}^{n+2})`$ whose projection to $`G(r,\text{IP}^{r+2})`$ is $`\mathrm{\Sigma }`$. Hence we have a commutative diagram
We have thus proved:
###### Theorem 5.2
Being $`n1`$, there exist a congruence $`\mathrm{\Sigma }_nG(n,n+2)`$ constructed in the next way: given $`L\text{IP}^{n+2}`$ an $`n`$–dimensional projective space and $`\mathrm{\Lambda }:\text{IP}^1H_LL^{}`$ a $`n`$–th Veronese embedding, $`\mathrm{\Sigma }_n:=_{HL}\mathrm{\Omega }_{\mathrm{\Lambda }(H),H}`$. $`\mathrm{\Sigma }_n`$ is a rational normal cone of degree $`n+1`$. Every congruence $`\mathrm{\Sigma }G(r,r+2)`$ of order $`1`$ and class $`n`$ without fixed points in the case III can be obtained cutting $`\mathrm{\Sigma }_n`$ with a suitable projective subspace $`\text{IP}^{r+2}`$, that is:
$$\mathrm{\Sigma }_nG(n,n+2)\stackrel{\rho _{\text{IP}^{r+2}}}{}\mathrm{\Sigma }G(r,r+2).$$
## 6 Smoothness of a congruence of order $`1`$
According to (1.11), in order to study the smoothness of congruences of order $`1`$ we only need to consider congruences without fixed points, translating the results to the general case. Summarizing (3.1), (4.5), (4.6) and (5.1), we have the next theorems:
###### Theorem 6.1
Let $`r,s`$ be two integers such that $`r1`$ and $`1sr1`$. If $`\mathrm{\Sigma }G(r,r+2)`$ is a surface of order $`1`$ and class $`n`$ such that $`dim(T(\mathrm{\Sigma }))=s`$, then $`nrs1`$. Furthermore:
* If $`s=r1`$, $`\mathrm{\Sigma }`$ is a plane.
* For every $`sr2`$, there exist singular surfaces $`\mathrm{\Sigma }G(r,r+2)`$ of bidegree $`(1,n)`$ such that $`dim(T(\mathrm{\Sigma }))=s`$.
If $`s=r2`$, then:
1. If $`n=1,2`$ and $`\mathrm{\Sigma }`$ is not a cone, $`\mathrm{\Sigma }`$ is smooth.
2. If $`n4`$, $`\mathrm{\Sigma }`$ is singular.
3. If $`n=3`$, there exist smooth and singular surfaces $`\mathrm{\Sigma }G(r,r+2)`$ of bidegree $`(1,3)`$ such that $`dimT(\mathrm{\Sigma })=s`$.
If $`1sr3`$, then:
* If $`n=rs1`$ or $`n=rs`$ and $`\mathrm{\Sigma }`$ is not a cone, $`\mathrm{\Sigma }`$ is smooth.
* If $`nrs+1`$, there exist smooth and singular surfaces $`\mathrm{\Sigma }G(r,r+2)`$ of bidegree $`(1,n)`$ such that $`dimT(\mathrm{\Sigma })=s`$.
Proof: Let $`\mathrm{\Sigma }G(r,r+2)`$ be a congruence of order $`1`$ and class $`n`$ such that the dimension of its fixed locus is $`dim(T(\mathrm{\Sigma }))=s`$; if $`s=r1`$, $`\mathrm{\Sigma }`$ is the family of $`\text{IP}^r`$’s containing $`T(\mathrm{\Sigma })`$, which is parametrized by a plane; in this case $`n=0=rs1`$; suppose now that $`sr2`$. Applying (1.11), $`\mathrm{\Sigma }`$ can be projected isomorphically to a congruence $`\mathrm{\Sigma }^{}G(rs1,rs+1)`$ without fixed points, having order $`1`$ and class $`n`$; if $`\mathrm{\Sigma }^{}`$ is in the case I, then $`n=3`$ and $`rs1\{1,2,3\}`$, and then $`nrs1`$; otherwise (2.10) shows that $`n=cl(\mathrm{\Sigma }^{})rs1`$. Moreover for every $`nrs1`$, there exist a congruence $`\mathrm{\Sigma }^{}G(rs1,rs+1)`$ of class $`n`$ in the case III without fixed points (see (5.2)), which is parametrized by a cone. Excepting those congruences, the rest of the theorem is consequence of (4.5) and (3.1) and (4.6).
###### Theorem 6.2
Being $`r1`$, the only smooth surfaces in $`G(r,r+2)`$ of bidegree $`(1,n)`$ are:
1. The plane parametrizing the family of $`\text{IP}^r`$’s containing a fixed $`\text{IP}^{r1}`$ ($`n=0`$).
2. The Veronese surface, described in (3.1), for $`r=1,2,3`$ ($`n=3`$); embedding $`G(1,3)`$, $`G(2,4)`$ or $`G(3,5)`$ through the map given in (1.11), it can be consider in every $`G(r,r+2)`$.
3. The rational ruled surfaces $`\mathrm{\Sigma }(\pi ,X)`$ described in Section 4 ($`n=\mathrm{deg}X`$), and verifying that $`\pi `$ contains at most one generator of $`X`$; this condition is always verified for $`n=r,r+1`$.
4. The surfaces described in 2. and 3. embedded in a higher dimension grassmannian $`G(r,r+2)`$ through the isomorphism exposed in (1.11).
Let us finally enumerate the smooth congruences of order $`1`$ in $`G(1,3)`$, recovering a Ziv Ran’s result given in . We remark that, for a congruence $`\mathrm{\Sigma }(\pi ,X)`$ in $`G(1,3)`$, $`\pi `$ is a line cutting the curve $`X`$ in $`n1`$ points; hence, imposing the smoothness condition of (4.5), we have $`n11`$, that is $`\mathrm{deg}X2`$.
###### Corollary 6.3
The only smooth congruences in $`G(1,3)`$ of order $`1`$ are:
1. The plane parametrizing the lines passing by a fixed point (bidegree (1,0)).
2. The Veronese surface $`\mathrm{\Sigma }(C)`$ parametrizing the secant lines to a rational normal cubic $`C`$ (bidegree (1,3)).
3. The quadric $`\mathrm{\Sigma }(\pi ,\pi ^{})`$ parametrizing the lines cutting two given lines $`\pi `$ and $`\pi ^{}`$ (bidegree (1,1)).
4. The rational normal cubic $`\mathrm{\Sigma }(\pi ,C)`$ parametrizing the lines cutting a given line $`\pi `$ and a given irreducible conic $`C`$ meeting in a point $`P=\pi C`$ (bidegree (1,2)).
> *Authors’ address:* Departamento de Algebra, Universidad de Santiago de Compostela, 15706 Santiago de Compostela, Galicia, Spain. Phone: 34-81563100-ext.13152. Fax: 34-81597054. e-mail: pedreira@zmat.usc.es
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# Recovering Coherence via Conditional Measurements
## I Introduction
Decoherence of non-classical states of a quantum system via coupling to a reservoir is of fundamental interest, as it constitutes the mechanism that yields the classical limit of the system dynamics . Recently it has also become a topic of great applied importance, because it determines the feasibility of quantum information storage, encoding (encrypting) and computing . In numerous current theoretical proposals, the irreversibility of decoherence processes in quantum computing is combatted by two generic means. One is the filtering out of the ensemble portion which has not decohered, i.e., has remained intact. This approach has been suggested for two-mode fields , but not for single-mode cavity fields. The other means is encoding the state (qubit) by means of several ancillas, decoding the result after a certain time, checking the ancillas for error syndromes and correcting them . Although the latter approach is in principle applicable to arbitrary errors, only extremely small error probabilities (per qubit or gate, per time step) can afford fault tolerant quantum computation . Instead of the “high level” unitary transformation approach to error correction in quantum computing—which involves substantial overhead in qubits and gates—the countering of decoherence of stored quantum information, e.g., in between computation steps, may be achieved by a “low level” approach: applying simple physical manipulations to the quantum storage device, which take advantage of its specific physical realization. Such approach has been advocated recently , relying on continuous monitoring of the dissipation channel for quantum jumps, with perfect photodetection efficiency, and on instantaneous feedback for the inversion of their effect.
The nature of quantum computing requires that decoherence be corrected without knowing which state is in error during the computation. There is, however, a simpler but still important problem: how to protect from decoherence the input states, prior to the onset of computation. Here we suggest a non-unitary approach to counter decoherence, which can be used to safely store quantum field states in dissipative cavities, in order to subsequently use them as apriori known input in information processing or in signal transmission. The basic idea is to restore the decohered field state by entangling it with an atom, and then projecting the entangled state onto a superposition of atomic eigenstates, whose phase and amplitudes are specifically tailored for the field state we wish to recreate. Such projection amounts to post-selection of the appropriate atomic state, i.e., to a conditional measurement (CM) . The specific scheme we put forward is based on modification of our optimized CM strategy for cavity-mode state preparation by resonant interaction with atoms, in the Jaynes-Cummings (JC) model, followed by projection onto selected atomic states. In the present problem we set the initial (unspoilt) superposition of zero-photon up to $`N`$-photon states as our target state, and work in Liouville space instead of Hilbert space, so as to account for the state decoherence. The results demonstrate that a few highly-probable CMs, in this simple model, can drastically reduce even a large error. One of our objectives is to find the optimal tradeoff between the CM probability and the error size, which grows in the course of dissipation.
The ability to approximately restore any mixture to any pure state (in our $`N`$+1-dimensional Hilbert space) is the advantage of our post-selection CM approach, compared to the non-selective measurement (tracing) approach: Mixed states can only evolve into the special “cotangent” and “tangent” pure states by a large number of JC interactions with atoms initially prepared in superposition states (under the atomic excitation-trapping condition) followed by tracing over the atomic states .
## II Decoherence Minimization by Conditional Measurements
We consider a single-mode cavity in which the quantized electromagnetic field is initially prepared in a finite superposition of Fock states,
$$|\psi (0)=\underset{n=0}{\overset{N}{}}c_n|n.$$
(1)
To model the effect of dissipation we assume the cavity field to be coupled to a zero-temperature heat bath. The master equation describing such coupling, in the interaction picture, is
$$\dot{\rho }__F=\gamma (2\widehat{a}\rho __F\widehat{a}^{}\widehat{a}^{}\widehat{a}\rho __F\rho __F\widehat{a}^{}\widehat{a}),$$
(2)
where $`\rho __F=\rho __F(t)`$ is the density matrix of the cavity field, $`\widehat{a}`$ and $`\widehat{a}^{}`$ are the annihilation and creation operators of the field, and $`\gamma `$ is the damping constant of the cavity.
The solution of Eq. (2) after dissipation over time $`\overline{t}>0`$ can be shown to have the form
$`\rho _{n,m}(\overline{t})`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}\rho _{n+k,m+k}(0)\sqrt{\left(\begin{array}{c}n+k\\ n\end{array}\right)(e^{2\gamma \overline{t}})^n(1e^{2\gamma \overline{t}})^k}`$ (4)
$`\times \sqrt{\left(\begin{array}{c}m+k\\ m\end{array}\right)(e^{2\gamma \overline{t}})^m(1e^{2\gamma \overline{t}})^k},`$
written here in Fock basis, $`\rho _{n,m}(t)=n|\rho __F(t)|m`$.
In order to recover the original state of the field we propose to apply an optimized CM (or a sequence thereof) to the cavity as follows: Using a classical field we prepare a two-level atom in a chosen superposition
$$|\varphi ^{(i)}=\alpha ^{(i)}|e+\beta ^{(i)}|g$$
(5)
of its ground $`|g`$ and excited $`|e`$ states, and let it interact with the field for a time $`\tau `$ by sending it through the cavity with controlled speed. The field-atom interaction is adequately described by the resonant Jaynes-Cummings (JC) model . We assume the field-atom interaction time $`\tau `$ to be much shorter than the cavity lifetime, $`\gamma \tau 1`$, so that we may neglect dissipation during each CM. Upon exiting the cavity the atom is conditionally measured, using a second classical field, to be in a state
$$|\varphi ^{(f)}=\alpha ^{(f)}|e+\beta ^{(f)}|g,$$
(6)
differing in general from the initial atomic state $`|\varphi ^{(i)}`$. This means that we post-select, using the same setup as in ref. , the atomic superposition state (6) which is correlated to a cavity field state that is as close as possible to the original state (23).
The effect the applied CM has on the cavity field is then calculated as follows: Initially, at the time the atom enters the cavity, the density matrix of the field-atom system is
$$\rho _{_{FA}}(\overline{t})=\rho __F(\overline{t})|\varphi ^{(i)}\varphi ^{(i)}|.$$
(7)
It then evolves unitarily by the JC interaction of duration $`\tau `$ into
$$\rho _{_{FA}}(\overline{t}+\tau )=\widehat{U}(\tau )\rho _{_{FA}}(\overline{t})\widehat{U}^{}(\tau ),$$
(8)
where $`\widehat{U}(\tau )`$ is the interaction picture evolution operator
$`\widehat{U}(\tau )|n|e`$ $`=`$ $`C_n|n|eiS_n|n+1|g`$ (9)
$`\widehat{U}(\tau )|n|g`$ $`=`$ $`C_{n1}|n|giS_{n1}|n1|e,`$ (10)
with $`C_n=\mathrm{cos}\left(\lambda \tau \sqrt{n+1}\right)`$ and $`S_n=\mathrm{sin}\left(\lambda \tau \sqrt{n+1}\right)`$, $`\lambda `$ being the field-atom coupling constant (known as the vacuum Rabi frequency). Finally, the conditional measurement of the atom in the state $`|\varphi ^{(f)}`$ results in a density matrix of the field given by
$$\rho __F(\overline{t}+\tau )=\mathrm{Tr}__A\left[\rho _{_{FA}}(\overline{t}+\tau )|\varphi ^{(f)}\varphi ^{(f)}|\right]/P,$$
(11)
where
$$P=\mathrm{Tr}__F\mathrm{Tr}__A\left[\rho _{_{FA}}(\overline{t}+\tau )|\varphi ^{(f)}\varphi ^{(f)}|\right]$$
(12)
is the success probability of the CM. The explicit expressions for $`\rho __F(\overline{t}+\tau )`$ and $`P`$ are given in the Appendix for an initial superposition of $`|0`$ and $`|1`$ states.
To nearly recover the original state of the field, we use the dependence of $`\rho __F(\overline{t}+\tau )`$ on the initial and final atomic states and the field-atom interaction time, choosing optimal parameters $`\alpha ^{(i)}`$, $`\beta ^{(i)}`$, $`\alpha ^{(f)}`$, $`\beta ^{(f)}`$ and $`\tau `$ such that
$$\rho __F(\overline{t}+\tau )\rho __F(0)$$
(13)
holds (see Appendix for an explicit form of this condition), along with high CM success probability (12). These optimal CM parameters are found by minimizing the cost function
$$G=\frac{d(\rho __F(\overline{t}+\tau ),\rho __F(0))}{P^r},$$
(14)
where $`d`$ is a distance function between two density matrices, defined as
$$d(\rho __F^{(1)},\rho __F^{(2)})=\sqrt{\underset{nm}{}(\rho _{nm}^{(1)}\rho _{nm}^{(2)})^2},$$
(15)
$`P`$ is the CM success probability (12), and the adjustable exponent $`r>0`$ determines the relative importance of the two factors in $`G`$. If this CM does not bring us as close to the original state as our experimental accuracy permits, we can repeat the process over and over again, as long as the distance to the original state keeps decreasing, while the CM success probability remains high. The atomic states (5) and (6) are determined by the minimization of (14) at each step. Let us note here that the application of each CM may introduce widening of the photon-number distribution by one photon, and yet the optimized CMs are capable of avoiding this widening and, moreover, of restoring the field to its initial pure state. Eqs. (73-77) in the Appendix exemplify the widening-avoidance requirements which are implicit in condition (13). These requirements amount to an effective control of a large Fock-state subspace.
## III Examples
We illustrate our approach with two examples below, using the $`Q`$-function $`Q_{\rho __F}(\alpha ,\alpha ^{})=\alpha |\rho __F|\alpha `$, $`|\alpha `$ being a coherent state of complex amplitude $`\alpha `$, to visualize the error-correction process:
1) Let us take as the original field state an equal-amplitude superposition of our basis states, e.g.,
$$|\psi (0)=(|0+e^{i\pi /3}|1)/\sqrt{2},$$
(16)
whose $`Q`$-function is shown in Fig. 1(a). Dissipation by $`\gamma \overline{t}=0.3`$ renders the error matrix $`\rho __F(\overline{t})\rho __F(0)`$ of considerable magnitude, as seen in Fig. 1(b). After the application of one CM ($`|\varphi ^{(i)}=\mathrm{cos}(3\pi /8)|e+\mathrm{sin}(3\pi /8)e^{i5\pi /4}|g`$, $`\lambda \tau =37.95`$, $`|\varphi ^{(f)}=\mathrm{cos}(3\pi /8)|e+\mathrm{sin}(3\pi /8)e^{i\pi /4}|g`$), optimized to yield high success probability ($`r=2`$), the remaining error matrix $`\rho __F(\overline{t}+\tau )\rho __F(0)`$ is roughly 2.5 times smaller than before the correction, as seen in Fig. 1(c). The success probability of the CM is a high 74%. Subsequent CMs can further reduce the distance to $`1/6`$ (one sixth) its original magnitude, with 62% success probability for the full CM sequence (Fig. 3). Stronger error reduction is obtainable at the expense of success probability: the application of 4 CMs optimized for $`r=1`$ (respectively $`r=0`$) yields an error reduction factor of 11 (respectively 28) with sequence probability of 33% (respectively 16%).
2) If the original field state is a strongly unequal superposition of the basis states, such as
$$|\psi (0)=10^1|0+e^{i\pi /3}\sqrt{110^2}|1$$
(17)
(Fig. 2(a)), the error matrix after dissipation by $`\gamma \overline{t}=0.3`$ is again significant (Fig. 2(b)). Successive application of 4 CMs, optimized for $`r=2`$, reduces this error by a factor of 30 (Fig. 2(c)), which means that the recovered state is practically indistinguishable from the original state. The success probability of the total CM sequence, 50%, is markedly high (Fig. 3). If we ignore success probability in Eq. (14) ($`r=0`$) we obtain a higher error-reduction factor of 75, with sequence probability of 28%.
In Fig. 3 we plot the distance $`d_K=d(\rho ^K,\rho __F(0))`$ (Eq. (15)) between the recovered state and the original state and the CM sequence probability $`P_{seq,K}=_{l=1}^KP_l`$, with $`P_l`$ given by (12), as a function of the number of CMs performed. It shows that the first CMs achieve a strong reduction of such a distance, whereas after a few successive CMs saturation sets on, in terms of both distance and success probability.
It is interesting to compare the success probability in our approach with the theoretical probability to find the original state in the dissipation-spoilt state, namely, $`\mathrm{Tr}__F[\rho __F(0)\rho __F(\overline{t})]`$, which we call the filtering probability. In Table I we list the success probability of a sequence of 4 CMs (optimized for $`r=2`$), $`P_{seq,K=4}`$, and the corresponding filtering probability for various values of the dissipation parameter $`\gamma \overline{t}`$, taking as the original state the state (17) of example 2. The probability $`\mathrm{Tr}__F[\rho __F(0)\rho __F^{K=4}]`$ of finding the original state in the recovered state is 0.99 or higher for all entries.
## IV Discussion
In conclusion, we have demonstrated here the effectiveness of simple JC-dynamics CMs as a means of reversing the effect of dissipation on coherent superpositions of Fock-states of a cavity field: the application of a small number of optimized CMs recovers the original state of the field with high success probability, which is comparable or even surpasses the filtering probability. The simplest tactics may employ a single highly-probable trial to achieve nearly-complete error correction. As noted above, although we have only five control parameters at our disposal for each CM, our optimization procedure is able to effectively control the amplitudes in a large Fock-state subspace.
Among the experimental imperfections that can degrade the effectiveness of any CM approach , realistic atomic velocity fluctuations (of 1%) and cavity-temperature effects (below $`1^{}`$K) are relatively unimportant, and especially so in the present scheme which makes use of a single or few CMs so that the effect of experimental imperfections is linear in the input errors. Only atomic detection efficiency is an experimental challenge . Although the detection efficiency is currently low, it is expected to rise considerably in the coming future.
Extensions of this approach to field-atom interaction Hamiltonians with more controllable degrees of freedom can make a single trial within this correction procedure effective for highly complicated states, encoding many qubits of information. Nevertheless, even in its present form the suggested approach has undoubted merits: (a) it can yield higher success probabilities than the filtering approach; (b) it is not limited to small errors as “high level” unitary-transformation approaches are; (c) it corrects errors after their occurrence, with no reliance on ideal continuous monitoring of the dissipation channel and on instantaneous feedback; and (d) it is realistic in that it can counter combined phase-amplitude errors which arise in cavity dissipation, and is of general applicability—not restricted to specific models of dissipation.
###### Acknowledgements.
The support of the German-Israeli Foundation (GIF) is acknowledged. M.F. thanks the European Economic Community (Human Capital and Mobility programme) for support.
##
The reduced density matrix of the field resulting from its interaction with the atom followed by the conditional measurement on the latter can be found using the formula (Eqs. (7-12))
$`\rho __F`$ $`(`$ $`\overline{t}+\tau )=`$ (19)
$`\mathrm{Tr}__A\left[\widehat{U}(\tau )\rho __F(\overline{t})|\varphi ^{(i)}\varphi ^{(i)}|\widehat{U}^{}(\tau )|\varphi ^{(f)}\varphi ^{(f)}|\right]/P,`$
where the normalization constant $`P`$ is the success probability of the conditional measurement and is given by
$`P=`$ $`\mathrm{Tr}__F\mathrm{Tr}__A\left[\widehat{U}(\tau )\rho __F(\overline{t})|\varphi ^{(i)}\varphi ^{(i)}|\widehat{U}^{}(\tau )|\varphi ^{(f)}\varphi ^{(f)}|\right].`$ (21)
In the simple case where the initial field state is a superposition of the vacuum and one-photon states
$$|\psi (0)=c_0|0+c_1|1,\rho __F(0)=|\psi (0)\psi (0)|,$$
(23)
the density matrix resulting from dissipation over time $`\overline{t}`$ is
$`\rho __F(\overline{t})`$ $`=`$ $`\rho _{00}(\overline{t})|00|+\rho _{01}(\overline{t})|01|+`$ (25)
$`+\rho _{10}(\overline{t})|10|+\rho _{11}(\overline{t})|11|,`$
with
$`\rho _{00}(\overline{t})`$ $`=`$ $`|c_0|^2+(1e^{2\gamma \overline{t}})|c_1|^2`$ (26)
$`\rho _{01}(\overline{t})`$ $`=`$ $`e^{\gamma \overline{t}}c_0c_1^{}`$ (27)
$`\rho _{10}(\overline{t})`$ $`=`$ $`e^{\gamma \overline{t}}c_1c_0^{}`$ (28)
$`\rho _{11}(\overline{t})`$ $`=`$ $`e^{2\gamma \overline{t}}|c_1|^2.`$ (29)
The explicit expressions for $`\rho __F(\overline{t}+\tau )`$ and the success probability $`P`$ are then
$`\rho __F`$ $`(\overline{t}+\tau )=P^1\{[|\alpha ^{(\mathrm{f})}|^2A+\alpha ^{(\mathrm{f})}\beta ^{(\mathrm{f})}M+\beta ^{(\mathrm{f})}\alpha ^{(\mathrm{f})}C`$ (30)
$`+`$ $`|\beta ^{(\mathrm{f})}|^2O]|00|+[|\alpha ^{(\mathrm{f})}|^2B+\alpha ^{(\mathrm{f})}\beta ^{(\mathrm{f})}N+\beta ^{(\mathrm{f})}\alpha ^{(\mathrm{f})}D`$ (31)
$`+`$ $`|\beta ^{(\mathrm{f})}|^2K]|01|+[|\alpha ^{(\mathrm{f})}|^2E+\alpha ^{(\mathrm{f})}\beta ^{(\mathrm{f})}R+\beta ^{(\mathrm{f})}\alpha ^{(\mathrm{f})}H`$ (32)
$`+`$ $`|\beta ^{(\mathrm{f})}|^2T]|10|+[|\alpha ^{(\mathrm{f})}|^2F+\alpha ^{(\mathrm{f})}\beta ^{(\mathrm{f})}S+\beta ^{(\mathrm{f})}\alpha ^{(\mathrm{f})}I`$ (33)
$`+`$ $`|\beta ^{(\mathrm{f})}|^2U]|11|+[\alpha ^{(\mathrm{f})}\beta ^{(\mathrm{f})}W+|\beta ^{(\mathrm{f})}|^2Y]|20|`$ (34)
$`+`$ $`[\beta ^{(\mathrm{f})}\alpha ^{(\mathrm{f})}G+|\beta ^{(\mathrm{f})}|^2Q]|02|+[\alpha ^{(\mathrm{f})}\beta ^{(\mathrm{f})}X`$ (35)
$`+`$ $`|\beta ^{(\mathrm{f})}|^2J]|21|+[\beta ^{(\mathrm{f})}\alpha ^{(\mathrm{f})}L+|\beta ^{(\mathrm{f})}|^2V]|12|`$ (36)
$`+`$ $`|\beta ^{(\mathrm{f})}|^2Z|22|\},`$ (37)
and
$`P`$ $`=`$ $`|\alpha ^{(\mathrm{f})}|^2(A+F)+\alpha ^{(\mathrm{f})}\beta ^{(\mathrm{f})}(M+S)`$ (39)
$`+\beta ^{(\mathrm{f})}\alpha ^{(\mathrm{f})}(C+I)+|\beta ^{(\mathrm{f})}|^2(O+U+Z).`$
The coefficients $`A,B,\mathrm{}`$ here are given by
$`A`$ $`=`$ $`\rho _{00}(\overline{t})|\alpha ^{(\mathrm{i})}|^2C_0^2+i\rho _{01}(\overline{t})\alpha ^{(\mathrm{i})}\beta ^{(\mathrm{i})}C_0S_0`$ (42)
$`+\rho _{11}(\overline{t})|\beta ^{(\mathrm{i})}|^2S_0^2i\rho _{10}(\overline{t})\alpha ^{(\mathrm{i})}\beta ^{(\mathrm{i})}S_0C_0,`$
$`B`$ $`=`$ $`\rho _{01}(\overline{t})|\alpha ^{(\mathrm{i})}|^2C_0C_1i\rho _{11}(\overline{t})\alpha ^{(\mathrm{i})}\beta ^{(\mathrm{i})}C_1S_0,`$ (43)
$`C`$ $`=`$ $`\rho _{00}(\overline{t})\alpha ^{(\mathrm{i})}\beta ^{(\mathrm{i})}C_0i\rho _{10}(\overline{t})|\beta ^{(\mathrm{i})}|^2S_0,`$ (44)
$`D`$ $`=`$ $`i\rho _{00}(\overline{t})|\alpha ^{(\mathrm{i})}|^2C_0S_0i\rho _{11}(\overline{t})|\beta ^{(\mathrm{i})}|^2S_0C_0`$ (46)
$`+\rho _{01}(\overline{t})\alpha ^{(\mathrm{i})}\beta ^{(\mathrm{i})}C_0^2+\rho _{10}(\overline{t})\alpha ^{(\mathrm{i})}\beta ^{(\mathrm{i})}S_0^2,`$
$`F`$ $`=`$ $`\rho _{11}(\overline{t})|\alpha ^{(\mathrm{i})}|^2C_1^2,`$ (47)
$`G`$ $`=`$ $`\rho _{11}(\overline{t})\alpha ^{(\mathrm{i})}\beta ^{(\mathrm{i})}S_0S_1+i\rho _{01}(\overline{t})|\alpha ^{(\mathrm{i})}|^2C_0S_1,`$ (48)
$`H`$ $`=`$ $`\rho _{10}(\overline{t})\alpha ^{(\mathrm{i})}\beta ^{(\mathrm{i})}C_1,`$ (49)
$`I`$ $`=`$ $`\rho _{11}(\overline{t})\alpha ^{(\mathrm{i})}\beta ^{(\mathrm{i})}C_1C_0+i\rho _{10}(\overline{t})|\alpha ^{(\mathrm{i})}|^2C_1S_0,`$ (50)
$`K`$ $`=`$ $`i\rho _{00}(\overline{t})\alpha ^{(\mathrm{i})}\beta ^{(\mathrm{i})}S_0+\rho _{01}(\overline{t})|\beta ^{(\mathrm{i})}|^2C_0,`$ (51)
$`L`$ $`=`$ $`i\rho _{11}(\overline{t})|\alpha ^{(\mathrm{i})}|^2C_1S_1,`$ (52)
$`O`$ $`=`$ $`\rho _{00}(\overline{t})|\beta ^{(\mathrm{i})}|^2,`$ (53)
$`Q`$ $`=`$ $`i\rho _{01}(\overline{t})\alpha ^{(\mathrm{i})}\beta ^{(\mathrm{i})}S_1,`$ (54)
$`U`$ $`=`$ $`\rho _{00}(\overline{t})|\alpha ^{(\mathrm{i})}|^2S_0^2i\rho _{01}(\overline{t})\alpha ^{(\mathrm{i})}\beta ^{(\mathrm{i})}S_0C_0`$ (56)
$`+\rho _{11}(\overline{t})|\beta ^{(\mathrm{i})}|^2C_0^2+i\rho _{10}(\overline{t})\alpha ^{(\mathrm{i})}\beta ^{(\mathrm{i})}C_0S_0,`$
$`V`$ $`=`$ $`i\rho _{11}(\overline{t})\alpha ^{(\mathrm{i})}\beta ^{(\mathrm{i})}C_0S_1+\rho _{01}(\overline{t})|\alpha ^{(\mathrm{i})}|^2S_0S_1,`$ (57)
$`Z`$ $`=`$ $`\rho _{11}(\overline{t})|\alpha ^{(\mathrm{i})}|^2S_1^2,`$ (58)
with the following relations holding between them
$`A`$ $`=`$ $`A^{},B=E^{},C=M^{},D=R^{}`$ (60)
$`F`$ $`=`$ $`F^{},G=W^{},H=N^{},I=S^{}`$ (61)
$`L`$ $`=`$ $`X^{},O=O^{},K=T^{},Q=Y^{}`$ (62)
$`U`$ $`=`$ $`U^{},V=J^{},Z=Z^{}.`$ (63)
(The coefficients $`K`$ and $`N`$ in this appendix bear no relation to $`K`$ and $`N`$ mentioned in the main text).
The explicit form of condition (13) for recovering the original field state is given by the following list of approximation relations:
$`\rho _{00}(\overline{t}+\tau )`$ $`=`$ $`P^1[|\alpha ^{(\mathrm{f})}|^2A+\alpha ^{(\mathrm{f})}\beta ^{(\mathrm{f})}M`$ (66)
$`+\beta ^{(\mathrm{f})}\alpha ^{(\mathrm{f})}C+|\beta ^{(\mathrm{f})}|^2O]|c_0|^2`$
$`\rho _{01}(\overline{t}+\tau )`$ $`=`$ $`P^1[|\alpha ^{(\mathrm{f})}|^2B+\alpha ^{(\mathrm{f})}\beta ^{(\mathrm{f})}N`$ (68)
$`+\beta ^{(\mathrm{f})}\alpha ^{(\mathrm{f})}D+|\beta ^{(\mathrm{f})}|^2K]c_0c_1^{}`$
$`\rho _{10}(\overline{t}+\tau )`$ $`=`$ $`P^1[|\alpha ^{(\mathrm{f})}|^2E+\alpha ^{(\mathrm{f})}\beta ^{(\mathrm{f})}R`$ (70)
$`+\beta ^{(\mathrm{f})}\alpha ^{(\mathrm{f})}H+|\beta ^{(\mathrm{f})}|^2T]c_1c_0^{}`$
$`\rho _{11}(\overline{t}+\tau )`$ $`=`$ $`P^1[|\alpha ^{(\mathrm{f})}|^2F+\alpha ^{(\mathrm{f})}\beta ^{(\mathrm{f})}S`$ (72)
$`+\beta ^{(\mathrm{f})}\alpha ^{(\mathrm{f})}I+|\beta ^{(\mathrm{f})}|^2U]|c_1|^2`$
$`\rho _{02}(\overline{t}+\tau )`$ $`=`$ $`P^1\left[\beta ^{(\mathrm{f})}\alpha ^{(\mathrm{f})}G+|\beta ^{(\mathrm{f})}|^2Q\right]0`$ (73)
$`\rho _{20}(\overline{t}+\tau )`$ $`=`$ $`P^1\left[\alpha ^{(\mathrm{f})}\beta ^{(\mathrm{f})}W+|\beta ^{(\mathrm{f})}|^2Y\right]0`$ (74)
$`\rho _{12}(\overline{t}+\tau )`$ $`=`$ $`P^1\left[\beta ^{(\mathrm{f})}\alpha ^{(\mathrm{f})}L+|\beta ^{(\mathrm{f})}|^2V\right]0`$ (75)
$`\rho _{21}(\overline{t}+\tau )`$ $`=`$ $`P^1\left[\alpha ^{(\mathrm{f})}\beta ^{(\mathrm{f})}X+|\beta ^{(\mathrm{f})}|^2J\right]0`$ (76)
$`\rho _{22}(\overline{t}+\tau )`$ $`=`$ $`P^1\left[|\beta ^{(\mathrm{f})}|^2Z\right]0.`$ (77)
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# Classical Radiation Processes in the Weizsäcker-Williams Approximation
## I The Weizsäcker-Williams Approximation
Following an earlier discussion by Fermi , Weizsäcker and Williams noted that the electromagnetic fields of an electron in uniform relativistic motion are predominantly transverse, with $`𝐄𝐁`$ (in Gaussian units). This is very much like the fields of a plane wave, so one is led to regard a fast electron as carrying with it a cloud of virtual photons that it can shed (radiate) if perturbed.
The key feature of the frequency spectrum of the fields can be estimated as follows. To an observer at rest at distance $`b`$ from the electron’s trajectory, the peak electric field is $`E=\gamma e/b^2`$, and the field remains above half this strength for time $`b/\gamma c`$, so the frequency spectrum of this pulse extends up to $`\omega _{\mathrm{max}}\gamma c/b`$. The total energy of the pulse (relevant to this observer) is $`UE^2\mathrm{Vol}\gamma ^2e^2/b^4b^2b/\gamma \gamma e^2/b`$.
If the electron radiates all of this energy, the energy spectrum would be
$$\frac{dU(\omega )}{d\omega }\frac{U}{\omega _{\mathrm{max}}}\frac{e^2}{c}.$$
(1)
This result does not depend on the choice of impact parameter $`b`$, and is indeed of general validity (to within a factor of $`\mathrm{ln}\gamma `$). The number of photons $`n_\omega `$ of frequency $`\omega `$ is thus
$$dn_\omega =\frac{dU(\omega )}{\mathrm{}\omega }\frac{e^2}{\mathrm{}c}\frac{d\omega }{\omega }=\alpha \frac{d\omega }{\omega },$$
(2)
where $`\alpha =e^2/\mathrm{}c1/137`$ is the fine structure constant.
The quick approximation (1)-(2) is not accurate at high frequencies. In general, additional physical arguments are needed to identify the maximum frequency of its validity, called the characteristic or critical frequency $`\omega _C`$, or equivalently, the minimum relevant impact parameter $`b_{\mathrm{min}}`$. A more detailed evaluation of the high-frequency tail of the virtual photon spectrum shows it to be
$$dn_\omega \alpha \frac{d\omega }{\omega }e^{2\omega b_{\mathrm{min}}/\gamma c}\text{(high frequency)}.$$
(3)
From this, we see the general relation between the critical frequency and the minimum impact parameter is
$$\omega _C\gamma \frac{c}{b_{\mathrm{min}}},b_{\mathrm{min}}\gamma \lambda _C.$$
(4)
The characteristic angular spread $`\theta _C`$ of the radiation pattern near the critical frequency can be estimated from eq. (4) by noting that the radiation is much like that of a beam of light with waist $`b_{\mathrm{min}}`$. Then, from the laws of diffraction we conclude that
$$\theta _C\frac{\lambda _C}{b_{\mathrm{min}}}\frac{1}{\gamma }.$$
(5)
This behavior is also expected in that a ray of light emitted in the electron’s rest frame at $`90^{}`$ appears at angle $`1/\gamma `$ to the laboratory direction of the electron.
### A The Formation Length
To complete an application of the Weizsäcker-Williams method, we must also know over what interval the virtual photon cloud is shaken off the electron to become the radiation detected in the laboratory. Intense (and hence, physically interesting) radiation processes are those in which the entire cloud of virtual photons is emitted as rapidly as possible. This is usefully described by the so-called formation time $`t_0`$ and the corresponding formation length $`L_0=vt_0`$ where $`vc`$ is the velocity of the relativistic electron.
The formation length (time) is the distance (time) the electron travels while a radiated wave advances one wavelength $`\lambda `$ ahead of the projection of the electron’s motion onto the direction of observation. The wave takes on the character of radiation that is no longer tied to its source only after the formation time has elapsed. That is,
$$\lambda =ct_0vt_0\mathrm{cos}\theta L_0(1\beta \mathrm{cos}\theta )L_0\left(\frac{1}{2\gamma ^2}+\frac{\theta ^2}{2}\right),$$
(6)
for radiation observed at angle $`\theta `$ to the electron’s trajectory. Thus, the formation length is given by
$$L_0\frac{2\lambda }{\theta ^2+1/\gamma ^2}$$
(7)
If the frequency of the radiation is near the critical frequency (4), then the radiated intensity is significant only for θ
<
θC1/γ
<
𝜃subscript𝜃𝐶1𝛾\theta\mathrel{\vbox{\kern 0.0pt\hbox{$<$} \kern 0.0pt\hbox{$\sim$} }}\theta_{C}\approx 1/\gamma, and the formation length is
$$L_0\gamma ^2\lambda (\lambda \lambda _C).$$
(8)
A good discussion of the formation length in both classical and quantum contexts has been given in ref. .
### B Summary of the Method
A relativistic electron carries with it a virtual photon spectrum of $`\alpha `$ photons per unit frequency interval. When radiation occurs, for whatever reason, the observed frequency spectrum will closely follow this virtual spectrum. In cases where the driving force for the radiation extends over many formation lengths, the spectrum of radiated photons per unit path length for intense processes is given by expressions (2)-(3), which describe the radiation emitted over one formation length, divided by the formation length (7):
$$\frac{dn_\omega }{dl}\frac{\alpha }{L_0(\omega )}\frac{d\omega }{\omega }\times \{\begin{array}{cccc}1\hfill & & & (\omega <\omega _C),\hfill \\ e^{\omega /\omega _C}\hfill & & & (\omega \omega _C).\hfill \end{array}$$
(9)
Synchrotron radiation, undulator radiation, transition radiation, and Čerenkov radiation are examples of processes which can be described within the context of classical electromagnetism, but for which the Weizsäcker-Williams approximation is also suitable. Čerenkov radiation and transition radiation are often thought of as rather weak processes, but the Weizsäcker-Williams viewpoint indicates that they are actually as intense as is possible for radiation by a single charge, in the sense that the entire virtual photon cloud is liberated over a formation length.
In this paper, we emphasize a simplified version of the Weizsäcker-Williams method with the goal of illustrating the principle qualitative features of various radiation processes. A more detailed analysis can reproduce the complete forms of the classical radiation, as has been demonstrated for synchrotron radiation by Lieu and Axford .
## II Synchrotron Radiation
Synchrotron radiation arises when a charge, usually an electron, is deflected by a magnetic field. For a large enough region of uniform magnetic field, the electron’s trajectory would be a complete circle. However, synchrotron radiation as described below occurs whenever the magnetic field region is longer than a formation length. The radiation observed when the magnetic field extends for less than a formation length has been discussed in refs. .
### A The Critical Frequency
An important fact about synchrotron radiation is that the frequency spectrum peaks near the critical frequency, $`\omega _C`$, which depends on the radius $`R`$ of curvature of the electron’s trajectory, and on the Lorentz factor $`\gamma `$ via
$$\omega _C\gamma ^3\frac{c}{R}.$$
(10)
Since $`\omega _0=c/R`$ is the angular velocity for particles with velocity near the speed of light, synchrotron radiation occurs at very high harmonics of this fundamental frequency. The wavelength at the critical frequency is then
$$\lambda _C\frac{R}{\gamma ^3}.$$
(11)
For completeness, we sketch a well-known argument leading to eq. (10). The characteristic frequency $`\omega _C`$ is the reciprocal of the pulselength of the radiation from a single electron according to an observer at rest in the lab. In the case of motion in a circle, the electron emits a cone of radiation of angular width $`\theta =1/\gamma `$ according to eq. (5) that rotates with angular velocity $`\omega =c/R`$. Light within this cone reaches the fixed observer during time interval $`\delta t^{}=\theta /\omega R/\gamma c`$. However, this time interval measures the retarded time $`t^{}`$ at the source, not the time $`t`$ at the observer. Both $`t`$ and $`t^{}`$ are measured in the lab frame, and are related by $`t^{}=tr/c`$ where $`r`$ is the distance between the source and observer. When the source is heading towards the observer, we have $`\delta r=v\delta t^{}`$, so $`\delta t=\delta t^{}(1v/c)\delta t^{}/2\gamma ^2R/\gamma ^3c`$, from which eq. (10) follows.
### B The Formation Length
The formation length $`L_0`$ introduced in eq. (7) applies for radiation processes during which the electron moves along a straight line, such as Čerenkov radiation and transition radiation. But, synchrotron radiation occurs when the electron moves in the arc of a circle of radius $`R`$. During the formation time, the electron moves by formation angle $`\theta _0=L_0/R`$ with respect to the center of the circle. We now reconsider the derivation of the formation time, noting that while the electron moves on the arc $`R\theta _0=vt_0`$ of the circle, the radiation moves on the chord $`2R\mathrm{sin}(\theta _0/2)R\theta _0R\theta _0^3/24`$. Hence,
$`\lambda `$ $`=`$ $`ct_0\text{chord}{\displaystyle \frac{cR\theta _0}{v}}R\theta _0+{\displaystyle \frac{R\theta _0^3}{24}}`$ (12)
$``$ $`R\theta _0(1\beta )+{\displaystyle \frac{R\theta _0^3}{24}}{\displaystyle \frac{R\theta _0}{2\gamma ^2}}+{\displaystyle \frac{R\theta _0^3}{24}},`$ (13)
for radiation observed at small angles to the chord.
For wavelengths longer than $`\lambda _C`$, the formation angle grows large compared to the characteristic angle $`\theta _C1/\gamma `$, and the first term of eq. (13) can be neglected compared to the second. In this case,
$$\theta _0\left(\frac{\lambda }{R}\right)^{1/3}\frac{1}{\gamma }\left(\frac{\lambda }{\lambda _C}\right)^{1/3}(\lambda \lambda _C),$$
(14)
and
$$L_0R^{2/3}\lambda ^{1/3}\gamma ^2\lambda _C\left(\frac{\lambda }{\lambda _C}\right)^{1/3}(\lambda \lambda _C),$$
(15)
using eq. (11).
The formation angle $`\theta _0(\lambda )`$ can also be interpreted as the characteristic angular width of the radiation pattern at this wavelength. A result not deducible from the simplified arguments given above is that for $`\lambda \lambda _C`$, the angular distribution of synchrotron radiation falls off exponentially: $`dU(\lambda )/d\mathrm{\Omega }e^{\theta ^2/2\theta _0^2}`$. See, for example, sec. 14.6 of .
For wavelengths much less than $`\lambda _C`$, the formation length is short, the formation angle is small, and the last term of eq. (13) can be neglected. Then, we find that
$$\theta _0\frac{\lambda }{\gamma \lambda _C},L_0\gamma ^2\lambda (\lambda \lambda _C),$$
(16)
the same as for motion along a straight line, eq. (8). In this limit, our approximation neglects the curvature of the particle’s trajectory, which is an essential aspect of synchrotron radiation, and we cannot expect our analysis to be very accurate. But for $`\lambda \lambda _C`$, the rate of radiation is negligible.
Of greater physical interest is the region $`\lambda \lambda _C`$ where the frequency spectrum begins to be exponentially damped but the rate is still reasonably high. The cubic equation (13) does not yield a simple analytic result in the region. So, we interpolate between the limiting results for $`\theta _0`$ at large and small wavelengths, eqs. (14) and (16), and estimate that
$$\theta _0\frac{1}{\gamma }\sqrt{\frac{\lambda }{\lambda _C}}(\lambda \lambda _C),$$
(17)
which agrees with a more detailed analysis . The corresponding formation length $`R\theta _0`$ is then
$$L_0\gamma ^2\sqrt{\lambda \lambda _C}(\lambda \lambda _C).$$
(18)
### C Transverse Coherence Length
The longitudinal origin of radiation is uncertain to within one formation length $`L_0`$. Over this length, the trajectory of the electron is curved, so there is an uncertainty in the transverse origin of the radiation as well. A measure of the transverse uncertainty is the sagitta $`L_0^2/8R`$, which we label $`w_0`$ anticipating a useful analogy with the common notation for the waist of a focused laser beam. For $`\lambda \lambda _C`$, we have from eq. (15),
$$w_0\frac{L_0^2}{R}R^{1/3}\lambda ^{2/3}\gamma \lambda _C\left(\frac{\lambda }{\lambda _C}\right)^{2/3}(\lambda \lambda _C).$$
(19)
The sagitta (19) is larger than the minimum transverse length (4), so we expect that the full virtual photon cloud is shaken off over one formation length.
For $`\lambda \lambda _C`$, the characteristic angular spread (14) of the radiation obeys
$$\theta _0\frac{\lambda }{w_0},$$
(20)
consistent with the laws of diffraction. Hence, the distance $`w_0`$ of eq. (19) can also be called the transverse coherence length of the source of synchrotron radiation.
The analogy with laser notation is also consistent with identifying the formation length $`L_0`$ with the Rayleigh range $`z_0=w_0/\theta _0`$, since we see that
$$L_0\frac{\lambda }{\theta _0^2}\frac{w_0}{\theta _0}.$$
(21)
A subtle difference between the radiation of a relativistic charge and a focused laser beam is that the laser beam has a Guoy phase shift between its waist and the far field, while radiation from a charge does not.
For $`\lambda \lambda _C`$, the sagitta is $`L_0^2/R\gamma ^2\lambda `$, using eq. (18). When $`\lambda <\lambda _C`$, the characteristic angle $`\theta _0`$ given by eq.(17) is less than $`\lambda `$/sagitta, and the sagitta is no longer a good measure of the transverse coherence length, which is better defined as $`\lambda /\theta _0\gamma \sqrt{\lambda \lambda _C}`$.
### D Frequency Spectrum
The number of photons radiated per unit path length $`l`$ during synchrotron radiation is obtained from the Weizsäcker-Williams spectrum (9) using eqs. (15) and (18) for the formation length:
dnωdl{αωC2/3dω/γ2cω2/3(λλC),αωC1/2eω/ωCdω/γ2cω1/2(λ
<
λC).𝑑subscript𝑛𝜔𝑑𝑙cases𝛼superscriptsubscript𝜔𝐶23𝑑𝜔superscript𝛾2𝑐superscript𝜔23missing-subexpressionmissing-subexpressionmuch-greater-than𝜆subscript𝜆𝐶𝛼superscriptsubscript𝜔𝐶12superscript𝑒𝜔subscript𝜔𝐶𝑑𝜔superscript𝛾2𝑐superscript𝜔12missing-subexpressionmissing-subexpression
<
𝜆subscript𝜆𝐶{dn_{\omega}\over dl}\approx\left\{\begin{array}[]{llll}\alpha\omega_{C}^{2/3}d\omega/\gamma^{2}c\omega^{2/3}&&&(\lambda\gg\lambda_{C}),\\
\alpha\omega_{C}^{1/2}e^{-\omega/\omega_{C}}d\omega/\gamma^{2}c\omega^{1/2}&&&(\lambda\mathrel{\vbox{\kern 0.0pt\hbox{$<$} \kern 0.0pt\hbox{$\sim$} }}\lambda_{C}).\end{array}\right. (22)
We multiply by $`\mathrm{}\omega `$ to recover the energy spectrum:
dU(ω)dl{e2ωC2/3ω1/3dω/γ2c2(λλC),e2ωC1/2ω1/2eω/ωCdω/γ2c2(λ
<
λC).𝑑𝑈𝜔𝑑𝑙casessuperscript𝑒2superscriptsubscript𝜔𝐶23superscript𝜔13𝑑𝜔superscript𝛾2superscript𝑐2missing-subexpressionmissing-subexpressionmuch-greater-than𝜆subscript𝜆𝐶superscript𝑒2superscriptsubscript𝜔𝐶12superscript𝜔12superscript𝑒𝜔subscript𝜔𝐶𝑑𝜔superscript𝛾2superscript𝑐2missing-subexpressionmissing-subexpression
<
𝜆subscript𝜆𝐶{dU(\omega)\over dl}\approx\left\{\begin{array}[]{llll}e^{2}\omega_{C}^{2/3}\omega^{1/3}d\omega/\gamma^{2}c^{2}&&&(\lambda\gg\lambda_{C}),\\
e^{2}\omega_{C}^{1/2}\omega^{1/2}e^{-\omega/\omega_{C}}d\omega/\gamma^{2}c^{2}&&&(\lambda\mathrel{\vbox{\kern 0.0pt\hbox{$<$} \kern 0.0pt\hbox{$\sim$} }}\lambda_{C}).\end{array}\right. (23)
Thus, the Weizsäcker-Williams method shows that the energy spectrum varies as $`\omega ^{1/3}`$ at low frequencies, and as $`\sqrt{\omega }e^{\omega /\omega _C}`$ at frequencies above the critical frequency $`\omega _C=\gamma ^3c/R`$.
The total radiated power is estimated from eq. (23) using $`\omega d\omega \omega _C\gamma ^3c/R`$, and multiplying by $`vc`$ to convert $`dl`$ to $`dt`$:
$$\frac{dU}{dt}\frac{e^2\gamma ^4c}{R^2}.$$
(24)
This well-known result is also obtained from the Larmor formula, $`dU/dt=2e^2a^2/3c^2`$, where the rest-frame acceleration is given by $`a^{}=\gamma ^2a\gamma ^2c^2/R`$ in terms of lab quantities.
## III Undulator Radiation
An undulator is a device that creates a region of transverse magnetic field that whose magnitude oscillates with spatial period $`\lambda _0`$. This field is constant in time, and is usually lies in a transverse plane (although helical undulators have been built, and are actually somewhat easily to analyze). As an electron with velocity $`v`$ traverses the undulator, its trajectory involves transverse oscillations with laboratory wavelength $`\lambda _0`$, and laboratory frequency $`\omega _0=c/\mathrm{\lambda ̄}_0`$. The oscillating electron then emits undulator radiation.
This radiation is usefully described by first transforming to the average rest frame of the electron, which is done by a Lorentz boost of $`\gamma =1/\sqrt{1(v/c)^2}`$ in the first approximation. The undulator wavelength in this frame is $`\lambda ^{}=\lambda _0/\gamma `$, and the frequency of the oscillator is $`\omega ^{}=\gamma \omega _0`$. The electron emits dipole radiation at this frequency in its average rest frame. The laboratory radiation is the transform of this radiation.
Thus, undulator radiation is readily discussed as the Lorentz transform of a Hertzian dipole oscillator, and the Weizsäcker-Williams approximation does not offer much practical advantage here. However, an analysis of undulator radiation can validate the Weizsäcker-Williams approximation, while also exploring the distinction between undulator radiation and wiggler radiation.
### A A First Estimate
The characteristic angle of undulator radiation in the laboratory is $`\theta _C1/\gamma `$, this being the transform of a ray at $`\theta ^{}=90^{}`$ to the electron’s lab velocity. The radiation is nearly monochromatic, with frequency
$$\omega _C2\gamma \omega ^{}=2\gamma ^2\omega _0,$$
(25)
and wavelength
$$\lambda _C\frac{\lambda _0}{2\gamma ^2}.$$
(26)
The formation length, defined as the distance over which radiation pulls one wavelength ahead of the electron, is $`L_0\gamma ^2\lambda \lambda _0`$, the undulator period. But when the electron advances one period, it continues to oscillate, and the amplitude of the radiation emitted during the second period is in phase with that of the first. Assuming that the radiation from successive period overlaps in space, there will be constructive interference which continues over the entire length of the undulator. In this case, the radiation is not clearly distinct from the near zone of the electron until it leaves the undulator. Hence, the formation length of undulator radiation is better defined as
$$L_0=N_0\lambda _0,$$
(27)
where $`N_0`$ is the number of periods in the undulator.
The frequency spread of undulator radiation narrows as the number of undulator periods increases, and
$$\frac{\mathrm{\Delta }\omega }{\omega _C}\frac{1}{N_0}$$
(28)
We now try to deduce the radiated photon spectrum from the Weizsäcker-Williams approximation (9). The constructive interference over the $`N_0`$ undulator periods implies that the radiated energy will be $`N_0^2`$ times that if there were only one period. So we multiply eq. (9) by $`N_0^2`$ to obtain
$$\frac{dn_\omega }{dl}\frac{N_0^2\alpha }{L_0}\frac{d\omega }{\omega }\frac{\alpha }{\lambda _0},$$
(29)
in the narrow band (28) around the characteristic frequency (25). The radiated power is $`v\mathrm{}\omega _Cc\mathrm{}\omega _C`$ times eq. (29):
$$\frac{dU}{dt}\frac{e^2c\gamma ^2}{\lambda _0^2},$$
(30)
using eq. (25).
This estimate proves to be reasonable only for that part of the range of undulator parameters. To clarify this, we need to examine the electron’s trajectory through the undulator in greater detail.
### B Details of the Electron’s Trajectory
A magnetic field changes the direction of the electron’s velocity, but not its magnitude. As a result of the transverse oscillation in the undulator, the electron’s average forward velocity $`\overline{v}`$ will be less than $`v`$. The boost to the average rest frame is described by $`\overline{\gamma }`$ rather than $`\gamma `$.
In the average rest frame, the electron is not at rest, but oscillates in the electric and magnetic fields $`\stackrel{~}{E}\stackrel{~}{B}=\overline{\gamma }B_0`$, where we use the symbol $`\stackrel{~}{}`$ to indicate quantities in the average rest frame. The case of a helical undulator is actually simpler than that of a linear one. For a helical undulator, the average-rest-frame fields are essentially those of circularly polarized light of frequency $`\stackrel{~}{\omega }=\overline{\gamma }\omega _0`$. The electron moves in a circle of radius $`R`$ at this frequency, in phase with the electric field $`\stackrel{~}{E}`$, and with velocity $`\stackrel{~}{v}`$ and associated Lorentz factor $`\stackrel{~}{\gamma }`$, all related by
$$\frac{\stackrel{~}{\gamma }m\stackrel{~}{v}^2}{R}=\stackrel{~}{\gamma }m\stackrel{~}{v}\stackrel{~}{\omega }=e\stackrel{~}{E}.$$
(31)
From this we learn that
$$\stackrel{~}{\gamma }\stackrel{~}{\beta }=\frac{e\stackrel{~}{E}}{m\stackrel{~}{\omega }c}\frac{eB_0}{m\omega _0c}\eta ,$$
(32)
and hence,
$$\stackrel{~}{\gamma }=\sqrt{1+\eta ^2},\stackrel{~}{\beta }=\frac{\eta }{\sqrt{1+\eta ^2}},$$
(33)
and
$$R=\frac{\stackrel{~}{\beta }c}{\stackrel{~}{\omega }}=\frac{\eta \stackrel{~}{\mathrm{\lambda ̄}}}{\sqrt{1+\eta ^2}}=\frac{\eta \mathrm{\lambda ̄}_0}{\overline{\gamma }\sqrt{1+\eta ^2}}$$
(34)
Thus, the dimensionless parameter $`\eta `$ describes many features of the transverse motion of an electron in an oscillatory field. It is a Lorentz invariant, being proportional to the magnitude of the 4-vector potential.
For a linear undulator, $`\eta `$ is usefully defined as
$$\eta =\frac{eB_{0,\mathrm{rms}}}{m\omega _0c},$$
(35)
where the root-mean-square (rms) average is taken over one period. With the definition (35), the rms values of $`\stackrel{~}{\beta }`$, $`\stackrel{~}{\gamma }`$ and $`R`$ for a linear undulator of strength $`\eta `$ are also given by eqs. (33)-(34).
We can now display a relation for $`\overline{\gamma }`$ by noting that in the average rest frame, the electron’s (average) energy is $`\stackrel{~}{\gamma }mc^2=m\sqrt{1+\eta ^2}c^2`$, while its average momentum is zero there. Hence, on transforming back to the lab frame, we have $`\gamma mc^2=\overline{\gamma }\stackrel{~}{\gamma }mc^2`$, and so
$$\overline{\gamma }=\frac{\gamma }{\sqrt{1+\eta ^2}}.$$
(36)
The transverse amplitude of the motion is obtained from eqs. (34) and (36):
$$R=\frac{\eta \mathrm{\lambda ̄}_0}{\gamma }=2\eta \gamma \mathrm{\lambda ̄}_C,$$
(37)
recalling eq. (26).
### C $`\eta >1`$: Wiggler Radiation
The pitch angle of the helical trajectory is
$$\theta \mathrm{tan}\theta =\frac{R}{\mathrm{\lambda ̄}_0}=\frac{\eta }{\gamma }.$$
(38)
Since the characteristic angle of the radiation is $`\theta _C1/\gamma `$, we see that the radiation from one period of the oscillation does not overlap the radiation from the next period unless
η
<
1.
<
𝜂1\eta\mathrel{\vbox{\kern 0.0pt\hbox{$<$} \kern 0.0pt\hbox{$\sim$} }}1. (39)
Hence, there is no constructive interference and consequent sharpening of the frequency spectrum unless condition (39) is satisfied.
For $`\eta >1`$, the radiation is essentially the sum of synchrotron radiation from $`N_0`$ separate magnets each $`\lambda _0`$ long, and this case is called wiggler radiation.
The laboratory frequency of the radiation is now
$$\omega _C2\overline{\gamma }^2\omega _0,$$
(40)
rather than eq. (25). However, in the regime of undulator radiation, (39), there is little difference between the two expressions.
Another aspect of wiggler radiation is that for η
>
1
>
𝜂1\eta\mathrel{\vbox{\kern 0.0pt\hbox{$>$} \kern 0.0pt\hbox{$\sim$} }}1 the motion of the electron in its average rest frame is relativistic, as can be see from eq. (33). In this case, higher multipole radiation becomes important, which appears at integer multiples of frequency $`\omega ^{}`$ in the average rest frame, and at the corresponding Lorentz transformed frequencies in the lab frame. The total radiated power is still given by eq. (30), so the amount of power radiated any particular frequency is less than when $`\eta 1`$.
### D $`\eta <1`$: Weak Undulators
The estimate (30) for the power of undulator radiation holds only if essentially the whole virtual photon cloud around the electron is shaken off. This can be expected to happen only if the amplitude of the electron’s transverse motion exceeds the minimum impact parameter $`b_{\mathrm{min}}\gamma \lambda _C`$ introduced in eq. (4). From eq. (37) we see that the transverse amplitude obeys
$$R\eta b_{\mathrm{min}}.$$
(41)
Thus, for $`\eta `$ less than one, the undulator radiation will be less than full strength. We readily expect that the intensity of weak radiation varies as the square of the amplitude of the motion, so the estimate (30) should be revised as
dUdtη2e2cγ2λ02,(η
<
1).𝑑𝑈𝑑𝑡superscript𝜂2superscript𝑒2𝑐superscript𝛾2superscriptsubscript𝜆02
<
𝜂1{dU\over dt}\approx{\eta^{2}e^{2}c\gamma^{2}\over\lambda_{0}^{2}},\qquad(\eta\mathrel{\vbox{\kern 0.0pt\hbox{$<$} \kern 0.0pt\hbox{$\sim$} }}1). (42)
The radiated power can be calculated exactly via the Larmor formula,
$$\frac{dU}{dt}=\frac{2e^2a^2}{3c^3},$$
(43)
where $`a^{}=eE^{}/m`$ is the acceleration of the electron in its instantaneous rest frame. The electron is moving in a helix with its velocity perpendicular to $`𝐁_0`$, so the electric field in the instantaneous rest frame is $`E^{}=\gamma \beta B_0\gamma B_0`$. Hence,
$$\frac{dU}{dt}\frac{2e^2\gamma ^2}{3c}\left(\frac{eB_0}{mc}\right)^2=\frac{2e^2c\gamma ^2\eta ^2}{3\mathrm{\lambda ̄}_0^2},$$
(44)
in agreement with the revised estimate (42).
In practice, $`\eta 1`$ is the region of greatest interest as it provides the maximum amount of radiation at the fundamental frequency $`\omega _C`$.
## IV Transition Radiation
As a charged particle crosses, for example, a vacuum/metal boundary, its interaction with charges in the material results in their acceleration and hence radiation, commonly called transition radiation. The formation zone extends outwards from each boundary, with formation length given by eq. (7). The number of photons emitted as the particle crosses each boundary is given by eq. (2) as $`\alpha `$ per unit frequency interval. If two boundaries are separated by less than a formation length, interference effects arise that will not be considered here.
The minimum relevant transverse scale, $`b_{\mathrm{min}}`$, is the plasma wavelength $`\mathrm{\lambda ̄}_p=c/\omega _p`$, so the critical frequency is $`\omega _C\gamma \omega _p`$, according to eq. (4). This is well into the x-ray regime. While the characteristic angle of transition radiation is $`1/\gamma `$, there is only a power-law falloff at larger angles, and the optical transition radiation from an intense beam of charged particles can be used to measure the spot size to accuracy of a few $`\lambda `$ .
## V Čerenkov Radiation
When a charged particle moves with velocity $`v`$ in a dielectric medium, Čerenkov radiation is emitted at those wavelengths for which $`v>c/n`$, where $`n(\lambda )`$ is the index of refraction. As the particle approaches lightspeed in a medium, its electric field is compressed into a “pancake”, which deforms into a cone when the the particle velocity exceeds lightspeed. The particle outruns its electric field, which is freed as the Čerenkov radiation. The fields become identifiable as radiation after the particle has moved a formation length, $`L_0=vt_0`$, which is the distance over which the electron pulls one wavelength ahead of the projection of the wave motion onto the electron’s direction. The Čerenkov angle $`\theta _C`$ is defined by the direction of the radiation, which is normal to the conical surface that contains the electric field. As usual, $`\mathrm{cos}\theta _C=c/nv=1/n\beta `$. The formation length is then
$$\lambda =vt_0\frac{c}{n}t_0\mathrm{cos}\theta _C=L_0\mathrm{sin}^2\theta _C.$$
(45)
Thus $`L_0=\lambda /\mathrm{sin}^2\theta _C`$, and the photon spectrum per unit path length from eq. (2) is
$$\frac{dn_\omega }{dl}\frac{\alpha }{L_0}\frac{d\omega }{\omega }\frac{\alpha \mathrm{sin}^2\theta _C}{\lambda }\frac{d\omega }{\omega }\alpha \mathrm{sin}^2\theta _C\frac{d\omega }{c},$$
(46)
as is well-known.
The characteristic angle $`\theta _C`$ of Čerenkov radiation is essentially independent of the Lorentz factor $`\gamma `$ of the charged particle, unlike that for the other radiation processes considered here. Correspondingly, the characteristic transverse length $`b`$ associated with Čerenkov radiation is also largely independent of $`\gamma `$. Rather, the region over which the Čerenkov radiation develops has radius roughly that of the Čerenkov cone after one formation length, i.e., $`b=L_0\mathrm{cos}\theta _C\mathrm{sin}\theta _C\lambda /\mathrm{tan}\theta _C`$. This is large only near the Čerenkov threshold where the radiated intensity is very small.
That the formation radius for Čerenkov radiation is of order $`\lambda `$ is supported by an analysis of a particle moving in vacuum along the axis of a tube inside a dielectric medium; the calculated Čerenkov radiation is negligible at wavelengths larger than the radius of the tube.
Čerenkov radiation is a form of energy loss for a particle passing through a medium, and is related to so-called ionization loss (see, for example, secs. 13.1-4 of ). The latter is important for frequencies higher than the ionization potential (divided by $`\mathrm{}`$) of the medium, for which the index of refraction is typically less than one. At frequencies that can cause ionization, the Čerenkov effect is insignificant. The transverse scale of the ionizing fields grows with $`\gamma `$ due to relativistic flattening, but shielding due to the induced dielectric polarization of the medium results in an effective transverse scale $`b\sqrt{\gamma }`$ for these fields when $`\gamma 1`$.
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# Computing with Highly Mixed States
## 1 Introduction
Ideally, a quantum computation is a sequence of local unitary transformations applied to a register of qubits which are initially in the state $`|0^n`$; followed by a measurement. Initializing the state of the quantum register is the biggest challenge in NMR quantum computing (which is perhaps the most advanced technology in terms of the scale of experiments performed to date ). The difficulty is that the register is actually initially in (approximately) the binomial distribution over pure states $`|x`$, in which each qubit is independently in the state $`|0`$ with probability $`\frac{1+ϵ}{2}`$; the currently achievable polarization $`ϵ`$ is quite small. There are currently two ways of implementing quantum computation in this technology. The first is used in current experiments , but does not scale beyond several qubits — the output signal decreases exponentially in the number of qubits in the quantum register.<sup>1</sup><sup>1</sup>1 The exponential decay in signal to noise ratio in any scheme that embeds virtual pure states on an $`n`$-qubit quantum computer with one clean qubit is unavoidable, due to the result . The second does scale, but is not feasible at the currently achievable values of polarization in liquid NMR . An intriguing third possibility was raised in . Suppose we start with one qubit in the pure state $`|0`$ in tensor product with $`n1`$ qubits in a maximally mixed state (i.e. in a uniform distribution over basis states $`|x`$). Is it possible to simulate general quantum computation by effecting a sequence of elementary quantum operations on this register? If the answer were affirmative, this would yield a procedure that would both scale and be currently feasible using the scalable initialization procedure to convert the initial binomial state to a state where the last $`n1`$ qubits are maximally mixed and the first bit has high polarization (the strength of the output signal is now proportional to this polarization). This is the question we focus on in this paper.
It is easy to see that if all $`n`$ qubits are in the maximally mixed state then no computation is possible. This is because applying any unitary transformation to this mixture leaves it invariant. This simple argument stands in striking contrast to the difficulty of the seemingly very similar case, in which just a single qubit is in a pure state, while all the others are maximally mixed. Since the initial state of the register is completely specified, the only real input in this model is the sequence of elementary quantum operations. So, given a quantum circuit $`C`$ which we would like to simulate on an input $`x`$, we wish to know whether there is a sequence of elementary quantum operations on the $`n`$ qubit register, which first prepares a quantum (mixed) state which encodes $`x`$, and then simulates $`C`$ on it. Of course, we will require that these mixed state encodings of basis states $`x`$ be distinguishable by some measurement with non-trivial probability.
Our main result shows that the above is impossible unless $`|x|O(\mathrm{log}n)`$, showing that the simulation is no more efficient than an exhaustive classical calculation. The technique used to show this uses some information about the representations of the symmetric group. The appendix 6 provides some necessary notions from representation theory.
We also show that using a 3-bit register it is possible to compute every language in NC1. This should give some indication of why the impossibility result is so much harder than for the case when all $`n`$ qubits are maximally mixed.
## 2 NC1
We begin by showing that in this model, even using a 3-bit register, we can compute every language in NC1.
Recall that the initial state of the register is a uniform distribution over the four $`3`$-bit strings starting with a $`0`$. In our simulation of NC1, all our operations will simply permute basis states.
###### Proposition 1
A $`3`$-bit quantum computer initialized with one clean qubit can recognize every language in NC1.
Our simulation is based on Barrington’s result that NC1 can be simulated by a width $`5`$ permutation branching program . The main idea is quite simple: let the $`5`$ states of the permutation branching program be represented by the states $`|000`$ through $`|100`$ of the $`3`$-qubit register. Without loss of generality assume that the permutation branching program accepts if the permutation it effects is the identity, and rejects if the permutation it effects is the transposition $`(000,100)`$. It is easy to simulate the permutation branching program by a sequence of elementary quantum operations. Now if we measure the first qubit in the register, then in the case that the permutation branching program accepts – i.e. the permutation effected is the identity – then measuring the first qubit in the register yields a $`0`$ with probability $`1`$. On the other hand, if the permutation branching program rejects, then measuring the first qubit in the register yields a $`1`$ with probability $`1/4`$. $`\mathrm{}`$
It is illuminating to try to extend this simulation to QNC1. First notice that in Barrington’s procedure for simulating NC1, each wire in the NC1 circuit is simulated at some stage in the branching program. In the case of a QNC1 circuit, the state of a wire is given by a qubit, which is, in general, entangled with the qubits carried by the other wires in the circuit. Therefore the state of this wire cannot be expressed in isolation, and there appears to be no alternative to creating that entangled state as part of any simulation. Thus the entire approach breaks down. One way to carry out such a construction, might be to apply a superposition of operations at each step: this extends the state space of the quantum computer and effectively provides many more clean qubits, making the model meaningless. Moreover all proposed implementations of quantum computation involve a classical, time-varying sequence of operations, applied to a quantum register. Since the control is classical, in any oblivious simulation the entangled quantum state of the simulated circuit must be encoded within the quantum register.
## 3 Limit on Computability
We are given a quantum computer with an $`n`$ qubit register, with one bit initialized to $`|0`$ and the rest of the $`n1`$ qubits in a maximally mixed state. We would like to simulate an $`m`$ qubit quantum circuit $`C`$ on input string $`x`$ using this model. If we wish to do an oblivious simulation, as sketched in the previous section, we must encode an arbitrary $`m`$ qubit state into the uninitialized $`n`$ qubit register. To do so, it is sufficient to consider the $`2^m`$ basis states of the $`m`$ wires, and encode them as distinguishable states of the uninitialized $`n`$ qubit register (for this to be an efficient encoding, we should have $`nO(poly(m))`$). The states must be distinguishable in the following sense: since we can prepare several copies of any state by repeating the simulation, we only require that there be a sequence of measurements on $`O(poly(n))`$ many copies of the state, that (with high probability) uniquely identify the state. Indeed, it is possible to do this with $`n=m`$, as follows: take the subspaces, spanned by the basis vectors in the sets $`A_b=\{x\{0,1\}^n:xb=0\mathrm{mod}\mathrm{\hspace{0.33em}2}\}`$ for $`b\{0,1\}^n`$.
However, to perform an oblivious simulation, the encoding must satisfy another property – permutability. The quantum circuit $`C`$ might carry out any unitary operation on its quantum state, and in particular an arbitrary permutation on its classical states. Again it is not hard to demonstrate an efficient encoding that satisfies this permutability condition, without distinguishability: take the subspaces spanned by the basis vectors in the sets $`A_b=\{x=(x_1\mathrm{}x_n)\{0,1\}^n:x_1=0\text{ or }(x_2\mathrm{}x_n)=b\}`$ for $`b\{0,1\}^{n1}`$.
However, it is not possible to construct an efficient encoding that satisfies both conditions simultaneously. This is the content of the following theorem.
Let $`M=2^m`$ be the total number of basis states of the ideal quantum computer which is being simulated. Note that each $`X𝒳`$ encoding one of these, is a subspace of dimension $`2^{n1}`$ within the Hilbert space $`^{2^n}`$ of the computer. If the computer has $`k`$ clean qubits, then $`X`$ is of dimension $`2^{nk}`$.
###### Theorem 2
Suppose that computations on $`m`$ qubits can be obliviously simulated in an $`n`$-qubit, $`k`$-clean-qubit computer in our model, and that $`dim(XY)/dim(X)<1\frac{1}{\text{poly}(m)}`$ for every pair of input encodings $`X,Y𝒳`$. Then $`m(2k+\mathrm{log}n)(1+o(1))`$.
This incidentally implies that the computation of an $`n`$-qubit, $`k`$-clean-qubit computer can be simulated by a classical computer with a $`\text{poly}(n2^k)`$ computational overhead.
It may be illuminating to consider a simpler, classical analogue of our problem. A classical circuit (taking inputs in $`\{0,1\}^n`$) composed of reversible gates executes a permutation of $`\{0,1\}^n`$. The analogous problem (just considering the case $`k=1`$) is that we can only represent inputs as uniform probability distributions over a set of half the elements of $`\{0,1\}^n`$. (In the quantum case this corresponds to axis-parallel subspaces of dimension $`2^{n1}`$.) The question is, what is the largest number of such subsets (probability distributions) which such a circuit can permute at will. It is also essential that the probability distributions be readily distinguishable by sampling, in other words the subsets must have small intersection. It is possible (though we omit it in this extended abstract) to provide a strictly combinatorial argument expressing the fact that this task is impossible for more than $`\text{poly}(n)`$ subsets, because of the tension between the two requirements (permutability and distinguishability). The large size of the subsets means that we have far more constraints than we have degrees of freedom. The combinatorial argument shows that if the requirement of full permutability is imposed, and we have a superpolynomial (in $`n`$) number of subsets, then the symmetric difference of every two sets must be a vanishing fraction of the size of the sets. The two types of sets $`\{A_b\}`$ described above, however, separately achieve distinguishability and permutability.
In the quantum case we have arbitrary subspaces in place of “subsets” (or correspondingly axis-parallel subspaces). And the circuit of course can perform not just permutations of the basis, but general unitary operations. In sharp contrast with the classical case, two subspaces of half the dimensionality of the space typically will not intersect. Nevertheless, the large dimension of the subspaces imposes strict constraints on an operator which must permute them; the difficulty is in formulating the incompatibility of these requirements when the number of subspaces is large and the subspaces are required to be very distinct.
Proof:
By assumption, there are unitary operators (each corresponding to some sequence of steps in the computer) permuting $`𝒳`$ in all ways. Let $`f_\pi `$ be the unitary operator corresponding to a permutation $`\pi S_M`$. If we have $`f_{\pi \sigma }=f_\pi f_\sigma `$ for all $`\pi `$ and $`\sigma S_M`$, then the operators $`f_\pi `$ form a representation of $`S_M`$ and we can apply the representation theory of the symmetric group.
Actually, the situation is slightly more complicated.
Let $`U`$ be the unitary group on $`^{2^n}`$. Let $`H`$ be the subgroup of $`U`$ acting on $`𝒳`$, i.e. carrying any $`X𝒳`$ to some $`Y𝒳`$. Let $`G`$ be the subgroup of $`H`$ that fixes all of $`𝒳`$; thus $`G`$ is normal in $`H`$.
It is apparent that $`H/GS_M`$, but although this means that $`H`$ can permute the subspaces $`𝒳`$ in arbitrary ways, it is different from saying that there is a subgroup of $`H`$ isomorphic to $`S_M`$ (or in other words that we can pick elements of $`H`$ so as to have these operators compose properly).
## 4 Proof of Theorem 2: the simple case
First, we show how to prove Theorem 2 if we can select transformations $`f_\pi `$ so that they form a representation ($`f_\pi f_\sigma =f_{\pi \sigma }`$). The more general case will be handled in the next section. Appendix (section 6) explains the notions of representation theory used in this and the next section.
We show that every pair $`X,Y\widehat{𝒳}`$ have a substantial intersection. Consider the decomposition of $`^{2^n}`$ into irreducible representations $`\rho _1\mathrm{}\rho _k`$. Let $`N=2^n`$.
###### Lemma 3
Either the first row or the first column of the Young diagram of each $`\rho _i`$ is of length more than $`Mcn`$.
Proof:
This results follows from a theorem by Rasala:
###### Theorem 4
\[10, pp.151-152\]
1. Let $`AM/2`$ and $`\rho `$ be an irreducible representation of $`S_M`$ such that the first row of the Young diagram of $`\rho `$ is of length exactly $`MA`$. Then,
$$dim\rho \phi _A(M)$$
where $`\phi _A(M)=\left(\genfrac{}{}{0pt}{}{M}{A}\right)\left(\genfrac{}{}{0pt}{}{M}{A1}\right)=\frac{M2A1}{MA1}\left(\genfrac{}{}{0pt}{}{M}{A}\right)`$ is the dimension of the irreducible representation corresponding to the partition $`(MA,A)`$.
2. If $`\rho `$ is an irreducible representation with both the first row and the first column of length at most $`M/2`$, then
$$dim\rho \phi _{M/2}(M).$$
This theorem means that any representation with the first row of the Young diagram having length at most $`Mk`$, $`kM/2`$ has dimension at least
$$\underset{B:ABM/2}{\mathrm{min}}\phi _B(M).$$
Simple algebra shows that this expression is minimized by $`B=A`$ if $`AM/2c\sqrt{M}`$ for some constant $`c`$ and $`B=M/2`$ if $`A>M/2c\sqrt{M}`$.<sup>2</sup><sup>2</sup>2In the second case, the lowest-dimensional representation actually has the first row less than $`MA`$. It is quite surprsing because, in most cases, removing a square from the first row of a Young diagram and adding a square somewhere else increases the dimension.
To deduce our lemma, assume that the Young diagram of an irreducible representation of $`S_M`$ has both first row and column of length at most $`MA`$. We show that $`N2^A`$. Consider two cases:
Case 1: $`AM/2c\sqrt{M}`$.
Notice that $`AM\sqrt{M}`$ because otherwise the Young diagram would fit into a square with a side less than $`\sqrt{M}`$ and area less than $`M`$. Theorem 4 implies that
$$dim\rho \phi _{M/2}(M)=\mathrm{\Omega }\left(\frac{2^M}{M\sqrt{M}}\right)=2^{M\frac{3\mathrm{log}M}{2}O(1)}>2^A.$$
Case 2: $`AM/2c\sqrt{M}`$. Then,
$$dim\rho \phi _A(M)=\frac{M2A+1}{MA+1}\left(\genfrac{}{}{0pt}{}{M}{A}\right)\frac{1}{M}\left(\frac{M}{A}\right)^A$$
$$=\frac{1}{M}\left(\frac{M}{2A}\right)^A2^A2^A.$$
$`\mathrm{}`$
Another lower bound on the longest row or column of a low dimension representation (for a different range of parameters) was given by Mischenko.
Next consider the stabilizer of $`X`$ in $`S_M`$, which is isomorphic to $`S_{M1}`$, and which we will denote $`S_{M1}^X`$. $`X`$ decomposes into irreducible representations $`V_1,\mathrm{}V_{\mathrm{}}`$ of $`S_{M1}^X`$. $`V_1`$ is carried by $`S_M`$ into each $`Y\widehat{𝒳}`$, and all these copies of $`V_1`$ are contained within some irreducible $`W`$ of $`S_M`$ in $`^{2^n}`$.
By the previous lemma, the Young diagram of $`W`$ has a long first column or row. We use the following fact from the representation theory of the symmetric group: when we restrict an irreducible representation0 $`\rho _\lambda `$ of $`S_M`$ of shape $`\lambda `$ to a subgroup $`S_{M1}S_M`$, it decomposes into irreducibles of $`S_{M1}`$ in the following way:
$$\rho =\underset{\lambda ^{}}{}\rho _\lambda ^{}$$
where $`\lambda ^{}`$ ranges over all shapes of size $`M1`$ that can be obtained by deleting an “inside corner” from $`\lambda `$. (An inside corner is simply a point of the shape whose deletion leaves a legal shape.)
We now use:
###### Lemma 5
Suppose the shape $`\lambda `$ has a first row (column) of length $`|\lambda |\mathrm{}`$ for $`\mathrm{}<|\lambda |/2`$. Let $`\lambda _1`$ denote the “$`\lambda ^{}`$” obtained by deleting the last element of the first row (column). Then $`dim(\rho _{\lambda _1})\frac{|\lambda |2\mathrm{}}{|\lambda |}dim(\rho _\lambda )`$.
Proof: Consider the ratio $`\frac{dim\rho _{\lambda _1}}{dim\rho _\lambda }=\frac{1}{|\lambda |}\frac{_{x\lambda }|x|}{_{x\lambda _1}|x|}`$. In the last ratio, points $`x`$ outside of the last row (column) appear identically in the numerator and denominator. Moreover for each $`x`$ in the first row (column) in the numerator other than the very last point (which contributes a factor of $`1`$ in the numerator and is absent in the denominator), the ratio between its contributions in the numerator and denominator is $`\frac{|x|}{|x|1}`$. Just examining the $`|\lambda |2\mathrm{}1`$ points of the first row furthest from the upper-left corner (and excepting the last point), we obtain a lower bound on these contributions of $`_{i=1}^{|\lambda |2\mathrm{}1}\frac{i+1}{i}=|\lambda |2\mathrm{}`$. Overall therefore $`\frac{dim\rho _{\lambda _1}}{dim\rho _\lambda }\frac{|\lambda |2\mathrm{}}{|\lambda |}`$. $`\mathrm{}`$
###### Lemma 6
Let $`f_\pi `$ be an $`N`$-dimensional representation of $`S_M`$ that acts as a permutation representation on a collection of $`M`$ subspaces $`\widehat{𝒳}`$. Then, for any $`X,Y\widehat{𝒳}`$,
$$dimXdimXY\frac{2cn}{M}N.$$
(We write $`n=\mathrm{lg}N`$ and $`m=\mathrm{lg}M`$.)
Proof: We decompose $`f_\pi `$ into irreducible representations. Let $`W`$ be one of these irreducible representations. We show that $`dimXWdimXYW\frac{2cn}{M}dimW`$.
We look at $`W`$ as a representation of $`S_{M1}^X`$. Let $`V`$ be the highest dimensional irreducible representation of $`S_{M1}^X`$ within $`W`$. If $`M<2cn`$ then the assertion is trivial. Otherwise $`Mcn>M/2`$ and the hypothesis of lemma 5 is satisfied, implying that $`dimV\frac{M2cn}{M}dimW`$. Consider two cases:
Case 1: $`VX`$.
Take $`\pi S_M`$ such that $`Y=f_\pi (X)`$. Then, $`YW=f_\pi (XW)`$. Therefore, $`dim(YW)=dim(XW)dimV\frac{M2cn}{M}dimW`$ and
$$dimXWdimXYW\frac{2cn}{M}dimW.$$
Case 2: $`VX`$.
Then, $`VX=0`$ because $`VX`$ is invariant under $`S_{M1}^X`$ and $`V`$ is irreducible. $`VX=0`$ implies
$$dimV+dimXWdimW.$$
Together with $`dimV\frac{M2cn}{M}dimW`$, this implies
$$dimXWdimXYWdimXW\frac{2cn}{M}dimW.$$
The lemma follows by summation over all irreducible $`W`$. $`\mathrm{}`$
If $`f_\pi `$ form a representation, Lemma 6 almost immediately implies Theorem 2. Namely, we have
$$\frac{dimXY}{dimX}=\frac{dimX(dimXdimXY)}{dimX}$$
$$\frac{2^{nk}\frac{2cn}{M}2^n}{2^{nk}}=1\frac{2^{k+1}cn}{M}.$$
If this is at most $`1\frac{1}{\text{poly}(m)}`$, then $`m(k+\mathrm{log}n)(1+o(1))`$. $`\mathrm{}`$
## 5 Proof of Theorem 2: the difficult case
### 5.1 Proof outline
Next, we deal with the case when $`f_\pi f_\sigma f_{\pi \sigma }`$ for some $`\pi `$ and $`\sigma S_M`$. Let $`G`$ be the group of transformations that map every subspace $`X\widehat{𝒳}`$ to itself. Then, $`f_\pi f_\sigma f_{\pi \sigma }^1`$ is an element of $`G`$ for any $`\pi ,\sigma S_M`$. We would like to modify $`f`$ so that this element becomes identity for all $`\pi `$ and $`\sigma S_M`$. Then, $`f_\pi f_\sigma =f_{\pi \sigma }`$, i.e., $`f_\pi `$ would form a representation of $`S_M`$ and we would be able to analyse this representation similarly to the previous section.
To achieve this, we look at $`\text{ }\mathrm{C}^{2^n}`$ as a representation of $`G`$ and express $`\text{ }\mathrm{C}^{2^n}`$ as $`V_1V_2\mathrm{}V_k`$, with $`V_i`$ corresponding to different types of irreducible representations of $`G`$.
Then, we compose each $`f_\pi `$ with an appropriate $`g_\pi G`$. The resulting transformation $`f_\pi ^{}=g_\pi f_\pi `$ still implements the same permutation $`\pi `$ of $`\widehat{𝒳}`$ because $`g_\pi `$ maps every $`X\widehat{𝒳}`$ to itself. We can choose the transformations $`g_\pi `$ so that, on every $`V_i`$, $`f_\pi ^{}f_\sigma ^{}`$ is the same as $`f_{\pi \sigma }^{}`$ up to a phase ($`f_{\pi \sigma }^{}=c_{\pi ,\sigma ,i}f_\pi ^{}f_\sigma ^{}`$ for some unit $`c_{\pi ,\sigma ,i}\text{ }\mathrm{C}`$).
The next step is eliminating the phase factors $`c_{\pi ,\sigma ,i}`$. This is done by considering a larger space $`V_1V_1^{}+\mathrm{}+V_kV_k^{}`$ and transformations $`f_\pi ^{\prime \prime }=f_\pi ^{}(f_\pi ^{})^{}`$ on this larger space. Then, the phase factors $`c_{\pi ,\sigma ,i}`$ (from $`f^{}`$) and $`c_{\pi ,\sigma ,i}^{}`$ (from $`(f^{})^{}`$) cancel out and we get $`f_{\pi \sigma }^{\prime \prime }=c_{\pi ,\sigma ,i}c_{\pi ,\sigma ,i}^{}f_\pi ^{\prime \prime }f_\sigma ^{\prime \prime }=f_\pi ^{\prime \prime }f_\sigma ^{\prime \prime }`$. Thus, $`f_\pi ^{\prime \prime }`$ form a representation of $`S_M`$ on the linear space $`V_1V_1^{}+\mathrm{}+V_kV_k^{}`$. This representation can be analysed similarly to section 4, obtaining lower bounds on intersections of invariant subspaces.
### 5.2 Representation up to phases $`c_{\pi ,\sigma ,i}`$
Let $`G`$ be the group of unitary transformations that fix every one of the subspaces $`X\widehat{𝒳}`$.
Then, $`^{2^n}`$ is a representation of $`G`$ and all $`h\widehat{𝒳}`$ are invariant subspaces. (They are fixed by every element of $`G`$ according to the definition of $`G`$.) These invariant subspaces decompose into irreducible invariant subspaces.
Consider all the irreducible invariant subspaces of $`^{2^n}`$. Split them into equivalence classes consisting of isomorphic irreducible subspaces. Let $`E_1,\mathrm{},E_k`$ be these equivalence classes. Let $`V_1`$ be the subspace of $`^{2^n}`$ spanned by all the irreducible subspaces in $`E_1`$ (i.e., the subspace spanned by all the vectors belonging to at least one subspace in $`E_1`$). Let $`V_2`$, $`\mathrm{}`$, $`V_k`$ be defined similarly.
###### Claim 1
If $`i,j\{1,\mathrm{},k\}`$ and $`ij`$, then $`V_iV_j`$.
Therefore, $`^{2^n}=V_1V_2\mathrm{}V_k`$. Next, we show that transformations $`f_\pi `$ map each $`V_i`$ to some (possibly different) $`V_i^{}`$.
###### Claim 2
Let $`V`$ be an invariant subspace. Then, $`f_\pi (V)`$ is invariant as well. If $`V`$ is irreducible, $`f_\pi (V)`$ is irreducible. Moreover, if $`V`$ and $`V^{}`$ are two isomorphic irreducible subspaces, $`f_\pi (V)`$ and $`f_\pi (V^{})`$ are isomorphic as well.
Proof: The map $`gf_\pi gf_\pi ^1`$ is an automorphism of $`G`$. If $`V`$ is invariant under the action of $`g`$, $`f_\pi (V)`$ is invariant under the action of $`f_\pi gf_\pi ^1`$. Therefore, if $`V`$ is invariant under $`G`$, so is $`f_\pi (V)`$.
If $`f_\pi (V)`$ is not irreducible, it decomposes into two or more invariant subspaces: $`f_\pi (V)=W_1W_2`$. Then, $`f_\pi ^1(W_1)`$ is invariant as well, implying that $`V`$ is not irreducible.
Finally, let $`h:VV^{}`$ be a $`G`$-isomorphism of $`V`$ and $`V^{}`$ (an isomorphism that commutes with the action of $`G`$). Let $`h^{}:f_\pi (V)f_\pi (V^{})`$ be defined by $`h^{}=f_\pi hf_\pi ^1`$. Then, for any $`g=f_\pi g^{}f_\pi ^1`$, we have
$$h^{}g=(f_\pi hf_\pi ^1)(f_\pi g^{}f_\pi ^1)=f_\pi hg^{}f_\pi ^1=f_\pi g^{}hf_\pi ^1=gh^{}$$
and every $`gG`$ can be expressed in the form $`f_\pi g^{}f_\pi ^1`$. Therefore, $`h^{}`$ is a $`G`$-isomorphism of $`f_\pi (V)`$ and $`f_\pi (V^{})`$. $`\mathrm{}`$
Remark. $`V`$ does not have to be isomorphic to $`f_\pi (V)`$ as a representation of $`G`$. $`f_\pi `$ establishes the isomorphism of $`g`$ on $`V`$ with $`f_\pi gf_\pi ^1`$ on $`f_\pi (V)`$, but $`f_\pi gf_\pi ^1`$ does not have to equal $`g`$ on $`f_\pi (V)`$.
###### Claim 3
For every $`i\{1,\mathrm{},k\}`$ there is an $`i^{}`$ such that $`f_\pi (V_i)=V_i^{}`$.
Proof: By Claim 2, every two isomorphic irreducible subspaces get mapped to isomorphic irreducible subspaces. Therefore, all subspaces in $`E_i`$ get mapped to subspaces in the same $`E_i^{}`$ and $`f_\pi (V_i)V_i^{}`$. Similar reasoning applied to $`f_\pi ^1`$ implies $`f_\pi ^1(V_i^{})V_i`$. $`\mathrm{}`$
For each $`i\{1,\mathrm{},k\}`$, $`V_i`$ is the direct sum of some number of isomorphic irreducible subspaces: $`V_i=V_{i1}V_{i2}\mathrm{}V_{ij_i}`$. We fix $`G`$-isomorphisms $`h_{ijj^{}}`$ between $`V_{ij}`$ and $`V_{ij^{}}`$ so that $`h_{ij^{}j^{\prime \prime }}h_{ijj^{}}=h_{ijj^{\prime \prime }}`$. (By Schur’s lemma, each of these isomorphisms is unique up to a multiplicative constant. The isomorphisms can be made to compose properly by adjusting these constants. Note of course that $`h_{ijj}`$ is the identity.)
###### Claim 4
$`WV_i`$ is an irreducible invariant subspace if and only if
$$W=\{a_jx+a_{j+1}h_{ij(j+1)}(x)+\mathrm{}+a_{j_i}h_{ijj_i}(x)|xV_{ij}\}$$
for some $`j\{1,\mathrm{},j_i\}`$ and $`a_j,\mathrm{},a_{j_i}\text{ }\mathrm{C}`$.
Proof: “If” part:
Invariance:
$$g(\underset{\mathrm{}=j}{\overset{j_i}{}}a_{\mathrm{}}h_{ij\mathrm{}}(x))=\underset{\mathrm{}=j}{\overset{j_i}{}}a_{\mathrm{}}g(h_{ij\mathrm{}}(x))=\underset{\mathrm{}=j}{\overset{j_i}{}}a_{\mathrm{}}h_{ij\mathrm{}}(g(x))$$
because each $`h_{ij\mathrm{}}`$ is a $`G`$-isomorphism.
$`W`$ is irreducible because, if $`W_1W`$ and $`W_1`$ is invariant, then
$$\{x|a_jx+a_{j+1}h_{ij(j+1)}(x)+\mathrm{}+a_{j_i}h_{ijj_i}(x)W_1\}$$
is an invariant subspace of $`V_{ij}`$; but $`V_{ij}`$ is irreducible and $`dim(W)=dim(V_{ij})`$.
“Only if” part:
Let $`W`$ be an irreducible invariant subspace of $`V_i`$. Let $`x^{}W`$. Then, we can write $`x^{}`$ as $`x_1^{}+\mathrm{}+x_{j_i}^{}`$, $`x_1^{}V_{i1}`$, $`\mathrm{}`$, $`x_{j_i}^{}V_{ij_i}`$. If $`x^{}x^{\prime \prime }W`$, then for any index $`j`$, $`x_j^{}x_j^{\prime \prime }`$ or $`x_j^{}=x_j^{\prime \prime }=0`$. (Otherwise, $`W_\mathrm{}jV_i\mathrm{}`$ is a nontrivial subspace of $`W`$. It is invariant because $`V_i`$ and $`V_{ij}`$ are invariant. Contradiction with the irreducibility of $`W`$.)
Let $`j`$ be the smallest index for which there is an $`x^{}W`$ with $`x_j^{}0`$. Then, for every $`xV_{ij}`$, there is an $`x^{}W`$ with $`x_j^{}=x`$. (For, if $`A`$ and $`B`$ are invariant subspaces of a unitary representation, the projection of $`A`$ onto $`B`$ is invariant. Apply this with $`A=W`$ and $`B=V_{ij}`$, then use the irreducibility of $`V_{ij}`$.)
The above considerations allow us to define the mapping $`h_{jj^{}}:V_{ij}V_{ij^{}}`$ by $`h_{jj^{}}(x_j^{})=x_j^{}^{}`$. By the definition, $`h_{jj^{}}(g(x_j^{}))=h_{jj^{}}((g(x^{}))_j)=(g(x^{}))_j^{}=g(x_j^{}^{})=g(h_{jj^{}}(x_j^{})`$, so $`h_{jj^{}}`$ is a $`G`$-isomorphism. By Schur’s lemma, this implies that $`h_{jj^{}}=a_j^{}h_{ijj^{}}`$ for some $`a_j^{}\text{ }\mathrm{C}`$. $`\mathrm{}`$
In general, if $`W`$ is any irreducible invariant subspace of $`^{2^n}`$, then $`W`$ must be in the form described by claim 4 for some $`i`$. ($`W`$ belongs to some equivalence class $`E_i`$ and therefore is contained in the corresponding $`V_i`$.)
For each $`V_{i1}`$ ($`i\{1,\mathrm{},k\}`$), we fix an orthonormal basis $`v_{i1},\mathrm{},v_{it}`$. This also fixes a related basis $`h_{i1j}(v_{i1})`$, $`\mathrm{}`$, $`h_{i1j}(v_{it})`$ for each $`V_{ij}`$. Moreover, we also get a similar basis
$$\underset{\mathrm{}=j}{\overset{j_i}{}}a_{\mathrm{}}h_{i1\mathrm{}}(v_{i1}),\mathrm{},\underset{\mathrm{}=j}{\overset{j_i}{}}a_{\mathrm{}}h_{i1\mathrm{}}(v_{it})$$
for every invariant irreducible $`WV_i`$ because any such $`W`$ can be written in the form given by claim 4. We call these bases designated.
This designated basis is exactly the basis for $`W`$ that can be obtained by applying the isomorphism between $`V_{i1}`$ and $`W`$ to the basis for $`V_{i1}`$. Moreover, if $`W,W^{}`$ are two isomorphic irreducible subspaces, the designated basis for $`W`$ is mapped to the designated basis for $`W^{}`$ by the isomorphism between $`W`$ and $`W^{}`$.
We are going to impose the following condition on $`f_\pi ^{}`$:
Condition. Let $`W`$ be an irreducible representation of $`G`$ and $`w_1,\mathrm{},w_l`$ be the designated basis of $`W`$. Let $`w_1^{},\mathrm{},w_l^{}`$ be the designated basis of $`f_\pi ^{}(W)`$. Then, there exists $`c`$, $`|c|=1`$ such that $`f_\pi ^{}(w_1)=cw_1^{}`$, $`\mathrm{}`$, $`f_\pi ^{}(w_l)=cw_l^{}`$.
Next, we show that this condition suffices to guarantee $`f_\pi ^{}f_\sigma ^{}=c_{\pi ,\sigma ,i}f_{\pi \sigma }^{}`$ on every $`V_i`$ and that any $`f_\pi `$ that permutes $`X\widehat{𝒳}`$ without satisfying this condition can be transformed into $`f_\pi ^{}`$ that satisfies the condition and still permutes the subspaces in the same way.
First, we show that it is enough to ensure that the designated basis of $`V_{i1}`$ is mapped correctly for every $`i\{1,\mathrm{},k\}`$.
###### Claim 5
Assume that the condition is true for $`W=V_{i1}`$. Then, it is also true for any irreducible $`WV_i`$.
Proof: Let $`h`$ be the isomorphism between $`V_{i1}`$, $`W`$. Note that $`h`$ maps the designated basis of $`V_{i1}`$ to the designated basis of $`W`$.
Then (by claim 2) $`f_\pi ^{}h(f_\pi ^{})^1`$ is an isomorphism between $`f_\pi ^{}(V_{i1})`$ and $`f_\pi ^{}(W)`$. We know that there is an isomorphism between these two irreducibles that maps the designated basis of one of them to the designated basis of the other. By Schur’s lemma, any two isomorphisms of irreducible subspaces can differ only by a multiplicative constant $`c`$. The unitarity of $`f_\pi ^{}h(f_\pi ^{})^1`$ implies that $`|c|=1`$.
Therefore, $`f_\pi ^{}h(f_\pi ^{})^1`$ maps the designated basis of $`f_\pi ^{}(V_{i1})`$ to $`c`$ times the designated basis of $`f_\pi ^{}h(f_\pi ^{})^1(f_\pi ^{}(V_{i1}))=f_\pi ^{}h(V_{i1})=f_\pi ^{}(W)`$. We know that $`(f_\pi ^{})^1`$ maps the designated basis of $`f_\pi ^{}(V_{i1})`$ to the designated basis of $`V_{i1}`$ and that $`h`$ maps the designated basis of $`V_{i1}`$ to the designated basis of $`W`$. This implies that $`f_\pi ^{}`$ maps the designated basis of $`W`$ to $`c`$ times the designated basis of $`f_\pi ^{}(W)`$. $`\mathrm{}`$
Next, we show how to transform $`f_\pi `$ into $`f_\pi ^{}`$ that performs the same permutation $`\pi `$ of $`\widehat{𝒳}`$ and maps the designated basis of every $`V_{i1}`$ as required.
Let $`W_1`$, $`\mathrm{}`$, $`W_k`$ be $`f_\pi (V_{11})`$, $`\mathrm{}`$, $`f_\pi (V_{k1})`$. Each of $`W_i`$ lies within one of $`V_1,\mathrm{},V_k`$. Denote this subspace $`V_i^{}`$. Then, for $`ij`$, $`V_i^{}V_j^{}`$. For each $`i\{1,\mathrm{},k\}`$, we define a unitary transformation $`g_{\pi ,i}`$ on $`V_i`$ such that $`g_{\pi ,i^{}}f_\pi `$ maps the designated basis of $`V_{i1}`$ to the designated basis of $`f_\pi (V_{i1})`$.
By Claim 4, the irreducible subspace $`W_i=f_\pi (V_{i1})`$ is just
$$\{a_{i^{}j}x+a_{i^{}(j+1)}h_{i^{}j(j+1)}(x)+\mathrm{}+a_{i^{}j_i^{}}h_{i^{}jj_i^{}}(x)|xV_{i^{}j}\}$$
for some $`j`$. Moreover, the mapping that maps each $`vW_i`$ to its $`V_{i^{}j}`$-component is an isomorphism of $`W_i`$ and $`V_{i^{}j}`$ w.r.t. $`G`$ (similarly to proof of Claim 4).
Let $`v_1,\mathrm{},v_l`$ be the designated basis of $`V_{i1}`$, $`v_1^{},\mathrm{},v_l^{}`$ be $`f_\pi (v_1),\mathrm{},f_\pi (v_l)`$ and $`v_1^{\prime \prime }`$, $`\mathrm{}`$, $`v_l^{\prime \prime }`$ be the $`V_{i^{}j}`$ components of $`v_1^{}`$, $`\mathrm{}`$, $`v_l^{}`$.
Let $`w_1`$, $`\mathrm{}`$, $`w_l`$ be the designated basis of $`V_{i^{}j}`$ and $`g_{\pi ,i^{}j}`$ be the unitary transformation on $`V_{i^{}j}`$ that maps $`v_1^{\prime \prime },\mathrm{},v_l^{\prime \prime }`$ to $`w_1`$, $`\mathrm{}`$, $`w_l`$. We define a unitary transformation $`g_{\pi ,i^{}j^{}}`$ (for every $`j^{}j`$) on $`V_{i^{}j^{}}`$ to be $`h_{i^{}jj^{}\mathrm{`}}g_{\pi ,i^{}j}h_{i^{}jj^{}}^1`$. Finally, we take the transformation $`g_{\pi ,i^{}}`$ of $`V_i^{}`$ that is equal to $`g_{\pi ,i^{}j}`$ on each $`V_{i^{}j}`$. Then, $`g_{\pi ,i^{}}`$ maps
$$v_1^{}=a_{i^{}j}v_1^{\prime \prime }+a_{i^{}(j+1)}h_{i^{}j(j+1)}(v_1^{\prime \prime })+\mathrm{}+a_{i^{}j_i^{}}h_{i^{}jj_i^{}}(v_1^{\prime \prime })$$
to
$$a_{i^{}j}g_{\pi ,i^{}j}(v_1^{\prime \prime })+a_{i^{}(j+1)}g_{\pi ,i^{}(j+1)}h_{i^{}j(j+1)}(v_1^{\prime \prime })+\mathrm{}=$$
$$a_{i^{}j}g_{\pi ,i^{}j}(v_1^{\prime \prime })+a_{i^{}(j+1)}h_{i^{}j(j+1)}g_{\pi ,i^{}j}(v_1^{\prime \prime })+\mathrm{}$$
$$=a_{i^{}j}w_1+a_{i^{}(j+1)}h_{i^{}j(j+1)}(w_1)+\mathrm{}$$
which is exactly the first vector of the designated basis for $`W_i`$. The same is true for $`v_2^{}`$, $`\mathrm{}`$, $`v_l^{}`$, implying that $`g_{\pi ,i^{}}f_\pi `$ maps the designated basis of $`V_{i1}`$ to the designated basis of $`W_i`$.
Now, we take $`g_\pi `$ that is equal to $`g_{\pi ,i}`$ on each $`V_i`$ and take $`f_\pi ^{}=g_\pi f_\pi `$.
###### Claim 6
$`g_\pi `$ preserves all $`X\widehat{𝒳}`$.
Proof: By definition, the restriction $`g_\pi |_{V_i}`$ is equal to $`g_{\pi ,i}`$, and $`g_{\pi ,i}`$ clearly preserves $`V_{i1},\mathrm{},V_{ij_i}`$. Moreover, $`g_{\pi ,i}`$ (and, hence, $`g_\pi `$) preserves any irreducible subspace $`WV_i`$ because any such subspace is in the form of claim 4.
Every $`X\widehat{𝒳}`$ is invariant under $`G`$. Therefore, it decomposes into a direct sum of irreducible subspaces. Each of these subspaces is in one of the classes $`E_1`$, $`\mathrm{}`$, $`E_k`$ and, therefore, lies in one of $`V_1`$, $`\mathrm{}`$, $`V_k`$. This means that it is preserved by $`g_\pi `$. Therefore, $`X`$ which is a direct sum of such irreducible subspaces is preserved by $`g_\pi `$ as well. $`\mathrm{}`$
Hence, $`f_\pi ^{}=g_\pi f_\pi `$ realizes the same permutation $`\pi `$ of $`X\widehat{𝒳}`$ as $`f_\pi `$.
###### Claim 7
On every $`V_i`$, $`f_\pi ^{}f_\sigma ^{}=c_{\pi ,\sigma ,i}f_{\pi \sigma }^{}`$ for some $`c_{\pi ,\sigma ,i}`$.
Proof: This is equivalent to showing that $`(f_{\pi \sigma }^{})^1f_\pi ^{}f_\sigma ^{}`$ is equal to $`c_{\pi ,\sigma ,i}`$ times the identity. To show that, notice that $`(f_{\pi \sigma }^{})^1f_\pi ^{}f_\sigma ^{}`$ maps every subspace $`X\widehat{𝒳}`$ to itself because $`(f_{\pi \sigma }^{})^1`$ performs the inverse of the permutation $`\pi \sigma `$ on $`\widehat{𝒳}`$. Therefore, $`(f_{\pi \sigma }^{})^1f_\pi ^{}f_\sigma ^{}G`$. This means that $`V_{ij}`$ are all preserved by $`(f_{\pi \sigma }^{})^1f_\pi ^{}f_\sigma ^{}`$.
Moreover, $`f_\sigma ^{}`$, $`f_\pi ^{}`$ and $`f_{\pi \sigma }^1`$ all map the designated bases to $`c`$-times designated bases (Claim 5). Therefore, $`(f_{\pi \sigma }^{})^1f_\pi ^{}f_\sigma ^{}`$ maps the designated basis of $`V_{ij}`$ to $`c`$ times the designated basis of $`(f_{\pi \sigma }^{})^1f_\pi ^{}f_\sigma ^{}(V_{ij})=V_{ij}`$.
It remains to show that $`c`$ is the same for all irreducible subspaces $`V_{ij}`$ contained in $`V_i`$. Let $`c_j`$ and $`c_j^{}`$ be the values of $`c`$ for $`V_{ij}`$ and $`V_{ij^{}}`$. Consider the subspace
$$W=\{x+h_{ijj^{}}(x)|xV_{ij}\}.$$
By Claim 4, this is an irreducible invariant subspace. Now, $`(f_{\pi \sigma }^{})^1f_\pi ^{}f_\sigma ^{}`$ maps it to
$$W^{}=\{c_jx+c_j^{}h_{ijj^{}}(x)|xV_{ij}\}=$$
$$\{x+\frac{c_j}{c_j^{}}h_{ijj^{}}(x)|xV_{ij}\}.$$
The invariance of $`W`$ means that $`W^{}=W`$ and $`c_j=c_j^{}`$.
Therefore, $`c_j`$ are all equal. This means that $`(f_{\pi \sigma }^{})^1(x)=c_jx`$ for all $`xV_i`$ because the designated bases of $`V_{ij}`$ together form a basis for entire subspace $`V_i`$.
Unfortunately, arguments of this type (composing $`f_\pi `$ with an appropriate transformation that fixes all $`U_i`$) cannot be used to eliminate phases $`c_{\pi ,\sigma ,i}`$.
The reason for this is that there exist so-called projective representations. A projective representation is a set of maps $`f_\pi `$ such that $`f_\pi f_\sigma =c_{\pi ,\sigma }f_{\pi \sigma }`$, $`c_{\pi ,\sigma }\text{ }\mathrm{C}`$. It is known that the symmetric group has projective representations which are not equivalent to any of the usual representations.
One possible solution would be to use the standard forms of projective representations which are quite well studied. However, to be able to use them, we would need to show that the multiplicative constants $`c_{\pi ,\sigma ,i}`$ are the same for all $`V_i`$ (or show that we can split all $`V_i`$ in several groups so that $`c_{\pi ,\sigma ,i}`$ is the same within one group) and we do not know if this is possible.
Our solution is to replace $`f_\pi ^{}`$ by transformations $`f_\pi ^{\prime \prime }`$ on a larger space $`V_1V_1^{}+\mathrm{}V_kV_k^{}`$ so that $`f_\pi ^{\prime \prime }f_\sigma ^{\prime \prime }=f_{\pi \sigma }^{\prime \prime }`$. Then, $`f_\pi ^{\prime \prime }`$ form a representation in the usual sense and we can analyse them similarly to section 4.
### 5.3 Solving the problem with phases
We split $`V_1,\mathrm{},V_k`$ into equivalence classes $`𝒱_1,\mathrm{}𝒱_l`$. $`V_i`$ and $`V_j`$ are in one class if there is a $`\pi S_M`$ such that $`f_\pi (V_i)=V_j`$. Let $`W_i`$ be the union of all $`V_j`$ that belong to $`𝒱_i`$. Then, $`f_\pi (W_i)=W_i`$ for any $`\pi S_M`$ (because $`f_\pi `$ maps every $`V_j𝒱_i`$ to some $`V_j^{}𝒱_i`$). Therefore, we can look at each $`W_i`$ separately.
###### Lemma 7
Let $`X,Y\widehat{𝒳}`$. Then, for any $`t`$,
$$dimXW_tdimXYW_t\sqrt{\frac{4cn}{M}}dimW_t.$$
Proof: To simplify the notation, assume that $`W_t=V_1V_2\mathrm{}V_l`$.
Consider the linear space $`W_t^{}=V_1V_1^{}\mathrm{}V_lV_l^{}`$ and the linear transformations $`f_\pi ^{\prime \prime }=f_\pi ^{}f_\pi `$. These linear transformations form a representation because
$$f_{\pi \sigma }^{}(f_{\pi \sigma }^{})^{}=c_{\pi ,\sigma ,i}f_\pi ^{}f_\sigma ^{}c_{\pi ,\sigma ,i}^{}(f_\pi ^{})^{}(f_\sigma ^{})^{}=$$
$$f_\pi ^{}f_\sigma ^{}(f_\pi ^{})^{}(f_\sigma ^{})^{}=f_\pi ^{\prime \prime }f_\sigma ^{\prime \prime }$$
on every $`V_iV_i^{}`$.
Let
$$X^{}=_{i=1}^l(XV_i)(XV_i)^{}$$
be the subspace of $`W_t^{}`$ corresponding to $`X`$. Then, $`f_\pi ^{}(X)=Y`$ implies $`f_\pi ^{\prime \prime }(X^{})=Y^{}`$. (To see this, consider one of $`(XV_i)(XV_i)^{}`$. Assume that $`f_\pi ^{}`$ maps $`V_i`$ to $`V_i^{}`$. Then, $`f_\pi ^{}(X)=Y`$ implies $`f_\pi ^{}(XV_i)=YV_i^{}`$ and
$$f_\pi ^{\prime \prime }((XV_i)(XV_i)^{})=(YV_i^{})(YV_i^{})^{}.$$
Combining these equalities for all $`V_i`$ gives $`f_\pi ^{\prime \prime }(X^{})=Y^{}`$.)
In particular, $`f_\pi ^{\prime \prime }(X^{})=Y^{}`$ means that $`X^{}`$ is invariant under all $`\pi S_M`$ satisfying $`\pi (X)=X`$. Therefore, by Lemma 6,
$$dimX^{}dimX^{}Y^{}\frac{4cn}{M}dimW_t^{}.$$
(1)
We use this inequality to derive a bound on $`dimXW_tdimXYW_t`$. To do this, we relate the dimensions of $`XV_i`$ and $`X^{}(V_iV_i^{})`$. First, notice that we have
$$XW_t=_{i=1}^l(XV_i)$$
(2)
because $`X`$ is invariant under $`G`$ and, therefore, can be written as a sum of irreducible invariant subspaces (and each of these irreducibles is contained in some $`V_i`$). The same is true about $`Y`$ and $`XY`$:
$$XYW_t=_{i=1}^l(XYV_i)$$
(3)
Let $`d_i`$ and $`d_i^{}`$ be the dimensions of $`XV_i`$ and $`XYV_i`$. Then, (2) and (3) imply that $`dimXW_t=_{i=1}^ld_i`$, $`dimXYW_t=_{i=1}^ld_i^{}`$ and
$$dimXW_tdimXYW_t=\underset{i=1}{\overset{l}{}}(d_id_i^{}).$$
If we look at $`V_iV_i^{}`$, then
$$X^{}(V_iV_i^{})=(XV_i)(XV_i)^{}.$$
This implies $`dimX^{}(V_iV_i^{})=d_i^2`$ and $`dimX^{}=_id_i^2`$. Similarly, $`dimX^{}Y^{}=_{i=1}^ld_{i}^{}{}_{}{}^{2}`$.
Let $`d`$ be the dimension of $`V_1`$. Then, the dimensions of $`V_2`$, $`\mathrm{}`$, $`V_l`$ are $`d`$ as well because, for every $`i\{2,\mathrm{},l\}`$, there is a unitary $`f_\pi `$ such that $`f_\pi (V_1)=V_i`$. Therefore, $`dimW_1=ld`$. Also, $`dimW_t^{}=ld^2`$ because $`dimV_iV_i^{}=d^2`$ for every $`i\{1,\mathrm{},l\}`$. Hence, we have
$$\frac{dimXW_tdimXYW_t}{dimW_t}=\frac{_{i=1}^m(d_id_i)}{md}=$$
$$\frac{1}{m}\underset{i=1}{\overset{m}{}}\sqrt{\frac{(d_id_i^{})^2}{d^2}}\frac{1}{m}\underset{i=1}{\overset{m}{}}\sqrt{\frac{(d_id_i^{})(d_i+d_i^{})}{d^2}}=$$
$$\frac{1}{m}\underset{i=1}{\overset{m}{}}\sqrt{\frac{(d_i^2d_i^{}_{}{}^{}2)}{d^2}}.$$
Convexity of the square root implies that this is at most
$$\sqrt{\frac{_{i=1}^m(d_i^2d_i^{}_{}{}^{}2)}{md^2}}=\sqrt{\frac{dimX^{}dimX^{}Y^{}}{dimW_t^{}}}.$$
Equation 1 implies that this is at most $`\sqrt{(4cn)/M}`$. This completes the proof of lemma. $`\mathrm{}`$
With Lemma 7, we can finish the proof similarly to the simple case (section 4). By summing over $`W_t`$’s, we get
$$dimXdimXY\sqrt{\frac{2cn}{M}}\underset{t}{}dimW_t=\sqrt{\frac{2cn}{M}}2^n.$$
Therefore,
$$\frac{dimXY}{dimX}\frac{2^{nk}\sqrt{\frac{2cn}{M}}2^n}{2^{nk}}=1\frac{2^{k+1}\sqrt{cn}}{\sqrt{M}}.$$
If $`\frac{2^{k+1}\sqrt{cn}}{\sqrt{M}}\frac{1}{\text{poly}(m)}`$, then $`m=(2k+\mathrm{log}n)(1+o(1))`$. This completes the proof of Theorem 2. $`\mathrm{}`$
## 6 Representation theory
A representation $`\rho `$ of a group $`G`$ is a homomorphism $`\rho `$ from $`G`$ to the group of linear transformations $`GL(V)`$ of a vector space $`V`$. This means that, for any $`g,hG`$, $`\rho (gh)=\rho (g)\rho (h)`$. If the mapping $`\rho `$ is clear from the context, we often call the space $`V`$ itself representation of $`G`$.
We say that a subspace $`W`$ is an *invariant* subspace of a representation $`\rho `$ if $`\rho (g)WW`$ for all $`gG`$. In order for $`W`$ to be an invariant subspace for $`\rho `$, it must be simultanously fixed under all $`\rho (g)`$. The zero subspace and the subspace $`V`$ are always invariant. If no nonzero proper subspaces are invariant, the representation is said to be *irreducible*.
Two representations $`\rho :GGL(V)`$ and $`\rho ^{}:GGL(W)`$ are isomorphic if there is a bijective linear map $`\phi :VW`$ such that $`\phi \rho (g)=\rho ^{}(g)\phi `$ for any $`gG`$.
If $`\rho `$ and $`\rho ^{}`$ are two irreducible representations and $`\phi `$ is an isomorphism between them, then any other isomorphism $`\phi ^{}`$ between $`\rho `$ and $`\rho ^{}`$ is $`c\phi `$ for some constant $`c\text{ }\mathrm{C}`$.
Schur’s lemma is usually stated for finite groups. However, if the representation is unitary (as in this paper), it is also true for infinite groups.
When a representation *does* have a nonzero proper invariant subspace $`V_1V`$, it is always possible to find a complementary subspace $`V_2`$ (so that $`V=V_1V_2`$) which is also invariant. Since $`\rho (g)`$ fixes $`V_1`$, we may let $`\rho _1(g)`$ be the linear map on $`V_1`$ given by $`\rho (g)`$. It is not hard to see that $`\rho _1:GGL(V_1)`$ is in fact a representation. Similarly define $`\rho _2(g)`$ to be $`\rho (g)`$ restricted to $`V_2`$. Since $`V=V_1V_2`$, the linear map $`\rho (g)`$ is completely determined by $`\rho _1(g)`$ and $`\rho _2(g)`$, and in this case we write $`\rho =\rho _1\rho _2`$.
Repeating the process described above, a representation $`\rho `$ may be written $`\rho =\rho _1\rho _2\mathrm{}\rho _k`$, where each $`\rho _i`$ is irreducible.
In this paper, we use representations of the symmetric group $`S_M`$. The irreducible representations of $`S_M`$ may be placed into one-to-one correspondence with the partitions of $`n`$. A *partition* of $`M`$ is a sequence $`(\lambda _1,\mathrm{},\lambda _k)`$ of positive integers, with $`\lambda _1\mathrm{}\lambda _k`$ for which $`\lambda _i=M`$. It is customary to identify the partition $`\lambda =(\lambda _1,\mathrm{},\lambda _k)`$ with a diagram consisting of $`k`$ rows of boxes, the $`i`$th row containing $`\lambda _i`$ boxes. We will let $`\lambda `$ stand for both the partition and the associated diagram. For example, the diagram corresponding to the partition $`\lambda =(4,4,2,1)`$ is shown in figure 1.
The irreducible representation associated with $`\lambda `$ is denoted $`\rho _\lambda `$. There is an explicit formula for the dimension of $`\rho _\lambda `$. This involves the notion of a *hook*: for a cell $`(i,j)`$ of a Young tableau $`\lambda `$, the $`(i,j)`$-hook $`h_{i,j}`$ is the collection of all cells of $`\lambda `$ which are beneath $`(i,j)`$ (but in the same column) or to the right of $`(i,j)`$ (but in the same row), including the cell $`(i,j)`$. The *length* of the hook $`\mathrm{}(h)`$ is the number of cells appearing in the hook. With this notation, the dimension of $`\rho _\lambda `$ may be expressed:
$$dim\rho _\lambda =\frac{n!}{_{i,j}\mathrm{}(h_{i,j})},$$
(4)
this product being taken over all hooks $`h`$ of $`\lambda `$. Figure 2 shows the hook lengths for the partition $`\lambda =(4,4,3,1)`$. Formula (4) implies that the dimension of corresponding representation is
$$\frac{11!}{75326423}=1320.$$
A representation $`\rho `$ of a group $`G`$ is also automatically a representation of any subgroup $`H`$. Note that even if a representation is irreducible over $`G`$, it may no longer be irreducible when restricted to $`H`$.
In particular, we will be considering the restrictions of irreducible representations of $`S_M`$ to $`S_{M1}`$. Let $`\lambda `$ be a partition of $`M`$ and $`\rho _\lambda `$ be the corresponding irreducible representation. Then, when we restrict to $`S_{M1}`$, $`\rho _\lambda `$ decomposes into irreducible representations of $`S_{M1}`$ in the following way:
$$\rho =\underset{\lambda ^{}}{}\rho _\lambda ^{}$$
where $`\lambda ^{}`$ ranges over all shapes of size $`M1`$ that can be obtained by deleting an “inside corner” from $`\lambda `$. (An inside corner is simply a point of the shape whose deletion leaves a legal shape.)
For example, the representation $`\rho _\lambda `$, $`\lambda =(4,4,2,1)`$ of $`S_{11}`$ decomposes into 3 irreducible representations of $`S_{10}`$. The Young diagrams of these representations are shown in Fig. 3.
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# Landau-Pomeranchuk-Migdal effect in thermal field theory
## 1 Introduction
The photon production rate is thought to be a quantity of phenomenological interest in heavy ions collisions, possibly enabling one to detect the formation of a quark-gluon plasma. Part of the interest in this electromagnetic observable comes from the fact that photons are relatively weakly coupled to nuclear matter ($`\alpha _{_{EM}}\alpha __S`$). Given the typical size of the system in such collisions (much smaller than the mean free path of a photon), they do not re-interact between their production and their observation. As a consequence, photons (real photons, or virtual photons decaying eventually into a lepton pair) can provide information on the state of the system at the time they were produced.
In order to calculate the photon yield from a hot quark-gluon plasma, thermal field theory is the tool of choice since its Feynman rules automatically take into account the presence of a thermal bath with the appropriate distributions of partons. Thermal gauge theories are however plagued by infrared singularities arising from the Bose-Einstein distribution functions which are singular at zero energy. An improvement over the bare Feynman rules is achieved by the resummation of one-loop thermal contributions known as hard thermal loops (HTL) . These thermal corrections make it possible to include in the propagators effects like Debye screening or Landau damping, and transform partons into massive quasi-particles. From a quantitative perspective, they provide important changes to the dynamics of soft modes (of momentum of order $`gT`$ or less). This resummation lacks however two features: it does not provide any Debye screening for static magnetic fields (such a screening is expected to arise non-perturbatively in QCD at the length scale $`(g^2T)^1`$), and its quasi-particles do not undergo collisions (their collisional<sup>1</sup><sup>1</sup>1In a plasma, it is very important to distinguish two contributions to the total width: the decay width made of the zero temperature contribution to the imaginary part of the quark self-energy, and the collisional width which exists only in a medium. The latter is also called “anomalous damping rate” (or just “damping rate” for short) and is of order $`g^2T\mathrm{ln}(1/g)`$, while in QED/QCD the former starts at the order $`g^4`$. The width $`\mathrm{\Gamma }`$ that we introduce in this paper is the collisional width, and its inverse is the mean free path of the quark in the plasma. width is also a sub-leading effect of order $`g^2T\mathrm{ln}(1/g)`$).
In the framework of thermal field theory, the photon/dilepton rate is obtained via the calculation of the imaginary part of the retarded photon polarization tensor . This object has been evaluated at one-loop in the HTL-improved perturbative expansion, both for virtual and real photons . More recently, new processes like bremsstrahlung were studied in detail in this framework and have been found to be dominant sources of low invariant mass photons . Despite the fact that this process arrives only in 2-loop contributions to the photon polarization tensor, it is always important because of a strong collinear enhancement. Indeed, it was found in that for a very small photon invariant mass, the corresponding diagram contains collinear singularities that, when regularized by a thermal quark mass of order $`gT`$, give an enhancement by a factor of order $`1/g^2`$ over naive estimates coming from power counting.
After finding that certain 2-loop terms are contributing at leading order, one may wonder whether this result is specific to this 2-loop contribution only or if, on the contrary, this is an indication of the breakdown of perturbative expansion (even improved with HTLs). In a recent paper , we studied what is the effect of loop corrections to this 2-loop diagram. Power counting indeed indicates a problem very similar to the problem raised by Linde for the calculation of the free energy, due to the lack of Debye screening for static magnetic modes. This problem is avoided for the production of massive enough photons, because some cancellations (occuring within any given topology, when one is summing over all the cuts contributing to the imaginary part) generate a kinematical cutoff large enough to prevent any sensitivity to the non-perturbative scale $`g^2T`$. Unfortunately, this cutoff is smaller than $`g^2T`$ whenever the invariant mass of the produced photon is too small (typically $`Q^2<g^2Tq_0`$ for $`q_0<T`$). In this low invariant mass region, the photon rate is therefore non perturbative: exchanged transverse gluons reach the scale $`g^2T`$ of the non-perturbative “magnetic mass”, and an infinite set of diagrams must be resummed.
In this paper, we present a completely different approach to higher order corrections, that completes the picture outlined in . The idea behind the present study is that a width on the quark propagator will act as a regulator in the collinear sector, because it moves the poles of the propagator away from the real energy axis. Such a collisional width is necessarily a higher-loop effect, because the hard thermal loop framework does not take into account the collisions of quasi-particles. Having in mind the fact that 2-loop contributions are important because of collinear enhancement, an important question to answer is how much of this enhancement is lost when an additional regulator like a width is taken into account. This is the question we want to address in this paper, by calculating the same 2-loop diagrams as in , in the presence of a quark width. To be more definite, and keep the model as well as the calculations simple, we use a momentum-independent width.
A word of caution is necessary here: the formulae found using this simple model should not be taken as an accurate quantitative account of what the effect of such a width will be on thermal photon production rates. Indeed, the constant width model is probably too naive to be realistic, and more importantly our calculation disregards the fact that a modification of the vertices should in principle accompany the modification of the quark propagators. Nevertheless, this simple approach is sufficient here for our purpose which is just to determine the region in which effects of a width of order $`g^2T\mathrm{ln}(1/g)`$ are to be expected, since that gives another handle on how and when higher order corrections may affect photon production by a quark-gluon plasma.
We find that the region where a width of order $`g^2T\mathrm{ln}(1/g)`$ is important is very similar to the one found in for the contribution of higher-loop topologies due to an IR sensitivity to the scale of the magnetic mass. Despite their similarity, these two non perturbative regions have different physical interpretations. In particular, we find that the sensitivity to the collisional width of the quarks is a manifestation of the Landau-Pomeranchuk-Migdal effect. Indeed, it occurs when the formation time of the photon is larger than the mean free path of the quark producing the photon. An interesting consequence of our study is that the LPM effect also modifies the spectrum of highly energetic photons.
The structure of the paper is as follows. Section 2 makes more precise our modeling of the mean free path for the quarks. In section 3, we start by computing the 1-loop contribution in presence of a width. Although not related to collinear singularities in any way, the purpose of this warmup exercise is two-fold: it illustrates the technology (how one does calculations with a width), and it shows in a simple way how collisions can open up the phase space. We also show that this 1-loop contribution is canceled by the resummation of vertex corrections, and is therefore not physical.
In section 4, we present the 2-loop calculation with a quark width, and obtain a rather simple generalization of the formulae of . These formulae are discussed extensively in section 5, in which we also determine the region where a width of order $`g^2T\mathrm{ln}(1/g)`$ plays an important role.
Section 6 is devoted to establishing the connection between the previous results and the LPM effect. In section 7, we study the process obtained from bremsstrahlung by crossing symmetry, which turns out to be important in the region of large photon energy. We show that this process is also affected by the LPM effect. The last section contains concluding remarks. In particular, we combine the present work on the LPM effect with previous results on infrared singularities in order to make a syntesis and present a reasonable physical picture of thermal photon production.
## 2 Model
Let us first make our framework more definite. The only modification compared to the context extensively described in is that the quark propagators are given a width, as outlined by the following substitution in the retarded and advanced propagators:
$$\mathrm{\Delta }_{_{R,A}}(P)\frac{1}{P^2M_{\mathrm{}}^2\pm ip_0ϵ}\Delta _{_{R,A}}(P)\frac{1}{(p_0\pm i\mathrm{\Gamma })^2p^2M_{\mathrm{}}^2},$$
(1)
where $`M_{\mathrm{}}gT`$ is the usual asymptotic thermal mass for hard quarks , and where $`\mathrm{\Gamma }`$ is a constant width. The expression in Eq. (1) is sufficient for the physics we consider, which is dominated by hard quark momenta close to the mass shell. Whenever we need to estimate the order of magnitude of a term containing $`\mathrm{\Gamma }`$ in the following, we assume that it is of order $`g^2T\mathrm{ln}(1/g)`$ (with $`g1`$). Note that we could as well have written
$$\Delta _{_{R,A}}(P)=\frac{1}{P^2M_{\mathrm{}}^2\pm 2ip_0\mathrm{\Gamma }}$$
(2)
since the two differ by the small correction $`\mathrm{\Gamma }^2M_{\mathrm{}}^2`$ to the real part of the denominator. The retarded (resp. advanced) propagator has two complex poles in the $`p_0`$ plane, located at $`p_0=\pm \omega _pi\mathrm{\Gamma }`$ (resp. $`p_0=\pm \omega _p+i\mathrm{\Gamma }`$), with $`\omega _p\sqrt{(p^2+M_{\mathrm{}}^2)}`$.
## 3 1-loop study
### 3.1 Calculation
Let us now consider the simplest calculation conceivable in this framework, namely the computation of the 1-loop contribution to the photon polarization tensor. Our purpose is simply to illustrate possible changes brought in when taking into account a width for the quarks running in the loop. We do not attempt a complete calculation including HTL vertices, as pictured on Fig. 1.
The contribution of this diagram to the imaginary part of the retarded polarization tensor is trivial to obtain:
$`\text{Im}\mathrm{\Pi }^\mu {}_{\mu }{}^{}(Q){\displaystyle \frac{1}{2}}e^2{\displaystyle }{\displaystyle \frac{d^4P}{(2\pi )^4}}[n__F(r_0)n__F(p_0)]\text{Tr}`$
$`\times [\Delta __R(P)\Delta __A(P)][\Delta __R(R)\Delta __A(R)],`$ (3)
where $`\mathrm{Tr}`$ denotes the Dirac’s trace associated to the quark loop:
$$\mathrm{Tr}=8PR.$$
(4)
The case with $`\mathrm{\Gamma }=0`$ is made simple by the fact that the differences $`\Delta __R\Delta __A`$ are $`\delta `$ functions, thereby enabling some of the integrals to be performed trivially. Let us remind that for the case $`Q^2<4M_{\mathrm{}}^2`$ which we assume throughout this paper, $`\mathrm{Im}\mathrm{\Pi }{}_{}{}^{\mu }{}_{\mu }{}^{}(Q)=0`$ if $`\mathrm{\Gamma }=0`$ due to incompatibilities between the $`\delta `$ functions. To evaluate Eq. (3) when $`\mathrm{\Gamma }0`$, we perform the integral over $`p_0`$ by closing the real axis in the complex energy plane and using the theorem of residues.
Note first that one can disregard the poles of the statistical weights $`n__F(p_0)n__F(r_0)`$. Indeed, these poles are the imaginary fermionic Matsubara frequencies and are to be plugged into differences like $`\Delta __R(P)\Delta __A(P)`$. The fact that $`\mathrm{\Gamma }T`$ makes these differences very small<sup>2</sup><sup>2</sup>2Strictly speaking, these contributions are needed to ensure that the final result is a real number. Indeed, when one picks poles like $`p_0=\omega _p+i\mathrm{\Gamma }`$ from the propagators and plugs them in the distribution functions, the latter become complex numbers. Their (small) imaginary part is canceled by the (small) contribution coming from the poles of $`n__F(r_0)n__F(p_0)`$. For this approximation to be consistent, one must also neglect $`i\mathrm{\Gamma }`$ in the hard argument of statistical weights.. In other words, the only important terms are those for which denominators like $`P^2M_{\mathrm{}}^2`$ are small, and comparable to $`p_0\mathrm{\Gamma }`$. This cannot happen if $`p_0`$ is an imaginary number of order $`T`$.
We are therefore left with the poles of the propagators themselves. At this point, we obtain the following result:
$`\text{Im}\mathrm{\Pi }^\mu {}_{\mu }{}^{}(Q){\displaystyle \frac{1}{2}}e^2{\displaystyle }{\displaystyle \frac{d^3p}{(2\pi )^3}}{\displaystyle \underset{\eta =\pm 1}{}}[n__F(q_0+\eta \omega _p)n__F(\eta \omega _p)]\text{Tr}(p_0=\eta \omega _p)`$
$`\times {\displaystyle \frac{\eta }{2\omega _p}}\left[{\displaystyle \frac{1}{(q_0+\eta \omega _p+2i\mathrm{\Gamma })^2\omega _r^2}}{\displaystyle \frac{1}{(q_0+\eta \omega _p2i\mathrm{\Gamma })^2\omega _r^2}}\right].`$ (5)
We can also note for later use that the result of this approximation could have been obtained by starting from an expression slightly different from Eq. (3):
$`\text{Im}\mathrm{\Pi }^\mu {}_{\mu }{}^{}(Q){\displaystyle \frac{1}{2}}e^2{\displaystyle }{\displaystyle \frac{d^4P}{(2\pi )^4}}[n__F(r_0)n__F(p_0)]\text{Tr}`$
$`\times 2\pi ϵ(p_0)\delta (P^2M_{\mathrm{}}^2)[\Delta __R^{2\mathrm{\Gamma }}(R)\Delta __A^{2\mathrm{\Gamma }}(R)],`$ (6)
where we denote
$$\Delta _{_{R,A}}^{2\mathrm{\Gamma }}(R)\frac{1}{(r_0\pm 2i\mathrm{\Gamma })^2r^2M_{\mathrm{}}^2}.$$
(7)
In other words, if $`\mathrm{\Gamma }T`$, then we can as well put twice the width on one of the two quark lines, and nothing in the other quark line<sup>3</sup><sup>3</sup>3This is where it is important to have a constant width. Indeed, some intermediate step relies on the cancellation $`\mathrm{\Gamma }(P)\mathrm{\Gamma }(R)=0`$..
After Eq. (5), the angular integration over $`d\mathrm{\Omega }_p`$ is trivial. In the case where $`\eta =+1`$ (process $`qq\gamma `$), we find:
$`\mathrm{Im}\mathrm{\Pi }^\mu {}_{\mu }{}^{}(Q){\displaystyle \frac{e^2\mathrm{\Gamma }}{2\pi ^2q_0}}{\displaystyle \underset{0}{\overset{+\mathrm{}}{}}}dp(q_0+p)[n__F(q_0+p)n__F(p)]`$
$`\times \mathrm{ln}\left({\displaystyle \frac{(\omega _p+p)^2q^2+\mathrm{\Gamma }^2(q_0+p)^2}{(\omega _pp)^2q^2+\mathrm{\Gamma }^2(q_0+p)^2}}\right).`$ (8)
In particular, assuming for the sake of simplicity that $`qq_0`$, it is trivial to obtain the following asymptotic behaviors for soft and hard photons:
$`\mathrm{If}q_0T,\mathrm{Im}\mathrm{\Pi }^\mu {}_{\mu }{}^{}(Q)e^2\mathrm{\Gamma }T\mathrm{ln}(1+4q_0^2/\mathrm{\Gamma }^2),`$ (9)
$`\mathrm{If}q_0T,\mathrm{Im}\mathrm{\Pi }^\mu {}_{\mu }{}^{}(Q)e^2\mathrm{\Gamma }T\mathrm{ln}(T^2/\mathrm{\Gamma }^2).`$ (10)
The main point is that these contributions are proportional to the width and vanish when $`\mathrm{\Gamma }=0`$, in agreement with a direct calculation. As a side remark, one may note that for $`\mathrm{\Gamma }g^2T\mathrm{ln}(1/g)`$ and soft photons ($`q_0gT`$), Eq. (9) is larger by a factor $`1/g`$ than the 1-loop HTL result, while in the regime of Eq. (10) it is of the same order. It is the collision partners of the quarks that open up the phase-space (hard quarks colliding in the medium can emit a photon, a process forbidden for collisionless quarks) and make these contributions so large.
Another property of this result is that there is a suppression if $`q_0T`$, such that $`\mathrm{Im}\mathrm{\Pi }{}_{}{}^{\mu }{}_{\mu }{}^{}(Q)`$ tends to 0 when $`q_00`$. Equations similar to Eqs. (9) and (10) can be obtained when considering the contribution $`\eta =1`$ to Eq. (6).
### 3.2 Cancellation with vertex corrections
This suppression was interpreted in as the manifestation of the LPM effect. However, despite the suppression at small $`q_0`$, the connection with the LPM effect is not clear in this context. Indeed, the LPM effect is expected when the photon formation time is larger than the quark mean free path, a condition which never appears in , nor in the above 1-loop calculation. In fact, following , we know that the propagator of a quark close to its mass shell on which a width is resummed can be evaluated by an eikonal resummation of the multiple collisions. By the same method, one can include in this resummation all the vertex corrections. Indeed, one can check that the photon polarization tensor at this level of approximation is proportional to (in space-time coordinates)
$$\mathrm{\Pi }{}_{}{}^{\mu }{}_{}{}^{\nu }(x,y)=e^2[dA^\mu ]e^{iS[A]}\mathrm{Tr}[\gamma ^\mu S(x,y|A)\gamma ^\nu S(y,x|A)],$$
(11)
where $`S[A]`$ is the action of the gauge fields<sup>4</sup><sup>4</sup>4This action does not play any role in the argument, and therefore can even include the HTL effective action for gluons., and where $`S(x,y|A)`$ is the propagator of a quark in the background field $`A`$. This formula includes only one quark loop (in addition to the quark loops that may have been resummed in the gluon propagators and vertices in $`S[A]`$), and all orders in the gluon fields. In the eikonal approximation, the quark propagator has a very simple dependence on the field $`A_\mu `$ (which is a matrix $`T^aA_\mu ^a`$ in QCD):
$$S(x,y|A)=S_0(x,y)𝒫\mathrm{exp}ig\underset{x^0}{\overset{y^0}{}}v_\mu A^\mu 𝑑t,$$
(12)
where $`v_\mu `$ is the 4-velocity of the quark, and $`S_0(x,y)`$ is the free propagator of the quark. Now, if one is looking at photons of very small invariant mass, the two quarks are nearly parallel in the collinear limit (even if the photon is hard), so that the two quarks have mostly the same $`v`$. As a consequence, the two path ordered exponentials cancel each other, and the product of the propagators under the functional integral is
$$\gamma ^\mu S(x,y|A)\gamma ^\nu S(y,x|A)\gamma ^\mu S_0(x,y)\gamma ^\nu S_0(y,x).$$
(13)
Therefore, at the level of approximation at which the resummed quark propagator is calculated, the sum of all the gluon loop corrections is vanishing:
$$\underset{\mathrm{all}\mathrm{gluon}\mathrm{corrections}}{}\text{}.$$
(14)
This result generalizes a result already known for QED established by for the subset of ladder vertex corrections, and extended to all abelian topologies in , to the case where gluons are exchanged (i.e. to non abelian topologies). Indeed, we see that what makes gluons specific, namely the fact that they can couple to each other, is hidden in the action $`S[A]`$, which plays a passive role in the argument.
This result is also closely related to the fact that there are no $`\gamma \gamma g\mathrm{}g`$ HTL vertices (see for an interpretation of the cancellation found in as a consequence of the absence of HTL vertices with $`n>2`$ photons), because the sum of all the eikonal contributions can also be written as (when $`Q`$ is soft):
$$\mathrm{\Pi }{}_{}{}^{\mu }{}_{}{}^{\nu }(Q)=\text{}\text{ },$$
(15)
where the sum is extended to all the possible ways to close the external gluonic legs of the $`\gamma \gamma g\mathrm{}g`$ vertices. If is then obvious to see that the sum reduces to its first term if these vertices are vanishing.
One is therefore left with the bare 1-loop diagram, which does not contributes to the imaginary part of the photon polarization tensor. In other words, the results derived in Eqs. (9) and (10) as a warmup and that were claimed to be related to the LPM effect in , are just artifacts with no physical meaning. In section 6 of the present paper, we show where the LPM effect appears in thermal field theory.
Despite the fact that this contribution is not physical, one learns two things from this calculation: (1) one has to keep terms beyond the eikonal approximation in order to circumvent this cancellation, and (2) the width may open the phase space to new processes.
## 4 2-loop calculation
The lesson from the previous section is that physical contributions must be looked for beyond the eikonal approximation. This implies inserting explicitly a gluon exchange with momentum $`L`$ in the diagram, and not assuming that $`L^22PL`$.
There are in principle two 2-loop topologies contributing to the photon polarization tensor. However, it was found in that only the topology correcting the $`q\overline{q}\gamma `$ vertex (see Fig. 2) is relevant in the region where the collinear enhancement takes place. Since our purpose it to study how this collinear enhancement is affected by the width $`\mathrm{\Gamma }`$, we limit the present study to the terms that were found important in .
There is though one difference with the case $`\mathrm{\Gamma }=0`$ that we must take into account: with a zero width, processes corresponding to cuts $`(a)`$ and $`(d)`$ are kinematically forbidden for small $`Q^2`$. As in the previous section, turning on the width opens those production channels, and we cannot disregard them a priori. This is not a problem though from a purely technical point of view, because it turns out that the sum of all the cuts has a simpler expression than each individual piece<sup>5</sup><sup>5</sup>5In the appendix A, we show how one can obtain the difference between the cuts $`(c)`$ and $`(d)`$. Having already their sum, we can therefore reconstruct the two contributions separately..
If we select only terms that have the large Bose-Einstein factor $`n__B(l_o)`$, the contribution of Fig. 2 to the imaginary part of the photon polarization tensor takes the following simple form
$`\text{Im}\mathrm{\Pi }^\mu {}_{\mu }{}^{}(Q){\displaystyle \frac{1}{2}}e^2g^2{\displaystyle }{\displaystyle \frac{d^4P}{(2\pi )^4}}[n__F(r_0)n__F(p_0)]`$
$`\times {\displaystyle }{\displaystyle \frac{d^4L}{(2\pi )^4}}n__B(l_0)\rho _{_{T,L}}(L)P_{_{T,L}}^{\rho \sigma }(L)\text{Tr}_{\rho \sigma }`$
$`\times [\Delta __R(P)\Delta __R(P+L)\Delta __A(P)\Delta __A(P+L)]`$
$`\times [\Delta __R(R)\Delta __R(R+L)\Delta __A(R)\Delta __A(R+L)],`$ (16)
where $`\rho _{_{T,L}}(L)`$ are the spectral functions of transverse and longitudinal gluons, $`P_{_{T,L}}^{\rho \sigma }(L)`$ are the corresponding projectors, and where $`\mathrm{Tr}_{\rho \sigma }`$ is the result of the Dirac’s trace for the quark loop. The result of the previous section requires that we keep in this Dirac’s trace only terms that do not appear in the eikonal approximation (i.e. which includes soft corrections to the hard loop momentum). In the collinear limit, the first non vanishing term (beyond eikonal approximation) is
$$\mathrm{Tr}_{\rho \sigma }8L^2(R_\rho R_\sigma +P_\rho P_\sigma ),$$
(17)
which turns out to be the same as the term found in . Its contractions with the projectors are given by
$`P__T^{\rho \sigma }(L)\mathrm{Tr}_{\rho \sigma }8L^2(r^2+p^2)(1\mathrm{cos}^2\theta ^{})`$
$`P__L^{\rho \sigma }(L)\mathrm{Tr}_{\rho \sigma }+8L^2(r^2+p^2)(1\mathrm{cos}^2\theta ^{}),`$ (18)
where $`\theta ^{}`$ is the angle between the 3-vectors $`p`$ and $`l`$. At this point, we have used the fact that we are looking at collinearly enhanced terms, and consider that $`p`$ and $`r`$ are parallel (the only place where we do not do this approximation is in the denominators which are very sensitive to the details of the collinear sector).
In order now to perform the integral over $`p_0`$, we follow the method of the previous section, and use the same approximations concerning the statistical weights. In addition, we compute only the contribution of cuts $`(c)+(d)`$, and multiply the result by an overall factor $`2`$ in order to take into account the contribution of the cuts $`(a)+(b)`$. Following the remark of the previous section, we can start directly from the expression:
$`\text{Im}\mathrm{\Pi }^\mu {}_{\mu }{}^{}(Q){\displaystyle \frac{1}{2}}e^2g^2{\displaystyle }{\displaystyle \frac{d^4P}{(2\pi )^4}}[n__F(r_0)n__F(p_0)]`$
$`\times {\displaystyle }{\displaystyle \frac{d^4L}{(2\pi )^4}}n__B(l_0)\rho _{_{T,L}}(L)P_{_{T,L}}^{\rho \sigma }(L)\text{Tr}_{\rho \sigma }`$
$`\times 2\pi ϵ(p_0)\delta (P^2M_{\mathrm{}}^2){\displaystyle \frac{1}{(P+L)^2M_{\mathrm{}}^2}}`$
$`\times \mathrm{Disc}\left[\Delta __R^{2\mathrm{\Gamma }}(R)\Delta __R^{2\mathrm{\Gamma }}(R+L)\right],`$ (19)
where we use the notation $`\mathrm{Disc}f(\mathrm{\Gamma })f(\mathrm{\Gamma })f(\mathrm{\Gamma })`$. We first do the $`p_0`$ integration for free thanks to the $`\delta (P^2M_{\mathrm{}}^2)`$. In this section, we consider only the case of $`p_0=+\omega _p`$ (bremsstrahlung of a quark) in order to keep the calculation compact.
The contribution of $`p_0<0`$ is identical if $`q_0T`$, but is different if $`q_0`$ is large. In the latter case, the corresponding process is a $`q^{}\overline{q}`$ annihilation ($`q^{}`$ denotes a quark placed off-shell by a scattering) instead of bremsstrahlung (see Fig. 3), and is considered in more detail in section 7.
Then, it happens that the angular integral over the direction $`\mathrm{\Omega }_l`$ of the 3-vector $`l`$ can be done analytically in a rather simple way. We have to perform an integral like<sup>6</sup><sup>6</sup>6Strictly speaking, we have also a factor of $`1\mathrm{cos}^2\theta ^{}`$ in the numerator that depends on $`\mathrm{\Omega }_l`$. This factor can be taken into account analytically in the angular integral, the price to pay being more cumbersome expressions. One can however make the following simplification: if the width is small and if we use the collinear approximation for the numerators, then this angle is approximately given by $`\mathrm{cos}\theta ^{}l_0/l`$, and does not play any role in the angular integrals.
$$I__L\frac{d\mathrm{\Omega }_l}{4\pi }\frac{1}{2LP+L^2}\frac{1}{2LR+L^2+2PQ+Q^2+4ir_0\mathrm{\Gamma }},$$
(20)
which can be rewritten as
$$I__L=\frac{d\mathrm{\Omega }_l}{4\pi }\frac{1}{2\widehat{L}A}\frac{1}{2\widehat{L}B},$$
(21)
where $`\widehat{L}(1,\widehat{l})`$ provided we introduce the fictitious “4-vectors”
$`A(p_0l_0+{\displaystyle \frac{L^2}{2}},lp)`$
$`B(r_0l_0+{\displaystyle \frac{L^2}{2}}+PQ+{\displaystyle \frac{Q^2}{2}}+2ir_0\mathrm{\Gamma },lr).`$ (22)
The advantage of rewriting $`I__L`$ like this lies in the fact that the last integral is known in closed form<sup>7</sup><sup>7</sup>7$`I__L`$ is an analytic function of its arguments. The two possible choices for the square root of the complex number $`\mathrm{\Delta }`$ lead to the same result.:
$$I__L=\frac{1}{8\sqrt{\mathrm{\Delta }}}\left[\mathrm{ln}(AB+\sqrt{\mathrm{\Delta }})\mathrm{ln}(AB\sqrt{\mathrm{\Delta }})\right],$$
(23)
where $`\mathrm{\Delta }(AB)^2A^2B^2`$. We need now to evaluate the three quantities $`A^2`$, $`B^2`$ and $`AB`$, for which we will also use approximations based on the fact that $`L`$ is soft while the other momenta are hard. We obtain first a rather simple, and by now very familiar , expression for $`\mathrm{\Delta }`$:
$$\mathrm{\Delta }p^4q_0^2l^2\left[\left(1\mathrm{cos}\theta +\frac{M_{\mathrm{eff}}^2}{2p^2}+\frac{L^2}{2p^2}\right)^2\frac{L^2}{p^2}\frac{M_{\mathrm{eff}}^2}{p^2}\right],$$
(24)
where $`\theta `$ is the angle between the 3-vectors $`p`$ and $`q`$, and with
$$M_{\mathrm{eff}}^2M_{\mathrm{}}^2+\frac{Q^2}{q_0^2}p(p+q_0)+4i\frac{\mathrm{\Gamma }}{q_0}p(p+q_0).$$
(25)
Therefore, apart from the fact that the effective mass $`M_{\mathrm{eff}}^2`$ now gets an imaginary part coming from the width (and has been extended to hold for hard $`q_0`$ as well), this quantity is nothing but the denominator appearing in Eq. (39) of .
Since $`ABprL^2\sqrt{\mathrm{\Delta }}`$ and since the bremsstrahlung lives in the region where $`L^2<0`$, the difference of the two logarithms in Eq. (23) is just the discontinuity of the logarithm across its branch cut, which gives
$$\mathrm{ln}(AB+\sqrt{\mathrm{\Delta }})\mathrm{ln}(AB\sqrt{\mathrm{\Delta }})2i\pi ϵ(\mathrm{\Gamma }).$$
(26)
Therefore, in presence of a width $`\mathrm{\Gamma }`$, the contribution of bremsstrahlung to the photon polarization tensor is
$`\mathrm{Im}{}_{}{}^{\mu }{}_{\mu }{}^{}(Q){\displaystyle \frac{e^2g^2}{32\pi ^4}}{\displaystyle \frac{T}{q_0^2}}{\displaystyle \underset{0}{\overset{+\mathrm{}}{}}}dp{\displaystyle \frac{p^2+(p+q_0)^2}{p^2}}[n__F(p+q_0)n__F(p)]`$
$`\times {\displaystyle }{\displaystyle \frac{dl_0}{l_0}}{\displaystyle }l^3dl[\rho __T(l_0,l)\rho __L(l_0,l)][1\left({\displaystyle \frac{l_0}{l}}\right)^2]`$
$`\times \mathrm{Disc}{\displaystyle \underset{0}{\overset{2}{}}}{\displaystyle \frac{ϵ(\mathrm{\Gamma })du}{\left[u+\frac{M_{\mathrm{eff}}^2}{2p^2}\right]\left[\left(u+\frac{M_{\mathrm{eff}}^2}{2p^2}+\frac{L^2}{2p^2}\right)^2\frac{L^2}{p^2}\frac{M_{\mathrm{eff}}^2}{p^2}\right]^{1/2}}},`$ (27)
where $`u1\mathrm{cos}\theta `$. At this stage, this expression is formally similar to the one found for $`\mathrm{\Gamma }=0`$ and reproduces known results in this limit. Indeed, when $`\mathrm{\Gamma }=0`$, then $`M_{\mathrm{eff}}^2`$ is a real number, and taking the discontinuity just gives a factor of $`2`$ (do not forget the $`ϵ(\mathrm{\Gamma })`$ in the numerator). It is then possible to use the very same sequence of changes of variables as in to transform the integral over $`u`$, and write $`\mathrm{Im}\mathrm{\Pi }{}_{}{}^{\mu }{}_{\mu }{}^{}(Q)`$ as<sup>8</sup><sup>8</sup>8We used the fact that the spectral functions $`\rho _{_{T,L}}`$ have a simple expression in the space-like region:
$$\rho _{_{T,L}}(L)=\frac{2\mathrm{Im}\mathrm{\Pi }_{_{T,L}}^{^{HTL}}(L)}{(L^2\mathrm{Re}\mathrm{\Pi }_{_{T,L}}^{^{HTL}}(L))^2+(\mathrm{Im}\mathrm{\Pi }_{_{T,L}}^{^{HTL}})^2}.$$
(28)
$`\mathrm{Im}\mathrm{\Pi }{}_{}{}^{\mu }{}_{\mu }{}^{}(Q){\displaystyle \frac{2e^2g^2}{\pi ^4}}{\displaystyle \frac{T}{q_0^2}}{\displaystyle \underset{0}{\overset{+\mathrm{}}{}}}dp{\displaystyle \frac{p^2+(p+q_0)^2}{2}}[n__F(p+q_0)n__F(p)]`$
$`\times {\displaystyle \underset{m=T,L}{}}{\displaystyle \underset{0}{\overset{1}{}}}{\displaystyle \frac{dx}{x}}{\displaystyle \underset{0}{\overset{+\mathrm{}}{}}}dw{\displaystyle \frac{\left|\stackrel{~}{I}_m\right|K(w,\widehat{\mathrm{\Gamma }})}{(w+\stackrel{~}{R}_m)^2+(\stackrel{~}{I}_m)^2}},`$ (29)
where we denote:
$`w{\displaystyle \frac{L^2}{\mathrm{Re}M_{\mathrm{eff}}^2}},x{\displaystyle \frac{l_0}{l}},\widehat{\mathrm{\Gamma }}{\displaystyle \frac{\mathrm{Im}M_{\mathrm{eff}}^2}{\mathrm{Re}M_{\mathrm{eff}}^2}}`$
$`\stackrel{~}{I}_{_{T,L}}{\displaystyle \frac{\mathrm{Im}\mathrm{\Pi }_{_{T,L}}^{^{HTL}}}{\mathrm{Re}M_{\mathrm{eff}}^2}},\stackrel{~}{R}_{_{T,L}}{\displaystyle \frac{\mathrm{Re}\mathrm{\Pi }_{_{T,L}}^{^{HTL}}}{\mathrm{Re}M_{\mathrm{eff}}^2}},`$ (30)
and where the function $`K`$ comes from the integral over $`u`$ and is defined as<sup>9</sup><sup>9</sup>9One can go from the variable $`u`$ to $`y`$ by the following transformations
$$2p^2uM_{\mathrm{eff}}^2[w(1z)+z^11],\mathrm{and}y4(zz^2).$$
(31)
$`K(w,\widehat{\mathrm{\Gamma }}){\displaystyle \frac{1}{2}}{\displaystyle \underset{0}{\overset{1}{}}}{\displaystyle \frac{dy}{\sqrt{1y}}}{\displaystyle \frac{y+4/w}{(y+4/w)^2+(4\widehat{\mathrm{\Gamma }}/w)^2}}`$
$`={\displaystyle \frac{1}{4(\alpha ^2+\beta ^2)}}\{\alpha \mathrm{ln}\left({\displaystyle \frac{(1+\alpha )^2+\beta ^2}{(1\alpha )^2+\beta ^2}}\right)`$
$`2\beta [\mathrm{arctan}\left({\displaystyle \frac{\alpha +1}{\beta }}\right)\mathrm{arctan}\left({\displaystyle \frac{\alpha 1}{\beta }}\right)]\}`$ (32)
with $`\alpha +i\beta `$ a square root<sup>10</sup><sup>10</sup>10Explicitly, we have:
$$\alpha ,\beta =\sqrt{\frac{1}{2}\left[\sqrt{((4+w)/w)^2+(4\widehat{\mathrm{\Gamma }}/w)^2}\pm ((4+w)/w)\right]}.$$
(33) of the complex number $`(4+w)/w+4i\widehat{\mathrm{\Gamma }}/w`$. This function $`K(w,\widehat{\mathrm{\Gamma }})`$ is the generalization to the case of a non vanishing width of the factor $`\sqrt{w/(w+4)}\mathrm{tanh}^1\sqrt{w/(w+4)}`$ appearing in Eq. (89) of . In the limit of vanishing width $`(\widehat{\mathrm{\Gamma }}0)`$, we recover the results of . Eq. (29) cannot be further simplified analytically (except in some limiting cases), and must be evaluated numerically.
## 5 Discussion
### 5.1 Modification of the emission spectrum
By inspecting our final expression given in Eq. (29), the first thing we notice is the very strong similarity with the same result in the absence of the width. Only the function $`K(w,\widehat{\mathrm{\Gamma }})`$ contains the width, in the form of the dimensionless ratio $`\widehat{\mathrm{\Gamma }}`$. When this ratio is small, the width has a negligeable effect while on the contrary if $`\widehat{\mathrm{\Gamma }}1`$ then the width plays a dominant role.
A first simple conclusion is obtained by noticing that the width arrives in $`M_{\mathrm{eff}}^2`$ via the combination $`\mathrm{\Gamma }p(p+q_0)/q_0`$ which becomes large at small $`q_0`$. Therefore, the effect of the width is more important at small $`q_0`$. This property is illustrated by the plot of Fig. 4, which shows $`\mathrm{Im}\mathrm{\Pi }{}_{}{}^{\mu }_\mu `$ (obtained numerically from Eq. (29)) as a function of the photon energy $`q_0`$ (the invariant mass $`Q^2`$ is kept zero in this plot), for different values of the width $`\mathrm{\Gamma }`$.
On this plot, one can see that the $`q_0^1`$ behavior of bremsstrahlung at small $`q_0`$ is modified in presence of a width. Instead of that, one reaches a plateau<sup>11</sup><sup>11</sup>11A similar behavior has already been observed in a very different calculation by Weldon in (see Fig. 1 of ). However, the resummations considered in this paper deal with the possibility to emit more than one photon, and affect the spectrum only below the scale $`q_0\alpha _{_{EM}}T`$, much smaller than the scale at which $`\mathrm{\Gamma }g^2T\mathrm{ln}(1/g)`$ starts playing a role. The effect we are considering here appears at much larger photon energies (even including hard photons if $`\mathrm{\Gamma }`$ is large enough), since it is related to resummations of gluons and thus involves the strong coupling constant $`\alpha __S\alpha _{_{EM}}`$. when $`q_00`$. By comparing the second and third curves, we see that the value of $`q_0`$ at which we reach this plateau varies by a factor $`10^2`$ when $`\mathrm{\Gamma }`$ varies by the same factor. This is a consequence of the fact that the width enters in the result via the ratio $`\mathrm{\Gamma }/q_0`$ for small $`q_0`$. Moreover, the value of the plateau is proportional to $`1/\mathrm{\Gamma }`$ (this can be deduced from the plot since $`\mathrm{\Gamma }`$ increases by a factor $`10^2`$ between two consecutive curves). This leads to the conjecture that the rate has a very simple $`\mathrm{\Gamma }`$ dependence in the region where the width is dominant. Additionally, when $`\mathrm{\Gamma }`$ increases, the extension of this plateau increases as well. In particular, for a large enough width, even the region of hard photons is modified.
### 5.2 Region where the width is important
If one wants to go beyond this qualitative statement, and get a better idea of the region (in the $`(q,q_0)`$ plane) where the width is important, one must study the ratio $`\widehat{\mathrm{\Gamma }}`$. Indeed, the curve defined by the equation $`\widehat{\mathrm{\Gamma }}=1`$ is precisely the curve on which the effect of the width is comparable to the effect of $`M_{\mathrm{}}^2`$ and $`Q^2`$ (see Fig. 5). Above this curve (region II of Fig. 5), the width is an irrelevant parameter, and below this curve (region I – low invariant mass photons) we have the region where the width is the dominant collinear regulator.
At first sight, the region where the width is important looks very similar to the region where higher loop corrections are IR-sensitive to the scale $`g^2T`$ and must be resummed, as found in . We recall that in that work we considered the infrared structure of higher loop diagrams contributing to thermal photon production. It was shown that compensations between different cuts cancel all unscreened infrared divergences, and that the remaining terms are sensitive to gluon momenta down to a $`Q`$-dependent cutoff. In addition, we compared this cutoff with the scale $`\mu g^2T`$ of the magnetic mass, and concluded that higher loop diagrams would be sensitive to the magnetic mass for a small enough photon invariant mass (see the figure 5 of ).
This similarity comes from the fact that line in the $`(q,q_0)`$ plane where $`\mathrm{\Gamma }`$ starts to be important is given by the equation:
$$2\mathrm{\Gamma }\frac{q_0}{2}\frac{\mathrm{Re}M_{\mathrm{eff}}^2}{p(p+q_0)}.$$
(34)
while the line on which the magnetic mass $`\mu `$ becomes important has the following equation:
$$\mu \frac{q_0}{2}\frac{\mathrm{Re}M_{\mathrm{eff}}^2}{p(p+q_0)}.$$
(35)
In other words, the two lines are defined by comparing a common momentum scale alternatively with the width of the quarks and with the magnetic mass. The physical interpretation of the momentum scale appearing in the right hand side of Eqs. (34) and (35) will be given in the next section. The two boundaries are therefore very similar in QCD because $`\mathrm{\Gamma }g^2T\mathrm{ln}(1/g)`$ and $`\mu g^2T`$ are not very different. But on the other hand, $`\mathrm{\Gamma }`$ and $`\mu `$ have very different meanings<sup>12</sup><sup>12</sup>12In particular, the fact that $`\mathrm{\Gamma }`$ and $`\mu `$ are very close is specific to QCD. In QED, one would have $`\mu =0`$ and $`\mathrm{\Gamma }e^2T\mathrm{ln}(1/e)`$., and we may expect rather different interpretations for the two conditions Eqs. (34) and (35). In the next section, we show that Eq. (34) is closely related to the LPM effect.
There is also a more technical argument suggesting the difference of Eqs. (34) and (35), which becomes clear by examining how the width $`\mathrm{\Gamma }`$ would appear in perturbation theory. For that purpose, let us insert a self-energy correction on the leg of momentum $`R`$ in the two-loop diagram of Fig. 2. Since we are dealing with the width perturbatively here, the quark propagators contain only the asymptotic mass $`M_{\mathrm{}}`$. By a crude power counting, we can estimate that this insertion brings the following extra factor<sup>13</sup><sup>13</sup>13A self-energy correction of the type considered here modifies the real-part of the pole of the propagator by a mass-shift $`\delta M_{\mathrm{}}^2`$, negligible compared to $`M_{\mathrm{}}^2`$ arising from HTL corrections.
$$\frac{ir^0\mathrm{\Gamma }}{R^2M_{\mathrm{}}^2}\frac{1}{\mathrm{Re}M_{\mathrm{eff}}^2}\frac{i\mathrm{\Gamma }p(p+q_0)}{q_0}.$$
(36)
We can already see that a self-energy insertion increases the strength of the potential collinear divergences by bringing an additional denominator $`R^2M_{\mathrm{}}^2`$. On the other hand, this insertion does not modify the infrared properties of this diagram. By summing over the number of such insertions from $`0`$ to $`+\mathrm{}`$, we get a factor
$$\frac{\mathrm{Re}M_{\mathrm{eff}}^2}{\mathrm{Re}M_{\mathrm{eff}}^2i\mathrm{\Gamma }p(p+q_0)/q_0}.$$
(37)
Therefore, the effect of such a resummation is to substitute $`\mathrm{Re}M_{\mathrm{eff}}^2`$ by $`\mathrm{Re}M_{\mathrm{eff}}^2i\mathrm{\Gamma }p(p+q_0)/q_0`$ in the factor $`T^2/M_{\mathrm{eff}}^2`$ of collinear enhancement. This is precisely what has been observed in the more rigorous calculation of section 4. In addition, this simple argument demonstrates clearly that the mode of action of width insertions is to affect the collinear sector, leaving unmodified the infrared sector. On the contrary, the possibility to have a sensitivity to $`\mu `$ found in is related to infrared singularities due exclusively to transverse gluons<sup>14</sup><sup>14</sup>14In particular, longitudinal as well as transverse gluons contribute to $`\mathrm{\Gamma }`$. There is no contradiction with which found only the transverse gluons to be important, since the statement of was about infrared singularities..
### 5.3 Limit of dominant width
In region I, where the width becomes the dominant regulator, the ratio $`\widehat{\mathrm{\Gamma }}`$ is large, which enables us to make some additional approximations. In particular, we can perform very simply the integrals over $`w`$ and $`y`$ in Eq. (29). Indeed, we can first note that $`\widehat{\mathrm{\Gamma }}`$ sets the order of magnitude of the variable $`w`$. As a consequence, typical values of $`w`$ are large, and we can neglect the corrections $`\stackrel{~}{R}_{_{T,L}}`$ and $`\stackrel{~}{I}_{_{T,L}}`$ in the denominator of Eq. (29), as well as $`1/w`$ in front of $`y`$. Therefore, this equation becomes
$`\mathrm{Im}\mathrm{\Pi }{}_{}{}^{\mu }{}_{\mu }{}^{}(Q){\displaystyle \frac{2e^2g^2}{\pi ^4}}{\displaystyle \frac{T}{q_0^2}}{\displaystyle \underset{0}{\overset{+\mathrm{}}{}}}dp{\displaystyle \frac{p^2+(p+q_0)^2}{2}}[n__F(p+q_0)n__F(p)]`$
$`\times \left[{\displaystyle \underset{m=T,L}{}}{\displaystyle \underset{0}{\overset{1}{}}}{\displaystyle \frac{dx}{x}}\left|\stackrel{~}{I}_m\right|\right]\left[{\displaystyle \underset{0}{\overset{+\mathrm{}}{}}}{\displaystyle \frac{dw}{w^2}}{\displaystyle \frac{1}{2}}{\displaystyle \underset{0}{\overset{1}{}}}{\displaystyle \frac{dy}{\sqrt{1y}}}{\displaystyle \frac{y}{y^2+(4\widehat{\mathrm{\Gamma }}/w)^2}}\right]`$
$`={\displaystyle \frac{2e^2g^2}{\pi ^4}}{\displaystyle \frac{T}{q_0^2}}{\displaystyle \underset{0}{\overset{+\mathrm{}}{}}}𝑑p{\displaystyle \frac{p^2+(p+q_0)^2}{2}}[n__F(p+q_0)n__F(p)]\times {\displaystyle \frac{3\pi m_\mathrm{g}^2}{2\mathrm{R}\mathrm{e}M_{\mathrm{eff}}^2}}{\displaystyle \frac{\pi }{4\widehat{\mathrm{\Gamma }}}}`$ (38)
where $`m_\mathrm{g}gT`$ is the gluon thermal mass coming from the prefactor of $`\stackrel{~}{I}_{_{T,L}}`$. We see that this expression is of the form
$$\mathrm{Im}\mathrm{\Pi }{}_{}{}^{\mu }{}_{\mu }{}^{}(Q)e^2g^4\frac{T^3}{\mathrm{\Gamma }}\mathrm{\Phi }\left(\frac{q_0}{T}\right),$$
(39)
where $`\mathrm{\Phi }`$ is a function independent of $`\mathrm{\Gamma }`$ giving the shape of the photon spectrum. Therefore, in the region where the width dominates, the spectrum scales as $`\mathrm{\Gamma }^1`$. One notes that the combination $`\widehat{\mathrm{\Gamma }}\mathrm{Re}M_{\mathrm{eff}}^2`$ is proportional to $`\mathrm{\Gamma }p(p+q_0)/q_0`$, and therefore the shape of the spectrum is independent of $`Q^2`$ and $`M_{\mathrm{}}^2`$. Furthermore, for $`q_0T`$, it is independent of $`q_0`$ (see Fig. 4).
It may be interesting to compare in two extreme cases ($`q_0`$ soft and $`q_0T`$) the result of Eq. (39) with those obtained with a vanishing width. In the soft $`q_0`$ case, the result when $`\mathrm{\Gamma }=0`$ was obtained in :
$$\mathrm{Im}\mathrm{\Pi }{}_{}{}^{\mu }{}_{\mu }{}^{}(Q)e^2g^4\frac{T^3}{q_0}\frac{1}{g^2},$$
(40)
where we have explicitly isolated the factor $`T^2/M_{\mathrm{}}^21/g^2`$ that comes from the collinear enhancement, while we have now:
$$\mathrm{Im}\mathrm{\Pi }{}_{}{}^{\mu }{}_{\mu }{}^{}(Q)e^2g^4\frac{T^3}{q_0}\frac{q_0}{\mathrm{\Gamma }}.$$
(41)
The factor $`q_0/\mathrm{\Gamma }`$ is in fact the new enhancement factor coming from $`T^2/\mathrm{Im}M_{\mathrm{eff}}^2`$. Therefore, the two results differ only by the nature of the factor of collinear enhancement. For hard photon bremsstrahlung, we had on the other hand
$$\mathrm{Im}\mathrm{\Pi }{}_{}{}^{\mu }{}_{\mu }{}^{}(Q)e^2g^4T^2\frac{1}{g^2},$$
(42)
which becomes now
$$\mathrm{Im}\mathrm{\Pi }{}_{}{}^{\mu }{}_{\mu }{}^{}(Q)e^2g^4T^2\frac{T}{\mathrm{\Gamma }}.$$
(43)
Again, the two expressions differ only by the collinear enhancement factor: $`T^2/M_{\mathrm{}}^21/g^2`$ if $`\mathrm{\Gamma }=0`$ instead of $`T^2/\mathrm{Im}M_{\mathrm{eff}}^2T/\mathrm{\Gamma }`$ when the width is dominant. The numerical prefactors not written in Eqs. (42) and (43) could however be quite different since they reflect a different physics.
The $`\mathrm{\Gamma }^1`$ scaling law in Eq. (39) has also a very simple physical interpretation related to the fact that $`\mathrm{\Gamma }^1`$ is the mean free path of the quark in the medium. This result just tells us that the photon rate is proportional to the mean free path of the quark that emits the photon. In other words, quarks colliding very frequently do not have enough time to emit photons. This picture is at the basis of the Landau-Pomeranchuk-Migdal effect, that we discuss in the next section.
## 6 Connection with the LPM effect
### 6.1 Generalities on the LPM effect
The non-perturbative region found in the present paper can also be interpreted in a much more physical way, which seems to indicate that the region I in Fig. 5 is the region where the Landau-Pomeranchuk-Migdal effect modifies photon production by a plasma.
Let us first recall the condition for the LPM effect, by using a very heuristic argument valid for real photon production. In the bremsstrahlung process of photon production, the gluon exchanged during the scattering has a minimal momentum given by
$$l_{\mathrm{min}}\frac{q_0}{2}\frac{M_{\mathrm{}}^2}{p(p+q_0)},$$
(44)
where $`p`$ is the momentum of the quark. The inverse of this momentum transfer defines the “length” on which the photon is emitted, and for this reason is called the “coherence length” (denoted $`\lambda _{\mathrm{coh}}`$ in the following) for the production of this photon. It is also interpreted as the “formation time” of the photon. If $`\lambda _{\mathrm{coh}}`$ is much smaller that the typical distance between two consecutive scatterings (proportional to $`1/\mathrm{\Gamma }`$), then the photon rate is not affected by multiple scatterings. In other words, successive scatterings can be considered as independent. On the contrary, if the coherence length is larger than the mean free path, then successive scatterings are not independent anymore. This is the LPM effect. It is also possible to recast the previous condition as a comparison between the “formation time” and the mean free path of the quark in a very suggestive way: the photon must be produced before the quark scatters off another parton.
### 6.2 LPM effect in thermal field theory
Let us now show that the previous discussion arises automatically in the thermal field theory approach. In fact, there is in the calculation presented before a quantity very similar to the coherence length discussed in the semi-classical treatment of the LPM effect, namely,
$$\lambda _{\mathrm{coh}}^1\frac{q_0}{2}\frac{\mathrm{Re}M_{\mathrm{eff}}^2}{p(p+q_0)}.$$
(45)
This value agrees with the one of when $`Q^2=0`$ since it comes from kinematics, and generalizes it to the case where the invariant mass $`Q^2`$ is non vanishing. Additionally, this quantity was found in to be the lower bound for the momentum of exchanged gluons.
It is now possible to reformulate the condition $`\widehat{\mathrm{\Gamma }}=1`$ in more physical terms. Indeed, this condition can be rewritten as
$$\lambda _{\mathrm{mean}}\mathrm{\Gamma }^1\lambda _{\mathrm{coh}},$$
(46)
where $`\lambda _{\mathrm{mean}}`$ is the mean free path of the quark in the plasma, and the region where the width dominates corresponds to the condition $`\lambda _{\mathrm{mean}}<\lambda _{\mathrm{coh}}`$. At this point, the connection with the LPM effect is obvious: the region we found to be non-perturbative due to the width of the quarks is also the region where the LPM effect matters.
This discussion on the condition for the LPM effect to occur is in fact summarized in the expression for $`M_{\mathrm{eff}}^2`$ introduced in Eq. (25). This formula can be written very elegantly as:
$$M_{\mathrm{eff}}^2=\frac{2p(p+q_0)}{q_0}\left[\frac{1}{\lambda _{\mathrm{coh}}}+\frac{i}{\lambda _{\mathrm{mean}}}\right],$$
(47)
and the physical condition for the LPM effect<sup>15</sup><sup>15</sup>15In was investigated the influence of the finite mean free path of nucleons on axion (or neutrino pair) production by a supernova. The connection between the width $`\mathrm{\Gamma }`$ and the LPM effect is also mentioned in this paper. However, the condition given by Eq. (46) does not come out from the formalism of (it seems that the approximations made in for the bremsstrahlung do not enable one to track the coherence length, so that only $`\mathrm{\Gamma }`$ appears in the final result). What we have shown in the present section is that a careful calculation of the bremsstrahlung in thermal field theory generates automatically both terms of the comparison, through the effective mass given in Eq. (47). to be relevant is mathematically expressed by the dominance of $`\mathrm{Im}M_{\mathrm{eff}}^2`$ over $`\mathrm{Re}M_{\mathrm{eff}}^2`$. Conversely, when $`\mathrm{Re}M_{\mathrm{eff}}^2>\mathrm{Im}M_{\mathrm{eff}}^2`$, the perturbative approach is valid. As a side remark, let us note that the quantity $`M_{\mathrm{eff}}^2`$, which controls the physics of bremsstrahlung, combines in a very simple way three soft scales of the problem: $`Q^2`$, $`M_{\mathrm{}}^2`$ and $`\mathrm{\Gamma }`$. It is also possible to give a very simple physical interpretation of the ratio $`\widehat{\mathrm{\Gamma }}`$ which contains all the $`\mathrm{\Gamma }`$ dependence of the final result (Eq. (29)). Indeed, one can rewrite this ratio as $`\widehat{\mathrm{\Gamma }}=\lambda _{\mathrm{coh}}/\lambda _{\mathrm{mean}}`$, which is nothing but the typical number of coherent scatterings necessary to produce a photon.
The analogy with the standard treatment of the LPM effect is only partial though. Indeed, our estimate of the effect of the finite mean free path of the quarks is too crude (because of the constant width, and because of the missing vertex corrections) to be reliable from a quantitative perspective. The trend found here (suppression of the photon spectrum at small $`q_0`$) is however expected to be a solid prediction, and is in agreement with what is usually found from the LPM effect.
The main difference with the usual calculation of the LPM effect comes from the treatment of scattering centers: one usually assumes a fast moving charged particle going through a medium of “cold” (static) scattering centers. In the case we are considering here, both the particle emitting the photon and the scattering centers are thermalized particles, of comparable momenta. The approximations made in do not apply to this situation, and it is not clear whether a quantitative agreement is to be expected at all. In particular, the static approximation selects only Debye shielded longitudinal gluon exchanges, while the width (the resummation of which introduces multiple scatterings in our approach) of the thermalized quark receives contributions from transverse gluons as well. Contrary to , transverse gluon exchanges are very important in thermal field theory, and matter for the LPM effect as well.
There is a major difference between the LPM effect from multiple longitudinal gluon exchanges and multiple transverse gluon exchanges, which can be readily seen when one compares the respective ranges of electric and magnetic fields to the mean free path of the quark. Indeed, since $`\lambda _{\mathrm{mean}}m__D^1(gT)^1`$, the mean free path is much larger than the range of Debye screened electric fields. As a consequence, successive exchanges of longitudinal gluons are independent: they correspond to scatterings off different partons. On the contrary, we have $`\mu ^1>\lambda _{\mathrm{mean}}`$, which implies that the range of magnetic fields can extend beyond the mean free path of the quark. Therefore, successive exchanges of transverse gluons may correspond to scatterings off the same parton. This interpretation is supported by the infrared study performed in , where we found that only transverse gluons are causing trouble in the infrared sector, and that this problem occurs when $`\lambda _{\mathrm{coh}}>\mu ^1`$ (see Eq. (35)). This condition means that the production of a single photon occurs on a distance larger than the correlation length of magnetic fields: this emission process is therefore able to probe the scale of the magnetic screening, and the rate is expected to become non perturbative. This qualitative difference of transverse gluon exchanges is also supported by the fact that the contribution of longitudinal gluons to $`\mathrm{\Gamma }`$ is perturbative (saturated at 1-loop), while the transverse contribution to $`\mathrm{\Gamma }`$ is non-perturbative .
### 6.3 Comparison with the approach of Cleymans et al.
The LPM effect has already been studied in the context of photon production by a quark-gluon plasma, although with a very different approach, in . In this paper, the study follows the semi-classical approach of , to first deduce the photon rate from a single quark trajectory, and then average over the possible trajectories. Finally, the authors of manage to rewrite the rate as a function of a quantity $`F(𝜽,\tau )`$ which is related to the probability for the quark trajectory to undergo an angular deviation $`𝜽`$ after a time $`\tau `$. Additionally, they show that this object satisfies the following Fokker-Planck equation, which would be in our notations<sup>16</sup><sup>16</sup>16In , only soft photons are considered. As a consequence, in this comparison, we assume always $`q_0p`$, and therefore do not distinguish $`p`$ and $`p+q_0`$.:
$$\frac{F}{\tau }+i\frac{q_0}{2}\left[𝜽^2+\frac{M_{\mathrm{}}^2}{p^2}+\frac{Q^2}{q_0^2}\right]F=\frac{𝜽__S^2}{4}\mathbf{}_\theta ^2F,$$
(48)
where $`𝜽__S^2`$ is the mean square of scattering angle per unit length.
In order to make the connection with our approach more transparent, we can factorize out the quark width in the following way (recalling the fact that $`\mathrm{\Gamma }^1`$ gives the mean free path):
$$𝜽__S^216\overline{𝜽^2}\mathrm{\Gamma },$$
(49)
where $`\overline{𝜽^2}`$ is the average square of the scattering angle per collision (up to a purely numerical factor) and where the prefactor $`16`$ has been chosen for later convenience. We can now rewrite the equation satisfied by $`F`$ in the following way
$$\frac{F}{\tau }+iq_0\left[\frac{𝜽^2}{2}+\frac{1}{2p^2}\left\{\mathrm{Re}M_{\mathrm{eff}}^2+i\mathrm{Im}M_{\mathrm{eff}}^2\overline{𝜽^2}\mathbf{}_\theta ^2\right\}\right]F=0.$$
(50)
We see that this equation is governed by a complex number whose real and imaginary part also appear in our approach. In particular, the condition for having LPM effect (discussed after Eq. (2.21) of ) is the same as ours, given the fact that typical scattering angles satisfy $`\overline{𝜽^2}\mathrm{Re}M_{\mathrm{eff}}^2/p^2`$. Note however that our calculation also includes the case of hard photon production, and shows that transverse gluon exchanges are equally important, both points which are not contained in Cleymans et al.’s approach.
## 7 $`q^{}\overline{q}`$ annihilation
### 7.1 Technical differences
In section 4, when we performed the integral over the quark energy $`p_0`$ using the $`\delta (P^2M_{\mathrm{}}^2)`$ in Eq. 19, we considered only the positive energy solution, which corresponds to photon emission by the bremsstrahlung of a quark. The calculation is very similar for the contribution of $`p_0=\omega _p`$, with some peculiarities that we are going to highlight in this section.
One of the most important differences is that, since we have $`q_0>0`$, this contribution contains in fact the sum of two processes (see Fig. 3): the bremsstrahlung of an antiquark if $`\omega _p>q_0`$, and the annihilation of an off-shell quark (a quark put off-shell by a scattering) with an antiquark if $`\omega _p<q_0`$. The latter process has been studied in the case of zero width and real photons in and was indeed found to be a very important contribution to hard photon production. The purpose of this section is not to present a complete calculation of this process (which would require the discussion of a few other technical issues, especially in the region where $`Q^2>4M_{\mathrm{}}^2`$), but only to reproduce in this case the discussion of section 5, in order to have a qualitative picture of the influence of the width on this process.
The calculation of the $`p_0<0`$ contribution can be performed by the same method as the $`p_0>0`$ one. In particular, the angular integral over $`d\mathrm{\Omega }_l`$ generates an effective mass $`M_{\mathrm{eff}}^{\mathrm{\hspace{0.17em}2}}`$, the expression of which is
$$M_{\mathrm{eff}}^{\mathrm{\hspace{0.17em}2}}M_{\mathrm{}}^2+\frac{Q^2}{q_0^2}p(pq_0)+4i\frac{\mathrm{\Gamma }}{q_0}p(pq_0).$$
(51)
In addition to this change, the angular integral over $`d\mathrm{\Omega }_p`$ is now dominated by values of $`\theta `$ around $`\pi `$ (because $`1\mathrm{cos}\theta `$ in the expression of $`\mathrm{\Delta }`$ becomes $`1+\mathrm{cos}\theta `$)<sup>17</sup><sup>17</sup>17This difference is also obvious when one compares the process on the left of Fig. 3 with the two processes on the right of the same figure, which we are dealing with now.. A practical consequence of this is that $`r|pq_0|`$ in the collinear limit. All the subsequent steps can be reproduced and one is lead to the following final formula:
$`\mathrm{Im}\mathrm{\Pi }{}_{}{}^{\mu }{}_{\mu }{}^{}(Q){\displaystyle \frac{2e^2g^2}{\pi ^4}}{\displaystyle \frac{T}{q_0^2}}{\displaystyle \underset{0}{\overset{+\mathrm{}}{}}}dp{\displaystyle \frac{p^2+(q_0p)^2}{2}}[n__F(q_0p)n__F(p)]`$
$`\times {\displaystyle \underset{m=T,L}{}}{\displaystyle \underset{0}{\overset{1}{}}}{\displaystyle \frac{dx}{x}}{\displaystyle \underset{0}{\overset{+\mathrm{}}{}}}dw{\displaystyle \frac{\left|\stackrel{~}{I}_m\right|K(w,\widehat{\mathrm{\Gamma }})}{(w+\stackrel{~}{R}_m)^2+(\stackrel{~}{I}_m)^2}},`$ (52)
where the notations are the same as in Eq. (30), with $`M_{\mathrm{eff}}^{\mathrm{\hspace{0.17em}2}}`$ in place of $`M_{\mathrm{eff}}^2`$.
Another very important technical difference appears when one performs the integral over the quark momentum $`p`$. Indeed, the range of this integral is controlled by the statistical factors $`n__F(r_0)n__F(p_0)n__F(p)n__F(pq_0)`$. This factor is of order $`1`$ in a domain going from $`p=0`$ to $`p\mathrm{Max}(q_0,T)`$. In other words, if $`q_0`$ is soft, this range is the usual $`[0,T]`$ interval. On the contrary, if $`q_0T`$, the integral extends to $`pq_0T`$. A practical consequence is that for $`q_0T`$, the contribution of $`p_0<0`$ comes mainly from the bremsstrahlung of an antiquark (since most of $`dp`$ integral satisfies $`q_0<pT`$), while for $`q_oT`$ it is dominated by the $`q^{}\overline{q}`$ annihilation (since we then have $`p<q_0`$). In the intermediate region $`q_0T`$, this contribution is a mixture of both processes.
### 7.2 Modifications due to the width
Again, we see that the term depending on $`\mathrm{\Gamma }`$ in $`M_{\mathrm{eff}}^{\mathrm{\hspace{0.17em}2}}`$ comes with $`\mathrm{\Gamma }/q_0`$, indicating that the width will modify significantly the rate when $`q_0`$ is soft. But, in addition, when $`q_0T`$, $`p`$ can itself be of order $`q_0`$, so that the imaginary part of $`M_{\mathrm{eff}}^{\mathrm{\hspace{0.17em}2}}`$ is of order $`\mathrm{\Gamma }q_0`$, which can be much larger than $`M_{\mathrm{}}^2`$ at very large $`q_0`$. It means that the width also changes dramatically the spectrum of very hard real photons (in the case of bremsstrahlung, the width modifies only marginally the rate of hard photons, since the imaginary part of $`M_{\mathrm{eff}}^2`$ is of the same order of magnitude as $`M_{\mathrm{}}^2`$ if $`q_0T`$ and $`\mathrm{\Gamma }g^2T`$) coming from the $`q^{}\overline{q}`$ annihilation.
This has been checked by evaluating numerically Eq. (52), and the results are displayed in Fig. 6, where we plot the imaginary part of the photon polarization tensor as a function of $`q_0/T`$ (with the same values of temperature and coupling constant as in Fig. 4) for a set of values of $`\mathrm{\Gamma }`$ identical to the one chosen in Fig. 4.
From this plot, it is obvious that the regions $`q_0T`$ and $`q_0T`$ are very different. In fact, in the region $`q_0T`$, we obtain results that exactly match those of Fig. 4: indeed, in this region, the $`p_0<0`$ contribution is dominated by the bremsstrahlung of an antiquark, which as expected contributes equally as the bremsstrahlung of a quark.
The new feature coming with $`p_0<0`$ appears in the hard photon region ($`q_0T`$), where now the $`q^{}\overline{q}`$ annihilation contributes. At very small $`\mathrm{\Gamma }`$ (upper curve), we recover the result of according to which the process $`q^{}\overline{q}\gamma `$ gives a contribution to $`\mathrm{Im}\mathrm{\Pi }{}_{}{}^{\mu }{}_{\mu }{}^{}(Q)`$ that increases like $`q_0`$. Then, as $`\mathrm{\Gamma }`$ increases, we see a saturation at large enough $`q_0`$, to a value that behaves like $`\mathrm{\Gamma }^1`$. This confirms the above qualitative statement based on $`M_{\mathrm{eff}}^{\mathrm{\hspace{0.17em}2}}`$. In this case also, it is possible to obtain a formula like Eq. (39) showing explicitly the $`\mathrm{\Gamma }`$ scaling law when $`\mathrm{\Gamma }`$ is the dominant collinear regulator.
The contribution of the process $`q^{}\overline{q}\gamma `$, first evaluated at leading order in , has been found to be phenomenologically important for the production of direct photons in the $`210`$ GeV range, because it dominates previous estimates by a factor of order $`5`$ . The inclusion of this contribution in heavy nuclei collisions simulations has lead to a good agreement with the measured rates from the WA98 experiment . These rates did not include the LPM suppression advocated here, but one must realize that the energy range where thermal photons are relevant is also the region of minimum sensitivity to the width $`\mathrm{\Gamma }`$ (around the minimum of the curves in figure 6), since for a temperature of a few hundred MeV, the GeV range corresponds to $`q_0/T10`$. It is too early to be more quantitative here given the fact that vertex corrections have been completely disregarded in the present work, but having a precise prediction for this process would definitely be of important phenomenological interest.
### 7.3 LPM effect in $`q^{}\overline{q}`$ annihilation
When $`p_0<0`$, the condition under which the width is the dominant collinear regulator can be rewritten as
$$2\mathrm{\Gamma }>\frac{q_0}{2}\frac{\mathrm{Re}M_{\mathrm{eff}}^{\mathrm{\hspace{0.17em}2}}}{p|q_0p|},$$
(53)
and the quantity in the right hand side is the minimal value $`l_{\mathrm{min}}`$ for the momentum $`L`$ of the gluon exchanged in the scattering. Therefore, its inverse is also the coherence length for the emission process, and the above condition is nothing but the criterion for having LPM effect.
As a consequence, we are lead to the conclusion that the LPM effect plays a role not only for the photon production by bremsstrahlung at low energy, but also for the $`q^{}\overline{q}`$ annihilation at high energy. This conclusion was also recently obtained in by a semi-classical method following the original method of Migdal . The authors of in fact studied the inverse of the process $`q^{}\overline{q}\gamma `$, namely the production of a fermion pair out of a photon (this process is made possible by the fact that one of the fermions is produced off-shell and then scatters in the medium). They found that the LPM effect is important for the fragmentation of very high energy photons. This is equivalent to our observation that the LPM effect modifies $`q^{}\overline{q}\gamma `$ for large $`q_0`$.
## 8 Conclusions
In this paper, we have studied the effect of the finite mean free path of quarks on photon production by a quark gluon plasma. The main result is that such a correction may have important effects by affecting the collinear sector.
We have first considered the effect of a quark width in the simplest 1-loop diagram, and shown that potentially large 1-loop contributions cancel when a more complete summation is performed. This cancellation is not new in fact, and has been noticed elsewhere. It has connections with the fact that there are no hard thermal loops for vertices with 2 photons and any number of gluons.
Two-loop contributions beyond the eikonal approximation escape this cancellation, and lead to a result which exhibits features of the LPM effect. It is indeed possible to interpret the region where the width of the quark is important as the region where the formation length of the photon is larger than the mean free path of the quark. Our study also shows that the LPM effect modifies the emission spectrum of very hard photons. Despite its nice interpretation, our present result is not complete because it does not consider the vertex corrections that should accompany the resummation of a width.
Multiple scatterings due to longitudinal gluon exchanges are independent, as in the semi-classical treatment of the LPM effect. On the contrary, the transverse sector displays features that are qualitatively different from what is usually accepted in the semi-classical treatment (which just disregards transverse gluon exchanges). Indeed, we find that multiple transverse gluon exchanges are very important also, and that they are not independent due to the long range of magnetic fields.
The various scales in the problem are summarized in Fig. 7. Three of these scales are intrinsic to the quark gluon plasma: the Debye screening length $`m__D^1`$, the mean free path of the quarks $`\lambda _{\mathrm{mean}}\mathrm{\Gamma }^1`$ and the magnetic screening length $`\mu ^1`$. The fourth scale, the coherence length $`\lambda _{\mathrm{coh}}`$ (or formation length of the photon), depends on the energy and invariant mass of the photon one wants to observe, and should be compared to the first three.
This leads to three zones which have simple physical interpretations. In addition, the nature (and complexity) of the dominant higher loop corrections depend on the coherence length.
$``$ Region A: $`0<\lambda _{\mathrm{coh}}<\lambda _{\mathrm{mean}}`$, and photon production is dominated by single scatterings. The large scale structure of magnetic fields is irrelevant. The only relevant diagram is the 2-loop diagram (with HTL resummed quarks) already considered in . The contributions of transverse and longitudinal gluons are comparable.
$``$ Region B: $`\lambda _{\mathrm{mean}}<\lambda _{\mathrm{coh}}<\mu ^1`$, which implies that photon production is affected by multiple scatterings (LPM effect). In this region, the emission process is not yet sensitive to the magnetic screening. Indeed, according to , this sensitivity comes in when $`\lambda _{\mathrm{coh}}>\mu ^1`$. It is sufficient to resum a width (saturated at 1-loop<sup>18</sup><sup>18</sup>18When the IR cutoff is $`\lambda _{\mathrm{coh}}^1\mu `$, corrections to $`\mathrm{\Gamma }`$ coming from topologies beyond 1-loop are suppressed by $`g^2T/\lambda _{\mathrm{coh}}^11`$.) on the quarks, and to consider vertex corrections involving both longitudinal and transverse gluons. So far, we have said nothing about the vertex corrections that should come with self-energy insertions. Nevertheless, it is reasonable to assume from Ward identities that 1-loop self-energy insertions talk to ladder corrections<sup>19</sup><sup>19</sup>19The mechanism which makes these ladder corrections important is the same as in , and has nothing to do with infrared singularities. The arguments of cannot exclude them, even if the infrared cutoff is $`\lambda _{\mathrm{coh}}^1\mu `$. This is why these corrections are due to both longitudinal and transverse gluons.. We can note also that the discussion of this region is somewhat academic, because its extension is only proportional to $`\mathrm{ln}(1/g)`$. Additionally, the new topologies of the region C are at most suppressed by powers of $`\mathrm{ln}(1/g)^1`$ if evaluated in region B.
$``$ Region C: $`\mu ^1<\lambda _{\mathrm{coh}}`$. The LPM effect still modifies photon production. In addition, the emission of a photon lasts long enough for the process to be sensitive to the magnetic screening. This is the physical meaning of the result of . Diagrams with an arbitrary number of ultrasoft transverse gluons, connected in all the possible ways, must be resummed. In addition, the width included on the quarks contains non perturbative contributions due to transverse gluons , and the effective gluon vertices must be corrected to hold for ultrasoft momenta .
A full study (including vertex corrections) of the corrections to photon production due to longitudinal gluons seems within the reach of perturbation theory, but the situation is very different in the transverse sector which seems far beyond the possibilities of perturbative methods. In this respect, photon production is not very different from the calculation of the quark damping rate . New tools, like functional methods, transport equations, and eventually lattice techiques, are presumably the way to go in this area. In particular, the picture emerging from Fig. 7 suggests an analogy with the successive resummations of , in which the scales $`T`$, $`gT`$ and $`g^2T\mathrm{ln}(1/g)`$ are successively integrated out. Indeed, going to larger and larger coherence lengths requires the inclusion of more and more complicated topologies, a procedure which amounts to integrate out degrees of freedom encountered at smaller length scales.
## Acknowledgements
F.G. would like to thank R. Pisarski, A. Peshier, R. Venugopalan, D. Bödeker and S. Jeon for very interesting comments and discussions. H.Z. thanks the Institute of Nuclear Theory at the University of Washington for its hospitality where part of this work has been done and Dam Thanh Son for fruitful discussions. The work of F.G. is supported by DOE under grant DE-AC02-98CH10886.
## Appendix A Separation of the various cuts
Equation (27) includes automatically the sum of the cuts $`(c)`$ and $`(d)`$ of Fig. 2. To see it explicitly, let us first denote
$`F(\mathrm{\Gamma }){\displaystyle \frac{1}{u+\frac{M_{\mathrm{eff}}^2}{2p^2}}}`$
$`G(\mathrm{\Gamma }){\displaystyle \frac{ϵ(\mathrm{\Gamma })}{\left[\left(u+\frac{M_{\mathrm{eff}}^2}{2p^2}+\frac{L^2}{2p^2}\right)^2\frac{L^2}{p^2}\frac{M_{\mathrm{eff}}^2}{p^2}\right]^{1/2}}}.`$ (54)
Then, the discontinuity on the last line of Eq. (27) equals:
$`\mathrm{Disc}F(\mathrm{\Gamma })G(\mathrm{\Gamma })=F(\mathrm{\Gamma })G(\mathrm{\Gamma })F(\mathrm{\Gamma })G(\mathrm{\Gamma })`$
$`={\displaystyle \frac{[G(\mathrm{\Gamma })+G(\mathrm{\Gamma })]}{2}}\mathrm{Disc}F(\mathrm{\Gamma })+{\displaystyle \frac{[F(\mathrm{\Gamma })+F(\mathrm{\Gamma })]}{2}}\mathrm{Disc}G(\mathrm{\Gamma }).`$ (55)
In the last line, the first term corresponds to cut $`(d)`$ and the second term gives cut $`(c)`$. The fact that both cuts contribute to the photon rate when $`\mathrm{\Gamma }0`$ was to be expected. Indeed, the cut $`(d)`$ which corresponds to the process $`qq\gamma `$ is kinematically forbidden if $`\mathrm{\Gamma }=0`$, but is now allowed by the very fact that the width $`\mathrm{\Gamma }`$ reflects the collisions of the quasi-quark (and the channel $`qq\gamma `$ is allowed in the presence of a medium).
The integrals to be calculated are more complicated if one wants the two cuts separately. But it turns out that the difference between the two cuts is also very simple. This difference is given by:
$`(c)(d){\displaystyle \frac{[F(\mathrm{\Gamma })+F(\mathrm{\Gamma })]}{2}}\mathrm{Disc}G(\mathrm{\Gamma }){\displaystyle \frac{[G(\mathrm{\Gamma })+G(\mathrm{\Gamma })]}{2}}\mathrm{Disc}F(\mathrm{\Gamma })`$
$`=F(\mathrm{\Gamma })G(\mathrm{\Gamma })F(\mathrm{\Gamma })G(\mathrm{\Gamma }).`$ (56)
Therefore, in order to get $`(c)(d)`$, one has to start from Eq. (27) in which the last line is substituted by $`\mathrm{Disc}F(\mathrm{\Gamma })G(\mathrm{\Gamma })`$. The same transformations can be applied to this integral, and one finally obtains an equation similar to Eq. (29), but in which the function $`K(w,\widehat{\mathrm{\Gamma }})`$ is replaced by
$$L(w,\widehat{\mathrm{\Gamma }})\frac{1}{2}\underset{0}{\overset{1}{}}\frac{dy}{\sqrt{1y}}\frac{y+4/w}{(y+4/w)^2+(1y)(4\widehat{\mathrm{\Gamma }}/w)^2}.$$
(57)
This new integral is also elementary and can be obtained in closed form if needed.
It is obvious that in the limit of zero width $`(\widehat{\mathrm{\Gamma }}0)`$, the two integrals $`K(w,\widehat{\mathrm{\Gamma }})`$ and $`L(w,\widehat{\mathrm{\Gamma }})`$ become equal. The only way for this to happen is that the contribution of cut $`(d)`$ goes to zero. Therefore, we recover the fact that the process $`qq\gamma `$ disappears when $`\mathrm{\Gamma }=0`$.
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# Thermokinetic approach of the generalized Landau-Lifshitz-Gilbert equation with spin polarized current
## I Conservation laws
Spin-dependent transport processes in layered structures are described on the basis of the following simple picture. After entering in the k-th layer $`\mathrm{\Sigma }^k`$ , the incident current (which was spin polarized along the axis described by the unite vector $`\pm \stackrel{}{u}_{k1}`$ in the layer $`\mathrm{\Sigma }^{k1}`$) first aligns along the axis $`\pm \stackrel{}{u}_k`$. Inside the ferromagnetic layer, the population of spin up ($`N_+`$) and spin down ($`N_{}`$) is then not conserved due to spin-flip scattering, and some of the down spins relax to the up direction.
In this picture, the states of the sub-system $`\mathrm{\Sigma }^k`$ are described by the variables
$$(S^k,\stackrel{}{M_0^k},N_+^k,N_{}^k,\dot{\stackrel{}{M_0}}^k)$$
(1)
where $`S^k`$ is the entropy, $`\stackrel{}{M_0^k}=M_s^k\stackrel{}{u}_k`$ is the magnetization of the layer k without current, $`N_\pm ^k`$ is the number of conduction electrons with spin parallel to the unit vector $`\pm \stackrel{}{u}_k`$ and $`\dot{\stackrel{}{M_0}}^k`$ is the time derivative of the magnetization in the layer k without current.
If $`\stackrel{}{M_0^k}`$ and $`N_\pm ^k`$ are independent variables, the conservation of the magnetic momentum reads
$$\frac{d\stackrel{}{M^k}}{dt}=\dot{\stackrel{}{M_0}}^k+g\mu _B(\dot{N}_+^k\dot{N}_{}^k)\stackrel{}{u}_{k1}$$
(2)
where $`\mu _B`$ is the Bohr magneton and $`g`$ is the Landé factor.
In order to write the conservation laws, the spin-flip scattering mechanism is described as a chemical reaction transforming a spin down into a spin up along the axis $`\pm \stackrel{}{u}_k`$. In this context, the reaction rate $`\dot{\mathrm{\Psi }}^k`$ is introduced as the number of chemical events per unit of time , . Let $`I_+^{kk+1}`$ and $`I_{}^{kk+1}`$ be the current of particles flowing from the layer $`\mathrm{\Sigma }^k`$ to the layer $`\mathrm{\Sigma }^{k+1}`$ due respectively to the electrons with spin in the direction $`\stackrel{}{u}_k`$ and to the electrons with spin in the direction $`\stackrel{}{u}_k`$. The conservation of the particles is then described by
$$\{\begin{array}{ccc}\frac{dN_+^k}{dt}\hfill & =\hfill & \alpha (k1;k)I_+^{k1k}+\left(1\alpha (k1;k)\right)I_{}^{k1k}I_+^{kk+1}\dot{\mathrm{\Psi }}^k\hfill \\ \frac{dN_{}^k}{dt}\hfill & =\hfill & \left(1\alpha (k1;k)\right)I_+^{k1k}+\alpha (k1;k)I_{}^{k1k}I_{}^{kk+1}+\dot{\mathrm{\Psi }}^k\hfill \end{array}$$
(3)
where $`\alpha `$ is the spin flip probability of the alignment process . In the case of ballistic alignment $`\alpha (k1;k)=cos^2\left(\frac{1}{2}\theta (k1;k)\right)`$ with $`\theta `$ the angle between $`\stackrel{}{u}_{k1}`$ and $`\stackrel{}{u}_k`$. Introducing the polarized current $`I_p`$, and the normal current $`I_N`$ defined by
$$\{\begin{array}{ccc}I_p^{k1k}\hfill & =\hfill & I_+^{k1k}I_{}^{k1k}\hfill \\ I_N^{k1k}\hfill & =\hfill & I_+^{k1k}+I_{}^{k1k}\hfill \end{array}$$
(4)
the Eq. (3) can be put into the form:
$$\{\begin{array}{ccc}\frac{dN_+^k}{dt}\hfill & =\hfill & I_+^{k1k}I_+^{kk+1}\left[\gamma (k1;k)I_p^{k1k}+\dot{\mathrm{\Psi }}^k\right]\hfill \\ \frac{dN_{}^k}{dt}\hfill & =\hfill & I_{}^{k1k}I_{}^{kk+1}+\left[\gamma (k1;k)I_p^{k1k}+\dot{\mathrm{\Psi }}^k\right]\hfill \end{array}$$
(5)
with $`\gamma (k1;k)=\mathrm{\hspace{0.17em}1}\alpha (k1;k)=sin^2\left(\frac{1}{2}\theta (k1;k)\right)`$
Inserting Eq. (5) into Eq. (2) the conservation of the magnetization reads:
$$\frac{d\stackrel{}{M^k}}{dt}=\stackrel{}{\dot{M_0}}^k+g\mu _B\stackrel{}{u}_{k1}\left(I_p^{k1k}I_p^{kk+1}\mathrm{\hspace{0.17em}2}\left[\gamma (k1;k)I_p^{k1k}+\dot{\mathrm{\Psi }}^k\right]\right)$$
(6)
The problem of the spin transfer between the polarized current and the magnetic layer is hence solved if the polarized current $`I_p`$ and the reaction rates $`\dot{\mathrm{\Psi }}`$ can be describe as functions of the experimentally accessible parameters, the current $`I_N`$, the electric field $`E_0`$ and the kinetic coefficients. This task is typically performed by the application of the first and second laws of the thermodynamics.
## II Kinetic equations
The system $`\mathrm{\Sigma }^k`$ is open to heat transfer, to chemical transfer, and to mechanical work due to the magnetization and magnetic fields. Let us define the heat and chemical power by $`P_\varphi `$ and the mechanical power by $`P_W`$. The first law of the thermodynamics applied to the layer $`\mathrm{\Sigma }^k`$ gives
$$\frac{dE^k}{dt}=P_\varphi ^{k1k}P_\varphi ^{kk+1}+P_W^{extk}$$
(7)
where $`P_W^{extk}=\stackrel{}{H}^{ext}\dot{\stackrel{}{M}^k}`$. Furthermore, with using the canonical definitions $`T^k=\frac{E^k}{S^k},\mu _\pm ^k=\frac{E^k}{N_\pm ^k},\stackrel{}{H}^k=\frac{E^k}{\stackrel{}{M}^k}`$ the energy variation is:
$$\frac{dE^k}{dt}=T^k\frac{dS^k}{dt}+\mu _+^k\frac{dN_+^k}{dt}+\mu _{}^k\frac{dN_{}^k}{dt}\stackrel{}{H}^k\frac{d\stackrel{}{M}^k}{dt}+\frac{E^k}{\dot{\stackrel{}{M}_0^k}}\frac{d\dot{\stackrel{}{M}_0^k}}{dt}$$
(8)
In the present work we limit our analysis to the isothermal case, $`T^k=T`$. The entropy variation of the sub-layer is deduced from the two last equations, after introducing the conservation laws:
$`T{\displaystyle \frac{dS^k}{dt}}`$ $`=`$ $`P_\varphi ^{k1k}P_\varphi ^{kk+1}+(\stackrel{}{H}^k\stackrel{}{H}^{ext})\dot{\stackrel{}{M_0}^k}`$ (11)
$`{\displaystyle \frac{1}{2}}\left(A^k\mathrm{\hspace{0.17em}2}g\mu _B(H^{k1}H^{ext,k1})\right)\left(I_p^{k1k}I_p^{kk+1}\mathrm{\hspace{0.17em}2}\gamma (k1;k)I_p^{k1k}\mathrm{\hspace{0.17em}2}\dot{\mathrm{\Psi }}^k\right)`$
$`{\displaystyle \frac{1}{2}}\mu _0(I_N^{k1k}I_N^{kk+1}){\displaystyle \frac{E^k}{\dot{\stackrel{}{M}_0^k}}}{\displaystyle \frac{d\dot{\stackrel{}{M}_0^k}}{dt}}`$
where the total chemical potential is $`\mu _0^k\mu _+^k+\mu _{}^k`$. We have furthermore defined $`H^{k1}\stackrel{}{H}^k\stackrel{}{u}_{k1}`$ and $`H^{ext,k}\stackrel{}{H}^{ext}\stackrel{}{u}_k`$ . The chemical affinity of the reaction, defined by $`A^k\frac{E^k}{\mathrm{\Psi }}=\mu _+^k\mu _{}^k`$ has also been introduced.
The entropy being an extensive variable, the total entropy variation of the system is obtained by summation over the layers 1 to $`\mathrm{\Omega }`$ where the layer 1 is in contact to the left reservoir $`R^l`$ and the layer $`\mathrm{\Omega }`$ is in contact to the right reservoir $`R^r`$. Letting
$$\stackrel{~}{A}^kA^k\mathrm{\hspace{0.17em}2}g\mu _B(H^{k1}H^{ext,k1})$$
(12)
the total entropy variation is:
$`T{\displaystyle \frac{dS}{dt}}`$ $`=`$ $`[\mathrm{}]^{R^l1}[\mathrm{}]^{\mathrm{\Omega }R^r}`$ (16)
$`+{\displaystyle \underset{k=1}{\overset{\mathrm{\Omega }}{}}}(\stackrel{}{H}^k\stackrel{}{H}^{ext})\stackrel{}{M_0}^k+{\displaystyle \underset{k=1}{\overset{\mathrm{\Omega }}{}}}\left({\displaystyle \frac{E^k}{\dot{\stackrel{}{M}_0^k}}}\right){\displaystyle \frac{d\dot{\stackrel{}{M}_0^k}}{dt}}`$
$`+{\displaystyle \underset{k=2}{\overset{\mathrm{\Omega }}{}}}{\displaystyle \frac{1}{2}}\left(\stackrel{~}{A}^{k1}\stackrel{~}{A}^k+\mathrm{\hspace{0.17em}2}\gamma (k1;k)\stackrel{~}{A}^k\right)I_p^{k1k}`$
$`+{\displaystyle \underset{k=2}{\overset{\mathrm{\Omega }}{}}}{\displaystyle \frac{1}{2}}(\mu _0^{k1}\mu _0^k)I_N^{k1k}+{\displaystyle \underset{k=1}{\overset{\mathrm{\Omega }}{}}}\stackrel{~}{A}^k\dot{\mathrm{\Psi }}^k`$
where the two first terms in the right hand side of the equality stand for the heat and chemical transfer from the reservoirs to the system $`\mathrm{\Sigma }`$.
The variation of entropy takes the form
$$T\frac{dS}{dt}=\underset{i}{}F_i\dot{X}^i+P^{ext}(t)$$
(17)
where $`F_i`$ are generalized forces and $`\dot{X}^i`$ are the conjugated generalized fluxes. The variation of entropy is composed by an external entropy variation $`P^{ext}(t)/T`$ and by an internal entropy variation $`dS^{int}/dt`$.
By applying the second law of thermodynamics $`dS^{int}/dt0`$ we are leading to introduce the kinetic coefficients $`l_{\alpha \beta }`$ such that $`dS^{int}/dt=_iF_i\left(_jl_{ij}F^j\right)`$. By identification with the expression (16), the kinetic equations are obtained:
$`\left[\begin{array}{c}I_N^{k1k}\\ I_p^{k1k}\\ \dot{\mathrm{\Psi }}^k\\ \dot{\stackrel{}{M_0}^k}\\ \frac{d\dot{\stackrel{}{M}_0^k}}{dt}\end{array}\right]=\left[\begin{array}{ccccc}l_{NN}& l_{Np}& l_{Nc}& l_{NM}& l_{N\dot{M}}\\ l_{pN}& l_{pp}& l_{pc}& l_{pM}& l_{p\dot{M}}\\ l_{cN}& l_{cp}& l_{cc}& l_{cM}& l_{c\dot{M}}\\ l_{MN}& l_{Mp}& l_{Mc}& l_{MM}& l_{M\dot{M}}\\ l_{\dot{M}N}& l_{\dot{M}p}& l_{\dot{M}c}& l_{\dot{M}M}& l_{\dot{M}\dot{M}}\end{array}\right]\left[\begin{array}{c}\frac{1}{2}(\mu _0^{k1}\mu _0^k)\\ \frac{1}{2}\left(\stackrel{~}{A}^{k1}\stackrel{~}{A}^k+\mathrm{\hspace{0.17em}2}\gamma (k1;k)\stackrel{~}{A}^k\right)\\ \stackrel{~}{A}^k\\ \stackrel{}{H}^k\stackrel{}{H}^{ext}\\ \left(\frac{E^k}{\dot{\stackrel{}{M}_0^k}}\right)\end{array}\right]`$ (33)
The indices $`N`$ and $`p`$ stand respectively for the normal and polarized transport processes (see section III), the indices $`c`$ stands for the spin-flip scattering chemical reaction and the indices $`M`$ and $`\dot{M}`$ account for the dynamics of the magnetization (see section IV).
The kinetic coefficients are state functions: $`l_{ij}=l_{ij}(S^k,\stackrel{}{M^k},N_+^k,N_{}^k,\dot{\stackrel{}{M}}_0^k)`$ and the symmetrized matrix is positive : $`\frac{1}{2}\left\{l_{ji}+l_{ij}\right\}_{\{ij\}}0`$. Furthermore, according to Onsager relations, the kinetic coefficients are symmetric or antisymmetric $`l_{ij}=\pm l_{ji}`$
We assume in the following that cross effects between electronic transport and ferromagnetic transport are negligible, so that $`l_{ij}=0`$ if $`i=\{N,p,c\}`$ and $`j=\{M,\dot{M}\}`$. Note that a polarized current is directly produced by a non-uniform magnetization state, through the coefficient $`\gamma (k1,k)`$ in Eq (33).
The physical meaning of the kinetic coefficients is described in the two following sections.
## III Giant magnetoresistance
We focus in this section on the first three equations of (33) which describes the electric transport with spin polarization. After performing the continuum limit, we have:
$`\left[\begin{array}{c}J_N\\ J_p\\ \dot{\mathrm{\Psi }}\end{array}\right]=\left[\begin{array}{ccc}L_{NN}& L_{Np}& L_{Nc}\\ L_{pN}& L_{pp}& L_{pc}\\ L_{Nc}& L_{pc}& L_{cc}\end{array}\right]\left[\begin{array}{c}\frac{\mu _0}{z}\\ \frac{A}{z}\mathrm{\hspace{0.17em}2}\stackrel{~}{\gamma }A\\ A\end{array}\right]`$ (43)
where $`\stackrel{~}{\gamma }=lim_{dz0}\frac{d\gamma }{dz}`$ and the choice of symmetric coefficients $`L_{pN}=L_{Np}`$ and $`L_{Nc}=L_{cN}`$ is motivated by the fact that we neglect in this work the direct effects of the magnetic field on the charge carriers .
In the framework of the two-channel approximation , the coupling between the two conduction bands is neglected (i.e. there is no cross effect between the currents $`J_+,J_{}`$, and $`\dot{\mathrm{\Psi }}`$):
$$L_{Nc}=L_{pc}=0;L_{NN}=L_{pp}=\frac{2\sigma _0}{2}$$
(44)
where $`\sigma _0>0`$ is the mean conductivity of the two spin channels. The conductivity asymmetry $`\beta `$ of the two channels is given by
$$L_{Np}\frac{2\sigma _0}{e}\beta $$
(45)
Equation Eq (43) leads then to the set of equations:
$`\left[\begin{array}{c}J_N\\ J_p\\ \dot{\mathrm{\Psi }}\end{array}\right]={\displaystyle \frac{\sigma _0}{e}}\left[\begin{array}{ccc}1& \beta & 0\\ \beta & 1& 0\\ 0& 0& \frac{e}{\sigma _0}L_{cc}\end{array}\right]\left[\begin{array}{c}\frac{\mu _0}{z}\\ \frac{A}{z}\mathrm{\hspace{0.17em}2}\stackrel{~}{\gamma }A\\ A\end{array}\right]`$ (55)
with $`1\beta ^2`$ and $`L_{cc}0`$.
In the stationary state, $`\frac{J_N}{z}=0`$, and assuming that $`\beta `$, $`\sigma _0`$ and $`\gamma `$ are approximately independent of z the diffusion equation of the chemical affinity is deduced (see appendix):
$$\frac{^2A}{z^2}=\left(\frac{1}{l_{sf}^2}+\frac{1}{l_{DW}^2}\right)A+k\frac{A}{z},$$
(56)
where the spin diffusion length $`l_{sf}`$ is given by
$$l_{sf}\sqrt{\frac{\sigma _0\left(1\beta ^2\right)}{2eL_{cc}}},$$
(57)
the “Domain Wall” diffusion length $`l_{DW}`$ is given by
$$l_{DW}\sqrt{\frac{\left(1\beta ^2\right)}{4\stackrel{~}{\gamma }^2}},$$
(58)
and the parameter k is given by
$$k\stackrel{~}{\gamma }\frac{2\beta ^2}{\left(1\beta ^2\right)}$$
(59)
Note that the chemical affinity A is equal to the difference of the chemical potentials of the two conduction bands $`A=\mu _+\mu _{}`$. If we assume no rotation of the spin polarization axis, $`\stackrel{~}{\gamma }=\mathrm{\hspace{0.17em}0}`$ (which implies antiparallel magnetic configuration at the interface), then equation Eq (56) is the well known diffusion equation describing the so-called “spin accumulation” or “spin depletion” effect responsible for the giant magnetoresistance , , , . A straightforward calculation (see appendix) leads to the giant magnetoresistance of the interface:
$$R^{GMR}=2\frac{\beta ^2}{\sigma _0(1\beta ^2)}l_{sf}$$
(60)
## IV Landau-Lifshitz-Gilbert (LLG) equation
In this section we focus on the magnetic transport equation without electric current: $`J_N=J_p=0`$. From Eq (11) the entropy variation reduces to:
$$T\frac{dS}{dt}=P_\varphi ^{extin}P_\varphi ^{inext}+(\stackrel{}{H}\stackrel{}{H}^{ext})\frac{d\stackrel{}{M_0}}{dt}+\left(\frac{E}{\dot{\stackrel{}{M}_0}}\right)\frac{d\dot{\stackrel{}{M}_0}}{dt}$$
(61)
So that the application of the second law of thermodynamics yields,
$`\{\begin{array}{cc}(\stackrel{}{H}\stackrel{}{H}^{ext})& =\stackrel{~}{l}_{MM}\frac{d\stackrel{}{M_0}}{dt}+\stackrel{~}{l}_{M\dot{M}}\frac{d\dot{\stackrel{}{M}_0}}{dt}\\ \left(\frac{E}{\dot{\stackrel{}{M}_0}}\right)& =\stackrel{~}{l}_{\dot{M}M}\frac{d\stackrel{}{M_0}}{dt}+\stackrel{~}{l}_{\dot{M}\dot{M}}\frac{d\dot{\stackrel{}{M}_0}}{dt}\end{array}`$ (64)
where the kinetic coefficients $`\stackrel{~}{l}_{\alpha \beta }`$ are the coefficients of the inverse matrix $`\{l_{\alpha \beta }\}_{\{\alpha \beta \}}^1`$.
Note that in adiabatically closed systems, $`\left(\frac{E}{\dot{\stackrel{}{M}_0}}\right)`$, $`H`$ and $`\frac{d\stackrel{}{M_0}}{dt}`$ are state functions (i.e depend only of the state variables (S, $`\stackrel{}{M}_0`$, $`\dot{\stackrel{}{M}}_0`$), and not on $`\stackrel{}{H}^{ext}`$). Since the kinetic coefficients are also state functions, the first equation in (64) shows hence that $`\frac{d\dot{\stackrel{}{M}_0}}{dt}`$ depends on $`\stackrel{}{H}^{ext}`$. We are then leading to impose $`\stackrel{~}{l}_{\dot{M}\dot{M}}=\mathrm{\hspace{0.17em}0}`$ in order to satisfy the second equation in (64), which gives the magnetic kinetic energy . The coefficient $`\stackrel{~}{l}_{\dot{M}M}`$ can be identified to the magnetic mass, and the first equation in (64) gives the total magnetic force $`\stackrel{}{F}^{mag}`$ acting on the system:
$$\stackrel{}{F}^{mag}\stackrel{~}{l}_{M\dot{M}}\frac{d^2\stackrel{}{M}_0}{dt^2}=(\stackrel{}{H}\stackrel{}{H}^{ext})\stackrel{~}{l}_{MM}\frac{d\stackrel{}{M_0}}{dt}$$
(65)
equation (65) rewrites:
$$\stackrel{}{F}^{mag}=\frac{}{\stackrel{}{M}_0}(E\stackrel{}{H}^{ext}.\stackrel{}{M}_0)\eta \frac{d\stackrel{}{M_0}}{dt}$$
(66)
where we have identified the Gilbert friction coefficient , to $`\eta =\stackrel{~}{l}_{MM}`$.
The theorem of the kinetic momentum gives the equation of the dynamics:
$$\frac{d\stackrel{}{M}_0}{dt}=\mathrm{\Gamma }\left(\stackrel{}{M}_0\times \stackrel{}{F}^{mag}\right)=\mathrm{\Gamma }\stackrel{}{M}_0\times \left\{\frac{V}{\stackrel{}{M}_0}\eta \frac{d\stackrel{}{M_0}}{dt}\right\}$$
(67)
where $`\mathrm{\Gamma }`$ is the gyromagnetic ratio and the magnetic Gibbs potential is defined by $`V=E+\stackrel{}{M}_0.\stackrel{}{H}^{ext}`$. Equation (67) is the well known Gilbert equation , , and can be put into the following Landau-Lifshitz form. In the case of uniform magnetization we have $`\stackrel{}{M}_0=M_s\stackrel{}{u}_0`$, where $`M_s`$ is the saturation magnetization, Eq (67) rewrites
$$\dot{u_0}=g^{}\left(\stackrel{}{u}_0\times \stackrel{}{}V\right)h^{}\stackrel{}{u}_0\times \left(\stackrel{}{u}_0\times \stackrel{}{}V\right)$$
(68)
where $`\stackrel{}{}`$ is here the gradient operator on the surface of a unite sphere. The phenomenological parameters h’ and g’ are linked to the gyromagnetic ratio $`\mathrm{\Gamma }`$ and the Gilbert damping coefficient $`\eta `$ by the relations
$`\{\begin{array}{ccc}h^{}& =& \frac{\mathrm{\Gamma }\alpha }{(1+\alpha ^2)M_s}\\ g^{}& =& \frac{\mathrm{\Gamma }}{(1+\alpha ^2)M_s}\\ \alpha & =& \eta \mathrm{\Gamma }M_s\end{array}`$ (72)
## V LLG equation with spin polarized current
Let us assume an interface composed by an incident current $`I_p^i`$ of conductivity asymmetry $`\beta `$, polarized in the direction $`\stackrel{}{e}_p`$ which enters in a ferromagnetic layer (F) polarized in the direction $`\stackrel{}{u}_0`$ with conductivity asymmetry $`\beta `$. The transfer of magnetic moments is describes by the term $`\frac{dN_p^F}{dt}=\dot{N}_+^F\dot{N}_{}^F`$. Equation (5) rewrites
$`{\displaystyle \frac{dN_p^F}{dt}}`$ $`=`$ $`I_p^iI_p^F\mathrm{\hspace{0.17em}2}\gamma I_p^F=\beta _{eff}I_N`$ (73)
where $`\beta _{eff}=\mathrm{\hspace{0.17em}2}\beta (12\gamma (1\gamma ))`$ and the expressions of $`I_p^F`$ and $`I_p^i`$ are derived in the appendix.
The change of the magnetic moment of the layer due to the polarized current is given by equation (6) and (68):
$`\dot{u}g^{}\left(\stackrel{}{u}_0\times \stackrel{}{}V\right)h^{}\stackrel{}{u}_0\times \left(\stackrel{}{u}_0\times \stackrel{}{}V\right)+{\displaystyle \frac{g\mu _B}{M_0}}\beta _{eff}I_N\stackrel{}{e}_p`$ (74)
where the first, second and third term in the right hand side are respectively the precession term (or transverse relaxation), the longitudinal relaxation term, and the spin transfer due to spin polarized conduction electrons.
In order to estimate the effect of the injection of spin polarized current, the equation (74) is applied to the case of monodomain ferromagnet with applied field oriented at the angle $`\theta `$ from a single anisotropy axis (see Fig. 1). If the vector $`\stackrel{}{u}`$ makes an angle $`\phi `$ from the anisotropy axis, the Gibbs energy density can be written in the following form :
$`V(\phi \psi )=KS`$ $`(cos^2\phi \mathrm{\hspace{0.17em}2}h(cos(\theta )cos(\phi )`$ (76)
$`+sin(\theta )sin(\phi )cos(\psi ))`$
where $`h=H^{ext}/H_a`$ is the reduced applied field defined with the anisotropy field $`H_a`$, $`K=H_aM_s`$ is the anisotropy constant, S is the section and $`\psi `$ is the out-of-plane coordinate of the vector $`\stackrel{}{u}`$. Due to the cylindrical geometry, $`\psi =0`$. Before injecting the current, the angle $`\phi _0`$ is given by the equilibrium condition $`\stackrel{}{}V=0`$.
The precessional term can be neglected in (74) (low frequency response and/or high damping limite ), and $`MM_0`$.
Experiments and samples are described in Ref. , , . Ni nanowires are obtained by the method of electrodeposition in track etched membrane templates. A micro-contact is realized, and the magnetoresistance of a single nanowire is measured. The wires are about 70 nm diameter and 6000 nm length and the magnetic energy is dominated by the Zeeman energy term and the shape anisotropy (or magnetostatic term), very close to that of an infinite cylinder. The anisotropy field is calculated to be $`\mu _0H_a0.3`$ T.
The effect of the spin-polarized current was evidenced experimentally by injecting a strong current of about $`2.10^7A/cm^2`$ at a fixed value of the external field $`h=h_{sw}\mathrm{\Delta }h`$ smaller than the field $`h_{sw}`$ where the switching occurs without current. The magnetization switch occurs at the angle $`\phi ^c(\theta )`$. The maximum distance $`\mathrm{\Delta }h`$ where the jump of the magnetization can still be observed corresponds then to the variation of the angle $`\mathrm{\Delta }\phi =\phi ^c\phi _0`$ needed to shift the magnetization up to the unstable state.
For steady states, inserting $`h=h_{sw}\mathrm{\Delta }h`$, $`\phi =\phi _c`$, Eq. (74) leads to
$$\mathrm{\Delta }h=h_{sw}(\theta )\frac{2cI_N(\stackrel{}{e}_p.\stackrel{}{v})sin(2\phi ^c)}{sin(\phi ^c\theta )}$$
(77)
where $`\stackrel{}{v}`$ is the polar vector perpendicular to $`\stackrel{}{u}`$. The parameter c is defined by the relation
$$c=\frac{\beta _{eff}\mathrm{}}{eKv_a\alpha }$$
(78)
where the activation volume $`v_a`$ of magnetization $`M_s`$ was estimated to be $`v_a\mathrm{\hspace{0.17em}10}^{22}m^3`$, $`K\mathrm{\hspace{0.17em}10}^5J/m^3`$ , and $`\beta _{eff}\beta \mathrm{\hspace{0.17em}0.3}`$ , $`\alpha \mathrm{\hspace{0.17em}0.15}`$ . We obtain $`c\mathrm{\hspace{0.17em}200}A^1`$.
All parameters in Eq. (77) are known if the magnetization reversal mode, which describes the irreversible jump, is known. In some few theoretical models of magnetization reversal , the functions $`H_{sw}(\theta )`$ and $`\phi _c(\theta )`$ are analytical. In the framework of the present empirical approach, the experimental data are fitted by the relation deduced from a curling reversal mode in an infinite cylinder , :
$$h_{sw}(\theta )=\frac{a(a+1)}{\sqrt{a^2+(2a+1)cos^2(\theta )}}$$
(79)
The single adjustable parameter $`a=k(R_0/r)^2`$ is defined by the geometrical parameter k , by the exchange length $`R_0=\mathrm{\hspace{0.17em}20}nm`$ , and by the radius of the wire r. The experimental points $`H_{sw}(\theta )`$ are fitted in Fig. 2.
We obtained $`a=0.15`$ (which corresponds to r of about 60 nm). The relation between the angle of the applied field $`\theta `$ and the angle of the magnetization $`\phi ^c`$ is:
$$tan(\theta )=\frac{a+1}{a}tan(\phi ^c)$$
(80)
The curve $`\mathrm{\Delta }h`$, evaluated from Eq. (77) by numerical resolution with a polarization in the direction of the wire axis $`\stackrel{}{e}_p.\stackrel{}{v}=sin(\phi _0),`$ is plotted in Fig. 3 and Fig. 4, together with the experimental data. A strong discrepancy from the linear curve of $`\mathrm{\Delta }h(I_N)`$ at small current pulses can be observed. Above a critical current corresponding to about $`10^7A/cm^2`$ the linear fit gives a parameter $`c=\mathrm{\hspace{0.17em}190}`$, which is in accordance with the rough evaluation of Eq. (78). This critical current below which the linear regime failed in Fig. 4 could be interpreted following Ref and Ref as the current needed in order to excite spin waves or other magnetization inhomogeneities . The curve given by Eq. (77) can then be plotted without adjustable parameter (Fig. 4). The divergence at $`90^o`$ is due to the numerical resolution of Eq. (77) (numerator and denominator tend to zero at $`\theta =\mathrm{\hspace{0.17em}90}^o`$ angle).
## VI Conclusion
A systematic thermokinetic description of a metallic ferromagnetic layer open to electronic spin polarized reservoirs has been performed. At constant temperature, assuming the two current approximation and neglecting direct action of the magnetic field on charge carriers, five coupled transport equations account for the complexity of the system. The approximation of the explicit uncoupling of the transport processes leads to the known results about GMR and Landau-Lifshitz-Gilbert equations for magnetization dynamics. Within this approximation and on the basis of the conservation equation of the magnetic moment, the description of both polarized current and magnetization dynamics leads to a generalized Landau-Lifshitz-Gilbert equation. The application of this model to experimental data about current-induced magnetization reversal is performed. The existence of a critical current indicates that the kinetics of magnetization inhomogeneity also plays an important role. However, the comparison with experimental data shows that the derived thermokinetic generalized Landau-Lifshitz-Gilbert equation provides a description of the basic mechanism responsible for the effect of polarized-current-induced-magnetization reversal.
## VII Acknowledgment
This work is directly inspired by the synthesis performed by Professor C. Gruber on thermokinetics, and I thank him for his corrections and comments. I thank Professor J.-Ph. Ansermet for permanent support. Special thanks to Derek Kelly for his participation to the experimental work.
## VIII Appendix
The appendix is structured in three parts. In the first part (A), the equation of for the difference of chemical potentials $`A=\mathrm{\Delta }\mu `$ is derived from equation (55) in the case of steady states. In part (B) the equation is applied to the simplified case of GMR or spin accumulation, where the polarization axis is assumed constant through the interface. The GMR of the interface is deduced. In part (C), the equation is applied in the framework of the experimental study of polarized-current-induced-magnetization-reversal, where an abrupt charge of the polarization axis occurs at the interface.
(A) Assuming that the kinetic coefficients which coupled the dynamics of the magnetization and the electric currents vanish, we obtained the following set of kinetic equations (55):
$`\left[\begin{array}{c}J_N\\ J_p\\ \dot{\mathrm{\Psi }}\end{array}\right]={\displaystyle \frac{\sigma _0}{e}}\left[\begin{array}{ccc}1& \beta & 0\\ \beta & 1& 0\\ 0& 0& \frac{e}{\sigma _0}L_{cc}\end{array}\right]\left[\begin{array}{c}\frac{\mu _0}{z}\\ \frac{A}{z}+\mathrm{\hspace{0.17em}2}\stackrel{~}{\gamma }A\\ A\end{array}\right]`$ (90)
In the stationary state $`\frac{}{z}J_N(z)=0`$, and assuming that $`\beta `$, $`\sigma _0`$ and $`\gamma `$ are approximately independent of z the diffusion equation of the chemical affinity is deduced
$$\frac{^2\mu _0}{z^2}=\beta \left(\frac{^2A}{z^2}2\frac{(\stackrel{~}{\gamma }A)}{z}\right)$$
(91)
Inserting in (90) yields
$$\frac{J_p}{z}=\frac{\sigma _0}{e}\left((\beta ^21)\frac{^2A}{z^2}+2\frac{(\stackrel{~}{\gamma }A)}{z}\right)$$
(92)
and by integration,
$$J_p=\frac{\sigma _0}{e}\left((\beta ^21)\frac{A}{z}+2\stackrel{~}{\gamma }A\right)$$
(93)
where we assumed that $`J_p(\mathrm{})=\mathrm{\hspace{0.17em}0}`$ .
On the other hand, from the conservation equations we have
$$\frac{d}{dt}(N_+^kN_{}^k)=I_p^{k1k}I_p^{kk+1}\mathrm{\hspace{0.17em}2}\dot{\mathrm{\Psi }}^k\mathrm{\hspace{0.17em}2}\gamma (k1,k)I_p$$
(94)
At the continuum limit, we obtain the following relation
$$\frac{dn_p}{dt}=\frac{J_p}{z}\mathrm{\hspace{0.17em}2}L_{cc}A2\stackrel{~}{\gamma }J_p$$
(95)
where $`n_p`$ is the density of spin polarized conduction electrons.
Equation (95) rewrites
$$\frac{J_p}{z}=\mathrm{\hspace{0.17em}2}L_{cc}A2\stackrel{~}{\gamma }J_p\frac{dn_p}{dt}$$
(96)
where $`\frac{dn_p}{dt}`$ is constant for steady states . Furthermore, inside the ferromagnet and far away from the interface, $`J_p`$ is constant, whence
$$\frac{dn_p}{dt}=\mathrm{\hspace{0.17em}2}L_{cc}^{\mathrm{}}A_{\mathrm{}}=\mathrm{\hspace{0.17em}0}$$
(97)
where we assumed for simplicity that $`L_{cc}^{\mathrm{}}=\mathrm{\hspace{0.17em}0}`$. Together with (96), (93) and (92) the differential equation for A(z) is obtained:
$$\frac{^2A}{z^2}=\left(\frac{1}{l_{sf}^2}+\frac{1}{l_{DW}^2}\right)A+k\frac{A}{z},$$
(98)
where the spin diffusion length $`l_{sf}`$ is given by
$$l_{sf}\sqrt{\frac{\sigma _0\left(1\beta ^2\right)}{2eL_{cc}}},$$
(99)
the “Domain Wall” diffusion length $`l_{DW}`$ is given by
$$l_{DW}\sqrt{\frac{\left(1\beta ^2\right)}{4\stackrel{~}{\gamma }^2}},$$
(100)
and the parameter k is given by
$$k\stackrel{~}{\gamma }\frac{2\beta ^2}{\left(1\beta ^2\right)}$$
(101)
(B) APPLICATION TO GMR.
Assuming that the polarization axis is the same for all sub-layers, we have $`l_{DW}=0`$ and the last term in (98) vanishes.
The chemical affinity obeys the diffusion equation:
$$\frac{^2A}{z^2}=\frac{1}{l_{sf}^2}A.$$
(102)
The chemical affinity and the total chemical potential are then
$$A(z)=ae^{\frac{z}{l_{sf}}}+be^{\frac{z}{l_{sf}}}$$
(103)
$$\mu _0(z)=d+cz+\beta A(z)$$
(104)
where $`a`$, $`b`$, $`c`$ and $`d`$ are constants. The electric field E(z) is defined by $`eE(z)\frac{\mu _0}{z}`$ so that $`c=eE(\mathrm{})=eJ_N/\sigma _0`$. Under the condition of continuity of the currents of the two spin channels at the interface (no surface scattering) : $`J_\pm (0^{})=J_\pm (0^+)`$ we have $`a=\frac{el_{sf}\beta }{\sigma _0(1\beta ^2)}J_N`$. The spin polarized current on the left side of a single interface (b=0) is deduced:
$$J_p(z)=\frac{\sigma _0}{e}\left(\frac{A}{z}\beta \frac{\mu _0}{z}\right)=J_N\beta \left(e^{\frac{|z|}{l_{sf}}}1\right)$$
(105)
The electric field $`\frac{\mu _0}{ez}E(z)`$ is:
$$E(z)=\frac{eJ_N}{\sigma _0}\left(1+\frac{\beta ^2}{1\beta ^2}e^{\frac{|z|}{l_{sf}}}\right)$$
(106)
and the supplementary potential due to the spin-polarized current is
$$\mathrm{\Delta }V=_{\mathrm{}}^+\mathrm{}\left(E(z)\frac{eJ_N}{\sigma _0}\right)𝑑z=\mathrm{\hspace{0.17em}2}\frac{\beta ^2}{1\beta ^2}l_{sf}\frac{eJ_N}{\sigma _0}$$
(107)
from which the GMR resistance (60) is deduced.
(C) APPLICATION TO SPIN TRANSFER
In the case of an interface composed by an incident current of conductivity asymmetry $`\beta `$, polarized in the direction $`\stackrel{}{e}_p`$, entering in a ferromagnetic layer polarized in the direction $`\stackrel{}{u}`$ with conductivity asymmetry $`\beta `$, the Equation (103) and (104) still hold in the left and right hand sides of the interface. However, the change of the polarization axis at the interface leads to modify the continuity equation of the current of the two spin channels $`J_\pm (0^{})=J_\pm (0^+)\gamma J_p`$. The integration constant now reads
$$a=\frac{el_{sf}\beta (1\gamma )}{\sigma _0(1\beta ^2)}J_N$$
(108)
and the expression of the polarized current is
$$J_p(z)=J_N\beta (1\gamma )\left(e^{\frac{|z|}{l_{sf}}}1\right)$$
(109)
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# Kinks from Dynamical Systems: Domain Walls in a Deformed 𝑂(𝑁) Linear Sigma Model
## 1 Introduction
We divide this Introduction into three parts: A. A brief history and the “state of the art”. B. New developments and results to be presented in this work . C. Scenarios of possible physical applications.
A.
Kinks are solitary (non-dispersive) waves arising in several one-dimensional physical systems. Here, we shall focus on the relativistic theory of $`N`$-interacting scalar fields built on a space-time that is the (1+1)-dimensional Minkowski space $`𝐑^{1,1}`$. In this context, kinks are finite energy solutions of the Euler-Lagrange equations, such that the time-dependence is dictated by the Lorentz invariance: $`\stackrel{}{\varphi }_K(x,t)=\stackrel{}{q}_K\left(\frac{xvt}{\sqrt{1v^2}}\right)`$. Thus, the search for kinks leads to the solving of a system of $`N`$-coupled non-linear ordinary differential equations and therefore becomes a very interesting problem in Mathematical Physics.
The study of topological defects began as an area of research in field theory by the mid-seventies; see . It was immediately recognized that defects of the kink type are in one-to-one correspondence with the separatrix trajectories between the bounded and unbounded motion of a mechanical system, for which the motion equations precisely form the non-linear system of differential equations mentioned above. The equivalent mechanical system is also Lagrangian and thus automatically integrable if $`N=1`$. For $`N2`$, complete integrability is generically non-guaranteed and the equivalence to a mechanical system is not useful. This circumstance has been emphasized by Rajaraman; see pp. 23-24, and partially circumvented by himself: the trial orbit method allows one to guesstimate particular types of kink trajectories.
There are, nevertheless, theories with $`N=2`$-coupled scalar fields such that the equivalent dynamical system is completely integrable. The prototype of this kind of system is the MSTB model: in this was proposed in the context of the search for non-topological solitons with stability provided by a $`U(1)`$ internal symmetry. In Reference , the model was considered as a classical continuum approximation to a 1D crystal with a two-component order parameter and it was shown that the search for kinks in this system requires that a completely integrable dynamical system be addressed. Ito, in a seminal paper , showed the Hamilton-Jacobi separability of the system of non-linear differential equations. He found all the kink trajectories and explained a very peculiar kink energy sum rule.
Very rich manifolds of kinks were discovered in two $`N=2`$ field theoretical models, close relatives to the MSTB system, in a recent research performed by the authors of the present work . The investigation of kink properties in these models requires the analysis of the separatrix trajectories in two related dynamical systems which are type I and III respectively in the classification of Liouville bidimensional $`(N=2)`$ completely integrable systems, see .
In fact, on choosing between the four types of Liouville dynamical systems those that meet appropriate critical point structure, one builds an enormous list of related $`N=2`$ field theoretical models exhibiting manifolds of kinks of growing complexity; see . The rôle of these models can be understood by noticing that the MSTB system is a deformation of the $`O(2)`$-linear sigma model. Instead of spontaneous symmetry breaking of $`O(2)`$ by a degenerated $`S^1`$ vacuum manifold, the $`O(2)`$ symmetry group is explicitly broken to $`𝐙_2\times 𝐙_2`$ by a mass term; only invariance under $`\varphi _a(1)^{\delta _{ab}}\varphi _a`$, for $`b=1,2`$, survives. From the point of view of quantum field theory, this deformation is very natural because in (1+1)-dimensions infrared divergences forbid the existence of Goldstone bosons, according to a theorem of Coleman . Even if it is absent in the classical action, a mass term will be generated by quantum corrections.
We interpret this as follows: in the parameter space of the $`N=2`$ relativistic scalar field theories invariant under the $`𝐙_2\times 𝐙_2`$ with generators mentioned above, and potential energy of the form
$$U(\stackrel{}{\varphi })=\frac{1}{2}\left(\alpha _1\varphi _1^2+\alpha _2\varphi _2^2+\frac{\beta _1}{2}\varphi _1^4+\beta _{12}\varphi _1^2\varphi _2^2+\frac{\beta _2}{2}\varphi _2^4\right)+C$$
there are at least two distinguished points. There is a choice of coupling constants such that there is explicit $`O(2)`$ symmetry, which is spontaneously broken. This is the linear $`O(2)`$-sigma model. The other interesting point is the MSTB model where the explicit $`𝐙_2\times 𝐙_2`$ symmetry generated by $`\varphi _a(1)^{\delta _{ab}}\varphi _a`$, for $`b=1,2`$, breaks spontaneously to the $`𝐙_2`$ sub-group generated by $`\varphi _2\varphi _2`$. The key observation is that the renormalization group flow induced by quantum corrections in the parameter space avoids the $`O(2)`$-sigma system and instead leads to the MSTB model, which also offers a variety of kinks. All the other field theoretical models exhibiting an abundant supply of kinks also correspond to deformations of $`O(2)`$-symmetric systems with potential energies that depend on higher powers of $`\varphi _1`$ and $`\varphi _2`$, , .
There are strong analogies with the Zamolodchikov $`c`$-theorem, : deformations in the space of (1+1)-dimensional field theories leading from conformal to integrable systems are the most interesting ones. We meet an analogous finite dimensional situation: replace the (infinite dimensional) conformal group by the $`O(2)`$ group and integrability of one system with infinite degrees of freedom by integrability of a bidimensional mechanical system.
B.
This paper is devoted to investigating the kink solitary waves of the deformation of the linear $`O(N)`$-sigma model that generalize the MSTB system to the case of $`N`$-interacting scalar fields. Non-linear waves in relativistic field theories with $`N3`$ scalar fields were sketchily described for the first time in Reference . In this work, we offer a detailed analysis of this issue. The following points merit emphasis:
* a) The dynamical system that encodes the solitary waves of the model as separatrix trajectories has $`N`$ first integrals in involution and hence is completely integrable. Passing from Cartesian to Jacobi elliptic coordinates in the “internal” space, $`𝐑^N`$, the dynamical system becomes Hamilton-Jacobi separable. All the kink trajectories, and hence all the solitary waves, are then found by a special choice of the separation constants.
* b) Deep insight into the structure of the kink manifold is gained by focusing on the $`N=3`$ case. There are three kinds of kinks: 1. A two-parameter family of topological kinks with three non-null components that are “generic”, i.e. they are not fixed under the action of the $`𝐙_2\times 𝐙_2\times 𝐙_2`$ group generated by $`\varphi _a(1)^{\delta _{ab}}\varphi _a`$, for $`b=1,2,3`$. 2. Four one-parameter families of “enveloping” non-topological kinks, also with three non-null components. The four families are related through the action of one $`𝐙_2\times 𝐙_2`$ sub-group and, together, form the envelop of the separatrix trajectories. 3. All the solitary waves of the $`N=2`$ MSTB model appear “embedded” twice; once in each plane containing the two ground states. Different $`𝐙_2`$ sub-groups leave these embedded kinks invariant.
* c) The structure of the kink manifold of the $`O(N)`$ system with both explicit and spontaneous symmetry breaking repeats the patterns shown in the $`N=2`$ (MSTB) model and its generalization for $`N=3`$. There are also generic, enveloping and embedded kinks, although when $`N`$ increases the complexity of the kink manifold also increases. For instance, the $`N1`$ kink manifold is embedded $`N1`$ times in the manifold of kinks of the deformed linear $`O(N)`$-sigma model.
* d) In a remarkable system obtained from the generalized MSTB model by also allowing asymmetries in the quartic terms of the potential, only the embedded and enveloping topological kinks living on singular edges survive as solitary wave solutions. In this system, proposed in Reference for the $`N=2`$ case, the energy of all the above topological kinks is exactly the same. Together with vacuum degeneration, there is therefore kink degeneration, a phenomenon that deserves further analysis.
C.
Solitary waves of the kind that we are to describe play an important rôle in condensed matter physics. Phase transitions characterized through order parameters of the vector type are understood in terms of the linear (or non linear) $`O(N)`$-sigma model. The order parameter is organized in the fundamental representation of $`O(N)`$ and the system becomes non-linear when this $`N`$-vector is forced to take its values in the coset space $`=O(N)/O(N1)`$. In $`(1+1)`$-dimensional space-time, kinks are accompanied by the fermion fractionization phenomenon ; this describes the continuous approximation to the bizarre behaviour of certain one-dimensional polymers such as poly-acetilene. When the spatial dimension is 3, as in the real world, kinks become domain walls which are thus related to theories involving spontaneous breaking of discrete symmetries. This happens in the hot Big Bang cosmology, where domain wall topological defects can be formed in a phase transition occurring in the expansion of the very early Universe; see . More recently, domain walls have been characterized as BPS states of SUSY gluodynamics and the Wess-Zumino model, . In all these cases there are sets of scalar fields, as in our system, that presents a variety of domain walls with different characteristics when seen from a 3-dimensional perspective.
In quantum field theory, the linear $`O(N)`$-sigma model describes systems with spontaneous symmetry breakdown to an $`O(N1)`$ sub-group and $`N1`$ Goldstone bosons in the particle spectrum. At the beginning of the sixties Gell-Mann and Lèvy analyzed low energy hadronic phenomenology by introducing an effective Lagrangian field theory of this type . Besides becoming the central element of current algebra, linear sigma models also enter fundamental physics in the Higgs sector of gauge theories for elementary particle physics, see report for a comprehensive review (of the perturbative sub-sector). For instance, the linear $`O(4)`$-sigma model corresponds to the Higgs sector of the electro-weak theory, while the $`O(24)\times O(5)`$ case provides the bosonic sector of the $`SU(5)`$ Grand Unified Theory.
Either considered on their own or forming part of Gauge theories, there are reasons to discuss deformations of the linear sigma model. In the phenomenological approach, pions are identified with the Goldstone bosons of the model; a deformation is then necessary to convert these massless excitations in pseudo-Goldstone particles, accounting for the pions light mass. Gauge theories are today found in the low energy limit of (fundamental) string theory. Even though deformations in the bosonic sector of gauge theories produced by small mass terms spoil renormalizability, the low energy features remain (almost) untouched and it is (almost) legitimate to trust them.
Here, we shall search for domain walls when these mild deformations are performed in the linear $`O(N)`$-sigma model. It is precisely in this kind of model where the cosmological problem of wall domination is avoided . Moreover, the system has a rich manifold of topological and non-topological solitons, allowing for topological defects with “internal” structure and leading to the existence of defects inside defects, a situation that generalizes a proposal of Morris .
The organization of the paper is as follows: In Section §2 we discuss the particle spectrum of the deformed linear $`O(N)`$-sigma model as well as the manifold of the solitary wave solutions of the system. Section §3 is devoted to the $`N=3`$ case, which is described in full detail. We describe the situation of the generalized MSTB model for any $`N`$ in Section §4 and briefly discuss the phenomenon of kink degeneration in the Bazeia system. Finally, some conclusions are drawn and some new prospects opened in Section §5. An appendix on elliptic Jacobi coordinates is also offered.
## 2 Kinks in the deformed linear $`O(N)`$-sigma model
In a generic sense we understand “kinks” as the solitary waves of a relativistic (1+1)-dimensional scalar field theory. We shall stick to the standard definition of solitary waves; see and :
A solitary wave is a non-singular solution of the non-linear coupled field equations of finite energy such that their energy density has a space-time dependence of the form:
$$\epsilon (\stackrel{}{x},t)=\epsilon (\stackrel{}{x}\stackrel{}{v}t)$$
where $`\stackrel{}{v}`$ is some velocity vector.
Given one $`N`$-component scalar field, which is a map from the $`𝐑^{1,1}`$ Minkowski space-time to $`𝐑^N`$, $`\stackrel{}{\chi }(x,t)(\chi _1(x,t),`$ $`\chi _2(x,t),\mathrm{},\chi _N(x,t))`$, the dynamics of the system is governed by the action:
$$S=d^2y\left\{\frac{1}{2}_\mu \stackrel{}{\chi }^\mu \stackrel{}{\chi }\overline{V}(\stackrel{}{\chi })\right\}$$
Here, $`\mu =0,1`$ are indices in the space-time and we shall use $`a=1,2,\mathrm{},N`$ to label components of the field in the “internal” $`𝐑^N`$ space in such a way that $`\stackrel{}{\chi }\stackrel{}{\chi }={\displaystyle \underset{a=1}{\overset{N}{}}}\chi _a\chi _a`$. In $`𝐑^{1,1}`$ we choose the metric as $`g=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)`$ and the Einstein convention will be used throughout the paper only for the indices in $`𝐑^{1,1}`$. The potential energy density is:
$$\overline{V}(\chi _1,\mathrm{},\chi _N)=\frac{\lambda ^2}{4}\left(\stackrel{}{\chi }\stackrel{}{\chi }\frac{m^2}{\lambda ^2}\right)^2+\underset{a=1}{\overset{N}{}}\frac{\beta _a^2}{4}\chi _a^2$$
where $`\lambda ,m`$ and $`\beta _a`$ are coupling constants of inverse length. The linear $`O(N)`$-sigma model corresponds to the case $`\beta _a=0`$, $`a`$, which exhibits maximum $`O(N)`$ symmetry. We shall focus on the deformation of this system, which is maximally non-isotropic in the harmonic terms, i.e. $`\beta _a\beta _b`$, $`ab`$. Somehow, the deformation is natural from a quantum field theoretical vantage point as we shall explain later and, moreover, we shall stick to the range $`\beta _a^2<m^2`$, $`a`$, in the parameter space because in this regime the structure of the kink manifold is richer.
Introducing non-dimensional variables $`\chi \frac{m}{\lambda }\varphi `$, $`y_\mu \frac{\sqrt{2}}{m}x_\mu `$ and $`\frac{\beta _a^2}{m^2}\sigma _a^2`$, we find our expression for the action to be:
$`S={\displaystyle \frac{m^2}{\lambda ^2}}{\displaystyle d^2x\left\{\frac{1}{2}_\mu \stackrel{}{\varphi }^\mu \stackrel{}{\varphi }V(\varphi _1,\mathrm{},\varphi _N)\right\}}`$
(1)
$`V(\varphi _1,\mathrm{},\varphi _N)={\displaystyle \frac{1}{2}}\left(\stackrel{}{\varphi }\stackrel{}{\varphi }1\right)^2+{\displaystyle \underset{a=1}{\overset{N}{}}}{\displaystyle \frac{1}{2}}\sigma _a^2\varphi _a^2`$
### 2.1 Configuration space and particle spectrum
The Cauchy problem for the field equations
$$\mathrm{}\varphi _a=\frac{V}{\varphi _a},a=1,2,\mathrm{},N$$
(2)
is fixed by choosing a “point” $`\stackrel{}{\varphi }(x,t_0)\mathrm{Maps}(𝐑,𝐑^N)`$ in the configuration space $`𝒞`$ , and its “tangent”, $`\dot{\stackrel{}{\varphi }}(x,t_0)T_\stackrel{}{\varphi }\mathrm{Maps}(𝐑,𝐑^N)`$, as initial conditions to solve the system (2) of non-linear PDE .
The configuration space itself is isomorphic to the space of finite energy static configurations; if $`\stackrel{}{\varphi }(x,t)=\stackrel{}{q}(x)`$, $`𝒞`$ is the set of continuous maps $`\stackrel{}{q}:𝐑𝐑^N`$ ($`\stackrel{}{q}(q_1,\mathrm{},q_N)`$) such that the (static) energy is finite:
$$E=\frac{m^3}{\lambda ^2\sqrt{2}}𝑑x\left\{\frac{1}{2}\frac{d\stackrel{}{q}}{dx}\frac{d\stackrel{}{q}}{dx}+V(\stackrel{}{q})\right\}<+\mathrm{};$$
(3)
thus, $`𝒞=\left\{\stackrel{}{q}(x)/E[\stackrel{}{q}]<+\mathrm{}\right\}`$. $`\stackrel{}{q}(x)𝒞`$ only if $`\stackrel{}{q}`$ satisfies the asymptotic conditions:
$$\underset{x\pm \mathrm{}}{lim}\frac{dq_a}{dx}=0,\underset{x\pm \mathrm{}}{lim}q_a(x)=v_a,a=1,\mathrm{},N$$
(4)
where $`\stackrel{}{v}(v_1,\mathrm{},v_N)`$ is a constant vector that belongs to the set $``$ of vectors annihilating $`V`$. We assume, without loss of generality, the following ordering in the space of parameters: $`\sigma _1=0<\sigma _2<\mathrm{}<\sigma _N<1`$. $``$ is thus formed by two vectors
$$=\left\{\stackrel{}{v}^\pm =(\pm 1,0,\mathrm{},0)\right\}$$
(5)
which are the absolute minima of $`V`$.
We refer to $``$ as the vacuum manifold because in the quantum version of the theory points in $``$ are the expectation values of the quantum field operators $`\widehat{\varphi }_a`$ at the ground states (“vacua”) of the system. The vacuum degeneration - i.e.the existence of more than one vector in $``$ \- is related to the breaking of symmetry. Besides two-dimensional Poincaré invariance, there is a “internal” symmetry with respect to the discrete group $`G=𝐙_2\times \stackrel{N}{\mathrm{}}\times 𝐙_2=𝐙_2^{\times N}`$ generated by $`\varphi _a(1)^{\delta _{ab}}\varphi _a`$, for $`b=1,2,\mathrm{},N`$, $`a=1,\mathrm{},N`$. The vacuum manifold is the orbit of one element by the group action
$$=G/H_{\stackrel{}{v}^\pm }=𝐙_2,H_{\stackrel{}{v}^\pm }\stackrel{}{v}^\pm =\stackrel{}{v}^\pm $$
$`H_{\stackrel{}{v}^\pm }=𝐙_2\times \stackrel{N1}{\mathrm{}}\times 𝐙_2`$ is the little group of the vacuum $`\stackrel{}{v}^\pm `$. The generators of $`H_{\stackrel{}{v}^\pm }`$ are the transformations $`\varphi _a(1)^{\delta _{ab}}\varphi _a`$, for $`b=1,2,\mathrm{},N`$, and $`a=2,3,\mathrm{},N`$, so that $`H_{\stackrel{}{v}^\pm }`$ survives as a symmetry group when quantizing around $`\stackrel{}{v}^\pm `$. We can understand the internal parity group $`G`$ as the discrete “gauge” symmetry: in (1+1)-dimensions no dynamical degrees of freedom related to gauge potentials appear.
Vectors in $``$ are critical points of $`V`$ satisfying $`\frac{V}{\varphi _a}|_{\stackrel{}{\varphi }=\stackrel{}{v}^\pm }=0`$ and therefore constant solutions of the field equations (2). The plane wave expansion around $`\stackrel{}{\varphi }^\pm (x,t)=\stackrel{}{v}^\pm `$
$$\varphi _{v_a^\pm }(x,t)=v_a^\pm +\underset{k}{}A_a^\pm (k)e^{i\omega tikx}$$
is a solution of (2) if the dispersion relation
$$\delta _{ab}\omega ^2=\delta _{ab}k^2+M_{ab}^2(\stackrel{}{v}^\pm ),M_{ab}^2=\frac{^2V}{\varphi _a\varphi _b}(\stackrel{}{v}^\pm )$$
(6)
holds.
In the quantum theory, these plane waves become the fundamental quanta with mass matrix $`M_{ab}^2(\stackrel{}{v}^\pm )`$ and one reads the particle spectrum at a chosen critical point of $`V`$ from (6). Because $`\stackrel{}{v}^\pm `$ are minima of $`V`$ there are no negative eigenvalues of $`M_{ab}^2(\stackrel{}{v}^\pm )`$ and the dependence on time of the plane waves around $`\stackrel{}{v}^\pm `$ is bounded: $`e^{i\omega t}`$. The choice of $`\stackrel{}{v}^\pm `$ as the starting point of the quantization procedure “spontaneously” breaks the symmetry $`G=𝐙_2^{\times N}`$ of the action to $`H_{\stackrel{}{v}^\pm }=𝐙_2^{\times (N1)}`$, which is the remaining one that survives in the particle spectrum.
In our model, we read the particle spectrum from:
$$M^2(\stackrel{}{v}^\pm )=\frac{m^2}{2}\left(\begin{array}{cccc}4& 0& \mathrm{}& 0\\ 0& \sigma _2^2& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& \mathrm{}& \sigma _N^2\end{array}\right)$$
(7)
Considering this system as a physical description of the continuum approximation to a one-dimensional crystal with an $`N`$-component order parameter, the particle spectrum describes a single phase with $`N`$ phonon branches. We see explicitly how the symmetry group $`G=𝐙_2^{\times N}`$ is “broken” by the choice of the $`\stackrel{}{v}^\pm `$ vacuum to the $`H_{\stackrel{}{v}^\pm }=𝐙_2^{\times (N1)}`$ sub-group: the $`N`$ phonon branches have different masses or “energy gaps”. From the point of view of particle physics we can say that there are no tachyons; only a pseudo-Goldstone particle becomes a Goldstone boson if the corresponding $`\sigma _a`$ goes to zero.
It is interesting to see the model as a member of the family characterized by the potential energy densities:
$$V=\frac{1}{2}\left(\underset{a,b=1}{\overset{N}{}}\alpha _{ab}\varphi _a\varphi _b\gamma ^2\right)^2+\underset{a,b=1}{\overset{N}{}}\frac{\sigma _{ab}}{2}\varphi _a\varphi _b$$
where $`\alpha _{ab},\sigma _{ab}`$ and $`\gamma ^2`$ are “bare” non-dimensional parameters. Ultraviolet divergences are controlled by normal ordering in the quantum theory, but the need arises to introduce a renormalization ‘point’ $`\mu ^2`$, and the dependence of the renormalized parameters on $`\mu ^2`$ is determined by the renormalization group equation. One special solution, a specific renormalization group flow, might lead to the “point”:
$$\alpha _{ab}^R(\mu ^2)=\delta _{ab};\sigma _{ab}^R(\mu ^2)=0;\gamma ^R(\mu ^2)=1$$
in the space of quantum field theory models in the family. This point is the linear $`O(N)`$-sigma model which has $`G=O(N)`$ as the (continuous) symmetry group. The vacuum orbit is, however, $`=O(N)/O(N1)=S^{N1}`$, the $`(N1)`$-dimensional sphere, and thus there is no unbroken symmetry left: there are $`N1`$ massless particles. If the only modification of the renormalized parameters is to allow for non-zero values of $`\sigma _{ab}^R(\mu ^2),a=b=r+1,\mathrm{},N`$ there are still $`r1`$ Goldstone bosons.
Coleman established that in $`(1+1)`$-dimensions the infrared asymptotics of the two-point Green functions of a quantum scalar field forbids poles at $`\omega ^2=k^2`$; there are no Goldstone bosons in $`(1+1)`$-dimensions. It is thus impossible to reach the $`O(N)`$-sigma model or its deformation with the $`O(r)`$ symmetry spontaneously broken to $`O(r1)`$ in the renormalization group flow. The closest admissible points are the models characterized by:
$$\alpha _{ab}^R(\mu ^2)=\delta _{ab},\gamma ^R(\mu ^2)=1,\sigma _{ab}^R=0,ab$$
$$\sigma _{11}^R(\mu ^2)=0<\sigma _{22}^R(\mu ^2)=\sigma _2^2\mathrm{}\sigma _{NN}^R(\mu ^2)=\sigma _N^2<1$$
In this paper we shall focus on the case of maximal explicit symmetry breaking; i.e. when strict inequalities in the parameter space occur. Nevertheless, we shall comment on the allowed situation characterized by
$$\sigma _1^2=0<\sigma _2^2=\mathrm{}=\sigma _{r_1}^2<\sigma _{r_1+1}^2=\mathrm{}=\sigma _{r_2}^2<\mathrm{}<\sigma _{r_k+1}^2=\mathrm{}=\sigma _N^2<1$$
when there is degeneration in the spectrum but no Goldstone bosons. Note that the generators of the $`O(r_11)\times O(r_2r_1)\times \mathrm{}\times O(Nr_k)`$ symmetry sub-group are in the little group of the vacuum. The symmetry group is $`G=𝐙_2\times O(r_11)\times O(r_2r_1)\times \mathrm{}\times O(Nr_k)`$, $`H_{\stackrel{}{v}^\pm }=O(r_11)\times O(r_2r_1)\times \mathrm{}\times O(Nr_k)`$ and the vacuum orbit is $`=G/H_{\stackrel{}{v}^\pm }=𝐙_2`$.
### 2.2 Configuration space topology: kinks and dynamical systems
The configuration space of the model is the union of topologically disconnected sectors: $`𝒞={\displaystyle \underset{\alpha ,\beta =1}{\overset{2}{}}}𝒞^{\alpha \beta }`$; thus, $`\pi _0(𝒞)=𝐙_2\times 𝐙_2`$ and $`|\pi _0(𝒞)|=4`$ are respectively the zeroth-order homotopy group of $`𝒞`$ and its order. This comes from the asymptotic conditions (4) and the continuity of the time evolution . There are topological charges defined for each configuration in $`𝒞`$ as:
$$Q_a^T=\frac{1}{2}_{\mathrm{}}^{\mathrm{}}𝑑x\frac{d\varphi _a}{dx}=\frac{1}{2}\left(\varphi _a(+\mathrm{},t)\varphi _a(\mathrm{},t)\right)$$
It should be noted that $`Q_a^T`$ is independent of $`t`$, $`a`$, and in our system equal to zero if $`a2`$. Therefore the four sectors $`𝒞^{\alpha \beta }`$ are labelled by the values $`\alpha ,\beta `$ of the fields at infinity compatible with finite energy and $`Q_1^T`$ determines the homotopy class in $`\pi _0(𝒞)=𝐙_2\times 𝐙_2`$.
The critical points of $`E`$ are time-independent finite-energy solutions of the field equations. If they are not spatially homogeneous, the critical points correspond to solitary waves that are therefore related to the topological structure of $`𝒞^{\alpha \beta }`$. Besides complying with (4), solitary waves satisfy the system of ordinary differential equations:
$$\frac{d^2q_a}{dx^2}=\frac{V}{q_a}$$
(8)
Recall that $`\varphi _a(x,t)=q_a(x)`$. Solving the system (8) is tantamount to finding the solutions of the Lagrangian dynamical system in which $`x=\tau `$ plays the rôle of time, the “particle” position is determined by $`q_a(\tau )`$, and the potential energy of the particle is $`U(\stackrel{}{q})=V(\stackrel{}{q})`$. From this perspective the static field energy $`E`$ is seen as the particle action:
$$E=J=𝑑\tau \left\{\frac{1}{2}\frac{d\stackrel{}{q}}{d\tau }\frac{d\stackrel{}{q}}{d\tau }U(\stackrel{}{q})\right\}$$
(9)
Trajectories that behaves asymptotically in the $`\tau `$-time as ruled by (4) have a finite action, $`J`$, in the mechanical problem and are in one-to-one correspondence with solitary waves/kinks that have energy $`E=J`$ in the field theoretical system.
The mechanical analogy is very helpful when one is dealing with a real scalar field theory because, then, a first integral is all that we need to find all the solutions. Vector scalar fields of $`N`$ components lead to $`N`$-dimensional dynamical systems which are seldom solvable. Magyari and Thomas realized that the two-dimensional dynamical system arising in connection with the MSTB model is a completely integrable one in the Liouville sense; there are two first integrals in involution. Moreover, Ito has shown that the mechanical system is Hamilton-Jacobi separable, finding all the trajectories and hence all the kinks of the MSTB model. In a recent publication , we have developed this procedure for two $`N=2`$ models with interesting features: the first system is a deformation of the (1+1)-dimensional scalar field theory, where the potential energy density is the Chern-Simons-Higgs potential arising in self-dual planar gauge theories. The second one is a deformation of the linear $`O(2)`$-sigma model, which is different from the MSTB model.
To extend this method of finding kinks to the linear $`O(N)`$-sigma model, $`N3`$, deformed in such a way that the $`O(N)`$ symmetry is explicitly broken to $`G=𝐙_2^{\times N}`$, we start from the “particle” action:
$$J=𝑑\tau \left\{\frac{1}{2}\frac{d\stackrel{}{q}}{d\tau }\frac{d\stackrel{}{q}}{d\tau }+\frac{1}{2}\left(\stackrel{}{q}\stackrel{}{q}1\right)^2+\frac{1}{2}\underset{a=1}{\overset{N}{}}\sigma _a^2q_a^2\right\}=𝑑\tau (\stackrel{}{q},\dot{\stackrel{}{q}})$$
The particle motion equations are:
$$\frac{d^2q_a}{d\tau ^2}=2q_a(\stackrel{}{q}\stackrel{}{q}1)+\sigma _a^2q_a,a=1,\mathrm{},N$$
(10)
which are mathematically identical to the field equations for static configurations. Finite action trajectories, kinks in the field theory, should also satisfy the asymptotic conditions:
$$\underset{\tau \pm \mathrm{}}{lim}\frac{dq_a}{d\tau }=0,\underset{\tau \pm \mathrm{}}{lim}q_a(\tau )=\pm \delta _{a1}$$
(11)
We shall use the Hamiltonian formalism to integrate the mechanical system. The canonical momenta $`p_a(\tau )=\frac{}{\dot{q}_a}=\frac{dq_a}{d\tau }(\tau )`$, together with the positions $`q_a(\tau )`$, form a system of local coordinates in phase space. We should bear in mind that $`p_a(\tau )=\frac{d\varphi _a}{dx}`$ when going back to the field theory. The mechanical Hamiltonian
$$I_1=\frac{1}{2}\stackrel{}{p}\stackrel{}{p}\frac{1}{2}\left(\stackrel{}{q}\stackrel{}{q}1\right)^2\underset{a=1}{\overset{N}{}}\frac{1}{2}\sigma _a^2q_a^2$$
(12)
leads to the system of canonical equations
$$\frac{dq_a}{d\tau }=\{I_1,q_a\},\frac{dp_a}{d\tau }=\{I_1,p_a\}$$
equivalent to (10). Given any two functions $`F(\stackrel{}{q},\stackrel{}{p})`$, $`G(\stackrel{}{q},\stackrel{}{p})`$ in phase space, the Poisson bracket is defined in the usual way:
$$\{F,G\}=\underset{a=1}{\overset{N}{}}\left(\frac{F}{q_a}\frac{G}{p_a}\frac{F}{p_a}\frac{G}{p_a}\right)$$
Obviously $`\frac{dI_1}{d\tau }=0`$, but our mechanical system is full of other invariants. In fact, as early as 1919 Garnier solved the motion equations and described periodic trajectories in terms of Theta functions: the kink trajectories of finite “action” correspond to a limiting case and are the separatrices between the periodic trajectories and unbounded motion. More recently Grosse, and other authors have shown that the functions:
$$K_a=\underset{b=1,ba}{\overset{N}{}}\frac{1}{\sigma _b^2\sigma _a^2}l_{ab}^2+p_a^2+(2\sigma _a^2)q_a^2q_a^2\underset{b=1}{\overset{N}{}}q_b^2$$
(13)
$$l_{ab}=p_aq_bp_bq_a$$
are first integrals in involution:
$$\{I_1,K_a\}=0\{K_a,K_b\}=0$$
There is a set of $`N+1`$ invariants in involution: $`I_1,K_1,K_2,\mathrm{},K_N`$. The dynamical system is not superintegrable, however, because there are only $`N`$-independent invariants: $`K_1+K_2+\mathrm{}+K_N=2I_1+1`$. According to the Liouville theorem, the $`N`$-dimensional mechanical system is completely integrable and all trajectories can be found, at least in principle.
At this point we pause to explain the singular nature of the deformation of the linear $`O(N)`$-sigma model chosen from among many possibilities. The $`\frac{\sigma _a^2}{2}\varphi _a^2`$ terms explicitly break the $`O(N)`$-symmetry of the linear sigma model; the case $`\sigma _a^2=0`$, $`a=1,2,\mathrm{},N`$. In the mechanical system the $`O(N)`$ internal transformations become ordinary rotations. The angular momentum components, $`l_{ab}`$, conserved in the limit $`\sigma _a^2=0`$, $`a`$, are no longer ‘time’-independent if $`\sigma _a^20`$. There are, however, $`N`$ invariants $`K_a`$, which in the limit $`\sigma _a=0`$, $`a`$, are given in terms of the $`O(N)`$-invariants: the $`r`$ Casimir invariants and the $`r`$ generators of the Cartan sub-algebra, where $`r=\frac{N}{2}`$ or $`\frac{N1}{2}`$ if $`N`$-even or -odd is the rank of the group. A warning: in the $`N=`$ odd case, the energy must be added to the other $`N1`$ invariants built from the Cartan sub-algebra and the Casimir invariants. For any $`N`$, the maximally asymmetric chosen deformation is special because it retains enough symmetry to solve the mechanical system. There is no Lie algebra associated with $`K_a`$ however; since the invariants are quadratic in $`q_a`$, $`p_a`$, the action of $`K_a`$ in the phase space, given by $`\{K_a,q_b\}`$ and $`\{K_a,p_b\}`$, is non-linear.
In (1+1)-dimensional field theory, the energy-momentum tensor:
$$T^{\mu \nu }=\frac{}{(_\mu \varphi _a)}^\nu \varphi _ag^{\mu \nu }$$
is divergenceless due to invariance under space-time translations. $`P^\mu =𝑑xT^{0\mu }`$ are thus conserved quantities whatever the values of $`\sigma _a^2`$. The $`O(N)`$ “isospin” currents however,
$$J_a^\mu =\underset{b,c=1}{\overset{N}{}}c_{abc}\varphi _b^\mu \varphi _c$$
are only divergenceless if $`\sigma _a=0`$, $`a`$. The $`c_{abc}`$ are the Lie $`O(N)`$ structure constants and the charges $`Q_a=𝑑xJ_a^0`$ are not conserved if there is no symmetry with respect to the transformation generated by them. For static configurations, we have
$$𝑑xT^{00}=E=J,T^{10}=T^{01}=0,T^{11}=I_1$$
$$J_a^0=0,J_a^1=\underset{b,c=1}{\overset{N}{}}c_{abc}l_{bc}$$
In terms of the ‘isospin’currents, the invariants $`K_a`$ can be written as:
$$K_a=\underset{b=1,ba}{\overset{N}{}}\frac{1}{\sigma _b^2\sigma _a^2}\left(\underset{c=1}{\overset{N}{}}c_{abc}J_c^1\right)\left(\underset{d=1}{\overset{N}{}}c_{abd}J_d^1\right)+\left(\frac{\varphi _a}{x}\right)^2+(2\sigma _a^2)\varphi _a^2\varphi _a^2\underset{b=1}{\overset{N}{}}\varphi _b^2$$
We expect that the time-evolution occurs in such a way that there is some equation of non-linear character
$$F(\frac{L_a}{t},\frac{K_a}{x})=0$$
between $`K_a`$ and
$$L_a=\underset{b=1,ba}{\overset{N}{}}\frac{1}{\sigma _b^2\sigma _a^2}\left(\underset{c=1}{\overset{N}{}}c_{abc}J_c^0\right)\left(\underset{d=1}{\overset{N}{}}c_{abd}J_d^0\right)+\left(\frac{\varphi _a}{t}\right)^2+(2\sigma _a^2)\varphi _a^2\varphi _a^2\underset{b=1}{\overset{N}{}}\varphi _b^2$$
which reduces to $`\frac{J_a^0}{t}=\frac{J_a^1}{x}`$ when $`\sigma _a=0`$, $`a`$. The situation is analogous to that occurring between conformal field theories and models with infinite-dimensional algebraic symmetry as in (1+1)-dimensional Toda field theories and Toda affine models . There are two differences: (1) the conformal group is infinite dimensional in (1+1)-dimensions. We have only one finite-dimensional group $`O(N)`$ and thus we can solve only the static limit of the field theory model. (2) Due to the non-linear character of the deformation of the $`O(N)`$ Lie generators, we do not even have a finite-dimensional Lie algebra.
### 2.3 The Hamilton-Jacobi equation and kink trajectories
The $`K_a`$ invariants defined in (13) are quadratic in the momenta, but not orthogonal (they contain terms in $`p_ap_b,ab`$). Therefore, the Stäckel theorem can not be applied to assure Hamilton-Jacobi separability. This problem is surpassed in our dynamical system with the choice of some suitable system of coordinates. The appropriate system is provided by elliptic Jacobi coordinates, with a choice of separation constants determined by the deformation parameters giving mass to the Goldstone bosons; $`\overline{\sigma }_a^2=1\sigma _a^2,a=1,2,\mathrm{},N`$. Thus we define:
$$q_a^2=\frac{{\displaystyle \underset{b=1}{\overset{N}{}}}(\overline{\sigma }_a^2\lambda _b)}{{\displaystyle \underset{b=1,ba}{\overset{N}{}}}(\overline{\sigma }_a^2\overline{\sigma }_b^2)}=\frac{\mathrm{\Lambda }(\overline{\sigma }_a^2)}{A^{}(\overline{\sigma }_a^2)}$$
(14)
ruling the change of coordinates from Cartesian, $`\stackrel{}{q}(q_1,\mathrm{},q_N)`$ to elliptic $`\stackrel{}{\lambda }(\lambda _1,\mathrm{},\lambda _N)`$. In the Appendix, it is explained how the elliptic variables are split:
$$\mathrm{}<\lambda _1<\overline{\sigma }_N^2<\lambda _2<\overline{\sigma }_{N1}^2<\mathrm{}<\overline{\sigma }_2^2<\lambda _N<1$$
(15)
Notice that formula (14) coincides with formula (71) in the Appendix if we change $`q_a`$ by $`q_{Na+1}`$ and choose $`r_{Na+1}=\overline{\sigma }_a^2`$.
Together with formula (14), this splitting means that the change of coordinates produces a map from a sub-space of $`𝐑^N`$, characterized as the set of points which are not invariants under the $`𝐙_2^{\times N}`$ group generated by $`q_a(1)^{\delta _{ab}}q_a,b=1,\mathrm{},N`$, to the interior of the infinite parallelepiped $`P_N(\mathrm{})`$ obtained by replacing the inequalities in (15) by equalities: $`\mathrm{}<\lambda _1\overline{\sigma }_N^2\mathrm{}\overline{\sigma }_2^2\lambda _N1`$. Notice that in this map $`2^N`$ regular points in $`𝐑^N`$ go to a single point in the interior of $`P_N(\mathrm{})`$; Singular points lie in the $`𝐑^m`$, $`m=0,1,\mathrm{},N1`$, sub-spaces that are invariant under the action of some non-trivial element of $`G=𝐙^{\times N}`$. These singular sub-spaces are mapped into the boundary of $`P_N(\mathrm{})`$.
The standard length of an interval in Euclidean space is expressed in elliptic coordinates in the form
$$ds^2=\underset{a=1}{\overset{N}{}}dq_adq_a=\underset{a=1}{\overset{N}{}}g_{aa}(\stackrel{}{\lambda })d\lambda _ad\lambda _a$$
because the metric $`g_{aa}(\stackrel{}{\lambda })`$, as derived in the Appendix, is:
$$g_{aa}(\stackrel{}{\lambda })=\frac{1}{4}\frac{f_a(\stackrel{}{\lambda })}{A(\lambda _a)},g_{ab}(\stackrel{}{\lambda })=0,ab$$
where $`A(\lambda _a)={\displaystyle \underset{b=1}{\overset{N}{}}}(\lambda _a\overline{\sigma }_b^2)`$, and <sup>1</sup><sup>1</sup>1The standard notation in the literature on elliptic Jacobi coordinates is $`\mathrm{\Lambda }^{}(\lambda _a)={\displaystyle \underset{\genfrac{}{}{0pt}{}{b=1}{ba}}{\overset{N}{}}}(\lambda _a\lambda _b)`$, see Appendix. We shall use $`f_a(\stackrel{}{\lambda })`$ in the main text instead of $`\mathrm{\Lambda }^{}(\lambda _a)`$, to stress the fact that this quantity depends on all the components of $`\stackrel{}{\lambda }`$
$$f_a(\stackrel{}{\lambda })=f_a(\lambda _1,\lambda _2,\mathrm{},\lambda _N)=\underset{\genfrac{}{}{0pt}{}{b=1}{ba}}{\overset{N}{}}(\lambda _a\lambda _b)$$
Therefore, the Lagrangian reads:
$`L`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{a=1}{\overset{N}{}}}\left({\displaystyle \frac{dq_a}{d\tau }}\right)^2U(q_1,\mathrm{},q_N)`$ (16)
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{a=1}{\overset{N}{}}}g_{aa}(\stackrel{}{\lambda })\left({\displaystyle \frac{d\lambda _a}{d\tau }}\right)^2U_\lambda (\lambda _1,\mathrm{},\lambda _N)`$
where the potential in elliptic coordinates is:
$$U_\lambda (\lambda _1,\mathrm{},\lambda _N)=\underset{a=1}{\overset{N}{}}\frac{1}{2}\frac{\lambda _a^{N+1}(\alpha 1)\lambda _a^N+(1\alpha +\beta )\lambda _a^{N1}}{f_a(\stackrel{}{\lambda })}$$
(17)
$$\alpha =\underset{a=1}{\overset{N}{}}\overline{\sigma }_a^2,\beta =\underset{a=1,a<b}{\overset{N}{}}\underset{b=2}{\overset{N}{}}\overline{\sigma }_a^2\overline{\sigma }_b^2$$
The computation of $`U_\lambda (\stackrel{}{\lambda })`$ is highly non-trivial and requires the use of formulas that follow the Jacobi Lemma, such as (76), (77), (78), etc.
The canonical momenta associated to the $`\lambda _a`$ variables are:
$$\pi _a=\underset{b=1}{\overset{N}{}}g_{ab}(\stackrel{}{\lambda })\frac{d\lambda _b}{d\tau }=g_{aa}(\stackrel{}{\lambda })\frac{d\lambda _a}{d\tau }$$
and, through the standard Legendre transformation, we write the Hamiltonian:
$$H=\frac{1}{2}\underset{a=1}{\overset{N}{}}\frac{4A(\lambda _a)}{f_a(\stackrel{}{\lambda })}\pi _a^2+U_\lambda (\stackrel{}{\lambda })$$
(18)
The key point is that $`H`$ can be written in Stäckel’s form:
$`H`$ $`=`$ $`{\displaystyle \underset{a=1}{\overset{N}{}}}{\displaystyle \frac{H_a}{f_a(\stackrel{}{\lambda })}}=`$
$`=`$ $`{\displaystyle \underset{a=1}{\overset{N}{}}}{\displaystyle \frac{\left[2A(\lambda _a)\pi _a^2\frac{1}{2}\left(\lambda _a^{N+1}(\alpha 1)\lambda _a^N+(1\alpha +\beta )\lambda _a^{N1}\right)\right]}{f_a(\stackrel{}{\lambda })}}`$
such that the Hamilton-Jacobi equation
$$\frac{𝒮}{\tau }+H(\frac{𝒮}{\lambda _1},\mathrm{},\frac{𝒮}{\lambda _N},\lambda _1,\mathrm{},\lambda _N)=0$$
(20)
is completely separable. We now prove this last statement.
Fixing $`H=I_1`$, the first integral of energy, and having in mind the expression (2.3) of $`H`$, we write the solution of (20) as:
$$𝒮=I_1\tau +\underset{a=1}{\overset{N}{}}S_a(\lambda _a).$$
(21)
Therefore, (20) reduces to :
$$I_1=H(\frac{dS_1}{d\lambda _1},\mathrm{},\frac{dS_N}{d\lambda _N},\lambda _1,\mathrm{},\lambda _N)=\underset{a=1}{\overset{N}{}}\frac{H_a}{f_a(\stackrel{}{\lambda })}$$
(22)
The Hamilton-Jacobi PDE equation (20) becomes equivalent to the system of non-coupled ordinary differential equations
$$H_a(\frac{dS_a}{d\lambda _a},\lambda _a)=\eta _1\lambda _a^{N1}+\eta _2\lambda _a^{N2}+\mathrm{}+\eta _{N1}\lambda _a+\eta _N$$
(23)
where
$$H_a(\frac{dS_a}{d\lambda _a},\lambda _a)=2A(\lambda _a)\left(\frac{dS_a}{d\lambda _a}\right)^2\frac{1}{2}\left(\lambda _a^{N+1}(\alpha 1)\lambda _a^N+(1\alpha +\beta )\lambda _a^{N1}\right),$$
(24)
due to the identity
$$I_1=I_1\underset{a=1}{\overset{N}{}}\frac{\lambda _a^{N1}}{f_a(\stackrel{}{\lambda })}+\eta _2\underset{a=1}{\overset{N}{}}\frac{\lambda _a^{N2}}{f_a(\stackrel{}{\lambda })}+\mathrm{}+\eta _N\underset{a=1}{\overset{N}{}}\frac{1}{f_a(\stackrel{}{\lambda })}$$
-observe that $`\eta _1=I_1`$\- which follows from the Jacobi Lemma and the subsequent relations (Appendix)
$$\underset{a=1}{\overset{N}{}}\frac{\lambda _a^{N1}}{f_a(\stackrel{}{\lambda })}=1,\underset{a=1}{\overset{N}{}}\frac{\lambda _a^{Ni}}{f_a(\stackrel{}{\lambda })}=0,i=2,\mathrm{},N$$
Alternatively, one could take a more direct approach to show formula (23). We start from (22), written explicitly as
$`{\displaystyle \frac{H_1(\lambda _1)}{(\lambda _1\lambda _2)(\lambda _1\lambda _3)\mathrm{}(\lambda _1\lambda _N)}}`$ $`+`$ $`{\displaystyle \frac{H_2(\lambda _2)}{(\lambda _2\lambda _1)(\lambda _2\lambda _3)\mathrm{}(\lambda _2\lambda _N)}}+\mathrm{}`$ (25)
$`+`$ $`{\displaystyle \frac{H_N(\lambda _N)}{(\lambda _N\lambda _1)(\lambda _N\lambda _2)\mathrm{}(\lambda _N\lambda _{N1})}}=I_1`$
Here, each $`H_a(\lambda _a)`$ is of the form (24) due to the ansatz (21). Multiplying (25) by $`(\lambda _1\lambda _2)`$ and setting $`\lambda _1=\lambda _2`$ one sees that $`H_1(\lambda )H_2(\lambda )`$, and hence by symmetry, all $`H_a(\lambda )`$, $`a=1,2,\mathrm{},N`$ are identical. It suffices therefore to find $`H_1(\lambda )`$. Multiplying (25) by $`f_1(\stackrel{}{\lambda })`$ one obtains
$$H_1(\lambda _1)+P_2(\stackrel{}{\lambda })H_2(\lambda _2)+\mathrm{}+P_N(\stackrel{}{\lambda })H_N(\lambda _N)=I_1(\lambda _1\lambda _2)(\lambda _1\lambda _3)\mathrm{}(\lambda _1\lambda _N),$$
(26)
where each $`P_a(\stackrel{}{\lambda })`$, $`a=2,3,\mathrm{},N`$, is a polynomial of degree $`N1`$ in $`\lambda _1`$. Differentiating (26) $`N`$ times with respect to $`\lambda _1`$ yields
$$\frac{d^NH_1(\lambda _1)}{d\lambda _1^N}=N!I_1$$
(27)
whence it follows that $`H_1(\lambda )`$ is a degree $`N`$ polynomial in $`\lambda `$ with leading coefficient $`I_1`$, in agreement with equation (23).
The set of separation constants $`\eta _i,i=1,2,\mathrm{},N`$ is another system of first integrals in involution. They can be expressed in the elliptic phase space $`T^{}P_N(\mathrm{})`$ as functions of $`\lambda _a`$ and $`\pi _a=\frac{dS_a}{d\lambda _a}`$ by solving the linear system of equations (23) in the unknown $`\eta _i`$, which is a Vandermonde system. The $`N1`$ roots of the polynomial
$$\lambda _a^{N1}+\frac{\eta _2}{I_1}\lambda _a^{N1}+\mathrm{}+\frac{\eta _N}{I_1}=(\lambda _aF_2)(\lambda _aF_3)\mathrm{}(\lambda _aF_N)$$
together with the energy $`F_1=I_1`$ form another system of invariants in involution. Both systems are related through the identities $`\eta _1=I_1`$ and:
$$\eta _a=(1)^aI_1\underset{i_1<i_2<\mathrm{}<i_a}{}F_{i_1}F_{i_2}\mathrm{}F_{i_a},i_\alpha =2,3,\mathrm{},N$$
Therefore, all the separation constants $`\eta _a`$ are proportional to the “particle” energy $`I_1`$.
Defining the polynomial $`B(\lambda _a)`$ in the form:
$$B(\lambda _a)=\lambda _a^{N+1}(\alpha 1)\lambda ^N+(1\alpha +\beta +2\eta _1)\lambda _a^{N1}+2\eta _2\lambda _a^{N2}+\mathrm{}+2\eta _N$$
the solution of the differential equation (23) is a quadrature:
$$S_a(\lambda _a)=\frac{1}{2}\mathrm{sign}\left(\frac{dS_a}{d\lambda _a}\right)\sqrt{\left|\frac{B(\lambda _a)}{A(\lambda _a)}\right|}𝑑\lambda _a$$
(28)
and the general solution of the Hamilton-Jacobi equation reads:
$$𝒮=\eta _1\tau +\underset{a=1}{\overset{N}{}}\frac{1}{2}\mathrm{sign}\left(\frac{dS_a}{d\lambda _a}\right)\sqrt{\left|\frac{B(\lambda _a)}{A(\lambda _a)}\right|}𝑑\lambda _a$$
(29)
The explicit integration of the quadratures in (28) requires the theory of Theta functions of genus depending on $`N`$. The action of the associated trajectories is infinite because they are either periodic or unbounded. The asymptotic conditions (4) that guarantee finite action to continuous trajectories satisfying them also require that the energy used by the particle in these trajectories should be zero. This is so because $`I_1|_{\tau =\pm \mathrm{}}=0`$ and, being an invariant of the evolution, $`I_1=0,\tau `$.
We recall that the trajectories of finite action in an evolution lasting an infinite time are the kinks of the parent field theory system: one just trades the finite action of the trajectory for finite energy of the non-linear wave. Therefore, the kinks are the trajectories obtained when all the separation constants in (23) are zero: $`\eta _a=0,a`$. These are the separatrices between bounded and unbounded motion and the integrals in (28) are easier to compute.
The explicit trajectories are also provided by the Hamilton-Jacobi principle, through the set of equations:
$$\gamma _a=\frac{𝒮}{\eta _a},a=1,2,\mathrm{},N$$
where the $`\gamma _a`$ are integration constants. In the hypersurface of the phase space determined by $`\eta _1=\mathrm{}=\eta _N=0`$, the first equation
$$\gamma _1=\tau +\underset{a=1}{\overset{N}{}}\frac{1}{2}\mathrm{sign}(\pi _a)\frac{\lambda _a^{N1}d\lambda _a}{\left|A(\lambda _a)\right|}\sqrt{\left|\frac{A(\lambda _a)}{\lambda _a^{N+1}(\alpha 1)\lambda _a^N+(1\alpha +\beta )\lambda _a^{N1}}\right|}$$
(30)
rules the time-dependence of the particle in its journey through the orbit. From the field theoretical point of view, it provides the kink form factor. The other $`N1`$ equations, $`i=2,3,\mathrm{},N`$,
$$\gamma _i=\underset{a=1}{\overset{N}{}}\frac{1}{2}\mathrm{sign}(\pi _a)\frac{\lambda _a^{Ni}d\lambda _a}{\left|A(\lambda _a)\right|}\sqrt{\left|\frac{A(\lambda _a)}{\lambda _a^{N+1}(\alpha 1)\lambda _a^N+(1\alpha +\beta )\lambda _a^{N1}}\right|}$$
(31)
determine the orbit in $`P_N(\mathrm{})`$, the intersection of $`N1`$ hypersurfaces in the configuration space. Therefore, there is a $`N1`$-dimensional family of kinks parametrized by the finite values of $`\gamma _i`$.
Although (30) and (31) identify all the separatrix trajectories of the mechanical system and henceforth all the kink solutions of the deformed linear $`O(N)`$-sigma model, an explicit description of such solitary waves is difficult for two reasons: (1). (30) and (31) form a system of transcendent equations of impossible analytical resolution. (2). Even if it were possible, expressing back the solution in Cartesian coordinates through (71) for $`N3`$ is another impossible task by analytical means.
## 3 $`N`$=3
To gain insight into the nature of the different kinks of the model, in this Section we shall address in full detail the $`N=3`$ case. We shall deal with a (1+1)-dimensional field theory including three scalar fields which transform according to a vector representation of the $`O(3)`$ group. The structure of the solitary wave solutions of the $`N=3`$ system is extremely rich from different points of view and shows the behavioural pattern of the general case with $`N`$-component fields.
### 3.1 The general solution of the Hamilton-Jacobi equation
The Hamiltonian of the underlying dynamical system reads:
$$H=\frac{1}{2}\left(p_1^2+p_2^2+p_3^2\right)\frac{1}{2}\left(q_1^2+q_2^2+q_3^21\right)^2\frac{\sigma _2^2}{2}q_2^2\frac{\sigma _3^2}{2}q_3^2$$
(32)
in Cartesian coordinates. To write the Hamiltonian in elliptic coordinates, note that for $`N=3`$ we have:
$$\alpha =1+\overline{\sigma }_2^2+\overline{\sigma }_3^2,\beta =\overline{\sigma }_2^2\overline{\sigma }_3^2+\overline{\sigma }_2^2+\overline{\sigma }_3^2$$
$$A(\lambda _a)=(\lambda _a1)(\lambda _a\overline{\sigma }_2^2)(\lambda _a\overline{\sigma }_3^2)$$
$$\mathrm{\Lambda }^{}(\lambda _1)=(\lambda _1\lambda _2)(\lambda _1\lambda _3);\mathrm{\Lambda }^{}(\lambda _2)=(\lambda _2\lambda _1)(\lambda _2\lambda _3);\mathrm{\Lambda }^{}(\lambda _3)=(\lambda _3\lambda _1)(\lambda _3\lambda _2)$$
Hence, $`H={\displaystyle \underset{a=1}{\overset{3}{}}}{\displaystyle \frac{1}{\mathrm{\Lambda }^{}(\lambda _a)}}H_a`$, where
$$H_a=2A(\lambda _a)\pi _a^2\frac{1}{2}\left[\lambda _a^2(\lambda _a\overline{\sigma }_2^2)(\lambda _a\overline{\sigma }_3^2)\right]$$
(33)
The separatrix trajectories, those in one-to-one correspondence with solitary waves of kink type in the encompassing field theory, are fully determined by the equations (31) restricted to the $`N=3`$ case:
$`C_2=`$ $`|{\displaystyle \frac{\sqrt{1\lambda _1}+\sigma _2}{\sqrt{1\lambda _1}\sigma _2}}|^{\sigma _3\mathrm{sign}(\pi _1)}|{\displaystyle \frac{\sqrt{1\lambda _1}\sigma _3}{\sqrt{1\lambda _1}+\sigma _3}}|^{\sigma _2\mathrm{sign}(\pi _1)}`$
$`|{\displaystyle \frac{\sqrt{1\lambda _2}+\sigma _2}{\sqrt{1\lambda _2}\sigma _2}}|^{\sigma _3\mathrm{sign}(\pi _2)}|{\displaystyle \frac{\sqrt{1\lambda _2}\sigma _3}{\sqrt{1\lambda _2}+\sigma _3}}|^{\sigma _2\mathrm{sign}(\pi _2)}`$
$`\left|{\displaystyle \frac{\sqrt{1\lambda _3}+\sigma _2}{\sqrt{1\lambda _3}\sigma _2}}\right|^{\sigma _3\mathrm{sign}(\pi _3)}\left|{\displaystyle \frac{\sqrt{1\lambda _3}\sigma _3}{\sqrt{1\lambda _3}+\sigma _3}}\right|^{\sigma _2\mathrm{sign}(\pi _3)}`$
where $`C_2=\mathrm{exp}\{2\gamma _2\sigma _2\sigma _3(\sigma _2^2\sigma _3^2)\}`$ is constant, and:
$`C_3=`$ $`|{\displaystyle \frac{\sqrt{1\lambda _1}1}{\sqrt{1\lambda _1}+1}}|^{\sigma _2\sigma _3(\sigma _2^2\sigma _3^2)\mathrm{sign}(\pi _1)}|{\displaystyle \frac{\sqrt{1\lambda _1}+\sigma _2}{\sqrt{1\lambda _1}\sigma _2}}|^{\sigma _3\overline{\sigma }_3^2\mathrm{sign}(\pi _1)}`$
$`|{\displaystyle \frac{\sqrt{1\lambda _1}\sigma _3}{\sqrt{1\lambda _1}+\sigma _3}}|^{\sigma _2\overline{\sigma }_2^2\mathrm{sign}(\pi _1)}|{\displaystyle \frac{\sqrt{1\lambda _2}1}{\sqrt{1\lambda _2}+1}}|^{\sigma _2\sigma _3(\sigma _2^2\sigma _3^2)\mathrm{sign}(\pi _2)}`$
$`|{\displaystyle \frac{\sqrt{1\lambda _2}+\sigma _2}{\sqrt{1\lambda _2}\sigma _2}}|^{\sigma _3\overline{\sigma }_3^2\mathrm{sign}(\pi _2)}|{\displaystyle \frac{\sqrt{1\lambda _2}\sigma _3}{\sqrt{1\lambda _2}+\sigma _3}}|^{\sigma _2\overline{\sigma }_2^2\mathrm{sign}(\pi _2)}`$
$`|{\displaystyle \frac{\sqrt{1\lambda _3}1}{\sqrt{1\lambda _3}+1}}|^{\sigma _2\sigma _3(\sigma _2^2\sigma _3^2)\mathrm{sign}(\pi _3)}|{\displaystyle \frac{\sqrt{1\lambda _3}+\sigma _2}{\sqrt{1\lambda _3}\sigma _2}}|^{\sigma _3\overline{\sigma }_3^2\mathrm{sign}(\pi _3)}`$
$`\left|{\displaystyle \frac{\sqrt{1\lambda _3}\sigma _3}{\sqrt{1\lambda _3}+\sigma _3}}\right|^{\sigma _2\overline{\sigma }_2^2\mathrm{sign}(\pi _3)}`$
with $`C_3=\mathrm{exp}\{2\gamma _3\sigma _2\sigma _3\overline{\sigma }_2^2\overline{\sigma }_3^2(\sigma _2^2\sigma _3^2)\}`$.
Integration of (30) in the $`N=3`$ case shows the time-table of the particle in each trajectory, or, the kink form factor:
$`C_1(\tau )=`$ $`|{\displaystyle \frac{\sqrt{1\lambda _1}\sigma _2}{\sqrt{1\lambda _1}+\sigma _2}}|^{\sigma _3\overline{\sigma }_2^2\mathrm{sign}(\pi _1)}|{\displaystyle \frac{\sqrt{1\lambda _1}+\sigma _3}{\sqrt{1\lambda _1}\sigma _3}}|^{\sigma _2\overline{\sigma }_3^2\mathrm{sign}(\pi _1)}`$
$`|{\displaystyle \frac{\sqrt{1\lambda _2}+\sigma _2}{\sqrt{1\lambda _2}\sigma _2}}|^{\sigma _3\overline{\sigma }_2^2\mathrm{sign}(\pi _2)}|{\displaystyle \frac{\sqrt{1\lambda _2}\sigma _3}{\sqrt{1\lambda _2}+\sigma _3}}|^{\sigma _2\overline{\sigma }_3^2\mathrm{sign}(\pi _2)}`$
$`\left|{\displaystyle \frac{\sqrt{1\lambda _3}\sigma _2}{\sqrt{1\lambda _3}+\sigma _2}}\right|^{\sigma _3\overline{\sigma }_2^2\mathrm{sign}(\pi _3)}\left|{\displaystyle \frac{\sqrt{1\lambda _3}+\sigma _3}{\sqrt{1\lambda _3}\sigma _3}}\right|^{\sigma _2\overline{\sigma }_3^2\mathrm{sign}(\pi _3)}`$
if $`C_1(\tau )=\mathrm{exp}\{2(\gamma _1+\tau )(\sigma _3^2\sigma _2^2)\sigma _2\sigma _3\}`$. Therefore, there is a family of kinks parametrized by the integration constants $`\gamma _2`$, $`\gamma _3`$: it corresponds to the family of curves in $`P_3(\mathrm{})`$ determined by the intersection of the surfaces defined by (3.1) and (3.1). The third constant, $`\gamma _1`$, fixes the center of the kink, the point where the energy density reaches its maximum value.
Better intuition of the kink shapes requires an interpretation of the solutions described by equations (3.1) and (3.1) in Cartesian coordinates. We shall describe how is this achieved in the next sub-sections, but before this it is convenient to note some details of the change of coordinates from Cartesian to elliptic in $`𝐑^3`$:
$`q_1^2`$ $`=`$ $`{\displaystyle \frac{1}{\sigma _2^2\sigma _3^2}}(1\lambda _1)(1\lambda _2)(1\lambda _3)`$
$`q_2^2`$ $`=`$ $`{\displaystyle \frac{1}{\sigma _2^2(\sigma _3^2\sigma _2^2)}}(\overline{\sigma }_2^2\lambda _1)(\overline{\sigma }_2^2\lambda _2)(\overline{\sigma }_2^2\lambda _3)`$ (37)
$`q_3^2`$ $`=`$ $`{\displaystyle \frac{1}{\sigma _3^2(\sigma _2^2\sigma _3^2)}}(\overline{\sigma }_3^2\lambda _1)(\overline{\sigma }_3^2\lambda _2)(\overline{\sigma }_3^2\lambda _3)`$
The change of coordinates is singular at the three $`𝐑^2`$ coordinate planes; $`q_1=0`$, $`q_2=0`$ and $`q_3=0`$. The image of the $`q_1=0`$ plane is a unique face, $`\lambda _3=1`$, of the $`P_3(\mathrm{})`$ parallelepiped:
$$\mathrm{}<\lambda _1\overline{\sigma }_3^2\lambda _2\overline{\sigma }_2^2\lambda _31$$
(38)
The $`q_2=0`$ plane, however, is mapped into faces $`\lambda _2=\overline{\sigma }_2^2`$ and $`\lambda _3=\overline{\sigma }_2^2`$, while the $`q_3=0`$ plane goes to faces $`\lambda _2=\overline{\sigma }_3^2`$ and $`\lambda _1=\overline{\sigma }_3^2`$ of $`P_3(\mathrm{})`$. Observe that $`g_{11}(\overline{\sigma }_3^2,\lambda _2,\lambda _3)=g_{22}(\lambda _1,\overline{\sigma }_3^2,\lambda _3)=g_{22}(\lambda _1,\overline{\sigma }_2^2,\lambda _3)=g_{33}(\lambda _1,\lambda _2,\overline{\sigma }_2^2)=g_{33}(\lambda _1,\lambda _2,1)=\mathrm{}`$. The whole $`𝐑^3`$ space is mapped in $`P_3(\mathrm{})`$. Due to the symmetry under the group $`G=𝐙_2\times 𝐙_2\times 𝐙_2`$ generated by $`q_aq_a`$, the mapping (37) is eight to one in regular points of $`𝐑^3`$: to any point in the interior of $`P_3(\mathrm{})`$ correspond eight points in $`𝐑^3`$ away from the coordinate planes. These planes are fixed loci of some subgroup of $`G`$.
The asymptotic conditions (4) in $`q_a`$ restrict the motion to the compact sub-space $`D^3`$ of $`𝐑^3`$ bounded by the tri-axial ellipsoid:
$$q_1^2+\frac{q_2^2}{\overline{\sigma }_2^2}+\frac{q_3^2}{\overline{\sigma }_3^2}=1$$
(39)
and are satisfied by finite action and zero energy trajectories. Elliptic coordinates are best suited for demonstrating such a restriction. In this coordinate system $`D^3`$ is mapped to the finite parallelepiped $`P_3(0)`$:
$$0\lambda _1\overline{\sigma }_3^2\lambda _2\overline{\sigma }_2^2\lambda _31$$
(40)
The unique non-singular face of $`P_3(0)`$ with respect to the change of coordinates is $`\lambda _1=0`$ and the inverse image of this face is the ellipsoid (39). The asymptotic conditions (4) force $`\eta _b=0,b`$, and thus, by (23), $`H_a=0,a`$, for the finite action solutions. $`\lambda _2`$ and $`\lambda _3`$ are bounded, see (40). Thus, we focus on,
$$H_1=0\frac{1}{2}\pi _1^2+\frac{1}{8}\frac{\lambda _1^2}{\lambda _11}=0$$
(41)
Equation (41) describes the motion of a particle with zero energy moving under the influence of a potential
$`𝒱(\lambda _1)`$ $`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle \frac{\lambda _1^2}{\lambda _11}},\mathrm{}<\lambda _1\overline{\sigma }_3^2`$
$`=`$ $`\mathrm{},\overline{\sigma }_3^2<\lambda _1<\mathrm{}`$
$`𝒱(\lambda _1)`$ has a maximum at $`\lambda _1=0`$ and goes to $`\mathrm{}`$ when $`\lambda _1`$ tends to $`\mathrm{}`$; therefore, bounded motion occurs only in the $`\lambda _1[0,\overline{\sigma }_3^2]`$ interval and the trajectories giving rise to kinks lie in $`P_3(0)`$, seen in elliptic coordinates, or $`D^3`$ in Cartesian space.
In Figure 1 the whole picture is depicted and we notice the following important elements of the dynamics:
\- Points: (1) the origin. This is a fixed point of $`G=𝐙^{\times 3}`$ and thus only one point O in $`D^3`$ is mapped to the vertex O in $`P_3(\mathrm{})`$. (2) Points B, C, D: these are the intersection points of the three distinguished ellipses , $`q_1^2+\frac{q_2^2}{\overline{\sigma }_2^2}=1`$, $`q_1^2+\frac{q_3^2}{\overline{\sigma }_3^2}=1`$ and $`\frac{q_2^2}{\overline{\sigma }_2^2}+\frac{q_3^2}{\overline{\sigma }_3^2}=1`$, in the ellipsoid (39). They are fixed points under the action of a sub-group $`𝐙_2^{\times 2}`$ of $`G`$ and thus, two points in the boundary of $`D^3`$ are mapped to a single point in the boundary of $`P_3(0)`$. D is the point where the two vacuum points $`\stackrel{}{v}^\pm `$ are mapped and hence it is a very important point of the dynamics: every finite action trajectory starts and ends at D. (3) Points F<sub>1</sub>, F<sub>2</sub>, F<sub>3</sub>, the foci of the above ellipses. Again two points in $`D^3`$ are mapped in a unique point in $`P_3(0)`$. (4) The umbilicus A of the ellipsoid $`\lambda _1=0`$ is another characteristic point; in Cartesian coordinates A corresponds to four points in the boundary of $`D^3`$ because they are invariant only under a $`𝐙_2`$ sub-group of $`G`$.
\- Curves:
* The ellipse with foci F<sub>2</sub>
$$\frac{q_1^2}{\sigma _3^2}+\frac{q_2^2}{\sigma _3^2\sigma _2^2}=1$$
(42)
in the $`q_3=0`$ plane passing through F<sub>3</sub> and F<sub>1</sub>. This is the edge $`\lambda _2=\overline{\sigma }_3^2=\lambda _1`$ in $`P_3(0)`$. Observe that four points on the ellipse (42) are mapped to one point in the edge of $`P_3(0)`$, because it is invariant under a $`𝐙_2`$ sub-group of $`G`$. The map leading to F<sub>3</sub> and F<sub>1</sub> is, however, two to one: the invariance group of these points is bigger, $`𝐙_2^{\times 2}G`$.
* The hyperbola
$$\frac{q_1^2}{\sigma _2^2}\frac{q_3^2}{\sigma _3^2\sigma _2^2}=1$$
(43)
in the $`q_2=0`$ plane passing through F<sub>2</sub> and A and having foci F<sub>3</sub> (the edge $`\lambda _2=\overline{\sigma }_2^2=\lambda _3`$ in $`P_3(0)`$).
The above points and curves play a special rôle in the definition of the elliptic coordinates and are also “critical loci” of the dynamics.
### 3.2 Generic Kinks
The generic kinks of the $`N=3`$ system are the trajectories given by the solutions of (3.1)-(3.1) for non-zero finite values of $`C_2`$ and $`C_3`$. The solutions of the implicit equations (3.1)-(3.1) cannot be graphically represented by means of the built-in functions of Mathematica. We use a numerical algorithm implemented in Mathematica to obtain the graphic portrait of the trajectories. The algorithm allows us to calculate an arbitrary number of points on the orbit. These points joined by straight segments provide a visualization of the trajectory. There is a special step and an iteration of routine steps in the procedure, which is based on the Newton-Raphson method.
First step. Identification of two points on the trajectory.
For given values of $`C_2`$, $`C_3`$, $`\sigma _2`$, $`\sigma _3`$ and a choice of signs, we set the first variable to the “point” $`\lambda _1=\overline{\lambda }_1`$. (3.1)-(3.1) becomes a system of two equations in two unknowns that can be solved by the Newton-Raphson method with starting values $`(\lambda _2^0,\lambda _3^0)`$. The outcome is a point $`P_1(\overline{\lambda }_1,\overline{\lambda }_2,\overline{\lambda }_3)`$ on the trajectory. We repeat this operation starting from $`\lambda _1=\overline{\lambda _1}+ϵ=\overline{\lambda }_1^{}`$ to find a second point $`P_2(\overline{\lambda }_1^{},\overline{\lambda }_2^{},\overline{\lambda }_3^{})`$ on the orbit.
$`\overline{\lambda }_1`$, $`\overline{\lambda }_1^{}`$, $`\lambda _2^0`$ and $`\lambda _3^0`$ are chosen at random; good convergence is attained if these points belong to the middle zones of the variation ranges of $`\lambda _1`$, $`\lambda _2`$ and $`\lambda _3`$ or, at least, they are far away from the singularities on the faces of $`P_3(0)`$
Successive steps.
$`P_1`$ and $`P_2`$ provide an approximation of the curve by the secant line joining them. For some small $`\delta 𝐑^+`$, we choose $`P_3^0=P_1+\delta (P_2P_1)`$ as the starting value of the Newton-Raphson procedure applied to the solution of equations (3.1) and (3.1); we thus obtain the point $`P_3`$ on the curve. $`P_2`$ and $`P_3`$ lead to guesstimate by the same token another value $`P_4^0`$ that produce the next point $`P_4`$ on the orbit and now the iteration is obvious. Replacing $`\delta `$ by $`\delta `$ we travel along the opposite sense on the trajectory. The algorithm stops when one of the three variables $`\lambda _1`$, $`\lambda _2`$, $`\lambda _3`$ reaches its extreme value; it is applied independently on each stage, determined by the signs of $`\pi _a`$ and the global trajectory is obtained by the demand for continuity.
We now describe the portrait of these orbits. Having fixed $`\gamma _2`$ and $`\gamma _3`$, the corresponding kink trajectory is a non-plane curve in the interior of $`P_3(0)`$ that starts from the vacuum point D, reaches the top face BCF<sub>1</sub>O and hits the edge AF<sub>2</sub>. It then goes to the edge F<sub>1</sub>F<sub>3</sub>, back again to the top face, hits the edge AF<sub>2</sub> a second time, the top face a third time and ends at D: see Figure 2 and Figure 3. Varying $`\gamma _2`$ and $`\gamma _3`$ in the range of finite real numbers, other similar trajectories are obtained that hit the edges AF<sub>2</sub> and F<sub>1</sub>F<sub>3</sub> at different points. Given a sense of time there therefore exists a two-parameter family of kink trajectories in one-to-one correspondence with the points in the interior of AF<sub>2</sub> and F<sub>1</sub>F<sub>3</sub>. It should be mentioned that a whole congruence of trajectories parametrized by the interior of F<sub>1</sub>F<sub>3</sub> converges at one single point in the interior of AF<sub>2</sub> and viceversa.
The translation of a generic kink trajectory to Cartesian coordinates is a delicate matter; due to the non-uniqueness of the mapping implied by the change of coordinates special care is necessary in the analysis of the trajectory near the special conics (42)-(43) where several options are a priori possible. Since it requires continuity and derivability to the trajectories in the interior of the ellipsoid (39), the interior of $`D^3`$, the behaviour of the curve is mostly fixed. Two choices, the specification of D$`=\stackrel{}{v}^{}`$ as the starting point and the location of the intersection of the kink trajectory with the $`q_1=0`$ plane in the quadrant characterized by $`q_2>0`$, $`q_3>0`$, completely fix the itinerary.
There is a crossroad where the “particle” touches the $`q_1>0`$ branch of the hyperbola (43) and turns back towards the ellipse (42). There, the movement enters the $`q_3<0`$ half-space and the kink trajectory reaches the other branch, $`q_1<0`$, of the critical hyperbola after crossing the $`q_2=0`$ plane. At this stage, the particle makes its way for a third crossing of the $`q_2=0`$ plane and, finally, the journey ends at D$`=\stackrel{}{v}^+`$. This kind of kink trajectory is therefore heteroclinic: it starts and ends at different unstable points, so that:
$$Q_1^T=\frac{1}{2}_{\mathrm{}}^{\mathrm{}}𝑑\tau \frac{dq_1}{d\tau }=\pm 1,Q_2^T=Q_3^T=0$$
(44)
We call these topological kinks TK3 because they have three non-null components:
$$q_1(\tau )=\varphi _1(x)0,q_2(\tau )=\varphi _2(x)0,q_3(\tau )=\varphi _3(x)0$$
It should be noted that a unique, apparently non-derivable, kink trajectory in elliptic coordinates corresponds to eight derivable trajectories in Cartesian coordinates: the choices of $`\stackrel{}{v}^{}`$ or $`\stackrel{}{v}^+`$ and $`q_3<0`$ or $`q_3>0`$, $`q_2<0`$ or $`q_2>0`$ as the starting point and initial quadrant give the eight possibilities.
The energy of a three-component topological kink is the action of the trajectory times $`\frac{m^3}{\lambda ^2\sqrt{2}}`$ and hence computable from formula (28) for the $`N=3`$ case:
$`{\displaystyle \frac{\lambda ^2\sqrt{2}}{m^3}}E_{\mathrm{TK3}}`$ $`=`$ $`{\displaystyle _0^{\overline{\sigma }_3^2}}{\displaystyle \frac{\lambda _1d\lambda _1}{\sqrt{1\lambda _1}}}+{\displaystyle _{\overline{\sigma }_3^2}^{\overline{\sigma }_2^2}}{\displaystyle \frac{2\lambda _2d\lambda _2}{\sqrt{1\lambda _2}}}+{\displaystyle _{\overline{\sigma }_2^2}^1}{\displaystyle \frac{3\lambda _3d\lambda _3}{\sqrt{1\lambda _3}}}`$
$`=`$ $`{\displaystyle \frac{4}{3}}+{\displaystyle \frac{2}{3}}\left[\sigma _3(3\sigma _3^2)+\sigma _2(3\sigma _2^2)\right]`$
It is independent of $`\gamma _2`$ and $`\gamma _3`$ and hence the same for every kink in the TK3 family.
### 3.3 Enveloping Kinks
There is another family of $`N=3`$ kinks living on the surface $`M_3\{(\lambda _1,\lambda _2,\lambda _3)/`$ $`\lambda _1=0\}`$, the unique face of $`P_3(0)`$ where the elliptic coordinates are not singular. In $`M_3`$, the Hamiltonian becomes:
$$H=\underset{a=2}{\overset{3}{}}\frac{H_a}{\mathrm{\Lambda }^{}(\lambda _a)}=\underset{a=2}{\overset{3}{}}\left\{\frac{2A(\lambda _a)}{\mathrm{\Lambda }^{}(\lambda _a)}\pi _a^2\frac{1}{2}\frac{\lambda _a^2(\lambda _a\overline{\sigma }_2^2)(\lambda _a\overline{\sigma }_3^2)}{\mathrm{\Lambda }^{}(\lambda _a)}\right\}$$
(46)
and therefore there is a two-dimensional system hidden inside the $`N=3`$ model which is Hamilton-Jacobi separable. The orbit equations,
$`C`$ $`=`$ $`|{\displaystyle \frac{\sqrt{1\lambda _2}\sigma _2}{\sqrt{1\lambda _2}+\sigma _2}}|^{\sigma _3\mathrm{sign}(\pi _2)}|{\displaystyle \frac{\sqrt{1\lambda _2}+\sigma _3}{\sqrt{1\lambda _2}\sigma _3}}|^{\sigma _2\mathrm{sign}(\pi _2)}`$ (47)
$`\left|{\displaystyle \frac{\sqrt{1\lambda _3}\sigma _2}{\sqrt{1\lambda _3}+\sigma _2}}\right|^{\sigma _3\mathrm{sign}(\pi _3)}\left|{\displaystyle \frac{\sqrt{1\lambda _3}+\sigma _3}{\sqrt{1\lambda _3}\sigma _3}}\right|^{\sigma _2\mathrm{sign}(\pi _3)}`$
are parametrized by only one real constant $`\gamma _2`$ ($`C=e^{\sigma _2\sigma _3(\sigma _3^2\sigma _2^2)\gamma _2}`$).
The Mathematica plot of these solutions is shown in Figure 4. Having fixed $`\gamma _2`$, the corresponding kink trajectory is a plane curve in $`M_3`$ that starts from the vacuum point D, reaches the top edge BC, goes to the umbilicus A and then back to the edge BC, to end finally in the vacuum point D. The value of $`\gamma _2`$ determines the points in BC where the trajectory bounces back and thus the one-parameter family of this kind of kink trajectories is in one-to-one correspondence with the points in the interior of BC.
In Cartesian coordinates the enveloping kinks are trajectories that unfold on the ellipsoid (39):
$$q_1^2+\frac{q_2^2}{\overline{\sigma }_2^2}+\frac{q_3^2}{\overline{\sigma }_3^2}=1$$
The starting point is either D$`=\stackrel{}{v}^+`$ or D$`=\stackrel{}{v}^{}`$ and the trajectories also end in either D$`=\stackrel{}{v}^+`$ or D$`=\stackrel{}{v}^{}`$. Associated with “homoclinic” trajectories, the corresponding kinks are “non-topological”: $`Q_1^T=Q_2^T=Q_3^T=0`$. The three Cartesian components $`q_a`$ differ from zero and the appropriate name for this kind of solitary wave is a non-topological kink of three components, NTK3 for short. Every NTK3 trajectory on its way from D$`=\stackrel{}{v}^\pm `$ to D$`=\stackrel{}{v}^\pm `$ crosses the umbilicus point of the ellipsoid. Note that, again, eight trajectories in the Cartesian space $`𝐑^3`$ correspond to one trajectory in $`P_3(0)`$: the particle has the freedom to choose the points $`\stackrel{}{v}^+`$ or $`\stackrel{}{v}^{}`$ as base points of the curve. Having fixed one of them, the trajectory may develop in the half-ellipsoids determined in (44) by $`q_30`$ or $`q_30`$ and, finally, there are two travelling senses in each orbit.
Also, the energy of a three-component non-topological kink is essentially the action of the NTK3 trajectory:
$$\frac{\lambda ^2\sqrt{2}}{m^3}E_{\mathrm{NTK3}}=_{\overline{\sigma }_3^2}^{\overline{\sigma }_2^2}\frac{\lambda _2d\lambda _2}{\sqrt{1\lambda _2}}+_{\overline{\sigma }_2^2}^1\frac{2\lambda _3d\lambda _3}{\sqrt{1\lambda _3}}=2\left(\sigma _2+\sigma _3\frac{\sigma _2^2+\sigma _3^2}{3}\right)$$
(48)
according to the Hamilton-Jacobi theory.
### 3.4 Embedded Kinks
Three-component topological and non-topological kinks arise as genuine solitary waves in the $`N=3`$ model. Restriction to the $`q_3=0`$ and/or $`q_2=0`$ planes shows that the $`N=2`$ system is included twice, once in each plane, in the $`N=3`$ model. Therefore, all the solitary waves of the $`N=2`$ model are embedded twice as kinks of the larger $`N=3`$ system. The embedded kinks live on the $`q_2=0`$ and $`q_3=0`$ planes, i.e. the faces of $`P_3(0)`$ where the elliptic coordinate system is singular.
#### I. Embedded kinks in the $`q_2=0`$ plane
Both $`\lambda _2=\overline{\sigma }_2^2`$ and $`\lambda _3=\overline{\sigma }_2^2`$ give $`q_3=0`$, see (37), and hence this coordinate plane in $`𝐑^3`$ is the union of the two faces, $`\lambda _2=\overline{\sigma }_2^2`$ and $`\lambda _3=\overline{\sigma }_2^2`$, of $`P_3(0)`$. Therefore, in
$$M_{2_{\sigma _3}}=\left\{(\lambda _1,\lambda _2,\lambda _3)/\lambda _3=\overline{\sigma }_2^2\right\}\left\{(\lambda _1,\lambda _2,\lambda _3)/\lambda _2=\overline{\sigma }_2^2\right\}=M_{2_{\sigma _3}}^1M_{2_{\sigma _3}}^2$$
we expect to find all the kinks of the $`N=2`$ case.
* In $`M_{2_{\sigma _3}}^1`$, $`\lambda _3=\overline{\sigma }_2^2`$, we are in the face of $`P_3(0)`$ such that $`0<\lambda _1<\overline{\sigma }_3^2<\lambda _2<\overline{\sigma }_2^2`$, and
$$q_1^2=\frac{1}{\sigma _3^2}(1\lambda _1)(1\lambda _2),q_3^2=\frac{1}{\sigma _3^2}(\overline{\sigma }_3^2\lambda _1)(\overline{\sigma }_3^2\lambda _2^2)$$
The Hamiltonian also reduces to the $`N=2`$ Hamiltonian
$`H`$ $`=`$ $`{\displaystyle \frac{2A(\lambda _1)}{\mathrm{\Lambda }^{}(\lambda _1)}}\pi _1^2{\displaystyle \frac{1}{2\mathrm{\Lambda }^{}(\lambda _1)}}\left(\lambda _1^2(\lambda _1\overline{\sigma }_2^2)(\lambda _1\overline{\sigma }_3^2)\right)`$
$`{\displaystyle \frac{2A(\lambda _2)}{\mathrm{\Lambda }^{}(\lambda _2)}}\pi _2^2{\displaystyle \frac{1}{2\mathrm{\Lambda }^{}(\lambda _2)}}\left(\lambda _2^2(\lambda _2\overline{\sigma }_2^2)(\lambda _2\overline{\sigma }_3^2)\right)`$
and the Hamilton-Jacobi method prescribes the equation
$`e^{2\sigma _3\overline{\sigma }_3^2\gamma _2}`$ $`=`$ $`\left(\right|{\displaystyle \frac{\sqrt{1\lambda _1}\sigma _3}{\sqrt{1\lambda _1}+\sigma _3}}||{\displaystyle \frac{\sqrt{1\lambda _1}+1}{\sqrt{1\lambda _1}1}}|^{\sigma _3})^{\mathrm{sign}(\pi _1)}`$ (49)
$`\left(\left|{\displaystyle \frac{\sqrt{1\lambda _2}\sigma _3}{\sqrt{1\lambda _2}+\sigma _3}}\right|\left|{\displaystyle \frac{\sqrt{1\lambda _2}+1}{\sqrt{1\lambda _2}1}}\right|^{\sigma _3}\right)^{\mathrm{sign}(\pi _2)}`$
as ruling the portion of the trajectories at this face, bounded by the edges AD, AF<sub>2</sub>, F<sub>2</sub>F<sub>3</sub> and F<sub>3</sub>D (see Figure 5).
* In $`M_{2_{\sigma _3}}^2`$, $`\lambda _2=\overline{\sigma }_2^2`$ and the face in the boundary of $`P_3(0)`$ is $`0<\lambda _1<\overline{\sigma }_3^2<\overline{\sigma }_2^2<\lambda _3<1`$. The trajectory equations at this face are:
$`e^{2\sigma _3\overline{\sigma }_3^2\gamma _2}`$ $`=`$ $`\left(\right|{\displaystyle \frac{\sqrt{1\lambda _1}\sigma _3}{\sqrt{1\lambda _1}+\sigma _3}}||{\displaystyle \frac{\sqrt{1\lambda _1}+1}{\sqrt{1\lambda _1}1}}|^{\sigma _3})^{\mathrm{sign}(\pi _1)}`$ (50)
$`\left(\left|{\displaystyle \frac{\sqrt{1\lambda _3}\sigma _3}{\sqrt{1\lambda _3}+\sigma _3}}\right|\left|{\displaystyle \frac{\sqrt{1\lambda _3}+1}{\sqrt{1\lambda _3}1}}\right|^{\sigma _3}\right)^{\mathrm{sign}(\pi _3)}`$
and the boundary is formed by the edges AF<sub>2</sub>, F<sub>2</sub>O, OB and BA.
For finite values of $`\gamma _2`$, some kink trajectories given by (49)-(50) are depicted in Figure 5. It may be observed that the trajectory starts at D and then runs through the face $`M_{2_{\sigma _3}}^1`$ until the edge AF<sub>2</sub>. From this point, the particle enters the $`M_{2_{\sigma _3}}^2`$ face (here the path is not derivable), reaches the BO edge and comes back to the AF<sub>2</sub> edge. This is the second point of non-differentiability re-entering the trajectory the $`M_{2_{\sigma _3}}^1`$ face. All the trajectories then meet at the vertex F<sub>3</sub> and come back in a symmetric way to end in the D point. In Cartesian coordinates, these kink trajectories start and end in either D=$`\stackrel{}{v}^+`$ or D$`=\stackrel{}{v}^{}`$, do not leave the $`q_2=0`$ plane, and cross either the focus ($`q_1=\sigma _3`$, $`q_3=0`$) or ($`q_1=\sigma _3`$, $`q_3=0`$). We therefore call them NTK2$`_{\sigma _3}`$ because they are two-component non-topological kinks, merely the family of NTK2 kinks of the $`N=2`$ model, embedded this way within the manifold of kinks of the $`N=3`$ system. There are four trajectories of this kind inside the ellipsoid (44) in $`𝐑^3`$ per trajectory in the boundary of $`P_3(0)`$: there is freedom to choose $`\stackrel{}{v}^+`$ or $`\stackrel{}{v}^{}`$ and the sense of travel in each orbit. The NTK2$`_{\sigma _3}`$ kinks are fixed points of the $`𝐙_2`$ sub-group of $`G=𝐙_2^{\times 3}`$ generated by $`q_2q_2`$, that, however, does not leave invariant the TK3 and the NTK3 trajectories . The energy of these solutions is:
$`{\displaystyle \frac{\lambda ^2\sqrt{2}}{m^3}}E_{\mathrm{NTK2}_{\sigma _3}}`$ $`=`$ $`{\displaystyle _0^{\overline{\sigma }_3^2}}{\displaystyle \frac{\lambda _1d\lambda _1}{\sqrt{1\lambda _1}}}+{\displaystyle _{\overline{\sigma }_3^2}^{\overline{\sigma }_2^2}}{\displaystyle \frac{2\lambda _2d\lambda _2}{\sqrt{1\lambda _2}}}+{\displaystyle _{\overline{\sigma }_2^2}^1}{\displaystyle \frac{2\lambda _3d\lambda _3}{\sqrt{1\lambda _3}}}`$ (51)
$`=`$ $`{\displaystyle \frac{4}{3}}+2\sigma _3\left(1{\displaystyle \frac{\sigma _3^2}{3}}\right)`$
There is a limiting case to this family of kinks: a trajectory along the DA and AB edges and back to D through the same way. Elliptic coordinates are even more singular on the edges, but the dynamical system reduces to a one-dimensional Hamiltonian system which can be integrated analytically. We have a two-step trajectory:
At the DA edge, $`\lambda _1=0`$ and $`\lambda _3=\overline{\sigma }_2^2`$, the canonical equations (after use of the first integral) reduce to:
$$\frac{d\lambda _2}{d\tau }=\pm 2(\lambda _2\overline{\sigma }_3^2)\sqrt{1\lambda _2}$$
(52)
with the solution
$$\lambda _1^{\mathrm{TK2}_{\sigma _3}}(\tau )=0,\lambda _2^{\mathrm{TK2}_{\sigma _3}}(\tau )=1\sigma _3^2\mathrm{tanh}^2(\sigma _3\tau ),\lambda _3^{\mathrm{TK2}_{\sigma _3}}(\tau )=\overline{\sigma }_2^2$$
(53)
for $`\tau (\mathrm{},\frac{1}{\sigma _3}arctanh\frac{\sigma _2}{\sigma _3}][\frac{1}{\sigma _3}arctanh\frac{\sigma _2}{\sigma _3},\mathrm{})`$. The second step occurs on the AB edge, where, again, the canonical equations reduce to a single differential equation: if $`\lambda _1=0`$ and $`\lambda _2=\overline{\sigma }_2^2`$,
$$\frac{d\lambda _3}{d\tau }=\pm 2(\lambda _3\overline{\sigma }_3^2)\sqrt{1\lambda _3}$$
(54)
has the solution
$$\lambda _1^{\mathrm{TK2}_{\sigma _3}}(\tau )=0,\lambda _2^{\mathrm{TK2}_{\sigma _3}}(\tau )=\overline{\sigma }_2^2,\lambda _3^{\mathrm{TK2}_{\sigma _3}}(\tau )=1\sigma _3^2\mathrm{tanh}^2(\sigma _3\tau )$$
(55)
for $`\tau [\frac{1}{\sigma _3}arctanh\frac{\sigma _2}{\sigma _3},\frac{1}{\sigma _3}arctanh\frac{\sigma _2}{\sigma _3}]`$. The corresponding kinks in Cartesian coordinates are TK2$`_{\sigma _3}`$ and TK2$`{}_{}{}^{}{}_{\sigma _3}{}^{}`$, the four two-component topological kinks of the $`N=2`$ model:
$$\left(\begin{array}{c}q_1^{\mathrm{TK2}_{\sigma _3}}(\tau )\\ q_2^{\mathrm{TK2}_{\sigma _3}}(\tau )\\ q_3^{\mathrm{TK2}_{\sigma _3}}(\tau )\end{array}\right)=\pm \left(\begin{array}{c}\mathrm{tanh}(\sigma _3\tau )\\ 0\\ \pm \overline{\sigma }_3sech(\sigma _3\tau )\end{array}\right)\left(\begin{array}{c}\varphi _1^{\mathrm{TK2}_{\sigma _3}}(x,t)\\ \varphi _2^{\mathrm{TK2}_{\sigma _3}}(x,t)\\ \varphi _3^{\mathrm{TK2}_{\sigma _3}}(x,t)\end{array}\right)=\pm \left(\begin{array}{c}\mathrm{tanh}(\sigma _3x)\\ 0\\ \pm \overline{\sigma }_3sech(\sigma _3x)\end{array}\right)$$
(56)
Thus, the enveloping kinks of the $`N=2`$ model are also embedded in the $`N=3`$ system. The energy for these solutions and their anti-kinks is:
$$\frac{\lambda ^2\sqrt{2}}{m^3}E_{\mathrm{TK2}_{\sigma _3}}=_{\overline{\sigma }_3^2}^{\overline{\sigma }_2^2}\frac{\lambda _2d\lambda _2}{\sqrt{1\lambda _2}}+_{\overline{\sigma }_2^2}^1\frac{\lambda _3d\lambda _3}{\sqrt{1\lambda _3}}=2\sigma _3\left(1\frac{\sigma _3^2}{3}\right)$$
(57)
In the $`q_2=0`$ plane there is still one trajectory that is even more singular: it is a three step path running on the edges DF<sub>3</sub>, F<sub>3</sub>F<sub>2</sub>, F<sub>2</sub>O and back to D through the same way. The canonical equations and its solutions in the three steps are:
1. $`\lambda _2=\overline{\sigma }_3^2`$ and $`\lambda _3=\overline{\sigma }_2^2`$.
$$\frac{d\lambda _1}{d\tau }=\pm 2\lambda _1\sqrt{1\lambda _1},\tau (\mathrm{},arctanh\sigma _3][arctanh\sigma _3,\mathrm{})$$
$$\lambda _1^{\mathrm{TK1}}(\tau )=1\mathrm{tanh}^2\tau ,\lambda _2^{\mathrm{TK1}}(\tau )=\overline{\sigma }_3^2,\lambda _3^{\mathrm{TK1}}(\tau )=\overline{\sigma }_2^2$$
2. $`\lambda _1=\overline{\sigma }_3^2`$ and $`\lambda _3=\overline{\sigma }_2^2`$.
$$\frac{d\lambda _2}{d\tau }=\pm 2\lambda _2\sqrt{1\lambda _2},\tau [arctanh\sigma _3,arctanh\sigma _2][arctanh\sigma _2,arctanh\sigma _3]$$
$$\lambda _1^{\mathrm{TK1}}(\tau )=\overline{\sigma }_3^2,\lambda _2^{\mathrm{TK1}}(\tau )=1\mathrm{tanh}^2\tau ,\lambda _3^{\mathrm{TK1}}(\tau )=\overline{\sigma }_2^2$$
3. $`\lambda _1=\overline{\sigma }_3^2`$ and $`\lambda _2=\overline{\sigma }_2^2`$.
$$\frac{d\lambda _3}{d\tau }=\pm 2\lambda _3\sqrt{1\lambda _3},\tau [arctanh\sigma _2,arctanh\sigma _2]$$
$$\lambda _1^{\mathrm{TK1}}(\tau )=\overline{\sigma }_3^2,\lambda _2^{\mathrm{TK1}}(\tau )=\overline{\sigma }_2^2,\lambda _3^{\mathrm{TK1}}(\tau )=1\mathrm{tanh}^2\tau $$
Only one Cartesian component is different from zero:
$$\left(\begin{array}{c}q_1^{\mathrm{TK1}}(\tau )\\ q_2^{\mathrm{TK1}}(\tau )\\ q_3^{\mathrm{TK1}}(\tau )\end{array}\right)=\left(\begin{array}{c}\pm \mathrm{tanh}\tau \\ 0\\ 0\end{array}\right)\left(\begin{array}{c}\varphi _1^{\mathrm{TK1}}(x,t)\\ \varphi _2^{\mathrm{TK1}}(x,t)\\ \varphi _3^{\mathrm{TK1}}(x,t)\end{array}\right)=\left(\begin{array}{c}\pm \mathrm{tanh}x\\ 0\\ 0\end{array}\right)$$
(58)
and hence the one-component topological kink of the $`N=1`$ model is embedded first in the manifold of kinks of the $`N=2`$ model, and then in the $`N=3`$ system. There are two kinks of this kind in Cartesian coordinates which are mapped in a unique trajectory in the boundary of $`P_3(0)`$. The TK1 trajectories are fixed points of the $`𝐙^{\times 2}`$ sub-group of $`G`$ generated by $`q_2q_2`$ and $`q_3q_3`$. The energy is
$$\frac{\lambda ^2\sqrt{2}}{m^3}E_{\mathrm{TK1}}=_0^{\overline{\sigma }_3^2}\frac{\lambda _1d\lambda _1}{\sqrt{1\lambda _1}}+_{\overline{\sigma }_3^2}^{\overline{\sigma }_2^2}\frac{\lambda _2d\lambda _2}{\sqrt{1\lambda _2}}+_{\overline{\sigma }_2^2}^1\frac{\lambda _3d\lambda _3}{\sqrt{1\lambda _3}}=\frac{4}{3}$$
(59)
#### Embedded Kinks in the $`q_3=0`$ plane
The DF<sub>3</sub>, F<sub>3</sub>F<sub>2</sub> and F<sub>2</sub>O edges form the intersection of the $`q_2=0`$ and $`q_3=0`$ planes. Therefore, the TK1 kinks also live in the $`q_3=0`$ plane. There is another maximally singular trajectory living on the “edge” in the $`q_3=0`$ plane:
At the DC edge, $`\lambda _2=\overline{\sigma }_3^2`$ and $`\lambda _1=0`$, the canonical equations for the finite action trajectories are:
$$\frac{d\lambda _3}{d\tau }=\pm 2(\lambda _3\overline{\sigma }_2^2)\sqrt{1\lambda _3}$$
(60)
The path
$$\lambda _1^{\mathrm{TK2}_{\sigma _2}}(\tau )=0,\lambda _2^{\mathrm{TK2}_{\sigma _2}}(\tau )=\overline{\sigma }_3^2,\lambda _3^{\mathrm{TK2}_{\sigma _2}}(\tau )=1\sigma _2^2\mathrm{tanh}^2(\sigma _2\tau )$$
(61)
solves (60) and runs when $`\tau `$ goes from $`\mathrm{}`$ to $`+\mathrm{}`$ from D to D passing through the vertex C at $`\tau =0`$. In Cartesian coordinates we recover the four two-component topological kinks of the $`N=2`$ model, now embedded in the $`q_3=0`$ plane:
$$\left(\begin{array}{c}q_1^{\mathrm{TK2}_{\sigma _2}}(\tau )\\ q_2^{\mathrm{TK2}_{\sigma _2}}(\tau )\\ q_3^{\mathrm{TK2}_{\sigma _2}}(\tau )\end{array}\right)=\pm \left(\begin{array}{c}\mathrm{tanh}(\sigma _2\tau )\\ \pm \overline{\sigma }_2^2sech(\sigma _2\tau )\\ 0\end{array}\right)\left(\begin{array}{c}\varphi _1^{\mathrm{TK2}_{\sigma _2}}(x,t)\\ \varphi _2^{\mathrm{TK2}_{\sigma _2}}(x,t)\\ \varphi _3^{\mathrm{TK2}_{\sigma _2}}(x,t)\end{array}\right)=\pm \left(\begin{array}{c}\mathrm{tanh}(\sigma _2x)\\ \pm \overline{\sigma }_2^2sech(\sigma _2x)\\ 0\end{array}\right)$$
(62)
These are heteroclinic trajectories that produce the TK2$`_{\sigma _2}`$ and TK2$`{}_{}{}^{}{}_{\sigma _2}{}^{}`$ topological kinks. The energy is:
$$\frac{\lambda ^2\sqrt{2}}{m^3}E_{\mathrm{TK2}_{\sigma _2}}=_{\overline{\sigma }_2^2}^1\frac{\lambda _3d\lambda _3}{\sqrt{1\lambda _3}}=2\sigma _2\left(1\frac{\sigma _2^2}{3}\right)$$
(63)
Of course,the full manifold of kinks of the $`N=2`$ model is embedded in the $`q_3=0`$ plane : the set of kinks of the $`N=3`$ system is completed by the two-component non-topological kinks living in the $`q_3=0`$ plane, see Figure 6. The $`q_3=0`$ plane is mapped to the union of two faces in the boundary of $`P_3(0)`$:
$$M_{2_{\sigma _2}}=\left\{(\lambda _1,\lambda _2,\lambda _3)/\lambda _1=\overline{\sigma }_3^2\right\}\left\{(\lambda _1,\lambda _2,\lambda _3)/\lambda _2=\overline{\sigma }_3^2\right\}=M_{2_{\sigma _2}}^1M_{2_{\sigma _2}}^2$$
* 1. In $`M_{2_{\sigma _2}}^1`$, $`\lambda _1=\overline{\sigma }_3^2`$ implies $`q_3=0`$. Therefore:
$`q_1^2`$ $`=`$ $`{\displaystyle \frac{1}{\sigma _2^2}}(1\lambda _2)(1\lambda _3)`$
$`q_2^2`$ $`=`$ $`{\displaystyle \frac{1}{\sigma _2^2}}(\overline{\sigma }_2^2\lambda _2)(\overline{\sigma }_2^2\lambda _3)`$
is a well defined change of coordinates in the range $`\overline{\sigma }_3^2<\lambda _2<\overline{\sigma }_2^2<\lambda _3<1`$. In this region, the interior of the ellipse (42), the trajectories providing kinks are given by the equations:
$`e^{2\sigma _2\overline{\sigma }_2^2\gamma _2}`$ $`=`$ $`\left(\right|{\displaystyle \frac{\sqrt{1\lambda _2}\sigma _2}{\sqrt{1\lambda _2}+\sigma _2}}||{\displaystyle \frac{\sqrt{1\lambda _2}+1}{\sqrt{1\lambda _2}1}}|^{\sigma _2})^{\mathrm{sign}(\pi _2)}`$ (64)
$`\left(\left|{\displaystyle \frac{\sqrt{1\lambda _3}\sigma _2}{\sqrt{1\lambda _3}+\sigma _2}}\right|\left|{\displaystyle \frac{\sqrt{1\lambda _3}+1}{\sqrt{1\lambda _2}1}}\right|^{\sigma _2}\right)^{\mathrm{sign}(\pi _3)}`$
* 2. In $`M_{2_{\sigma _2}}^2`$, $`\lambda _2=\overline{\sigma }_3^2`$ also implies $`q_3=0`$. In the range $`0<\lambda _1<\overline{\sigma }_3^2<\overline{\sigma }_2^2<\lambda _3<1`$ the change of coordinates is defined as
$`q_1^2`$ $`=`$ $`{\displaystyle \frac{1}{\sigma _2^2}}(1\lambda _1)(1\lambda _3)`$
$`q_1^2`$ $`=`$ $`{\displaystyle \frac{1}{\sigma _2^2}}(\overline{\sigma }_2^2\lambda _1)(\overline{\sigma }_2^2\lambda _3)`$
The kink trajectories satisfy the equations:
$`e^{2\sigma _2\overline{\sigma }_2^2\gamma _2}`$ $`=`$ $`\left(\right|{\displaystyle \frac{\sqrt{1\lambda _1}\sigma _2}{\sqrt{1\lambda _1}+\sigma _2}}||{\displaystyle \frac{\sqrt{1\lambda _1}+1}{\sqrt{1\lambda _1}1}}|^{\sigma _2})^{\mathrm{sign}(\pi _1)}`$ (65)
$`\left(\left|{\displaystyle \frac{\sqrt{1\lambda _3}\sigma _2}{\sqrt{1\lambda _3}+\sigma _2}}\right|\left|{\displaystyle \frac{\sqrt{1\lambda _3}+1}{\sqrt{1\lambda _2}1}}\right|^{\sigma _2}\right)^{\mathrm{sign}(\pi _3)}`$
The features of this kind of kinks are identical to the characteristics of the two-component non-topological kinks that exist in the $`q_2=0`$ plane. The only difference is that they have support in the faces $`M_{2_{\sigma _2}}^1`$ and $`M_{2_{\sigma _2}}^2`$ instead of $`M_{2_{\sigma _3}}^1`$ and $`M_{2_{\sigma _3}}^2`$ and we therefore call them NTK2$`_{\sigma _2}`$. They meet at the vertex F<sub>2</sub>, and therefore at the foci $`(q_1=\pm \sigma _2,0,0)`$ in $`𝐑^3`$; see Figure 6. The energy is:
$`{\displaystyle \frac{\lambda ^2\sqrt{2}}{m^3}}E_{\mathrm{NTK2}_{\sigma _2}}`$ $`=`$ $`2\left[{\displaystyle _0^{\overline{\sigma }_3^2}}{\displaystyle \frac{\lambda _1d\lambda _1}{\sqrt{1\lambda _1}}}+{\displaystyle _{\overline{\sigma }_3^2}^{\overline{\sigma }_2^2}}{\displaystyle \frac{\lambda _2d\lambda _2}{\sqrt{1\lambda _2}}}+{\displaystyle _{\overline{\sigma }_2^2}^1}{\displaystyle \frac{\lambda _3d\lambda _3}{\sqrt{1\lambda _3}}}\right]=`$ (66)
$`=`$ $`{\displaystyle \frac{4}{3}}+2\sigma _2\left(1{\displaystyle \frac{\sigma _2^2}{3}}\right)`$
In sum: the manifold of kinks of the $`N=2`$ model is embedded twice in the $`N=3`$ system, once in the $`q_2=0`$ plane and other in the $`q_3=0`$ plane. They are sewn togehter by the common TK1, embedded from the $`N=1`$ model . The embedded kinks fill the gaps left by the TK3 families of kinks in the interior of the ellipsoid (44) and also develop through the curves left by the NTK3 families on the boundary of $`D^3`$. $`D^3`$ is thus a “totally” geodesic manifold with respect to the separatrices between bounded and unbounded motion in the $`N=3`$ dynamical system. The NTK3 family form the envelop of the separatrices and the NTK3 kinks are themselves separatrices.
## 4 Further Comments
We now infer the general structure of the kink manifold of the linear $`O(N)`$-sigma model from the pattern shown by the $`O(2)`$\- and $`O(3)`$-sigma models. We can safely state that all the kink trajectories live in the sub-manifold $`D^N𝐑^N`$ determined by the inequality:
$$q_1^2+\frac{q_2^2}{\overline{\sigma }_2^2}+\mathrm{}+\frac{q_2^2}{\overline{\sigma }_N^2}1$$
There are three categories:
1. Generic Kinks
* There exists a family of generic kinks parametrized by $`N1`$ real constants that live in the interior of $`D^N`$. The intersection loci of the generic kinks are the singular quadrics:
$$\begin{array}{ccccc}\frac{q_1^2}{\overline{\sigma }_N^2}+\frac{q_2^2}{\overline{\sigma }_N^2\overline{\sigma }_2^2}+\hfill & \mathrm{}& \mathrm{}+\frac{q_{N1}^2}{\overline{\sigma }_N^2\overline{\sigma }_{N1}^2}& =1,& q_N=0;\lambda _1=\lambda _2=\overline{\sigma }_N^2\hfill \\ \frac{q_1^2}{\overline{\sigma }_{N1}^2}+\frac{q_2^2}{\overline{\sigma }_{N1}^2\overline{\sigma }_2^2}+\hfill & \mathrm{}& +\frac{q_{N2}^2}{\overline{\sigma }_{N1}^2\overline{\sigma }_{N2}^2}\frac{q_N^2}{\overline{\sigma }_N^2\overline{\sigma }_{N1}^2}& =1,& q_{N1}=0;\lambda _2=\lambda _3=\overline{\sigma }_{N1}^2\hfill \\ & \mathrm{}& \mathrm{}\mathrm{}& & \mathrm{}\hfill \\ \frac{q_1^2}{\overline{\sigma }_2^2}\frac{q_3^2}{\overline{\sigma }_3^2\overline{\sigma }_2^2}\hfill & \mathrm{}& \mathrm{}\frac{q_N^2}{\overline{\sigma }_N^2\overline{\sigma }_2^2}& =1,& q_2=0;\lambda _{N1}=\lambda _N=\overline{\sigma }_2^2\hfill \end{array}$$
* Generic kinks are non-topological- hence NTKN- if $`N`$ even, and topological- hence TKN- if N is odd.
2. Enveloping kinks.
* The restriction of the dynamical system to the boundary $`D^N`$ of $`D^N`$, the hyper-ellipsoid:
$$q_1^2+\frac{q_2^2}{\overline{\sigma }_2^2}+\mathrm{}+\frac{q_2^2}{\overline{\sigma }_N^2}=1$$
provides a family of enveloping kinks parametrized by $`N2`$ real constants. Recalling that in elliptic coordinates $`D^N`$ is characterized by the equation $`\lambda _1=0`$, the intersection loci of this congruence are the umbilical sub-manifolds:
$$\lambda _1=0,\lambda _a=\overline{\sigma }_{Na+1}^2=\lambda _{a+1},a=2,3,\mathrm{},N1$$
of dimension $`N3`$ of the hyper-ellipsoid $`D^N`$.
* Enveloping kinks are topological, hence TKN, if $`N`$ is even, and non-topological, hence NTKN, if $`N`$ is odd.
3. Embedded Kinks.
On the $`N1`$ $`𝐑^{N1}`$ sub-manifolds determined by the conditions $`q_a=0`$, if $`a=2`$ or $`3`$ or $`\mathrm{}`$ or $`N`$, the dynamical system reduces to the mechanical system that arises in the linear $`O(N1)`$-sigma model. Thus, the kink manifold of the $`N1`$ case is included $`N1`$ times in the $`O(N)`$-model, filling the holes left in the interior of $`D^N`$ by the generic kinks, and also covering in $`D^N`$ the sub-spaces which are not covereded by the enveloping kinks. Each $`N1`$ kink sub-manifold is not, however, included $`(N1)\times (N2)`$ times in the $`O(N)`$-model because the $`𝐑^{N2}`$ sub-spaces are intersections of the $`N1`$ $`𝐑^{N1}`$, defined above. The $`N1`$ kink manifolds are not separated but sewn togehter through the $`N2`$ kink sub-manifolds. This is a iterative process in such a way that the kink manifold of the $`O(Nr)`$-sigma model is included $`\left(\begin{array}{cc}N1& \\ r& \end{array}\right)`$ times in the kink manifold of the $`O(N)`$ system.
One can ask what happens if a continuous sub-group $`O(r)`$ of $`O(N)`$ survives as symmetry group of the system. This happens if the deformation is chosen in such a way that $`0<\sigma _2^2=\sigma _3^2=\mathrm{}=\sigma _r^2<\sigma _{r+1}^2<\mathrm{}<\sigma _N^2<1`$. In this case we obtain a sub-manifold of kinks from $`O(r)`$ rotations around the $`q_1`$ axis of the kink manifold of the $`N=2`$ system that lives in the $`q_1:q_2`$ plane. The remaining kinks correspond to the solitary waves of the $`N=r1`$ system defined in the orthogonal $`𝐑^{Nr+1}`$ sub-space. Also the deformations where $`1<\sigma _{r+1}^2<\mathrm{}<\sigma _N^2`$ are easy to understand. Finite action trajectories spread out in the domain in $`𝐑^N`$ bounded by the hyper-hyperboloid:
$$q_1^2+\frac{q_2^2}{\overline{\sigma }_2^2}+\mathrm{}+\frac{q_r^2}{\overline{\sigma }_r^2}\frac{q_{r+1}^2}{|\overline{\sigma }_{r+1}^2|}\mathrm{}\frac{q_n^2}{|\overline{\sigma }_N^2|}=1.$$
The kink manifold of this system is the kink manifold of the $`N=r`$ model defined in the sub-space $`𝐑^r𝐑^N`$ such that $`q_{r+1}=\mathrm{}=q_N=0`$.
Finally we consider a mild deformation of our model by introducing asymmetries in the non-harmonic terms of the potential energy and also adapting the quadratic terms in a suitable manner:
$`V`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(\varphi _1^2+(1+\epsilon _2)\varphi _2^2+\mathrm{}+(1+\epsilon _N)\varphi _N^21\right)^2+`$
$`+{\displaystyle \frac{\sigma _2^2}{2}}\varphi _2^2+\mathrm{}+{\displaystyle \frac{\sigma _N^2}{2}}\varphi _N^2+{\displaystyle \frac{\delta _2}{4}}\varphi _2^4+\mathrm{}+{\displaystyle \frac{\delta _N}{4}}\varphi _N^4`$
The new non-dimensional constants $`\epsilon _a`$ and $`\delta _a`$ are defined in terms of the old $`\sigma _a`$’s through:
$$1+\epsilon _a=\frac{\sigma _a(\sigma _a+1)}{2},\delta _a=\frac{\sigma _a^3(2+\sigma _a)}{2},a=2,3,\mathrm{},N$$
Among the kinks of the deformed linear $`O(N)`$-sigma model only the following survive as solitary wave solutions of this perturbed system:
1. The TK1.
$$\varphi _1=\mathrm{tanh}x,\varphi _2=\mathrm{}=\varphi _N=0$$
2. All the TK2 kinks. On the ellipse,
$$\varphi _1^2+\frac{1+\epsilon _a}{1\sigma _a^2}\varphi _a^2=1$$
the TK2$`\sigma _a`$ and TK2$`{}_{}{}^{}\sigma _{a}^{}`$ configurations,
$$\varphi _1=\mathrm{tanh}\sigma _ax,\varphi _a=\pm \sqrt{\frac{1\sigma _a^2}{1+\epsilon _a}}\mathrm{sech}\sigma _ax,\varphi _b=0,ba,b1$$
are solutions of the field equations. The amazing fact is that in this deformation of the $`O(N)`$-linear sigma model, discussed by Bazeia et al. if $`N=2`$ , the energy of all these kinks is the same:
$$E_{\mathrm{TK1}}=E_{\mathrm{TK2}\sigma _2}=\mathrm{}=E_{\mathrm{TK2}\sigma _N}=\frac{4}{3\sqrt{2}}\frac{m^3}{\lambda ^2}$$
On one hand, we have a deformation of the linear $`O(N)`$-sigma model that exhibits a complex variety of kinks; on the other hand, another deformation of the $`O(N)`$-model rejects almost every kink but the simplest ones as solutions, and all of the surviving kinks are degenerated in energy.
## 5 Outlook
The developments disclosed in this paper suggest a general strategy in the search for kinks in two space-time dimensional field theories. When the fields have $`N`$ components assembled in a vector representation of the $`O(N)`$ group, we focus on systems with symmetry breaking to a discrete sub-group of $`O(N)`$ which has more than one element. If the dynamical system that determines the localized static solutions is completely integrable, all the solitary waves can be found, at least in principle. Particularly interesting is the situation where the $`N1`$ invariants in involution with the mechanical energy act non-trivially on the manifold of localized solutions and the orbit is a continuous space. One can then perturb such a system, loosing in the perturbation many of the solitary wave solutions: only few of the localized static solutions survive as kinks of the perturbed (more realistic) model.
We finally list several interesting questions that will be postponed for future research:
* study of the structure of the kink manifold of the deformed linear $`O(N)`$-sigma model as a moduli space seems to be worthwhile.
* A detailed analysis of the sum rules between the energies of the different kinds of kinks is necessary to fix the structure mentioned above.
* A treatment à la Bogomolny is also possible. This allows for a supersymmetric extension of the model in such a way that the kinks become BPS states.
* The difficult problem remains of determining the stability of the different kinds of kinks. Application of the Morse index theorem helps in finding the stability properties,which in turn provide information about the quantization of these topological defects.
## 6 Acknowledgements
The authors are grateful to Askold Perelomov for teaching them the magic of the elliptic Jacobi coordinates and their relationship to dynamical problems on ellipsoids.
## Appendix: Elliptic coordinates
Given any set of $`N`$ real positive numbers such that $`0<r_1<r_2<\mathrm{}<r_N`$, let us consider the equation:
$$\underset{a=1}{\overset{N}{}}\frac{q_a^2}{r_a\lambda }=1$$
(67)
The left-hand member $`Q_\lambda (\stackrel{}{q})={\displaystyle \underset{a=1}{\overset{N}{}}}{\displaystyle \frac{q_a^2}{r_a\lambda }}`$ of this equation can be understood either as a function of $`𝐑^N`$, for fixed $`\lambda 𝐂`$, or as a function of the complex variable $`\lambda `$, for fixed $`\stackrel{}{q}𝐑^N`$. From (67) one immediately deduces:
$$1Q_\lambda (\stackrel{}{q})=1+\underset{a=1}{\overset{N}{}}\frac{q_a^2}{\lambda r_a}=\frac{{\displaystyle \underset{a=1}{\overset{N}{}}}(\lambda \lambda _a)}{{\displaystyle \underset{a=1}{\overset{N}{}}}(\lambda r_a)}=0$$
(68)
Therefore, the $`N`$ roots $`\lambda _a`$ of the polynomial in the numerator of $`1Q_\lambda (\stackrel{}{q})`$, a rational function of $`\lambda `$, are the roots of equation (67). The roots $`\lambda _a`$ are also real numbers and $`1Q_\lambda (\stackrel{}{q})`$ is a rational function such that $`\lambda _1<r_1<\lambda _2<\mathrm{}<r_{N1}<\lambda _N<r_N`$, see Figure 8. To prove this point one needs to study $`1Q_\lambda (\stackrel{}{\mathrm{q}})`$ along the $`\lambda `$-real axis, near the poles $`\lambda =r_a`$, using Bolzano’s theorem.
Definition. The elliptic coordinates of the point $`\stackrel{}{q}(q_1,\mathrm{},q_N)𝐑^N`$ are the roots $`\stackrel{}{\lambda }_E(\lambda _1,\mathrm{},\lambda _N)P^N(\mathrm{})`$ of $`Q_\lambda (\stackrel{}{q})=1`$.
$`P^N(\mathrm{})𝐑^N`$ is the open sub-space of $`𝐑^N`$ given by: $`\mathrm{}<\lambda _1<r_1`$, $`r_1<\lambda _2<r_2`$, $`\mathrm{}`$, $`r_{N1}<\lambda _N<r_N`$. The solution of (67) for $`\lambda =\lambda _1`$constant $`(\mathrm{},r_1)`$ is, geometrically, a quadric surface, a hyperellipsoid, when $`\stackrel{}{q}`$ varies in $`𝐑^N`$. The family of quadrics obtained by taking $`\lambda =\lambda _a`$ constant $`(r_{a1},r_a)`$, $`a2`$, corresponds to a family of hyper-hyperboloids of every possible signature in $`𝐑^N`$, if $`Q_\lambda (\stackrel{}{q})`$ is considered as a function of $`\stackrel{}{q}`$.
It is convenient to denote the products in (68) as:
$$\mathrm{\Lambda }(\lambda )=\underset{a=1}{\overset{N}{}}(\lambda \lambda _a),A(\lambda )=\underset{a=1}{\overset{N}{}}(\lambda r_a)$$
so that
$$1+\underset{a=1}{\overset{N}{}}\frac{q_a^2}{\lambda r_a}=\frac{{\displaystyle \underset{a=1}{\overset{N}{}}}(\lambda \lambda _a)}{{\displaystyle \underset{a=1}{\overset{N}{}}}(\lambda r_a)}\frac{\mathrm{\Lambda }(\lambda )}{A(\lambda )}$$
(69)
An explicit formula for defining $`q_a^2`$ as a function of the $`\lambda `$’s , $`a`$, is obtained by applying the residue theorem to both members of the equation (69):
$`\mathrm{Res}\left(1Q_\lambda (\stackrel{}{\mathrm{q}})\right)|_{\lambda =r_a}=q_a^2`$
(70)
$`\mathrm{Res}\left({\displaystyle \frac{\mathrm{\Lambda }(\lambda )}{A(\lambda )}}\right)|_{\lambda =r_a}={\displaystyle \frac{\mathrm{\Lambda }(r_a)}{A^{}(r_a)}},A^{}(r_a)={\displaystyle \frac{dA(\lambda )}{d\lambda }}|_{\lambda =r_a}`$
Therefore:
$$q_a^2=\frac{\mathrm{\Lambda }(r_a)}{A^{}(r_a)}=\frac{{\displaystyle \underset{b=1}{\overset{N}{}}}(r_a\lambda _b)}{{\displaystyle \underset{b=1,ba}{\overset{N}{}}}(r_ar_b)}$$
(71)
and we see that the transformation $`(q_1,\mathrm{},q_N)(\lambda _1,\mathrm{},\lambda _N)`$ is $`2^N`$ to 1.
Inverting (71) to express $`\lambda _a`$ as a function of the $`q`$’s, $`a`$, requires one to solve an algebraic equation in $`\lambda _a`$ with powers up to $`\lambda _a^N`$. This is easy for $`N=2`$, possible, but very difficult for $`N=3,4`$ using Cardano’s formulas, and impossible if $`N5`$. For this reason, another derivation of (71) is useful, which in passing allows one to show identities between Cartesian and elliptic coordinates that make practical computations possible. To do this, notice that (69) implies:
$$\mathrm{\Lambda }(\lambda )=\underset{a=1}{\overset{N}{}}(\lambda r_a)+\underset{a=1}{\overset{N}{}}q_a^2\underset{b=1,ba}{\overset{N}{}}(\lambda r_b)$$
(72)
Setting $`\lambda =r_c`$ in (72) one immediately derives (71). More important, expanding the two members of (72) in a power series in $`\lambda `$ , we obtain:
$$\lambda ^N+\lambda ^{N1}\left(\underset{a=1}{\overset{N}{}}r_a+\underset{a=1}{\overset{N}{}}q_a^2\right)+\lambda ^{N2}\left(\underset{b<a}{\overset{N}{}}r_ar_b\underset{a=1}{\overset{N}{}}q_a^2\underset{b=1,ba}{\overset{N}{}}r_b\right)+\mathrm{}$$
$$=\lambda ^N\left(\underset{a=1}{\overset{N}{}}\lambda _a\right)\lambda ^{N1}+\left(\underset{b<a}{\overset{N}{}}\lambda _a\lambda _b\right)\lambda ^{N2}+\mathrm{}+(1)^N\underset{a=1}{\overset{N}{}}\lambda _a$$
Equalizing the coefficients of the terms with the same power of $`\lambda `$ in the last equation we have $`N`$ non trivial identities. We shall use the equalities between the coefficients of $`\lambda ^{N1}`$ and $`\lambda ^{N2}`$:
$$\underset{a=1}{\overset{N}{}}\lambda _a=\underset{a=1}{\overset{N}{}}r_a\underset{a=1}{\overset{N}{}}q_a^2$$
(73)
$$\underset{b<a}{\overset{N}{}}\lambda _a\lambda _b=\underset{a=1}{\overset{N}{}}r_aq_a^2\underset{a=1}{\overset{N}{}}q_a^2\underset{b=1}{\overset{N}{}}r_b+\underset{b<a}{\overset{N}{}}r_br_a$$
(74)
Another important tool using elliptic coordinates is the Jacobi lemma:
Lemma. The expression
$$\underset{a=1}{\overset{N}{}}\frac{\alpha _a^s}{(\alpha _a\alpha _1)(\alpha _a\alpha _2)\mathrm{}\stackrel{(a)}{\mathrm{}}\mathrm{}(\alpha _a\alpha _N)}$$
where $`(\alpha _1,\mathrm{},\alpha _N)`$ are $`N`$ real numbers such that $`\alpha _1<\alpha _2<\mathrm{}<\alpha _N`$ is equal to 0 if $`sN2`$ and $`1`$ for $`s=N1`$.
Proof:
Consider the function
$$f_s(z)=\frac{z^s}{{\displaystyle \underset{a=1}{\overset{N}{}}}(z\alpha _a)}$$
which has $`N`$ poles in the complex plane at $`z=\alpha _a`$, $`a`$, and another pole at infinity. If $`\gamma `$ is a closed curve which is the boundary of a region $`D`$ of the complex plane containing all the finite poles of $`f_s(z)`$, the residue theorem tells us that:
$$\frac{1}{2\pi i}_\gamma f_s(z)𝑑z=\underset{a=1}{\overset{N}{}}\mathrm{Res}(f_s)(\alpha _a)=\mathrm{Res}(f_s)(\mathrm{})=\frac{1}{2\pi i}_\gamma f_s(z)𝑑z$$
Also,
$`{\displaystyle \underset{a=1}{\overset{N}{}}}\mathrm{Res}(f_s)(\alpha _a)={\displaystyle \underset{a=1}{\overset{N}{}}}{\displaystyle \frac{\alpha _a^s}{(\alpha _a\alpha _1)(\alpha _a\alpha _2)\mathrm{}\stackrel{(a)}{\mathrm{}}\mathrm{}(\alpha _a\alpha _N)}}`$
$`\mathrm{Res}(f_s)(\mathrm{})=\{\begin{array}{c}0,s<N1\hfill \\ 1,s=N1\hfill \end{array}`$
and the lemma is proved.
Applying this result to the choice $`\alpha _a=\lambda _a`$, $`a=1,\mathrm{},N`$, we obtain new identities
$`{\displaystyle \underset{a=1}{\overset{N}{}}}{\displaystyle \frac{\lambda _a^s}{\mathrm{\Lambda }^{}(\lambda _a)}}=0,s<N1`$ (76)
$`{\displaystyle \underset{a=1}{\overset{N}{}}}{\displaystyle \frac{\lambda _a^{N1}}{\mathrm{\Lambda }^{}(\lambda _a)}}=1`$ (77)
because $`\mathrm{\Lambda }^{}(\lambda _a)=(\lambda _a\lambda _1)(\lambda _a\lambda _2)\mathrm{}\stackrel{(a)}{\mathrm{}}\mathrm{}(\lambda _a\lambda _N)`$. Alternatively if we take $`\alpha _a=r_a`$, the lemma implies that:
$$\underset{a=1}{\overset{N}{}}\frac{r_a^s}{A^{}(r_a)}=0,s<N1;\underset{a=1}{\overset{N}{}}\frac{r_a^{N1}}{A^{}(r_a)}=1$$
(78)
An important identity obtained from the lemma is:
$$\underset{a=1}{\overset{N}{}}\frac{q_a^2}{(r_a\lambda _b)(r_a\lambda _c)}=0,b,c$$
(79)
The normal vectors to the family of quadrics (67)
$$Q_\lambda (\stackrel{}{q})=\underset{a=1}{\overset{N}{}}\frac{q_a^2}{r_a\lambda }=1$$
at the point $`\stackrel{}{q}(q_1,\mathrm{},q_N)`$, are
$$\stackrel{}{n}(\lambda )(n_1(\lambda ),\mathrm{},n_N(\lambda ))=(\frac{q_1}{r_1\lambda },\mathrm{},\frac{q_N}{r_N\lambda })$$
Observe that (79) implies
$$\underset{a=1}{\overset{N}{}}n_a(\lambda _b)n_a(\lambda _c)=0,a,b$$
Therefore, all the quadrics are orthogonal with each other and the elliptic coordinates form an orthogonal system. The standard Euclidean metric in Cartesian coordinates can be expressed in elliptic coordinates in the form:
$$ds^2=\underset{a=1}{\overset{N}{}}dq_a^2=\underset{a=1}{\overset{N}{}}\underset{b=1}{\overset{N}{}}g_{ab}(\stackrel{}{\lambda }_E)d\lambda _ad\lambda _b$$
Derivation of the two members of equation (71) leads to:
$$\frac{2dq_a}{q_a}=(1)^N\underset{b=1}{\overset{N}{}}\frac{d\lambda _b}{r_a\lambda _b}$$
and, using the Jacobi Lemma,
$$4dq_a^2=q_a^2\underset{b=1}{\overset{N}{}}\frac{d\lambda _b^2}{r_a\lambda _b}$$
Finally, we have:
$$g_{aa}=\frac{1}{4}\frac{\mathrm{\Lambda }^{}(\lambda _a)}{A(\lambda _a)}=\frac{1}{4}\frac{{\displaystyle \underset{ba,b=1}{\overset{N}{}}}(\lambda _a\lambda _b)}{{\displaystyle \underset{b=1}{\overset{N}{}}}(\lambda _ar_b)},g_{ab}=0,ab$$
(80)
The kinetic energy of a natural dynamical system in elliptic coordinates is;
$$T=\frac{1}{2}\underset{a=1}{\overset{N}{}}\dot{q}_a^2=\frac{1}{2}\underset{a=1}{\overset{N}{}}g_{aa}\dot{\lambda }_a^2$$
(81)
In terms of the canonical momentum $`\pi _a=\frac{T}{\dot{\lambda }_a}`$, $`T`$ reads:
$$T=\underset{a=1}{\overset{N}{}}\pi _a\dot{\lambda }_aT=\frac{1}{2}\underset{a=1}{\overset{N}{}}\frac{1}{g_{aa}}\pi _a^2=2\underset{a=1}{\overset{N}{}}\frac{A(\lambda _a)}{\mathrm{\Lambda }^{}(\lambda _a)}\pi _a^2$$
(82)
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# SPECTRAL PROPERTIES OF NON-SELF-ADJOINT OPERATORS IN THE SEMI-CLASSICAL REGIME
## 1 Introduction
This work was motivated by the paper of Shkalikov concerning the analysis of the semi-classical Airy operator, together with our computer simulations of the associated discrete problem using the numerical package Matlab. Specifically, we examined the operator $`H_{\delta ,h}`$ given formally by
$$H_{\delta ,h}:=h^2\frac{\mathrm{d}^2}{\mathrm{d}x^2}+V_\delta \text{ on }L^2(1,1)$$
(1)
where
$$V_\delta (x):=\{\begin{array}{cc}i(x+\delta )\hfill & \text{ for }x>0\text{ }\hfill \\ i(x\delta )\hfill & \text{ for }x<0\text{ }\hfill \end{array}$$
and both $`\delta >0`$, $`h>0`$ are small. In it is shown rigorously that as $`h0`$ the spectrum of $`H_{0,h}`$ becomes dense inside an arbitrarily small neighbourhood of the Y-shaped subset of $`𝐂`$ defined by
$$[i,1/\sqrt{3}],[i,1/\sqrt{3}]\text{ and }[1/\sqrt{3},\mathrm{}),$$
where we use $`[\alpha ,\beta ]`$ to denote the line segment joining $`\alpha ,\beta 𝐂`$ (see figure 2).
This paper confirms our surprising numerical results, which suggested that when an arbitrarily small jump discontinuity is inserted into the otherwise linear (purely imaginary) potential, and then the semi-classical parameter $`h`$ allowed to go to $`0`$, the asymptotic spectrum of $`H_{\delta ,h}`$ turns out to be completely different from that of $`H_{0,h}`$ (see figure 1 and Corollary 4).
Several papers have been written about the operator $`H_{0,h}`$ \- a major motivation being that it is a model operator for the ‘Orr-Sommerfeld’ problem . The operator also defines the ‘Squire model for the Couette flow’ in hydrodynamics; and in its own right, defines the semigroup which is the solution of the so-called ‘Torrey equation’ , related to the diffusion of magnetic fields. Thus, although the spectrum of this non-self-adjoint operator displays sometimes strange and singular behaviour, it must not be dismissed as a ‘pathological’ example from pure mathematics since it has important applications - for example, in magnetic resonance imaging devices.
It is well known that a basis for solutions of the so-called ‘Airy equation’
$$f^{\prime \prime }(z)+zf(z)=0$$
is given by any two of the Airy functions $`Ai(z)`$, $`Ai(\mathrm{e}^{2\pi i/3}z)`$ and $`Ai(\mathrm{e}^{2\pi i/3}z)`$ (see ). Thus, the analytic investigation of the eigenvalues will involve examining the asymptotic behaviour of certain Airy functions, where the expressions we use are WKB-type approximations. We will show that the eigenvalues lie inside a certain subset of the complex plane, which is intimately related to the $`Stokes^{}`$ $`lines`$ (or $`principal`$ $`curves`$) of the problem (see ). The proof will depend upon showing that for all $`\lambda `$ outside this subset, the eigenvalue problem
$$H_hf(x)=\lambda f(x)$$
has a well-defined Green’s function. As in , our analysis uses the concept of the $`characteristic`$ $`determinant`$ which we will describe in Section 2.
In Section 4 we analyse the spectral behaviour in the semi-classical limit $`h0`$ for a general complex-valued, piecewise linear potential. The surprising result is that each linear segment of the potential gives rise to a characteristic ‘Y-shaped’ set of eigenvalues; and the spectrum of the operator is contained within the superposition of these ‘Y-shaped’ sets. From the point of view of applications, this means that whilst the asymptotic spectrum is theoretically computable for an idealised linear potential; in practice, any $`arbitrarily`$ small deviation from the ideal can $`completely`$ change the spectrum. One way of expressing this for an operator $`H`$ is in terms of the $`pseudospectral`$ $`sets`$:
$$\mathrm{Spec}_ϵ(H):=\mathrm{Spec}(H)\{z𝐂:(Hz)^1ϵ^1\},$$
i.e. the contour sets of the resolvent norm, with the convention that $`z\mathrm{Spec}(H)`$ implies
$$(Hz)^1:=\mathrm{}.$$
For many non-self-adjoint operators it has been demonstrated that the pseudospectral sets become very large as some parameter varies, even though $`z`$ may be far from the spectrum of the operator. This is equivalent to saying that the spectrum is computationally unstable. Our aim in this paper is to demonstrate for a relatively transparent case, the mechanism behind this phenomenon; we believe the results to be capable of extension to a more general class of piecewise analytic potentials.
In Section 5 we provide an analysis of the simultaneous limit as $`h0`$ and $`\delta 0`$ together.
## 2 The Characteristic Determinant
In this section we describe the characteristic determinant of the operator $`H_h`$ defined by
$$H_hf(x):=h^2\frac{\mathrm{d}^2f(x)}{\mathrm{d}x^2}+V(x)f(x)$$
acting on $`L^2(1,1)`$ with Dirichlet boundary conditions, $`h>0`$ small, and $`V(x)`$ the complex valued, $`n`$-times piecewise linear function
$$V(x):=\{\begin{array}{cc}m_1x+l_1\hfill & x_0x<x_1\hfill \\ m_2x+l_2\hfill & x_1<x<x_2\hfill \\ \mathrm{}\hfill & \mathrm{}\hfill \\ m_nx+l_n\hfill & x_{n1}<xx_n\hfill \end{array}$$
with $`1=x_0<x_1<\mathrm{}<x_n=1`$, and the $`m_i`$, $`l_i`$ $`i=1,\mathrm{},n`$ complex constants. The domain of the operator is given precisely by
$$\mathrm{Dom}(H_h)=\{fC[1,1]:f(1)=f(1)=0,f^{}C[1,1]\text{ and }f^{\prime \prime }L^2(1,1)\}$$
(2)
where the primes denote differentiation with respect to $`x`$, and $`f^{\prime \prime }`$ is initially to be interpreted in the distributional sense. A direct substitution shows that a basis of solutions for the differential equation
$$h^2f^{\prime \prime }(x)+(V(x)\lambda )f(x)=0$$
where $`V(x):=mx+l`$; and $`l`$, $`m`$ are complex constants, is given by any two of the Airy functions $`Ai(w)`$ and $`Ai(\mathrm{e}^{\pm 2\pi i/3}w)`$, where
$$w:=h^{2/3}m^{2/3}(V(x)\lambda ).$$
It follows that, in order to construct an eigenfunction of the operator $`H_h`$, we seek constants $`\alpha _{i1},\alpha _{i2}`$ $`i=1,\mathrm{},n`$ such that the function
$$f(x):=\{\begin{array}{cc}\alpha _{11}u_{11}(x)+\alpha _{12}u_{12}(x)\hfill & x_0x<x_1\hfill \\ \alpha _{21}u_{21}(x)+\alpha _{22}u_{22}(x)\hfill & x_1<x<x_2\hfill \\ \mathrm{}\hfill & \mathrm{}\hfill \\ \alpha _{n1}u_{n1}(x)+\alpha _{n2}u_{n2}(x)\hfill & x_{n1}<xx_n\hfill \end{array}$$
satisfies all of the domain conditions (2), where
$$u_{i1}(x):=Ai(\mathrm{e}^{2k\pi i/3}h^{2/3}m_i^{2/3}((m_ix+l_i)\lambda ))$$
(3)
with $`k\{1,0,1\}`$. For each $`i=1,\mathrm{},n`$, the functions $`u_{i2}`$ are defined similarly, except that a different choice of $`k`$ is to be taken from $`\{1,0,1\}`$.
In addition to satisfying the boundary conditions $`f(1)=f(1)=0`$, $`f`$ must also be continuously differentiable, even at the points $`x_i`$. From the power series definition \[2, p.54\], it is clear that the Airy functions $`Ai`$ are analytic on the whole of $`𝐂`$, and so the requirement that $`f`$ be continuously differentiable reduces to the $`2(n1)`$ simultaneous ‘matching’ conditions
$$\alpha _{i1}u_{i1}(x_i)+\alpha _{i2}u_{i2}(x_i)\alpha _{(i+1)1}u_{(i+1)1}(x_i+)\alpha _{(i+1)2}u_{(i+1)2}(x_i+)=0$$
and
$$\alpha _{i1}u_{i1}^{}(x_i)+\alpha _{i2}u_{i2}^{}(x_i)\alpha _{(i+1)1}u_{(i+1)1}^{}(x_i+)\alpha _{(i+1)2}u_{(i+1)2}^{}(x_i+)=0.$$
The boundary conditions $`f(1)=f(1)=0`$ demand that
$$\alpha _{11}u_{11}(1)+\alpha _{12}u_{12}(1)=0$$
and
$$\alpha _{n1}u_{n1}(1)+\alpha _{n2}u_{n2}(1)=0.$$
Thus finding a solution of the differential equation
$$h^2f^{\prime \prime }(x)+V(x)f(x)=\lambda f(x)$$
which satisfies all of the domain conditions (2), involves solving the matrix equation
$$\left(\begin{array}{ccccccc}u_{11}(1)& u_{12}(1)& 0& 0& 0& \mathrm{}& 0\\ u_{11}(x_1)& u_{12}(x_1)& u_{21}(x_1)& u_{22}(x_1)& 0& \mathrm{}& 0\\ u_{11}^{}(x_1)& u_{12}^{}(x_1)& u_{21}^{}(x_1)& u_{22}^{}(x_1)& 0& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& \mathrm{}& 0& u_{(n1)1}(x_{n1})& u_{(n1)2}(x_{n1})& u_{n1}(x_{n1})& u_{n2}(x_{n1})\\ 0& \mathrm{}& 0& u_{(n1)1}^{}(x_{n1})& u_{(n1)2}^{}(x_{n1})& u_{n1}^{}(x_{n1})& u_{n2}^{}(x_{n1})\\ 0& \mathrm{}& 0& 0& 0& u_{n1}(1)& u_{n2}(1)\end{array}\right)$$
$$\times \left(\begin{array}{c}\alpha _{11}\\ \alpha _{12}\\ \alpha _{21}\\ \alpha _{22}\\ \mathrm{}\\ \alpha _{n1}\\ \alpha _{n2}\end{array}\right)=\left(\begin{array}{c}0\\ 0\\ 0\\ 0\\ \mathrm{}\\ 0\\ 0\end{array}\right)$$
(4)
for the constants $`\alpha _{i1},\alpha _{i2}.`$ Note that we are taking the left- and right-hand limits at the points $`x_i`$, although here and subsequently we will abuse the notation in order to add clarity, and simply write $`u_{i1}(x_i)`$ etc. It is the determinant of the matrix on the left-hand side of (4) that we shall call the characteristic determinant of the eigenvalue problem defined by $`H_h`$.
## 3 Airy functions and Stokes’ lines
For our proofs in the next section, we will need some notation and basic properties of the Airy functions. In all that follows we let the argument function $`\mathrm{Arg}`$ take principal values i.e.
$$\mathrm{Arg}:𝐂(\pi ,\pi ].$$
If $`\mathrm{Arg}(\beta \alpha ):=\theta `$, the subsets $`Y(\alpha ,\beta )`$ of $`𝐂`$ are to be constructed as follows: take the lines
$$\alpha +r\mathrm{e}^{2\theta i/3}\text{ and }\{\begin{array}{cc}\beta +r\mathrm{e}^{2\theta i/3+2\pi i/3}\hfill & \text{ for }\theta <0\hfill \\ \beta +r\mathrm{e}^{2\theta i/32\pi i/3}\hfill & \text{ for }\theta 0\hfill \end{array}$$
as $`r`$ ranges in $`[0,\mathrm{})`$, to their point of intersection, $`\mathrm{\Gamma }`$ say. Then, from $`\mathrm{\Gamma }`$ extend the infinite line defined by the set of $`\lambda 𝐂`$ such that
$$\mathrm{Re}\{(\mathrm{e}^{2\theta i/3}(\alpha \lambda ))^{3/2}\}=\mathrm{Re}\{(\mathrm{e}^{2\theta i/3}(\beta \lambda ))^{3/2}\}.$$
(5)
The motivation for this set will become clear during our proofs; in fact, it will be seen to comprise a curve asymptotic to the line
$$\{z𝐂:\mathrm{Im}(z)=\frac{\mathrm{Im}(\alpha )+\mathrm{Im}(\beta )}{2}\}.$$
Note for now, however, that (5) is $`h`$-independent. The $`\epsilon `$-neighbourhood of any subset $`T`$ of $`𝐂`$ will be defined by
$$\mathrm{Nhd}(T;\epsilon ):=\{t+z:tT\text{ and }\left|z\right|<\epsilon \}.$$
We will use the well-known asymptotic expansion of the Airy function $`Ai`$, giving the WKB-type approximation:
$$Ai(z)=\frac{z^{1/4}}{2\sqrt{\pi }}\mathrm{exp}\left(\frac{2}{3}z^{3/2}\right)(1+O(z^{3/2}))$$
(6)
as $`\left|z\right|\mathrm{}`$, valid for $`\left|\mathrm{Arg}(z)\right|<\pi `$; and where the principal value of $`z^{3/2}`$ is taken. Following the notation of Olver (see \[2, p.413\]), we define
$$S_0:=\{z:\left|\mathrm{Arg}(z)\right|<\pi /3\}$$
$$S_1:=\{z:\pi /3<\mathrm{Arg}(z)<\pi \}$$
$$S_1:=\{z:\pi /3>\mathrm{Arg}(z)>\pi \}$$
(suffixes enumerated modulo 3). One can check that for all complex $`z`$ (and taking principal values), we have
$$\mathrm{Re}\{(\mathrm{e}^{2\pi i/3}z)^{3/2}\}=\{\begin{array}{cc}\mathrm{Re}\{(\mathrm{e}^{2\pi i/3}z)^{3/2}\}\hfill & \text{ for }zS_1S_1\hfill \\ \mathrm{Re}\{(\mathrm{e}^{2\pi i/3}z)^{3/2}\}\hfill & \text{ for }zS_0\hfill \end{array}$$
(7)
$$\mathrm{Re}\{(\mathrm{e}^{2\pi i/3}z)^{3/2}\}=\{\begin{array}{cc}\mathrm{Re}\{(z)^{3/2}\}\hfill & \text{ for }zS_1S_0\hfill \\ \mathrm{Re}\{(z)^{3/2}\}\hfill & \text{ for }zS_1\hfill \end{array}$$
(8)
and
$$\mathrm{Re}\{(\mathrm{e}^{2\pi i/3}z)^{3/2}\}=\{\begin{array}{cc}\mathrm{Re}\{(z)^{3/2}\}\hfill & \text{ for }zS_0S_1\hfill \\ \mathrm{Re}\{(z)^{3/2}\}\hfill & \text{ for }zS_1.\hfill \end{array}$$
(9)
Then, putting
$$Ai_k(z):=Ai(\mathrm{e}^{2k\pi i/3}z)$$
(10)
the asymptotics (6) show that as $`\left|z\right|\mathrm{}`$, $`\left|Ai_k(z)\right|`$ decreases exponentially for $`zS_k`$, and increases exponentially for $`zS_{k1}S_{k+1}`$. The boundaries of the sectors $`S_k`$ i.e. the rays $`t\mathrm{e}^{\pm \pi i/3}`$ and $`t\mathrm{e}^{\pi i}`$ for $`t[0,\mathrm{})`$, are known as the Stokes’ lines (or principal curves) of the problem \[2, p.503\]. Indeed, for the Airy equation
$$f^{\prime \prime }(z)+zf(z)=0$$
the Stokes’ lines are defined to be the values of $`z`$ such that
$$\mathrm{Re}_0^z\sqrt{t}dt=\mathrm{Re}\left\{\frac{2}{3}z^{3/2}\right\}=0,$$
and denote the boundaries of the $`principal`$ $`subdomains`$ $`S_1`$ etc., inside of which the asymptotic expression (6) is valid for each $`k`$.
We will call the suffix $`k`$ ‘allowable’ for any given $`z𝐂`$, if
$$\left|\mathrm{Arg}(\mathrm{e}^{2k\pi i/3}z)\right|<\pi .$$
(11)
## 4 The Semi-Classical Limit
The following is a generalisation of the argument used in for the potential $`V(x)=ix`$, and will form a lemma for our main theorem.
###### Proposition 1
Let $`V`$ be the complex valued linear potential given by
$$V(x)=mx+lx[1,1]$$
where $`m`$ and $`l`$ are complex constants; $`u_{11}(x)`$, $`u_{12}(x)`$ are as defined in (3), and $`a,b[1,1]`$, $`a<b`$.
Let $`\epsilon >0`$ be given and $`\lambda 𝐂`$. If
$$\lambda \mathrm{Nhd}(Y(V(a),V(b));\epsilon ),$$
then
$$u_{11}(b)u_{12}(a)=o(u_{11}(a)u_{12}(b))$$
(12)
as $`h0`$.
Proof A simple scaling and translation of the operator $`H_h`$ allows us, without loss of generality, to assume that $`a:=1`$, $`b:=1`$ and $`l=0`$. That is, we assume
$$V(x):=x\mathrm{e}^{i\theta },\text{ where }\theta :=\mathrm{Arg}(m).$$
By elementary trigonometry, one can check that we then have
$$\mathrm{\Gamma }=\mathrm{e}^{i\theta }+\frac{4}{\sqrt{3}}\mathrm{sin}\left(\frac{\left|\theta \right|}{3}\right)\mathrm{e}^{2(\theta \pi )i/3}$$
or
$$\mathrm{\Gamma }=\mathrm{e}^{i\theta }+\frac{4}{\sqrt{3}}\mathrm{sin}\left(\frac{2\pi }{3}+\frac{\left|\theta \right|}{3}\right)\mathrm{e}^{2\theta i/3}.$$
Recalling our definition of the Airy functions $`u_{11}(x)`$ and $`u_{12}(x)`$ (3), we put
$$z(h,\lambda ,x):=h^{2/3}\mathrm{e}^{2\theta i/3}(x\mathrm{e}^{\theta i}\lambda ),$$
and can rewrite (6) explicitly in terms of $`h`$. Then, taking the modulus we obtain
$$\left|Ai(z(h,\lambda ,x))\right|$$
$$=h^{1/6}\frac{\left|x\mathrm{e}^{\theta i}\lambda \right|}{2\sqrt{\pi }}^{1/4}\mathrm{exp}\left(\frac{2}{3h}\mathrm{Re}(\mathrm{e}^{2\theta i/3}(x\mathrm{e}^{\theta i}\lambda ))^{3/2}\right)(1+O(h))$$
(13)
as $`h0`$, valid for $`\left|\mathrm{Arg}(z(h,\lambda ,x))\right|<\pi `$. Therefore, in order to estimate the moduli of the Airy functions $`Ai_k(z(h,\lambda ,x))`$ in the limit $`h0`$, it is sufficient to examine the behaviour of the functions
$$x\mathrm{Re}\{(\mathrm{e}^{2k\pi i/3}z(h,\lambda ,x))^{3/2}\},x𝐑$$
for $`k=1,0,1.`$
The basic idea of our proof is to show that for $`\lambda `$ outside an arbitrarily small $`\epsilon `$-neighbourhood of $`Y(V(1),V(1))`$, one can assign allowable values of $`j`$ and $`k`$ from $`\{1,0,1\}`$ (in the sense of (11)) such that the inequalities
$$\mathrm{Re}\{(\mathrm{e}^{2j\pi i/3}z(h,\lambda ,1)^{3/2}\}<\mathrm{Re}\{(\mathrm{e}^{2j\pi i/3}z(h,\lambda ,1)^{3/2}\}$$
(14)
and
$$\mathrm{Re}\{(\mathrm{e}^{2k\pi i/3}z(h,\lambda ,1)^{3/2}\}<\mathrm{Re}\{(\mathrm{e}^{2k\pi i/3}z(h,\lambda ,1)^{3/2}\}$$
(15)
hold in the limit $`h0`$. This will then be enough, by (13), to ensure that $`u_{11}(1)u_{12}(1)`$ and $`u_{11}(1)u_{12}(1)`$ are of different orders of magnitude as $`h0`$, thus implying (12).
Using the statements of the previous section; for all values of $`\lambda `$ such that
$$z(h,\lambda ,\pm 1):=h^{2/3}\mathrm{e}^{2\theta i/3}(\pm \mathrm{e}^{\theta i}\lambda )$$
does not lie within an $`\epsilon `$-neighbourhood of any of the Stokes’ lines, and $`z(h,\lambda ,1)`$, $`z(h,\lambda ,1)`$ lie in different principal domains, one can always obtain (14) and (15), and the asymptotics (13) will be valid. However, for $`\lambda `$ lying in the sector having its apex at $`\mathrm{\Gamma }`$, and bounded by the rays
$$\mathrm{\Gamma }+r\mathrm{e}^{2\theta i/3}\text{ and }\{\begin{array}{cc}\mathrm{\Gamma }+r\mathrm{e}^{2\theta i/3+2\pi i/3}\hfill & \text{ for }\theta <0\hfill \\ \mathrm{\Gamma }+r\mathrm{e}^{2\theta i/32\pi i/3}\hfill & \text{ for }\theta 0\hfill \end{array}$$
$`r[0,\mathrm{})`$, it is easy to check that $`\mathrm{e}^{2k\pi i/3}z(h,\lambda ,\pm 1)`$ both lie in the same principal domain, for each $`k\{1,0,1\}`$. Then it is also straightforward to check that as $`x`$ ranges from $`1`$ to $`1`$, the function
$$x\mathrm{Re}\{(\mathrm{e}^{2k\pi i/3}z(\epsilon ,\lambda ,x))^{3/2}$$
which has been at the heart of our analysis, has a single maximum/minimum. Together with the identities (7), (8) and (9), this means that there $`will`$ be values of $`\lambda `$ such that equality holds in both (14) and (15) - no matter what choices of $`j`$ and $`k`$ are made. Thus, (and without loss of generality assuming $`j=k=0`$,) the set of $`\lambda `$ satisfying
$$\mathrm{Re}\{(\mathrm{e}^{2\theta i/3}(\mathrm{e}^{\theta i}\lambda ))^{3/2}\}=\mathrm{Re}\{(\mathrm{e}^{2\theta i/3}(\mathrm{e}^{\theta i}\lambda ))^{3/2}\}$$
(16)
lies in $`Y(V(1),V(1))`$. We now examine this set in more detail. Expanding the Taylor series, we have
$$(\mathrm{e}^{2\theta i/3}(\mathrm{e}^{\theta i}\lambda ))^{3/2}=i\lambda ^{3/2}\mathrm{e}^{\theta i}+\frac{3}{2}i\lambda ^{1/2}\frac{3}{8}i\lambda ^{1/2}\mathrm{e}^{\theta i}\frac{1}{16}i\lambda ^{3/2}\mathrm{e}^{2\theta i}+O(\lambda ^{5/2})$$
and
$$(\mathrm{e}^{2\theta i/3}(\mathrm{e}^{\theta i}\lambda ))^{3/2}=i\lambda ^{3/2}\mathrm{e}^{\theta i}\frac{3}{2}i\lambda ^{1/2}\frac{3}{8}i\lambda ^{1/2}\mathrm{e}^{\theta i}+\frac{1}{16}i\lambda ^{3/2}\mathrm{e}^{2\theta i}+O(\lambda ^{5/2})$$
as $`\left|\lambda \right|\mathrm{}`$. Dividing through by $`i`$, this means that (16) will hold if and only if
$$\mathrm{Im}\left\{\frac{3}{2}\lambda ^{1/2}\frac{\lambda ^{3/2}\mathrm{e}^{2\theta i}}{16}\right\}=O(\lambda ^{5/2})$$
as $`\left|\lambda \right|\mathrm{}`$. Putting $`\lambda :=\rho \mathrm{e}^{\varphi i}`$, this is equivalent to the requirement
$$\mathrm{sin}\left(\frac{\varphi }{2}\right)\frac{1}{24\rho ^2}\mathrm{sin}\left(2\theta \frac{3\varphi }{2}\right)=O(\rho ^6)$$
as $`\rho \mathrm{}`$. But then
$`\mathrm{Im}(\lambda )`$ $`=\rho \mathrm{sin}(\varphi )`$
$`=2\rho \mathrm{sin}\left({\displaystyle \frac{\varphi }{2}}\right)\mathrm{cos}\left({\displaystyle \frac{\varphi }{2}}\right)`$
$`=2\rho \mathrm{cos}\left({\displaystyle \frac{\varphi }{2}}\right)\left\{{\displaystyle \frac{1}{24\rho ^2}}\mathrm{sin}\left(2\theta {\displaystyle \frac{3\theta }{2}}\right)+O(\rho ^6)\right\}`$
$`=O(\rho ^1)`$
as $`\rho \mathrm{}`$. By our definition of $`\mathrm{\Gamma }`$, $`z(h,\mathrm{\Gamma },\pm 1)`$ lies at the intersection of two Stokes’ lines, and so
$$\mathrm{Re}\{(z(h,\mathrm{\Gamma },\pm 1))^{3/2}\}=0$$
showing that $`\mathrm{\Gamma }`$ certainly lies in the set of $`\lambda `$ satisfying (16). Therefore, we deduce that $`Y(V(1),V(1))`$ contains a curve from $`\mathrm{\Gamma }`$ asymptotic to the positive real-axis.
Finally, we must examine what happens when $`z`$ does lie on one of the Stokes’ lines.
Firstly, suppose $`\mathrm{Arg}(z(h,\lambda ,1))=\pi /3`$, corresponding to $`\lambda `$ lying on the ray centred at $`\mathrm{e}^{\theta i}`$ and passing through $`\mathrm{\Gamma }`$. Then $`k=0,1`$ are allowable, and one checks that if $`\lambda `$ lies on the segment $`[\mathrm{e}^{\theta i},\mathrm{\Gamma })`$, we have $`z(h,\lambda ,1)S_1`$. It follows by (8) that
$$\mathrm{Re}\{(\mathrm{e}^{2\pi i/3}z(h,\lambda ,1))^{3/2}\}=\mathrm{Re}\{(z(h,\lambda ,1))^{3/2}\},$$
and so (14) and (15) cannot hold. However, if $`\lambda `$ lies on that part of the ray which extends past $`\mathrm{\Gamma }`$ (but not $`\lambda =\mathrm{\Gamma }`$ itself,) then $`z(h,\lambda ,1)S_1`$, and
$$\mathrm{Re}\{(\mathrm{e}^{2\pi i/3}z(h,\lambda ,1))^{3/2}\}=\mathrm{Re}\{(z(h,\lambda ,1))^{3/2}\}$$
causing (14) and (15) to hold for $`j=0`$, $`k=1`$.
An entirely similar argument holds when $`\mathrm{Arg}(z(h,\lambda ,1))=\pi `$, corresponding to $`\lambda `$ lying on the ray centred at $`\mathrm{e}^{\theta i}`$ and passing through $`\mathrm{\Gamma }`$ using (7), with $`j=1`$ and $`k=1`$.
Finally, the case where $`\mathrm{Arg}(z(h,\lambda ,\pm 1))=\pi /3`$ is taken care of using (9), which shows that we may use allowable values $`1`$ and $`0`$ to obtain (14) and (15).
This completes the proof.
In the case $`\theta =\pi /2`$; $`\mathrm{\Gamma }=1/\sqrt{3}`$ lies on the real-axis, and the figure $`Y(i,i)`$ has three $`linear`$ ‘arms’. When $`\theta =0`$, the symmetric case, $`Y(1,1)`$ is the semi-infinite interval $`[1,\mathrm{})`$, as is well-known from the theory of self-adjoint operators.
We now give our main result.
###### Theorem 2
Let
$$H_hf(x):=h^2\frac{\mathrm{d}^2f(x)}{\mathrm{d}x^2}+V(x)f(x)$$
act on $`L^2(1,1)`$ with Dirichlet boundary conditions, where $`h>0`$ is small, and $`V(x)`$ is the complex valued n-times piecewise linear function
$$V(x):=\{\begin{array}{cc}m_1x+l_1\hfill & x_0x<x_1\hfill \\ m_2x+l_2\hfill & x_1<x<x_2\hfill \\ \mathrm{}\hfill & \mathrm{}\hfill \\ m_nx+l_n\hfill & x_{n1}<xx_n\hfill \end{array}$$
with $`1=x_0<x_1<\mathrm{}<x_n=1`$ and the $`m_i`$, $`l_i`$, $`i=1,\mathrm{},n`$ complex constants. We assume for each $`i`$ that if $`m_ix_i+l_i=m_{i+1}x_i+l_{i+1}`$, then $`m_im_{i+1}`$. Put $`\theta _i:=\mathrm{Arg}(m_i)`$, and, using our earlier notation
$$T:=\underset{i=1}{\overset{n}{}}Y(V(x_i),V(x_{i+1})).$$
Let $`\epsilon >0`$ and $`N𝐙^+`$ be given. Then
$$\mathrm{Spec}(H_h)\{z:\left|z\right|N\}\mathrm{Nhd}(T;\epsilon )$$
for all small enough $`h>0`$.
Proof Our proof involves an analysis of the behaviour of the characteristic-determinant, i.e. the left-hand side of (4), as $`h0`$. We give a proof for the case $`n=3`$; the general case follows by a similar argument. For $`n=3`$, the characteristic-determinant is given by
$$\left|\begin{array}{cccccc}u_{11}(1)& u_{12}(1)& 0& 0& 0& 0\\ u_{11}(x_1)& u_{12}(x_1)& u_{21}(x_1)& u_{22}(x_1)& 0& 0\\ u_{11}^{}(x_1)& u_{12}^{}(x_1)& u_{21}^{}(x_1)& u_{22}^{}(x_1)& 0& 0\\ 0& 0& u_{21}(x_2)& u_{22}(x_2)& u_{31}(x_2)& u_{32}(x_2)\\ 0& 0& u_{21}^{}(x_2)& u_{22}^{}(x_2)& u_{31}^{}(x_2)& u_{32}^{}(x_2)\\ 0& 0& 0& 0& u_{31}(1)& u_{32}(1)\end{array}\right|$$
(17)
and we must prove that for certain values of $`\lambda 𝐂`$, this determinant does not vanish as $`h0`$. Expanding (17), one obtains
$$\{(u_{11}(1)u_{12}(x_1)u_{12}(1)u_{11}(x_1))(u_{22}^{}(x_1)u_{21}(x_2)u_{21}^{}(x_1)u_{22}(x_2))\times $$
$$\times (u_{31}(1)u_{32}^{}(x_2)u_{31}^{}(x_2)u_{32}(1))\}$$
$$\{(u_{11}(1)u_{12}(x_1)u_{12}(1)u_{11}(x_1))(u_{22}^{}(x_1)u_{21}^{}(x_2)u_{21}^{}(x_1)u_{22}^{}(x_2))\times $$
$$\times (u_{31}(1)u_{32}(x_2)u_{31}(x_2)u_{32}(1))\}+$$
$$+\{(u_{11}(1)u_{12}^{}(x_1)u_{12}(1)u_{11}^{}(x_1))(u_{22}(x_1)u_{21}^{}(x_2)u_{21}(x_1)u_{22}^{}(x_2))\times $$
$$\times (u_{31}(1)u_{32}(x_2)u_{31}(x_2)u_{32}(1))\}$$
$$\{(u_{11}(1)u_{12}^{}(x_1)u_{12}(1)u_{11}^{}(x_1))(u_{22}(x_1)u_{21}(x_2)u_{21}(x_1)u_{22}(x_2))\times $$
$$\times (u_{31}(1)u_{32}^{}(x_2)u_{31}^{}(x_2)u_{32}(1))\}$$
(18)
where so far, no asymptotics are involved.
Now, let $`\epsilon >0`$ and $`N𝐙^+`$ be as given in the statement of the theorem. Taking any
$$\lambda \{z:\left|z\right|N\}\backslash \mathrm{Nhd}(T;\epsilon ),$$
we can use the results of Proposition 1 to show that (17) is non-zero in the limit as $`h0`$. Indeed, by the proof of Proposition 1, we can ensure that the asymptotic estimates
$$u_{12}(1)u_{11}(x_1)=o(u_{11}(1)u_{12}(x_1)),$$
$$u_{21}(x_1)u_{22}(x_2)=o(u_{22}(x_1)u_{21}(x_2))$$
and
$$u_{31}(x_2)u_{32}(1)=o(u_{32}(x_2)u_{31}(1))$$
hold, as $`h0`$. Then, using the standard asymptotic expansions of the Airy functions , which give
$$Ai(z)=\frac{z^{1/4}}{2\sqrt{\pi }}\mathrm{exp}\left(\frac{2}{3}z^{3/2}\right)(1+O(z^{3/2}))$$
(19)
and
$$Ai^{}(z)=\frac{z^{1/4}}{2\sqrt{\pi }}\mathrm{exp}\left(\frac{2}{3}z^{3/2}\right)(1+O(z^{3/2}))$$
(20)
as $`\left|z\right|\mathrm{}`$, valid for all $`z`$ such that $`\left|\mathrm{Arg}(z)\right|<\pi `$; we see that, if
$$z(h,\lambda ,x_i):=h^{2/3}m_i^{2/3}((m_ix_i+l_i)\lambda ),$$
then
$`{\displaystyle \frac{\mathrm{d}}{\mathrm{d}x}}Ai(z)`$ $`={\displaystyle \frac{\mathrm{d}z}{\mathrm{d}x}}Ai^{}(z)`$
$`={\displaystyle \frac{\mathrm{d}z}{\mathrm{d}x}}{\displaystyle \frac{z^{1/4}}{2\sqrt{\pi }}}\mathrm{exp}\left({\displaystyle \frac{2}{3}}z^{3/2}\right)(1+O(z^{3/2}))`$
$`={\displaystyle \frac{\mathrm{d}z}{\mathrm{d}x}}z^{1/2}{\displaystyle \frac{z^{1/4}}{2\sqrt{\pi }}}\mathrm{exp}\left({\displaystyle \frac{2}{3}}z^{3/2}\right)(1+O(z^{3/2}))`$
as $`\left|z\right|\mathrm{}`$. Comparing this last expression with (19), and using $`fg`$ to mean that
$$\frac{f(h)}{g(h)}1\text{ as }h0,$$
we obtain
$$\frac{\mathrm{d}}{\mathrm{d}x}Ai(z(h,\lambda ,x_i))h^1((m_ix_i+l_i)\lambda )^{1/2}Ai(z(h,\lambda ,x_i))\text{ as }h0.$$
Moreover, similar calculations show that
$$\frac{\mathrm{d}}{\mathrm{d}x}Ai(\mathrm{e}^{\pm 2\pi i/3}z(h,\lambda ,x_i))h^1((m_ix_i+l_i)\lambda ))^{1/2}Ai(z(h,\lambda ,x_i))\text{ as }h0.$$
Reverting to our notation of (3), we will write
$$u_{i1}^{}(x_i)h^1c_{i1}(x_i)u_{i1}(x_i)\text{ etc.}$$
(21)
as $`h0`$, where it is important to note that the $`c_{ij}(x_i)`$, $`i=1,\mathrm{},(n1)`$, $`j=1,2`$ are independent of $`h`$. Then, since the constant terms $`c_{ij}(x_i)`$ are negligible in magnitude compared with the exponential terms $`u_{ij}(x_i)`$ as $`h0`$, the relations (21) imply that we also have the estimates
$$u_{12}(1)u_{11}^{}(x_1)=o(u_{11}(1)u_{12}^{}(x_1)),$$
$$u_{21}^{}(x_1)u_{22}(x_2)=o(u_{22}^{}(x_1)u_{21}(x_2))$$
$$u_{31}^{}(x_2)u_{32}(1)=o(u_{32}^{}(x_2)u_{31}(1))$$
$$u_{21}^{}(x_1)u_{22}^{}(x_2)=o(u_{22}^{}(x_1)u_{21}^{}(x_2))$$
and
$$u_{21}(x_1)u_{22}^{}(x_2)=o(u_{22}(x_1)u_{21}^{}(x_2))$$
as $`h0`$. Returning to (18), we first use the above estimates (since we may ignore the sub-dominant term in each round-bracketed expression), and then the relations (21) again, to obtain the asymptotic estimate on the first of the curly-bracketed terms:
$$\{(u_{11}(1)u_{12}(x_1)u_{12}(1)u_{11}(x_1))(u_{22}^{}(x_1)u_{21}(x_2)u_{21}^{}(x_1)u_{22}(x_2))\times $$
$$\times (u_{31}(1)u_{32}^{}(x_2)u_{31}^{}(x_2)u_{32}(1))\}$$
$`u_{11}(1)u_{12}(x_1)u_{22}^{}(x_1)u_{21}(x_2)u_{31}(1)u_{32}^{}(x_2)`$
$`u_{11}(1)u_{12}(x_1)\epsilon ^{1/2}c_{22}(x_1)u_{22}(x_1)u_{21}(x_2)u_{31}(1)\epsilon ^{1/2}c_{32}(x_2)u_{32}(x_2)`$
$`=h^2[c_{22}(x_1)c_{32}(x_2)](u_{11}(1)u_{12}(x_1)u_{22}(x_1)u_{21}(x_2)u_{31}(1)u_{32}(x_2))`$
as $`h0`$. Similar estimates apply to each of the remaining three terms in (18) i.e.
$$\{(u_{11}(1)u_{12}(x_1)u_{12}(1)u_{11}(x_1))(u_{22}^{}(x_1)u_{21}^{}(x_2)u_{21}^{}(x_1)u_{22}^{}(x_2))\times $$
$$\times (u_{31}(1)u_{32}(x_2)u_{31}(x_2)u_{32}(1))\}$$
$$h^2[c_{22}(x_1)c_{21}(x_2)](u_{11}(1)u_{12}(x_1)u_{22}(x_1)u_{21}(x_2)u_{31}(1)u_{32}(x_2)),$$
$$\{(u_{11}(1)u_{12}^{}(x_1)u_{12}(1)u_{11}^{}(x_1))(u_{22}(x_1)u_{21}^{}(x_2)u_{21}(x_1)u_{22}^{}(x_2))\times $$
$$\times (u_{31}(1)u_{32}(x_2)u_{31}(x_2)u_{32}(1))\}$$
$$h^2[c_{12}(x_1)c_{21}(x_2)](u_{11}(1)u_{12}(x_1)u_{22}(x_1)u_{21}(x_2)u_{31}(1)u_{32}(x_2))$$
and
$$\{(u_{11}(1)u_{12}^{}(x_1)u_{12}(1)u_{11}^{}(x_1))(u_{22}(x_1)u_{21}(x_2)u_{21}(x_1)u_{22}(x_2))\times $$
$$\times (u_{31}(1)u_{32}^{}(x_2)u_{31}^{}(x_2)u_{32}(1))\}$$
$$h^2[c_{12}(x_1)c_{32}(x_2)](u_{11}(1)u_{12}(x_1)u_{22}(x_1)u_{21}(x_2)u_{31}(1)u_{32}(x_2))$$
as $`h0`$. Collecting these estimates together, we see that the characteristic determinant (17) tends asymptotically towards
$$h^2\{(c_{22}(x_1)c_{12}(x_1))(c_{32}(x_2)c_{21}(x_2))\}(u_{11}(1)u_{12}(x_1)u_{22}(x_1)u_{21}(x_2)u_{31}(1)u_{32}(x_2))$$
as $`h0`$. The Airy functions $`Ai(z)`$ have countably many negative real zeros, ; and so by our choice of $`\lambda `$ outside $`\mathrm{Nhd}(T;\epsilon )`$ together with the proof of Proposition 1, we are assured that none of the Airy functions $`u_{ij}(x_i)`$ vanishes. Therefore, the determinant (17) does not vanish in the limit as $`h0`$, provided the ‘constant’ terms
$$c_{22}(x_1)c_{12}(x_1)\text{ and }c_{32}(x_2)c_{21}(x_2).$$
(22)
Our choice of $`\lambda `$ ensures that each of the individual constant terms $`c_{ij}(x_i)`$ is non-zero. Moreover, reviewing the proof of Proposition 1 and the identities (7)-(9), we see that the choices for $`j`$ and $`k`$ are not uniquely determined. Therefore, it is always possible to ensure that (22) holds, even when $`V`$ is continuous at some or all of the $`x_is`$. For example, if it happens that $`V(x_1)=V(x_1+)`$, then we choose $`j`$ and $`k`$ so that the constants $`c_{12}(x_1)`$ and $`c_{22}(x_1)`$ take different signs (by the calculations immediately above (21)). Thus, we deduce that such $`\lambda `$ cannot be an eigenvalue.
It now just requires the following compactness argument to complete the proof. Let $`B(z;\epsilon )`$ denote the open ball centred at $`z`$, with radius $`\epsilon `$. Our argument so far shows that for any
$$\lambda \{z𝐂:\left|z\right|N\}$$
such that
$$B(\lambda ;\epsilon )T=\mathrm{}$$
we have
$$B(\lambda ;\epsilon )\mathrm{Spec}(H_h)=\mathrm{}$$
for all $`0<h<E_\lambda `$, where $`E_\lambda `$ is some positive constant dependent upon $`\lambda `$. Let
$$M:=\{z𝐂:\left|z\right|N,\text{ and }\mathrm{dist}(z,T)2\epsilon \},$$
so that $`M`$ is compact. Then for all $`\lambda M`$
$$B(\lambda ;\epsilon )\mathrm{Nhd}(T;\epsilon )=\mathrm{}$$
and so
$$M\underset{\lambda M}{}B(\lambda ;\epsilon ).$$
But by compactness this means that there exists a finite sub-covering
$$M\underset{r=1}{\overset{n}{}}B(\lambda _r;\epsilon _{\lambda _r}).$$
Taking $`E`$ to be $`\mathrm{min}(E_{\lambda _1},\mathrm{},E_{\lambda _n})>0`$, we deduce that for all $`0<h<E`$ we have
$$\mathrm{Spec}(H_h)M=\mathrm{}$$
and this is equivalent to the statement of the theorem.
###### Remark 3
An important but subtle point, to note is that the zeros of
$$(u_{11}(1)u_{12}(x_1)u_{22}(x_1)u_{21}(x_2)u_{31}(1)u_{32}(x_2))$$
(23)
as a function of $`\lambda `$, are $`not`$ the same as the zeros of (17). However, by a similar argument to that of Shkalikov (i.e. using (19)), one can readily show that along each of the bounded arms of the Y-shaped figures making up $`T`$, the zeros (eigenvalues) do converge as $`h0`$ to form a dense set. Finding an asymptotic expression for the density of spectral points along the infinite lines (in the direction of the positive real-axis) appears to be a much more difficult problem; and we have no results yet in that direction.
To illustrate Theorem 2, in figure 3 we show a Matlab plot of the discretised version of the problem where the potential
$$V(x):=\{\begin{array}{cc}2ix+i\hfill & \text{ for }1x<0\text{ }\hfill \\ (i+1)x\hfill & \text{ for }0<x1\hfill \end{array}$$
We now return to the operator $`H_{\delta ,h}`$ defined in the first section.
###### Corollary 4
Let $`H_{\delta ,h}`$ be the non-self-adjoint operator defined by
$$H_{\delta ,h}:=h^2\frac{\mathrm{d}^2}{\mathrm{d}x^2}+V_\delta (x)$$
acting on $`L^2(1,1)`$ with Dirichlet boundary conditions, $`h>0`$ , and
$$V_\delta (x):=\{\begin{array}{cc}i(x\delta )\hfill & \text{ for }x<0\text{ }\hfill \\ i(x+\delta )\hfill & \text{ for }x>0\hfill \end{array}$$
with $`\delta >0`$. Define $`S𝐂`$ to be the double Y-shaped figure given by the line segments
$$[i\delta ,1/2\sqrt{3}+i(1+2\delta )/2]$$
$$[i(1+\delta ),1/2\sqrt{3}+i(1+2\delta )/2]$$
together with
$$[1/2\sqrt{3}+i(1+2\delta )/2,+\mathrm{}),$$
and
$$[i\delta ,1/2\sqrt{3}i(1+2\delta )/2]$$
$$[i(1+\delta ),1/2\sqrt{3}i(1+2\delta )/2]$$
together with
$$[1/2\sqrt{3}i(1+2\delta )/2,+\mathrm{}).$$
Then, given any $`\epsilon >0`$ and $`N𝐙^+`$, we have
$$\mathrm{Spec}(H_{\delta ,h})\{z:\left|z\right|n\}\mathrm{Nhd}(S;\epsilon )$$
for small enough $`h>0`$ (see figure 1).
By analyticity, however, for fixed $`h>0`$ we have
$$\underset{\delta 0}{lim}\mathrm{Spec}(H_{\delta ,h})=\mathrm{Spec}(H_{0,h}).$$
Hence, $`lim_{h0}lim_{\delta 0}\mathrm{Spec}(H_{\delta ,h})`$ is contained within an arbitrarily small neighbourhood of the line segments
$$[i,1/\sqrt{3}],[i,1/\sqrt{3}]\text{ and }[1/\sqrt{3},\mathrm{})$$
(see figure 2). Thus, the operations of taking limits do not commute, in the sense that as sets
$$\underset{h0}{lim}\underset{\delta 0}{lim}\mathrm{Spec}(H_{\delta ,h})\underset{\delta 0}{lim}\underset{h0}{lim}\mathrm{Spec}(H_{\delta ,h}).$$
## 5 Simultaneous limits for $`H_{\delta ,h}`$
In response to questions which we have been asked, we give an analysis for the situation in which $`\delta `$ and $`h`$ of Corollary 4 are not independent of each other.
###### Proposition 5
Defining the operator $`H_{\delta ,h}`$ as above, and putting
$$\delta :=h^{1/p}$$
we have
$$\underset{h0}{lim}\mathrm{Spec}(H_{h^{1/p},h})=\underset{h0}{lim}\mathrm{Spec}(H_{0,h})\text{ if }0<p<1$$
and
$$\underset{h0}{lim}\mathrm{Spec}(H_{h^{1/p},h})=\underset{\delta 0}{lim}\underset{h0}{lim}\mathrm{Spec}(H_{\delta ,h})\text{ if }p1.$$
Proof Referring to (4) and expanding, we see that the characteristic determinant of $`H_{\delta ,h}`$ is given by
$$\left\{(u_{11}(1)u_{12}(0)u_{12}(1)u_{11}(0))(u_{22}^{}(0)u_{21}(1)u_{21}^{}(0)u_{22}(1))\right\}$$
$$\left\{(u_{11}(1)u_{12}^{}(0)u_{12}(1)u_{11}^{}(0))(u_{22}(0)u_{21}(1)u_{21}(0)u_{22}(1))\right\},$$
(24)
whereas the characteristic determinant of $`H_{0,h}`$ is given by
$$u_{21}(1)u_{12}(1)u_{22}(1)u_{11}(1).$$
(25)
Now, putting
$`u_{12}(0):=Ai_k(h^{2/3}\mathrm{e}^{\pi i/3}(i\delta \lambda ))`$ and $`u_{22}(0):=Ai_k(h^{2/3}\mathrm{e}^{\pi i/3}(i\delta \lambda ))`$
it is clear by analyticity, that
$$u_{12}(0)u_{22}(0)\text{ and }u_{11}(0)u_{21}(0)$$
(26)
as $`\delta 0`$. Moreover, the calculations preceding (21) show that
$$u_{i1}^{}(0)\delta ^p(\lambda )^{1/2}u_{i1}(0)\text{ and }u_{i2}^{}(0)\delta ^p(\lambda )^{1/2}u_{i2}(0)$$
(27)
as $`\delta 0`$, for $`i=1,2`$. Therefore, using first (27) and then (26), the characteristic determinant (24) tends asymptotically towards
$$2\delta ^p(\lambda )^{1/2}\left(u_{11}(1)u_{12}(0)u_{21}(0)u_{22}(1)u_{12}(1)u_{11}(0)u_{22}(0)u_{21}(1)\right)$$
$$2\delta ^p(\lambda )^{1/2}\left(u_{22}(1)u_{11}(1)u_{21}(1)u_{12}(1)\right)$$
as $`\delta 0`$. Then the zeros of (24) tend asymptotically toward the zeros of (25) by $`\mathrm{Rouch}\stackrel{´}{\mathrm{e}}`$’s theorem, explaining the behaviour of
$$\underset{h0}{lim}\underset{\delta 0}{lim}\mathrm{Spec}(H_{\delta ,h}).$$
Substituting $`\delta =h^{1/p}`$, the character of $`lim_{h0}\mathrm{Spec}(H_{h^{1/p},h})`$ therefore depends upon the range of $`p`$ for which
$$\frac{u_{22}(0)}{u_{12}(0)}1\text{ and }\frac{u_{21}(0)}{u_{11}(0)}1$$
as $`h0`$. Now, without loss of generality, and using our earlier notation, let
$$\frac{u_{22}(0)}{u_{12}(0)}:=\frac{Ai_1(z_1)}{Ai_1(z_2)}$$
where
$$z_1:=h^{2/3}\mathrm{e}^{\pi i/3}(ih^{1/p}\lambda )\text{ and }z_2:=h^{2/3}\mathrm{e}^{\pi i/3}(ih^{1/p}\lambda )$$
so that, using the standard asymptotics (6)
$`{\displaystyle \frac{u_{22}(0)}{u_{12}(0)}}`$ $`={\displaystyle \frac{z_1^{1/4}\mathrm{exp}\left(\frac{2}{3}z_1^{3/2}\right)(1+O(z_1^{3/2}))}{z_2^{1/4}\mathrm{exp}\left(\frac{2}{3}z_2^{3/2}\right)(1+O(z_2^{3/2}))}}`$
$`=\left({\displaystyle \frac{z_1}{z_2}}\right)^{1/4}\mathrm{exp}\left({\displaystyle \frac{2}{3}}\left[z_1^{3/2}z_2^{3/2}\right]\right)(1+O(z_1^{3/2}))`$
$`\mathrm{exp}\left({\displaystyle \frac{2}{3}}\left[z_1^{3/2}z_2^{3/2}\right]\right)`$
as $`h0`$. But
$`z_1^{3/2}z_2^{3/2}`$ $`=h^1\mathrm{e}^{\pi i/3}\left\{(ih^{1/p}\lambda )^{3/2}(ih^{1/p}\lambda )^{3/2}\right\}`$
$`=h^1\mathrm{e}^{\pi i/3}(\lambda )^{3/2}\left\{(1ih^{1/p}/\lambda )^{3/2}(1+ih^{1/p}/\lambda )^{3/2}\right\}`$
$`=h^1\mathrm{e}^{\pi i/3}(\lambda )^{3/2}\left\{(13ih^{1/p}/2\lambda +\mathrm{})(1+3ih^{1/p}/2\lambda +\mathrm{})\right\}`$
$`=h^1\mathrm{e}^{\pi i/3}(\lambda )^{3/2}(3ih^{1/p}/\lambda +\mathrm{})`$
$`0`$
as $`h0`$ if and only if $`0<p<1`$. So, provided $`0<p<1`$
$$\frac{u_{22}(0)}{u_{12}(0)}1\text{ as }h0$$
and a similar calculation shows that we then also have
$$\frac{u_{21}(0)}{u_{11}(0)}1\text{ as }h0,$$
as required.
Acknowledgements I would like to thank E. B. Davies for suggesting this problem, and his guidance in the course of solving it.
Department of Mathematics
King’s College
Strand
London WC2R 2LS
England
e-mail:Redparth@mth.kcl.ac.uk
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# Charge on the quantum dot in the presence of tunneling current
## Abstract
The calculation of the charge present in central region of the double barrier structure at non-equilibrium conditions is discussed. We propose here a simple method to calculate non equilibrium Green’s functions which allows consistent calculations of retarded and distribution functions. To illustrate the approach we calculate the charge on the quantum dot coupled via tunnel barriers to two external leads having different chemical potentials $`\mu _L`$ and $`\mu _R`$. The obtained results have been compared with other approaches existing in the literature. They all agree in the equilibrium situation and the departures grow with increasing the difference $`\mu _L\mu _R`$.
Keywords: A. nanostructures; D. electron - electron interactions; D. tunneling
Recent advancements of nanotechnology allow the design and study of devices in which quantum phenomena play primary role . Typical device consists of a small central region coupled via tunnel barriers to two external electrodes (leads). In such quantum-dot devices Coulomb blockade, resonant tunneling and Kondo effects have been observed. Theoretical description of electron transport phenomena in quantum dot often requires (self-consistent) calculations of the charge accumulated in the central region . If the system is in an equilibrium state (i.e. the chemical potentials of the left and right leads are the same $`\mu _L=\mu _R=\mu `$) the calculation of charge on the dot is easy. One finds the (local) density of states for the dot $`N_d(\epsilon )`$ and integrates it with appropriate distribution function, which in equilibrium coincides with Fermi-Dirac one; $`f(\epsilon )=(e^{\beta (\epsilon \mu )}+1)^1`$, with $`\beta =1/k_BT`$ and gets
$$n=\underset{\mathrm{}}{\overset{\mathrm{}}{}}𝑑\epsilon N(\epsilon )f(\epsilon ).$$
(1)
The easiest way to derive (1) is to find the (equilibrium) Green’s function of the dot and apply spectral theorem to get $`n`$. In non equilibrium situation ($`\mu _L\mu _R`$) one can not use the equilibrium technique. Physically the problem is much more complicated as the tunneling current flows across the system and there is no a priori simple prescription as how to calculate the charge on the central region of the device (dot). Some additional complications arise if time dependent voltage or other time dependent fields are applied to the system. Here, however, we shall consider non-equilibrium but steady state. At nonequilibrium one can think about at least two ways of generalizing Eq. (1). The first is to replace the chemical potential $`\mu `$ by suitably chosen average chemical potential $`\overline{\mu }`$. The other is to replace $`f(\epsilon )`$ by the average distribution function $`\overline{f}(\epsilon )`$. The simplest two possibilities, which in fact have been used in literature to discuss various issues connected with transport across the system at hand are the following:
$$\overline{\mu }=\frac{\mu _L+\mu _R}{2}$$
(2)
or
$$\overline{f}(\epsilon )=\frac{f_L(\epsilon )+f_R(\epsilon )}{2},$$
(3)
where $`f_{L/R}=(e^{\beta (\epsilon \mu _{L/R})}+1)^1`$, is the equilibrium distribution function for the left/right lead. The immediate objection against the above proposals is that both give slightly different results, except at equilibrium and none of them is well justified. Equal weights assigned to $`\mu _{L(R)}`$ or $`f_{L(R)}(\epsilon )`$ may possibly be also justified for symmetrically coupled systems only.
Below we shall present an approach which allows the correct calculation of the charge $`n`$ in non equilibrium systems. To make the presentation of it as simple as possible we shall resort to the specific model describing the quantum dot coupled to two external leads. We use here the Anderson-Hubbard model
$$H=\underset{\lambda k\sigma }{}(\epsilon _{\lambda k}\mu _\lambda )c_{\lambda k\sigma }^+c_{\lambda k\sigma }+E_d\underset{\sigma }{}d_\sigma ^+d_\sigma +Un_{}n_{}+\underset{\lambda k\sigma }{}(V_{\lambda k}c_{\lambda k\sigma }^+d_\sigma +\mathrm{h}.\mathrm{c}.).$$
(4)
Here $`\lambda =R,L`$ denote the right ($`R`$) or left ($`L`$) lead in the system. The parameters have the following meaning: $`c_{\lambda k\sigma }^+(c_{\lambda k\sigma })`$ denote creation (annihilation) operator for a conduction electron with wave vector $`\stackrel{}{k}`$, spin $`\sigma `$ in the lead $`\lambda `$, $`V_{\lambda k}`$ is the hybridization matrix element between conduction electron of energy $`\epsilon _{\lambda k}`$ in a lead $`\lambda `$ with chemical potential $`\mu _\lambda `$ and localized electron on the dot. $`E_d`$ is the single particle energy of electrons in the dot. $`n_{}=d_{}^+d_{}`$ is the number operator for electrons with spin up localized on the dot and $`U`$ is the (repulsive) interaction energy between two electrons on the dot.
The correct way to calculate the charge on the dot under non equilibrium conditions is to use non equilibrium Green’s function (GF) of Keldysh . In this technique the average charge on dot at time $`t`$ is given by
$$n(t)=\underset{\sigma }{}d_\sigma ^+(t)d_\sigma (t)=i\underset{\sigma }{}G_\sigma ^<(t,t).$$
(5)
where $`G_\sigma ^<(t,t)`$ is diagonal (in time indices) element of the Keldysh ”lesser” GF
$$G_\sigma ^<(t,t^{})=id_\sigma ^+(t)d_\sigma (t^{}).$$
(6)
This function is the (1, 2) matrix element of the following contour ordered Green’s function matrix
$$\widehat{G}(t,t^{})=iT_cd(\tau )d^+(\tau ^{})=\left(\begin{array}{cc}G_c,\hfill & G^<\hfill \\ G^>,\hfill & G_{\overline{c}}\hfill \end{array}\right),$$
(7)
where $`T_c`$ is a complex time contour ordering operator. The time contour starts above time axis at $`t_0=\mathrm{}`$ passes through $`t`$ and $`t^{}`$ and returns back to $`t_0=\mathrm{}`$ below time axis. $`G_c=iTd(t)d^+(t^{})`$ is the usual time ordered (causal) GF for both times on the upper branch of the contour, while $`G_{\overline{c}}(t,t^{})`$ is antitime ordered GF with both time arguments on the lower branch. Greater ($`G^>`$) and lesser ($`G^<`$) GFs are distribution functions: $`G^>(t,t^{})=id(t)d^+(t^{})`$. Not all this functions are independent.
It turns out that the calculation of the ”lesser” GF is easy for noninteracting quantum dot i.e. for $`U=0`$. It can be found directly in the following way. Because, the Keldysh contour ordered GF, possesses the same perturbation expansion as the corresponding equilibrium GF it is possible to write down the matrix Dyson equation
$$\widehat{G}=\widehat{G}_0+\widehat{G}_0\widehat{\mathrm{\Sigma }}\widehat{G}.$$
(8)
The theorem due to Langreth allows the transition from contour ordered to real axis functions and one gets
$$G^<=(1+G^r\mathrm{\Sigma }^r)G_0^<(1+\mathrm{\Sigma }^aG^a)+G^r\mathrm{\Sigma }^<G^a$$
(9)
where $`G^{r(a)}`$ denotes retarded (advanced) GF.
For noninteracting quantum dot ($`U=0`$) exact expression for the retarded dot GF reads
$$G_\sigma ^r(\omega )=G_{0\sigma }^r(\omega )+G_{0\sigma }^r(\omega )\mathrm{\Sigma }^r(\omega )G_\sigma ^r(\omega )$$
(10)
with
$$\mathrm{\Sigma }^r(\omega )=\underset{\lambda k}{}|V_{\lambda k}|^2G_{0\lambda k}^r(\omega ),$$
(11)
where $`G_{0\lambda k}^r(\omega )=(\omega \epsilon _{\lambda k}+i0)^1`$ is the Green’s function of the lead $`\lambda `$. From this one immediately finds
$`\mathrm{\Sigma }^<(\omega )`$ $`=`$ $`{\displaystyle \underset{\lambda k}{}}V_{\lambda k}^{}{}_{}{}^{2}G_{0\lambda k}^<(\omega )=i{\displaystyle \underset{\lambda k}{}}|V_{\lambda k}|^2\delta (\omega \epsilon _{\lambda k})f_\lambda (\omega )=`$ (12)
$`=`$ $`2\pi i{\displaystyle \underset{\lambda }{}}\mathrm{\Gamma }_\lambda (\omega )f_\lambda (\omega ),`$
where we have used symbol $`\mathrm{\Gamma }_\lambda (\omega )`$ to denote the average coupling of the dot to lead $`\lambda `$; $`\mathrm{\Gamma }_\lambda (\omega )=2\pi _k|V_{\lambda k}|^2\delta (\omega \epsilon _{\lambda k})`$. Noting that $`G^<(t,t)=\frac{d\omega }{2\pi }G^<(\omega )`$ and using formula (9) we get
$$n=\underset{\sigma }{}𝑑\omega G^r(\omega )\mathrm{\Sigma }^<(\omega )G^a(\omega ),$$
(13)
which with help of Eqs. (10) and (12) we further rewrite as
$$n=\underset{\sigma }{}𝑑\omega \frac{\underset{\lambda }{}\mathrm{\Gamma }_\lambda (\omega )f_\lambda (\omega )}{\underset{\lambda }{}\mathrm{\Gamma }_\lambda (\omega )}\left(\frac{1}{\pi }\right)\mathrm{Im}G_\sigma ^r(\omega ).$$
(14)
Note that for symmetric coupling $`\mathrm{\Gamma }_L(\omega )=\mathrm{\Gamma }_R(\omega )`$ this formula reduces to the form (1) with average distribution function $`\overline{f}(\epsilon )`$ as given by (3). Such an expression for dot occupation has been used in for an interacting dot. We stress, however, that it is valid for noninteracting and symmetrically coupled dot only.
As one can see from the above the main difficulty in calculation of the charge on the dot is connected with calculation of the ”lesser” self energy $`\mathrm{\Sigma }^<`$. In the noninteracting case the knowledge of the exact expression for self energy (11) allowed easy calculation of exact $`\mathrm{\Sigma }^<`$ and exact expression for the charge. For interacting dot there is no way to get self energy exactly and one has to approximate it appropriately. One such approximation has been introduced by Ng and will be discussed latter on. Here we propose to use recently derived equation of motion method and calculate $`G^<(\omega )`$ with the desired accuracy and with approximations consistent with those made in calculation of the current or other response functions of the system.
To make the algebra easy we consider $`U=\mathrm{}`$ limit in the Hamiltonian (4) and use slave boson technique to handle it. We thus rewrite (4) with help of slave boson operators $`b,b^+`$ and use the commutation rules of LeGuillou and Ragoucy to evaluate the necessary quantum brackets. These rules treat exactly the local constraint which for $`U=\mathrm{}`$ prevents double occupancy of the dot. The procedure is simple. One rewrites the Hamiltonian in the form
$$H^{SB}=\underset{\lambda k\sigma }{}(\epsilon _{\lambda k}\mu _\lambda )c_{\lambda k\sigma }^+c_{\lambda k\sigma }+\epsilon _d\underset{\sigma }{}f_\sigma ^+f_\sigma +\underset{\lambda k\sigma }{}V_{\lambda k}(c_{\lambda k\sigma }^+b^+f_\sigma +f_\sigma ^+bc_{\lambda k\sigma })$$
(15)
and calculates the on-dot Green’s function $`D_\sigma ^<(\omega )=b^+f_\sigma |f_\sigma ^+b_\omega ^<`$ using equation (c.f. equation (28b) of )
$$A|B=g^<(\omega )[A,B]_\pm +g^r(\omega )[A,H_I]|B_\omega ^<+g^<(\omega )[A,H_I]|B_\omega ^a,$$
(16)
with the third term of $`H^{SB}`$ taken as a interaction part $`H_I`$ and lower case GFs being the corresponding GFs of the free Hamiltonian $`H_0`$ (consisting of first and second terms of $`H^{SB}`$). The higher order GFs appearing at this stage have been calculated in similar way. In the process we have used the same kind of factorization which one uses calculating the on-dot retarded GF necessary to get Kondo effect. Explicitly we neglected the GFs like $`<<c_{\lambda k\sigma }c_{\lambda ^{}k^{}\sigma }f_\sigma ^+b|f_\sigma ^+b>>^a`$ as they describe higher order spin correlations in the leads and performed the decoupling
$$<<c_{\lambda ^{}k^{}\sigma }^+c_{\lambda k\sigma }b^+f_\sigma |f_\sigma ^+b>>^af(\epsilon _{\lambda k})<<b^+f_\sigma |f_\sigma ^+b>>^a\delta _{kk^{}}\delta _{\lambda \lambda ^{}}.$$
(17)
The resulting ”lesser” self-energy takes on simple form
$$\mathrm{\Sigma }^<(\omega )=2\pi i\underset{\lambda }{}\mathrm{\Gamma }_\lambda (\omega )(1+f_\lambda (\omega ))f_\lambda (\omega ).$$
(18)
This is our final expression for $`\mathrm{\Sigma }^<(\omega )`$ to be used in calculations of the charge $`n`$.
Now let us rederive the expression for $`\mathrm{\Sigma }^<(\omega )`$ using the earlier mentioned approximation introduced by Ng and often used in the literature . This author proposes to assume that $`\mathrm{\Sigma }^<(\omega )=A\mathrm{\Sigma }_0^<(\omega )`$ and $`\mathrm{\Sigma }^>(\omega )=A\mathrm{\Sigma }_0^>`$ where $`A`$ is an unknown function and $`\mathrm{\Sigma }_0^<(\omega )`$ is noninteracting (and thus known) self-energy. For the present example it is given explicitly by (12). This together with the exact relations: $`\mathrm{\Sigma }^>\mathrm{\Sigma }^<=\mathrm{\Sigma }^r\mathrm{\Sigma }^a`$ and $`\mathrm{\Sigma }_0^>\mathrm{\Sigma }_0^<=\mathrm{\Sigma }_0^r\mathrm{\Sigma }_0^a`$ leads to
$$\mathrm{\Sigma }_{\mathrm{Ng}}^<(\omega )=\frac{\mathrm{\Sigma }^r\mathrm{\Sigma }^a}{\mathrm{\Sigma }_0^r\mathrm{\Sigma }_0^a}\mathrm{\Sigma }_0^<(\omega )$$
(19)
or in explicit form
$$\mathrm{\Sigma }_{\mathrm{Ng}}^<(\omega )=2\pi i\frac{\underset{\lambda }{}\mathrm{\Gamma }_\lambda (\omega )f_\lambda (\omega )}{\underset{\lambda }{}\mathrm{\Gamma }_\lambda (\omega )}\underset{\lambda }{}\mathrm{\Gamma }_\lambda (1+f_\lambda (\omega )),$$
(20)
which again contains the characteristic average distribution function, which for the symmetric coupling reduces to $`\overline{f}(\epsilon )`$. Note that this formula agrees with (18) only in equilibrium situation $`\mu _L=\mu _R=\mu `$.
It turns out that our formula, Eq. (18) can be rederived by suitably generalizing the Ng’s ansatz. To see this note that noninteracting self energy $`\mathrm{\Sigma }_0^<(\omega )`$ above can be written as a sum of pieces, each of which is connected with one of the leads i.e. $`\mathrm{\Sigma }_0^<(\omega )=\mathrm{\Sigma }_\lambda \mathrm{\Sigma }_{0\lambda }^<(\omega )`$. One then expects that due to locality of $`U`$ term in Hamiltonian the interacting self-energy will also be a sum of contributions from different leads $`\mathrm{\Sigma }^<(\omega )=_\lambda \mathrm{\Sigma }_\lambda ^<(\omega )`$. Generalising the ansatz and writing $`\mathrm{\Sigma }_\lambda ^<(\omega )=A_\lambda \mathrm{\Sigma }_{0\lambda }^<`$ independently for each $`\lambda `$ one immediately reproduces result (18). This shows that the application of equation of motion technique is a correct way to to derive the ”lesser” self energy. This technique, contrary to approximate schemes, allows mutually consistent calculations of both retarded and distribution GFs in Keldysh technique.
We have numerically calculated the charge $`n`$ on the interacting quantum dot (with $`U=\mathrm{}`$) coupled to leads with different chemical potentials using four formulae discussed in this paper. The results are presented in figures (1) and (2). Figure (1) shows the charge on the dot with $`E_d=1meV`$ as a function of the position of right lead chemical potential $`\mu _R`$ with $`\mu _L=0`$ and for symmetric couplings $`\mathrm{\Gamma }_L(\omega )=\mathrm{\Gamma }_R(\omega )=\mathrm{\Gamma }=17meV`$. Note tha for symmetric coupling the Ng’s approximation give the same results as average distribution function $`\overline{f}(\epsilon )`$ . In more general case of asymmetric coupling these two approaches lead to different results, but we abandoned this complication here. Solid line show results obtained with help of our formula (18), dashed one is calculated with help of (20), while dotted from (1) with $`\overline{\mu }`$ given by (2). The crossing of the lines for $`\mu _L=\mu _R=0`$ corresponds to equilibrium situation. The differences between the curves get larger for increasing voltage $`V=(\mu _L\mu _R)/e`$ applied to the system. The large (small) differences between various curves for negative (positive) values of $`\mu _R`$ is due to large (small) values of the on-dot density of states in the energy window between $`\mu _L`$ and $`\mu _R`$.
Figure (2) shows the charge on the dot calculated for constant difference $`\mu _L\mu _R=1meV`$ as a function of the parameter $`E_d`$. For $`E_d`$ much higher or much lower than both chemical potentials the differences are small as expected because the density of states around $`\mu _L`$ and $`\mu _R`$ is small. The largest differences of the calculated charge are found for such values of $`E_d`$ for which the density of states of electrons on the dot possesses the Kondo resonance.
We conclude by stressing that the calculation of charge on the dot is a necessary step towards self consistent calculations of the current flowing across quantum dot devices. In some parameter ranges all formulae give only slightly differing results. Generally, however, the differences between values of $`n`$ calculated by means of different formulae may be as large as 20-30%. The equation of motion method allows calculation of the lesser self energy consistent with retarded one. This is important as the current flowing in the system is, for proportionate couplings $`\mathrm{\Gamma }_L(\omega )=const\mathrm{\Gamma }_R(\omega )`$, expressed by the retarded Green’s function which in turn depends on the charge $`n`$ on the dot. For general couplings the consistency of approximations is even more important as the current depends on both retarded and lesser GFs and to be consistent one has to treat both self energies on equal footing.
This work has been partially supported by the Committee for Scientific Research under grant 2P03B 106 17.
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# Untitled Document
MAHLER MEASURE, LINKS AND HOMOLOGY GROWTH
Daniel S. Silver and Susan G. Williams
Both authors partially supported by NSF grant DMS-9704399.2000 Mathematics Subject Classification. Primary 57M25; secondary 37B10.
ABSTRACT. Let $`l`$ be an oriented link of $`d`$ components with nonzero Alexander polynomial $`\mathrm{\Delta }(u_1,\mathrm{},u_d)`$. Let $`\mathrm{\Lambda }`$ be a finite-index subgroup of $`H_1(S^3l)𝐙^d`$, and let $`M_\mathrm{\Lambda }`$ be the corresponding abelian cover of $`S^3`$ branched along $`l`$. The growth rate of the order of the torsion subgroup of $`H_1(M_\mathrm{\Lambda })`$, as a suitable measure of $`\mathrm{\Lambda }`$ approaches infinity, is equal to the Mahler measure of $`\mathrm{\Delta }`$.
1. Introduction. Associated to any knot $`kS^3`$ is a sequence of Alexander polynomials $`\mathrm{\Delta }_i,i1`$, in a single variable such that $`\mathrm{\Delta }_{i+1}`$ divides $`\mathrm{\Delta }_i`$. Likewise, for any oriented link of $`d`$ components there is a sequence of Alexander polynomials in $`d`$ variables. Following the usual custom, we refer to the first Alexander polynomial of a knot or a link as the Alexander polynomial, and we denote it simply by $`\mathrm{\Delta }`$.
In \[Go\] C. McA. Gordon examined the homology groups of $`r`$-fold cyclic covers $`M_r`$ of $`S^3`$ branched over a knot $`k`$. He proved that when each zero of the Alexander polynomial $`\mathrm{\Delta }`$ of $`k`$ has modulus one (and hence is a root of unity), the finite values of $`|H_1(M_r)|`$ are periodic in $`r`$. Gordon conjectured that when some zero of $`\mathrm{\Delta }`$ is not a root of unity, the finite values of $`|H_1(M_r)|`$ grow exponentially.
More than fifteen years later two independent proofs of Gordon’s conjecture, one by R. Riley \[Ri\] and another by F. Gonzaléz-Acuña and H. Short \[GoSh\], appeared. Both employed the Gel’fond-Baker theory of linear forms in the logarithms of algebraic integers \[Ba\], \[Ge\].
We extend the above results for knots, replacing the term “finite values of $`|H_1(M_r)|`$” with “order of the torsion subgroup of $`H_1(M_r)`$,” while at the same time proving a general result for links in $`S^3`$. Our proof, which is motivated by \[SiWi2\], identifies the torsion subgroup of the homology of a finite abelian branched cover with the connected components of periodic points in an associated algebraic dynamical system. Theorem 21.1 of \[Sc\], an enhanced version of a theorem of D. Lind, K. Schmidt and T. Ward \[LiScWa\], then completes our argument.
Recognizing that relatively few topologists are familiar with algebraic dynamical systems, we have endeavored to make this paper self-contained. The reader who desires to know more about such dynamical systems is encouraged to consult the extraordinary monograph \[Sc\].
We thank Jonathan Hillman, Douglas Lind and Klaus Schmidt for valuable discussions. Also, we are grateful to the University of Washington for its hospitality when part of this work was completed.
2. Statement of results. Let $`l=l_1\mathrm{}l_d`$ be an oriented link of $`d`$ components with exterior $`E=S^3\mathrm{int}N(l)`$, where $`N(l)`$ is a regular neighborhood of $`l`$. The meridianal generators of the link group $`G_l=\pi _1(S^3l)`$ represent distinguished generators $`u_1,\mathrm{},u_d`$ for the abelianization $`G_l/G_l^{}𝐙^d`$. We identify these generators with the standard basis of $`𝐙^d`$.
Given a finite-index subgroup $`\mathrm{\Lambda }𝐙^d`$ there exists a covering $`E_\mathrm{\Lambda }`$ of the link exterior corresponding to the epimorphism $`G_l𝐙^d/\mathrm{\Lambda }`$, the abelianization map composed with the canonical quotient map. By attaching solid tori to $`E_\mathrm{\Lambda }`$ so that meridians of the tori cover meridians of $`l`$ while the collection of cores map to the link, we obtain a cover $`M_\mathrm{\Lambda }`$ of $`S^3`$ branched over $`l`$.
Let $`_d`$ denote the ring $`𝐙[u_1^{\pm 1},\mathrm{},u_d^{\pm 1}]𝐙[𝐙^d]`$ of Laurent polynomials with integer coefficients. The Mahler measure of a nonzero polynomial $`f_d`$ is defined by
$$𝐌(f)=\mathrm{exp}(_{𝐒^d}\mathrm{log}|f(𝐬)|d𝐬),$$
where $`d𝐬`$ indicates integration with respect to normalized Haar measure, and $`𝐒^d`$ is the multiplicative subgroup of $`d`$-dimensional complex space $`𝐂^d`$ consisting of all vectors $`(s_1,\mathrm{},s_d)`$ with $`|s_1|=\mathrm{}=|s_d|=1`$. Clearly, Mahler measure is multiplicative, and the measure of any unit is $`1`$. It is known that $`M(f)=1`$ if and only if $`f`$ is equal up to a unit factor to the product of cyclotomic polynomials in a single variable evaluated at monomials (see \[Sc, Lemma 19.1\]).
The quantity $`𝐌(f)`$, which is the geometric mean of $`|f|`$ over the $`d`$-torus $`𝐒^d`$, was introduced by K. Mahler in \[Ma1\] and \[Ma2\]. It is a consequence of Jensen’s formula \[Al, p. 208\] that if $`f`$ is a nonzero polynomial $`c_nu^n+\mathrm{}+c_1u+c_0(c_n0)`$ in one variable, then
$$𝐌(f)=|c_n|\underset{j=1}{\overset{n}{}}\mathrm{max}(|r_j|,1),$$
where $`r_1,\mathrm{},r_n`$ are the zeros of $`f`$. A short proof can be found in either \[EvWa\] or \[Sc\].
For any finite-index subgroup $`\mathrm{\Lambda }`$ we let
$$\mathrm{\Lambda }=\mathrm{min}\{|v|:v\mathrm{\Lambda }0\},$$
where $`||`$ denotes the Euclidean metric.
Since $`M_\mathrm{\Lambda }`$ is a compact manifold, the homology group $`H_1(M_\mathrm{\Lambda })`$ is finitely generated. (All homology groups in this paper have integer coefficients.) We decompose $`H_1(M_\mathrm{\Lambda })`$ as the direct sum of a free abelian group of some rank $`\beta _\mathrm{\Lambda }`$ and a torsion subgroup $`TH_1(M_\mathrm{\Lambda })`$. We denote the order of $`TH_1(M_\mathrm{\Lambda })`$ by $`b_\mathrm{\Lambda }`$.
Theorem 2.1. Let $`l=l_1\mathrm{}l_d`$ be an oriented link of $`d`$ components having nonzero Alexander polynomial $`\mathrm{\Delta }=\mathrm{\Delta }(u_1,\mathrm{},u_d)`$. Then
$$\underset{\mathrm{\Lambda }\mathrm{}}{lim\; sup}\frac{1}{|𝐙^d/\mathrm{\Lambda }|}\mathrm{log}b_\mathrm{\Lambda }=\mathrm{log}𝐌(\mathrm{\Delta }).$$
When $`d=1`$ the $`lim\; sup`$ can be replaced by an ordinary limit.
As a consequence of the proof of Theorem 2.1 we obtain a new proof of a theorem of Gordon. Recall that for any knot, $`\mathrm{\Delta }_i`$ denotes the $`i`$th Alexander polynomial.
Corollary 2.2. \[Go\] Let $`k`$ be a knot in $`S^3`$. If $`\mathrm{\Delta }_1/\mathrm{\Delta }_2`$ divides $`t^N1`$ for some $`N`$, then $`H_1(M_r)H_1(M_{r+N})`$ for all $`r`$.
The Mahler measure of $`u_1^2u_1+1`$, the Alexander polynomial of the trefoil knot $`3_1`$, is $`1`$ since both zeros of the polynomial have unit modulus. On the other hand, $`u_1^23u_1+1`$, the Alexander polynomial of the figure eight knot $`4_1`$, has zeros $`(1\pm \sqrt{5})/2`$ and hence it has Mahler measure $`(1+\sqrt{5})/21.618`$.
Scanning the table of $`2`$-component links in \[Ro\] we find that $`6_2^2`$ is the first link with nonzero Alexander polynomial having Mahler measure greater than $`1`$. The polynomial is $`u_1+u_21+u_1^1+u_2^1`$, which has Mahler measure approximately equal to $`1.285`$.
The next link in the table, $`6_2^3`$, has Alexander polynomial $`2u_1u_2+2u_1u_2`$, which can be rewritten as $`(2u_1)+u_1u_2(2u_11)`$. Using Lemma 19.8 of \[Sc\] and an easy change of basis (replacing $`u_1u_2`$ with a new variable $`u_2^{}`$) we see that the Mahler measure of this Alexander polynomial is precisely $`2`$.
No polynomial with integer coefficients is known that has Mahler measure greater than 1 but less than that of $`\mathrm{\Delta }(x)=x^{10}+x^9x^7x^6x^5x^4x^3+x+1`$, which has Mahler measure approximately equal to $`1.176`$. (Only one of the nine zeros of $`\mathrm{\Delta }`$ lies outside the unit circle.) Deciding whether or not such a numerical gap truly exists is known as Lehmer’s problem, and it remains a vexing open question. (see \[Le\], \[EvWa\]). It is a provocative fact that $`\mathrm{\Delta }`$ is the Alexander polynomial of a knot. In fact, there are infinitely many, including infinitely many with complements that fiber over the circle.
For more calculations of Mahler measures of Alexander polynomials of links and further discussion of Lehmer’s question see \[SiWi4\].
3. Z<sup>d</sup>-shifts and link colorings. A $`𝐙^d`$-action by automorphisms on a topological group $`X`$ is a homomorphism $`\sigma :𝐦\sigma _𝐦`$ from $`𝐙^d`$ to $`\mathrm{Aut}(X)`$. Two $`𝐙^d`$-actions $`\sigma `$ and $`\sigma ^{}`$ on $`X`$ and $`X^{}`$, respectively, are algebraically conjugate if there exists a continuous group isomorphism $`\varphi :XX^{}`$ such that $`\varphi \sigma _𝐦=\sigma _𝐦^{}\varphi `$, for every $`𝐦𝐙^d`$.
$`_d`$-modules are an important source of $`𝐙^d`$-shifts via Pontryagin duality. Let $`𝐓`$ denote the additive circle group $`𝐑/𝐙`$. For any $`_d`$-module $`L`$, the Pontryagin dual $`L^{}=\mathrm{Hom}(L,𝐓)`$ is a group under pointwise addition. Here $`L`$ is given the discrete topology, and $`L^{}`$ is endowed with the compact-open topology; $`L^{}`$ is a compact abelian group. For $`𝐦𝐙^d`$, scalar multiplication $`a𝐦a`$ in $`L`$ induces $`\sigma _𝐦\mathrm{Aut}(L^{})`$ via its adjoint action. In this way we have a $`𝐙^d`$-action on $`L^{}`$. From a purely algebraic point of view, $`L^{}`$ is a $`_d`$-module. In the case that $`L`$ is a free $`_d`$-module of rank $`N`$, we obtain the compact group $`L^{}=(𝐓^N)^{𝐙^d}`$. The automorphism $`\sigma _𝐦`$ is the shift map given by $`\sigma _𝐦(\alpha _𝐧)=(\alpha _{𝐧+𝐦})`$ for $`\alpha =(\alpha _𝐧)(𝐓^N)^{𝐙^d}`$. The automorphisms $`\sigma _{u_1},\mathrm{},\sigma _{u_d}`$ will be denoted by $`\sigma _1,\mathrm{},\sigma _d`$ for notational ease.
Given a $`𝐙^d`$-action $`\sigma `$ on $`X`$, we say that a point $`xX`$ is periodic under $`\sigma `$ if its orbit $`\{\sigma _𝐦x|𝐦𝐙^d\}`$ is finite. We will be particularly interested in periodic point sets
$$\mathrm{Fix}_\mathrm{\Lambda }(\sigma )=\{xX|\sigma _𝐦x=x𝐦\mathrm{\Lambda }\},$$
where $`\mathrm{\Lambda }`$ is a subgroup of finite index in $`𝐙^d`$. For algebraically conjugate actions such sets clearly correspond under the isomorphism $`\varphi `$.
Definition 3.1. Assume that $`D`$ is a diagram of an oriented link $`l=l_1\mathrm{}l_d`$ of $`d`$ components. A $`𝐓^{𝐙^d}`$-coloring of $`D`$ is an assignment of elements (colors) $`\alpha ,\beta ,\mathrm{}𝐓^{𝐙^d}`$ to the arcs of $`D`$ such that the condition
$$\alpha +\sigma _t\beta =\gamma +\sigma _t^{}\alpha $$
$`(3.1)`$
is satisfied at any crossing. Here $`\alpha `$ corresponds to an overcrossing arc of the $`t`$th component of $`l`$, while $`\beta `$ and $`\gamma `$ correspond to undercrossing arcs of the $`t^{}`$th component. We encounter $`\beta `$ as we travel in the preferred direction along the arc labeled by $`\alpha `$, turning left at the crossing (see Figure 1). The terminology is motivated by the concept of Fox coloring for knots \[Fo\], which was generalized in \[SiWi1\], \[SiWi2\].
Figure 1: $`𝐓^{𝐙^d}`$-Coloring rule
If $`D`$ consists of $`N`$ arcs, then the set $`\mathrm{Col}_{𝐓,𝐙^d}(D)`$ of all $`𝐓^{𝐙^d}`$-colorings of $`D`$ is a closed subgroup of $`[𝐓^{𝐙^d}]^N[𝐓^N]^{𝐙^d}`$ that is invariant under $`\sigma _𝐦`$ for each $`𝐦𝐙^d`$. It will follow from observations in Section 4 that if $`D^{}`$ is a another diagram for $`l`$, then the $`𝐙^d`$-action on $`\mathrm{Col}_{𝐓,𝐙^d}(D)`$ is algebraically conjugate to the action on $`\mathrm{Col}_{𝐓,𝐙^d}(D^{})`$. In anticipation we make the following definition.
Definition 3.2. (Cf. \[SiWi1\]) Let $`l`$ be a $`d`$-component oriented link with diagram $`D`$. The color $`𝐙^d`$-shift $`\mathrm{Col}_{𝐓,𝐙^d}(l)`$ is the compact abelian group $`\mathrm{Col}_{𝐓,𝐙^d}(D)`$ together with the $`𝐙^d`$-action $`\sigma `$.
4. Alexander module and periodic points. This section contains the proofs of our main results. An example that illustrates the main ideas is given at the end.
A diagram $`D`$ for $`l`$ yields a finite Wirtinger presentation $`x_1,x_2\mathrm{},x_Nr_1,\mathrm{},r_N`$ for $`G_l=\pi _1(S^3l)`$. Let $`P`$ be the canonical $`2`$-complex with $`\pi _1PG_l`$, constructed with a single vertex $`v`$, directed edges labeled $`x_1,\mathrm{},x_N`$, and oriented $`2`$-cells $`c_1,\mathrm{},c_N`$ with each boundary $`c_i`$ attached to $`1`$-cells according to $`r_i`$ (see Chapter 11 of \[Li\]).
Let $`\stackrel{~}{P}`$ be the maximal abelian cover of $`P`$; that is, the cover corresponding to the abelianization map $`G_lG_l/G_l^{}𝐙^d`$. As usual each cell $`v,x_i,c_j`$ of $`P`$ lifts to a family $`𝐦\stackrel{~}{v},𝐦\stackrel{~}{x}_i,𝐦\stackrel{~}{c}_j`$ of oriented cells indexed by $`𝐙^d`$. By standard construction the chain complex $`0\stackrel{~}{C}_2\stackrel{_2}{}\stackrel{~}{C}_1\stackrel{_1}{}\stackrel{~}{C}_00`$ admits a quotient $`0\stackrel{~}{C}_2\stackrel{_2}{}\stackrel{~}{C}_10`$ that determines the relative homology group $`H_1(\stackrel{~}{P},\stackrel{~}{P}^0)`$, where $`\stackrel{~}{P}^0`$ is the $`0`$-skeleton of $`\stackrel{~}{P}`$. This $`_d`$-module is the Alexander module of the link, denoted here by $`A`$.
By the universal coefficient theorem \[Sp, p. 243\] the cohomology group $`H^1(\stackrel{~}{P},\stackrel{~}{P}^0;𝐓)`$ is naturally isomorphic to the dual group of the Alexander module. It is a closed subgroup of $`\mathrm{Hom}(\stackrel{~}{C}_1,𝐓)=[𝐓^N]^{𝐙^d}`$, the kernel of the coboundary operator
$$\mathrm{Hom}(\stackrel{~}{C}_1,𝐓)\stackrel{\mathrm{Hom}(_2,1)}{}\mathrm{Hom}(\stackrel{~}{C}_2,𝐓),$$
and hence it inherits a $`𝐙^d`$-action from $`\mathrm{Hom}(\stackrel{~}{C}_1,𝐓)`$ (see Section 3).
We observe that $`_2\stackrel{~}{c}_1,\mathrm{},_2\stackrel{~}{c}_N`$ closely resemble the relations of the coloring rule (3.1). The Wirtinger relator at the crossing in Figure 1 has the form $`x_ix_{j_1}x_i^1x_{j_2}^1`$, and the lifted loop in the cover that begins at $`\stackrel{~}{v}`$ determines the $`1`$-cycle $`\stackrel{~}{x}_i+u_t\stackrel{~}{x}_{j_1}u_t^{}\stackrel{~}{x}_i\stackrel{~}{x}_{j_2}`$, which induces the homology relation $`\stackrel{~}{x}_i+u_t\stackrel{~}{x}_{j_1}=\stackrel{~}{x}_{j_2}+u_t^{}\stackrel{~}{x}_i`$. Lifts that begin at other points of the cover are simply translates by elements of $`𝐙^d`$. We can regard the assignment of $`\alpha =(\alpha _𝐦)𝐓^{𝐙^d}`$ to an arc $`x_i`$ as an assignment of $`\alpha _𝐦𝐓`$ to the $`1`$-chain $`𝐦\stackrel{~}{x}_i`$. Then clearly $`\mathrm{Col}_{𝐓,𝐙^d}(D)`$ and $`H^1(\stackrel{~}{P},\stackrel{~}{P}^0;𝐓)`$ are described by identical subsets of $`[𝐓^N]^{𝐙^d}`$. Since the Alexander module is a link invariant, it follows that the algebraic conjugacy class of $`\mathrm{Col}_{𝐓,𝐙^d}(D)`$ is independent of the diagram for $`l`$.
Consider a finite-index subgroup $`\mathrm{\Lambda }`$ of $`𝐙^d`$. The unbranched cover $`E_\mathrm{\Lambda }`$ has the same homology as the quotient complex $`\stackrel{~}{P}/\mathrm{\Lambda }`$. A $`2`$-complex $`Q`$ with the same first homology group as the branched cover $`M_\mathrm{\Lambda }`$ is obtained from $`\stackrel{~}{P}/\mathrm{\Lambda }`$ by attaching additional $`2`$-cells as follows. Each Wirtinger generator $`x_i,1iN`$, maps to some $`u_{t(i)}`$ under abelianization. Assume that $`u_{t(i)}`$ represents an element of order $`n(i)`$ in $`𝐙^d/\mathrm{\Lambda }`$. Then $`\stackrel{~}{x}_i+u_{t(i)}\stackrel{~}{x}_i+\mathrm{}+u_{t(i)}^{n(i)1}\stackrel{~}{x}_i`$ is a $`1`$-cycle in $`\stackrel{~}{P}/\mathrm{\Lambda }`$, and we attach a $`2`$-cell along it and each of its translates. The additional cells added to the complex in this way correspond to the meridinal disks of tori that we attach to $`E_\mathrm{\Lambda }`$ when constructing $`M_\mathrm{\Lambda }`$. Consequently, $`Q`$ has the same fundamental group and hence the same first homology group as $`M_\mathrm{\Lambda }`$.
Elements of $`H^1(Q,Q^0;𝐓)`$ correspond to $`𝐙^d`$-colorings in $`\mathrm{Fix}_\mathrm{\Lambda }(\sigma )`$ such that if $`\alpha `$ is assigned to the $`i`$th arc of $`D`$, then
$$\alpha +\sigma _{t(i)}\alpha +\mathrm{}+\sigma _{t(i)}^{n(i)1}\alpha =0.$$
$`(4.1)`$
Such $`𝐙^d`$-colorings comprise a subgroup $`\mathrm{SFix}_\mathrm{\Lambda }(\sigma )\mathrm{Fix}_\mathrm{\Lambda }(\sigma )`$ of special periodic points.
The universal coefficient theorem implies that $`H^1(Q;𝐓)`$ is isomorphic to the dual group of $`H_1(Q)H_1(M_\mathrm{\Lambda })`$. By decomposing $`H_1(M_\mathrm{\Lambda })`$ as $`TH_1(M_\mathrm{\Lambda })𝐙^{\beta _\mathrm{\Lambda }}`$ (see Section 2) and recalling that $`A^{}A`$ for any finite abelian group $`A`$, we have
$$H^1(Q;𝐓)[TH_1(M_\mathrm{\Lambda })𝐙^{\beta _\mathrm{\Lambda }}]^{}TH_1(M_\mathrm{\Lambda })𝐓^{\beta _\mathrm{\Lambda }}.$$
$`(4.2)`$
Consider now the portion of the cohomology long exact sequence:
$$H^0(Q^0;𝐓)\stackrel{\delta }{}H^1(Q,Q^0;𝐓)H^1(Q;𝐓)0.$$
$`(4.3)`$
Lemma 4.1. The image of $`\delta `$ is a direct summand of $`H^1(Q,Q^0;𝐓)`$ isomorphic to $`𝐓^r`$, where $`r=|𝐙^d/\mathrm{\Lambda }|1`$.
Proof. The $`0`$-skeleton $`Q^0`$ consists of vertices indexed by elements of $`𝐙^d/\mathrm{\Lambda }`$. Elements of $`H^0(Q^0;𝐓)`$ can be regarded as functions $`f:𝐙^d/\mathrm{\Lambda }𝐓`$. The image $`\delta (f)`$ is an edge-labeling of $`Q`$, assigning $`f(𝐦^{})f(𝐦)`$ to an edge from $`𝐦\stackrel{~}{v}`$ to $`𝐦^{}\stackrel{~}{v}`$. Select a maximal tree $`T`$ in the $`1`$-skeleton of $`Q`$. It is clear that $`\delta (f)`$ is uniquely determined by its values on $`T`$, and such values can be prescribed arbitrarily. Since $`T`$ has $`r`$ edges, the image of $`\delta `$ is isomorphic to $`𝐓^r`$.
The map $`H^1(Q,Q^0;𝐓)\stackrel{ϵ}{}𝐓^r`$ given by restricting any cocycle to the edges of the maximal tree $`T`$ is an epimorphism, by what we have said above. We construct a right inverse $`\eta `$ for $`ϵ`$ as follows. Given a function $`g:T𝐓`$, choose an element $`fH^0(Q^0;𝐓)`$ such that $`\delta (f)`$ agrees with $`g`$ on $`T`$. Define $`\eta (g)=\delta (f)`$. Hence the image of $`\delta `$ is a direct summand of $`H^1(Q,Q^0;𝐓)`$.
Corollary 4.2. Let $`\mathrm{SFix}_\mathrm{\Lambda }^0(\sigma )`$ be the connected component of the identity in $`\mathrm{SFix}_\mathrm{\Lambda }(\sigma )`$. Then $`\mathrm{SFix}_\mathrm{\Lambda }(\sigma )/\mathrm{SFix}_\mathrm{\Lambda }^0(\sigma )TH_1(M_\mathrm{\Lambda })`$.
Proof. By equation (4.2) we have $`TH_1(M_\mathrm{\Lambda })𝐓^{\beta _\mathrm{\Lambda }}H^1(Q;𝐓)`$, and by the long exact sequence (4.3) the latter module is isomorphic to $`H^1(Q,Q^0;𝐓)/\mathrm{im}(\delta )`$. Recall that $`H^1(Q,Q^0;𝐓)`$ is isomorphic to $`\mathrm{SFix}_\mathrm{\Lambda }(\sigma )`$. By Lemma 4.1 the image of $`\delta `$ is connected and hence contained in $`\mathrm{SFix}_\mathrm{\Lambda }^0(\sigma )`$. Thus $`\mathrm{SFix}_\mathrm{\Lambda }(\sigma )/\mathrm{SFix}_\mathrm{\Lambda }^0(\sigma )TH_1(M_\mathrm{\Lambda })`$.
Corollary 4.3. The first Betti number $`\beta _\mathrm{\Lambda }`$ of $`M_\mathrm{\Lambda }`$ is equal to $`\mathrm{dim}\mathrm{SFix}_\mathrm{\Lambda }(\sigma )|𝐙^d/\mathrm{\Lambda }|+1.`$
Proof. By equation (4.2) the Betti number $`\beta _\mathrm{\Lambda }`$ is equal to the dimension of $`H^1(Q;𝐓)`$, and by the long exact sequence (4.3) the latter is $`\mathrm{dim}\mathrm{SFix}_\mathrm{\Lambda }(\sigma )\mathrm{dim}\mathrm{im}(\delta )`$. Lemma 4.1 completes the argument.
The quantity $`\mathrm{dim}\mathrm{SFix}_\mathrm{\Lambda }(\sigma )`$ can be generally computed as the nullity of a certain matrix. The computation is similar to that which results from the formula of M. Sakuma \[Sa, Theorem 1.1(2)\]. However, the dynamical systems perspective here is new.
Proof of Theorem 2.1. It follows from Corollary 4.2 that $`b_\mathrm{\Lambda }=|TH_1(M_\mathrm{\Lambda })|`$ is equal to the number of connected components of $`\mathrm{SFix}_\mathrm{\Lambda }(\sigma )`$. Now we apply techniques of symbolic dynamics to count the number of components and determine their exponential growth rate.
We denote the connected component of the identity in $`\mathrm{Fix}_\mathrm{\Lambda }(\tau )`$ by the symbol $`\mathrm{Fix}_\mathrm{\Lambda }^0(\tau )`$. For any $`𝐙^d`$-action $`\tau `$ associated to the dual group of a $`_d`$-module, Theorem 21.1 of \[Sc\] implies that the exponential growth rate of $`|\mathrm{Fix}_\mathrm{\Lambda }(\tau )/\mathrm{Fix}_\mathrm{\Lambda }^0(\tau )|`$ as $`\mathrm{\Lambda }`$ approaches infinity is equal to the topological entropy of $`\tau `$, provided that the topological entropy of $`\tau `$ is finite. The topological entropy of the $`𝐙^d`$-action $`\sigma `$ above is always infinite. However, we will show that there is a related $`𝐙^d`$-shift $`\sigma ^{}`$ such that (1) $`\sigma ^{}`$ has finite topological entropy equal to $`\mathrm{log}𝐌(\mathrm{\Delta })`$; and (2) $`|\mathrm{Fix}_\mathrm{\Lambda }(\sigma ^{})/\mathrm{Fix}_\mathrm{\Lambda }^0(\sigma ^{})|=|\mathrm{SFix}_\mathrm{\Lambda }(\sigma )/\mathrm{SFix}_\mathrm{\Lambda }^0(\sigma )|`$. The main conclusion of Theorem 2.1 follows from these assertions.
A $`𝐓^{𝐙^d}`$-coloring of $`D`$ that assigns $`0𝐓^{𝐙^d}`$ to some arc, say the arc corresponding to Wirtinger generator $`x_1`$, will be called a based $`𝐓^{𝐙^d}`$-coloring. The collection of based $`𝐓^{𝐙^d}`$-colorings is a closed shift-invariant subgroup of $`\mathrm{Col}_{𝐓,𝐙^d}(D)`$, independent of the choice of arc. It is the dual of $`B=A/\stackrel{~}{x}_1`$, the quotient of the Alexander module by the submodule generated by $`\stackrel{~}{x}_1`$, which one might call the based Alexander module of the link. We denote the $`𝐙^d`$-action on $`B`$ by $`\sigma ^B`$.
It is clear from the discussion above that for any knot the based Alexander module is isomorphic to the first homology group of the infinite cyclic cover of the knot exterior. Since all generators of a Wirtinger presentation for a knot group are conjugate, it follows that periodic points of $`\sigma ^B`$ are always special periodic points. (See the paragraph preceding Lemma 2.7 \[SiWi3\] for details.) Hence in this case $`\mathrm{Fix}((\sigma ^B)^r)`$ is isomorphic to the dual of $`H_1(M_r)`$, for any $`r`$.
In the general case we define special periodic points of $`\sigma ^B`$ just as we defined them for $`\sigma `$. We claim that
$$\mathrm{SFix}_\mathrm{\Lambda }(\sigma ^B)/\mathrm{SFix}_\mathrm{\Lambda }^0(\sigma ^B)\mathrm{SFix}_\mathrm{\Lambda }(\sigma )/\mathrm{SFix}_\mathrm{\Lambda }^0(\sigma ),$$
and hence the number of connected components of $`\mathrm{SFix}_\mathrm{\Lambda }(\sigma ^B)`$ is equal to that of $`\mathrm{SFix}_\mathrm{\Lambda }(\sigma )`$. One way to see this is by constructing a maximal tree $`T`$ for the $`1`$-skeleton of $`Q`$, selecting first a maximal number of edges of the form $`𝐦\stackrel{~}{x}_1`$. Recall that $`\mathrm{SFix}_\mathrm{\Lambda }(\sigma )`$ can be identified with $`H^1(Q,Q^0;𝐓)`$, and by the proof of Lemma 4.1 this group has a direct summand $`𝐓^r`$, the image of the map $`\eta `$ defined above. We pass from $`\mathrm{SFix}_\mathrm{\Lambda }(\sigma )`$ to $`\mathrm{SFix}_\mathrm{\Lambda }(\sigma )/\mathrm{SFix}_\mathrm{\Lambda }^0(\sigma )`$ in two stages. First we crush $`\eta (𝐓^r)`$; the quotient group is isomorphic to the subgroup of $`H^1(Q,Q^0;𝐓)`$ consisting of cocyles that vanish on the edges of the maximal tree $`T`$. However, in view of equation (4.1), any cocycle that vanishes on those edges vanishes on all of the edges $`𝐦\stackrel{~}{x}_1`$. Thus $`\mathrm{SFix}_\mathrm{\Lambda }(\sigma )/\eta (T^r)\mathrm{SFix}_\mathrm{\Lambda }(\sigma ^B)`$, the subgroup of $`H^1(Q,Q^0;𝐓)`$ consisting of based $`𝐓^{𝐙^d}`$-colorings. The connected component of the identity in this group is also a torus. Crushing it we find that $`\mathrm{SFix}_\mathrm{\Lambda }(\sigma )/\mathrm{SFix}_\mathrm{\Lambda }^0(\sigma )\mathrm{SFix}_\mathrm{\Lambda }(\sigma ^B)/\mathrm{SFix}_\mathrm{\Lambda }^0(\sigma ^B)`$.
The based Alexander module $`B`$ has an $`(N1)\times (N1)`$ presentation matrix $`R`$, which can be obtained from the matrix for the Alexander module $`A`$ by deleting the first row and column. Then $`\mathrm{\Delta }(u_1,\mathrm{},u_d)=(u_11)\mathrm{det}R`$ (see \[Li\], pp. 119 – 120). Since the Mahler measure of $`u_11`$ is equal to $`1`$, the determinant of $`R`$ has the same Mahler measure as $`\mathrm{\Delta }`$. By \[LiScWa, p. 611\] the topological entropy of $`\sigma ^B`$ is equal to the log of the Mahler measure of $`\mathrm{\Delta }`$.
The $`𝐙^d`$-action $`\sigma ^{}`$ that we need is a modification of $`\sigma ^B`$. Consider based $`𝐙^d`$-colorings of $`D`$, but replace any color $`\beta `$ by a pair $`(\beta ,\zeta )`$ of colors, and require in addition to the basic coloring rule (3.1) the condition:
$$\sigma _t\zeta =\zeta +\beta ,$$
$`(4.4)`$
where $`t`$ is the index of the component of $`l`$ containing the arc colored by $`(\beta ,\zeta )`$. Denote the associated $`𝐙^d`$-action by $`\sigma ^{}`$.
The $`𝐙^d`$-action $`\sigma ^{}`$ is on the dual group of a module $`B^{}`$ that we obtain from a presentation for $`B`$ by adding new generators $`z_2,\mathrm{},z_N`$ and relations $`u_{t(i)}z_i=z_i+x_i`$ ($`2iN`$), where $`t(i)`$ is the index of the component of $`l`$ containing the arc $`x_i`$. The determinant of the new relation matrix $`R^{}`$ is $`\mathrm{\Delta }(u_1,\mathrm{},u_d)`$ times a product of polynomials of the form $`u_t1`$. As before, since the Mahler measure of each $`u_t1`$ is equal to $`1`$, we have $`𝐌(\mathrm{det}(R^{}))=𝐌(\mathrm{\Delta })`$, and by \[LiScWa, p. 611\] the topological entropy of $`\sigma ^{}`$ is equal to this value. Hence assertion (1) above holds.
By a straightforward recursion argument we find that $`\mathrm{Fix}_\mathrm{\Lambda }(\sigma ^{})\mathrm{SFix}_\mathrm{\Lambda }(\sigma ^B)𝐓^s`$. Here $`s`$ is the number of second coordinates $`\zeta `$ that can be freely assigned: Assume that $`u_t`$ represents an element of order $`n`$ in $`𝐙^d/\mathrm{\Lambda }`$. Condition (4.4) implies
$$\zeta _{𝐦+u_t}=\zeta _𝐦+\beta _𝐦$$
$$\zeta _{𝐦+2u_t}=\zeta _𝐦+\beta _𝐦+\beta _{𝐦+u_t}$$
$$\zeta _{𝐦+nu_t}=\zeta _𝐦+\beta _𝐦+\beta _{𝐦+u_t}+\mathrm{}+\beta _{𝐦+(n1)u_t}.$$
Clearly, $`\zeta _{𝐦+nu_t}=\zeta _𝐦`$ if and only if $`\beta _{𝐦+u_t}+\mathrm{}+\beta _{𝐦+(n1)u_t}=0.`$ Moreover, the coordinates $`\zeta _{u_t},\mathrm{},\zeta _{(n1)u_t}`$ are uniquely determined from $`\zeta _\mathrm{𝟎}`$ and coordinates of $`\beta `$. When $`t^{}`$ is different from $`t`$, condition (4.4) imposes no new requirement; in such a case the coordinates $`\zeta _\mathrm{𝟎},\zeta _{u_i^{}},\mathrm{},\zeta _{(n^{}1)u_t^{}}`$ can be chosen arbitrarily, where $`u_t^{}`$ represents an element of order $`n^{}`$ in $`𝐙^d/\mathrm{\Lambda }`$. Assertion (2) is immediate, and the proof of the theorem is complete.
Example 4.4. We will illustrate the ideas and terminology used in the proof of Theorem 2.1 with an example.
The diagram $`D`$ for the link $`l=5_1^2`$, shown in Figure 2, yields a Wirtinger presentation for $`G_l`$:
$$x_1,x_2,x_3,x_4,x_5x_1x_3=x_5x_1,x_3x_2=x_1x_3,x_5x_4=x_3x_5,x_4x_2=x_1x_4,x_2x_4=x_5x_2.$$
Figure 2: Diagram for $`5_1^2`$
We assume that under abelianization $`x_1`$ and $`x_2`$ map to $`u_1`$ while the remaining generators are sent to $`u_2`$. A portion of the maximal abelian cover $`\stackrel{~}{P}`$ is shown in Figure 3. The $`2`$-cells are not shown.
Figure 3: Portion of $`\stackrel{~}{P}`$
We consider the subgroup $`\mathrm{\Lambda }`$ of $`𝐙^2`$ generated by $`u_1^3`$ and $`u_2^2`$. The $`1`$-skeleton of $`Q`$ (which is the same as the $`1`$-skeleton of $`\stackrel{~}{P}/\mathrm{\Lambda }`$) is shown in Figure 4. Since $`\mathrm{\Lambda }`$ has generators parallel to $`u_1,u_2`$, it is easy to visualize the additional $`2`$-cells that must be attached to $`\stackrel{~}{P}/\mathrm{\Lambda }`$ in order to build $`Q`$. In more general examples, the boundaries of the new cells might wind around the graph several times.
Figure 4: $`1`$-Skeleton of $`Q`$
An element $`x\mathrm{Fix}_\mathrm{\Lambda }(\sigma ^B)`$ can be represented by $`3\times 2`$ matrices $`\beta ,\gamma ,\delta ,ϵ`$ assigned to arcs corresponding to $`x_2,x_3,x_4,x_5`$, respectively. These matrices have entries in $`𝐓`$, and are the restrictions of the elements of $`𝐓^{𝐙^d}`$ to a fundamental region of $`𝐙^d/\mathrm{\Lambda }`$. The element $`x`$ is in $`\mathrm{SFix}_\mathrm{\Lambda }(\sigma ^B)`$ if the column sums of $`\beta `$ and the row sums of $`\gamma ,\delta ,ϵ`$ are all zero.
An element of $`\mathrm{Fix}_\mathrm{\Lambda }(\sigma ^{})`$ assigns additional $`3\times 2`$ matrices $`\zeta ^\beta ,\zeta ^\gamma ,\zeta ^d,\zeta ^ϵ`$. We can prescribe the coordinates $`\zeta _{0,0}^\beta ,\zeta _{0,1}^\beta `$ arbitrarily; the other coordinates of $`\zeta ^\beta `$ are uniquely determined by these and $`\beta `$. Similarly, the coordinates $`\zeta _{0,0}^\gamma ,\zeta _{1,0}^\gamma ,\zeta _{2,0}^\gamma ,\zeta _{0,0}^\delta ,\zeta _{1,0}^\delta ,\zeta _{2,0}^\delta ,\zeta _{0,0}^ϵ,\zeta _{1,0}^ϵ,\zeta _{2,0}^ϵ`$ are arbitrary. We have $`\mathrm{Fix}_\mathrm{\Lambda }(\sigma ^{})\mathrm{SFix}_\mathrm{\Lambda }(\sigma )𝐓^{11}`$.
We remark that for this link the Alexander polynomial is $`(1u_1)(1u_2)`$. Since the Mahler measure of the polynomial is $`1`$, the orders $`b_\mathrm{\Lambda }`$ have zero exponential growth rate.
Proof of Corollary 2.2. Since $`\mathrm{\Delta }_1/\mathrm{\Delta }_2`$ annihilates the Alexander module of any knot \[Cr1\], it follows that $`\sigma ^N=\mathrm{id}.`$ From this we have $`\mathrm{Fix}(\sigma ^{r+N})=\mathrm{Fix}(\sigma ^r)`$. Recall that for any knot, $`\mathrm{Fix}(\sigma ^r)`$ is isomorphic to the dual of $`H_1(M_r)`$. Hence $`H_1(M_{r+N})H_1(M_r)`$, for every $`r1`$.
5. Coloring with nonabelian groups. The coloring rule (3.1) generalizes in a natural way, allowing one to replace $`𝐓`$ with an arbitrary topological group $`\mathrm{\Sigma }`$.
Definition 5.1. Assume that $`D`$ is a diagram of an oriented link $`l=l_1\mathrm{}l_d`$ of $`d`$ components. A $`\mathrm{\Sigma }^{𝐙^d}`$-coloring of $`D`$ is an assignment of elements (colors) $`\alpha ,\beta ,\mathrm{}\mathrm{\Sigma }^{𝐙^d}`$ to the arcs of $`D`$ such that the condition
$$\alpha \sigma _t\beta =\gamma \sigma _t^{}\alpha $$
$`(5.1)`$
is satisfied at any crossing. The colors $`\alpha ,\beta ,\gamma `$ correspond to arcs that are described as in Definition 3.1.
As before, if $`D`$ consists of $`N`$ arcs, then the set $`\mathrm{Col}_{\mathrm{\Sigma },𝐙^d}(D)`$ of all $`\mathrm{\Sigma }^{𝐙^d}`$-colorings of $`D`$ is a closed subspace of $`[\mathrm{\Sigma }^N]^{𝐙^d}`$ that is invariant under $`\sigma _𝐦`$ for each $`𝐦𝐙^d`$. That $`\mathrm{Col}_{\mathrm{\Sigma },𝐙^d}(D)`$ does not depend on the choice of diagram for $`l`$ follows immediately from the following. Let $`\stackrel{~}{E}`$ denote the maximal abelian cover of the link exterior with projection $`p:\stackrel{~}{E}E`$, and let $``$ be a point of $`E`$. Let $`\stackrel{~}{}`$ denote a fixed lift of $``$. Any covering automorphism of $`\stackrel{~}{E}`$ induces a homeomorphism of the quotient space $`\stackrel{~}{E}/p^1()`$, and hence induces an automorphism of $`\pi _1(\stackrel{~}{E}/p^1(),\stackrel{~}{})`$. By considering the adjoint action we obtain a homeomorphism of the representation space $`\mathrm{Hom}[\pi _1(E/p^1(),\stackrel{~}{}),\mathrm{\Sigma }]`$. In this way we obtain a $`𝐙^d`$-action $`\sigma ^{}`$ on $`\mathrm{Hom}[\pi _1(E/p^1(),\stackrel{~}{}),\mathrm{\Sigma }]`$. A $`𝐙^d`$-action on a topological space is defined just as for $`𝐙^d`$-action on a topological group, eliminating the requirement of a group structure on the space; two $`𝐙^d`$-actions, $`\sigma `$ acting on $`X`$ and $`\sigma ^{}`$ acting on $`X^{}`$, are topologically conjugate if there is a homeomorphism $`\varphi :XX^{}`$ such that $`\varphi \sigma _𝐦=\sigma _𝐦^{}\varphi ,`$ for each $`𝐦𝐙^d`$.
Proposition 5.2. The $`𝐙^d`$-actions on $`\mathrm{Hom}[\pi _1(\stackrel{~}{E}/p^1(),\stackrel{~}{}),\mathrm{\Sigma }]`$ and $`\mathrm{Col}_{\mathrm{\Sigma },𝐙^d}(D)`$ are topologically conjugate.
Proof. The quotient space $`\stackrel{~}{E}/p^1()`$ has the same fundamental group as the quotient complex $`\stackrel{~}{P}/\stackrel{~}{P}^0`$, where $`\stackrel{~}{P}`$ is defined in Section 4 and $`\stackrel{~}{P}^0`$ is the $`1`$-skeleton. There is a group presentation for $`\pi _1(\stackrel{~}{P}/\stackrel{~}{P}^0)`$ in which the generators correspond to the edges of $`\stackrel{~}{P}`$; lifts in $`\stackrel{~}{P}`$ of closed paths representing Wirtinger relators are closed paths in $`\stackrel{~}{P}/\stackrel{~}{P}^0`$ representing the relators. Assignments of colors to the arcs of $`D`$, or equivalently to the edges of $`P`$, such that the condition (5.1) holds at each crossing then correspond to homomorphisms from $`\pi _1(\stackrel{~}{P}/\stackrel{~}{P}^0)`$ to $`\mathrm{\Sigma }`$. The correspondance defines a homeomorphism $`\varphi :\mathrm{Col}_{\mathrm{\Sigma },𝐙^d}(D)\mathrm{Hom}[\pi _1(\stackrel{~}{E}/p^1(),\stackrel{~}{}),\mathrm{\Sigma }]`$ such that $`\varphi \sigma _𝐦=\sigma _𝐦^{}\varphi `$, for each $`𝐦𝐙^d`$.
In view of this proposition we call $`\mathrm{Col}_{\mathrm{\Sigma },𝐙^d}(D)`$ the color $`\mathrm{\Sigma }^d`$-shift of the link $`l`$, and we denote it by $`\mathrm{Col}_{\mathrm{\Sigma },𝐙^d}(l)`$.
The abelianization of $`\pi _1(\stackrel{~}{E}/p^1())`$ is $`H_1(\stackrel{~}{P},\stackrel{~}{P}^0)`$. It is isomorphic to $`H_1(\stackrel{~}{E},p^1())`$, the Alexander module of the link. Hence we propose that $`\pi _1(\stackrel{~}{E}/p^1())`$ be called the Alexander group of the link; we denote the group by $`𝒜_l`$. (Shortly after completing the first draft of this paper, the authors discovered that the Alexander group is a special case of the derived group of a permutation representation, introduced by R. Crowell \[Cr2\].)
It is much easier to write a presentation for the Alexander group than for the commutator subgroup of $`\pi _1(S^3l)`$. One begins with families of generators $`a_𝐦,b_𝐦,c_𝐦,\mathrm{}`$ $`(𝐦𝐙^d)`$ corresponding to the arcs of the diagram. Each crossing gives rise to a family of relations: a crossing such as in Figure 1 imposes the relation $`a_𝐦b_{𝐦+u_t}=c_𝐦a_{𝐦+u_t}`$. In the case of a knot, when $`d=1`$, presentations of this sort are well known; J. C. Hausmann and M. Kervaire \[HaKe\] termed them $`𝐙`$-dynamic. A presentation for the Alexander group of a link such as the one we have described might be called $`𝐙^d`$-dynamic.
The next proposition describes the relationship between the Alexander group $`𝒜_l`$ of a link $`l`$ and the commutator subgroup $`G_l^{}`$. We use the terminology of Section 4. Recall that $`T`$ is a maximal tree in the $`1`$-skeleton of the cover $`\stackrel{~}{P}`$.
Proposition 5.3. Let $`l`$ be an oriented link of $`d`$ components. The generators of $`𝒜_l`$ corresponding to the edges of $`T`$ freely generate a subgroup $`F(E_T)`$ of $`𝒜_l`$. Moreover, $`𝒜_l`$ is isomorphic to the free product $`G_l^{}F(E_T)`$.
Proof. Let $`C\stackrel{~}{P}^0`$ denote the cone on the $`0`$-skeleton of $`\stackrel{~}{P}`$. The fundamental group of $`X=\stackrel{~}{P}_{\stackrel{~}{P}^0}C\stackrel{~}{P}^0`$ is isomorphic to $`𝒜_l`$. We can regard $`X`$ as the union of $`\stackrel{~}{P}`$ and $`T_{\stackrel{~}{P}^0}C\stackrel{~}{P}^0`$, which have contractible intersection $`T`$. An application of the Seifert van-Kampen theorem completes the argument.
Example 5.4. We examine two examples. Both are simple, but they highlight some of the advantages of working with the Alexander group rather the commutator subgroup of a link.
(i) Let $`l`$ be the trivial 2-component link. The group $`G_l`$ is free on two generators. Choosing a diagram without crossings, we find that the Alexander group $`𝒜_l`$ is free on generators $`a_{i,j},b_{i,j}`$, where $`i,j`$ range over $`𝐙`$. The commutator subgroup $`G^{}`$ is also free, but it does not admit a natural $`𝐙^2`$-action by automorphisms as does $`𝒜_l`$. (ii) Next consider the link $`l=2_1^2`$, a Hopf link. The group $`G_l`$ is free abelian of rank 2. The Alexander group $`𝒜_l`$ has presentation $`a_{i,j},b_{i,j}a_{i,j}b_{i+1,j}=b_{i,j}a_{i,j+1}`$. In this example the commutator subgroup $`G_l^{}`$ is trivial.
6. Conclusion. A possible direction for further inquiry involves links with zero Alexander polynomial.
The homology growth rate in Theorem 2.1 was computed as the topological entropy of a $`𝐙^d`$-action $`\sigma ^{}`$. When the Alexander polynomial of the link is zero, the entropy of $`\sigma ^{}`$ can be shown to be infinite; in such a case we obtained no information. However, we offer the following
Conjecture 6.1. If $`l`$ is an oriented link of $`d`$ components, then
$$\underset{\mathrm{\Lambda }\mathrm{}}{lim}\frac{1}{|𝐙^d/\mathrm{\Lambda }|}\mathrm{log}|TH_1(M_\mathrm{\Lambda }(l))|=\mathrm{log}𝐌(\mathrm{\Delta }_i),$$
where $`\mathrm{\Delta }_i`$ is the first nonzero Alexander polynomial of the link.
References.
\[Ah\] L. V. Ahlfors, Complex analysis, 2nd edition, McGraw-Hill, New York, 1966.
\[Ba\] A. Baker, “The theory of linear forms in logarithms,” in Transcendence Theory, Advances and Applications, Academic Press, 1977.
\[Cr1\] R. H.Crowell, “The annihilator of a knot module,” Proc. Amer. Math. Soc. 15 (1964), 696 – 700.
\[Cr2\] R. H.Crowell, “The derived group of a permutation representation,” Advances in Math. 53(1984), 99 – 124.
\[EvWa\] G. Everest, T. Ward, Heights of polynomials and entropy in algebraic dynamics, Springer-Verlag, London, 1999.
\[Fo\] R. H. Fox, “A quick trip through knot theory,” in Topology of $`3`$-Manifolds and Related Topics (edited by M. K. Fort), Prentice-Hall, NJ (1961), 120–167.
\[Ge\] A. O. Gel’fond, “On the approximation of transcendental numbers by algebraic integers,” Dokl. Akad. Navk. SSSR 2 (1935), 177 – 182.
\[Go\] C. McA. Gordon, “Knots whose branched coverings have periodic homology,” Trans. Amer. Math. Soc. 168 (1972), 357 – 370.
\[GoSh\] F. González-Acuña and H. Short, “Cyclic branched coverings of knots and homology spheres,” Revista Math. 4 (1991), 97 – 120.
\[HaKe\] J. C. Hausmann and M. Kervaire, “Sous-groupes dérivés des groupes des noeuds,” L’Enseign. Math. 24 (1978), 111 – 123.
\[Le\] D. H. Lehmer, “Factorization of certain cyclotomic functions,” Annals of Math. 34 (1933), 461 – 479.
\[Li\] W. B. Lickorish, An Introduction to Knot Theory, Springer-Verlag, Berlin, 1997.
\[LiScWa\] D. Lind, K. Schmidt and T. Ward, “Mahler measure and entropy for commuting automorphisms of compact groups,” Invent. Math. 101 (1990), 593–629.
\[Ma1\] K. Mahler, “An application of Jensen’s formula to polynomials,” Mathematika 7 (1960), 98 – 100.
\[Ma2\] K. Mahler, “On some inequalities for polynomials in several variables,” J. London Math. Soc. 37(1962), 341 – 344.
\[Ri\] R. Riley, “Growth of order of homology of cyclic branched covers of knots,” Bull. London Math. Soc. 22 (1990), 287 – 297.
\[Ro\] D. Rolfsen, Knots and Links, Publish or Perish, Berkeley, CA, 1976.
\[Sa\] M. Sakuma, “Homology of abelian coverings of links and spatial graphs,” Canad. J. Math. 47 (1995), 201 – 224.
\[Sc\] K. Schmidt, Dynamical Systems of Algebraic Origin, Birkhäuser Verlag, Basel, 1995.
\[SiWi1\] D. S. Silver and S. G. Williams, “Generalized $`n`$-colorings of links,” in Knot Theory, Banach Center Publications 42, Warsaw 1998, 381–394.
\[SiWi2\] D. S. Silver and S. G. Williams, “Coloring link diagrams with a continuous palette,” Topology 39 (2000), 1225 – 1237.
\[SiWi3\] D. S. Silver and S. G. Williams, “Knot invariants from symbolic dynamical systems,” Trans. Amer. Math. Soc. 351 (1999), 3243–3265.
\[SiWi4\] D. S. Silver and S. G. Williams, “Mahler measure of Alexander polynomials,” preprint, 2000.
\[Sp\] E. H. Spanier, Algebraic Topology, McGraw-Hill Book Co., New York, 1966.
Dept. of Mathematics and Statistics, Univ. of South Alabama, Mobile, AL 36688-0002 e-mail: silver@mathstat.usouthal.edu, williams@mathstat.usouthal.edu
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# 1 The International Gravitational Event Collaboration
## 1 The International Gravitational Event Collaboration
One of the most relevant scientific objectives for resonant detectors of gravitational waves (gw) is the search for short gw bursts, emitted during the gravitational collapse of stars or the final evolution of coalescing binaries.<sup>?</sup> To ensure the confidence of a detection, it is necessary to compare the observations made by multiple detectors with uncorrelated noise. This has already been done in the past years among pairs of cryogenic detectors with common time periods of observation of about a semester.<sup>?</sup><sup>,</sup><sup>?</sup> A few days of observation have been reported also for three simultaneously operating detectors.<sup>?</sup> In these attempts, the gw search consisted of a time coincidence analysis among the candidate signals reported by the different detectors and no statistically significant excess of coincidences was found.
An increase of the number of cryogenic resonant detectors in simultaneous operation in recent years has greatly improved the chance of making a confident detection of gw bursts. In fact, the International Gravitational Event Collaboration, IGEC, currently consists of the research groups operating the five cryogenic bar detectors ALLEGRO,<sup>?</sup> AURIGA,<sup>?</sup> EXPLORER,<sup>?</sup> NAUTILUS<sup>?</sup> and NIOBE.<sup>?</sup> The IGEC was established in July 1997 with an agreement<sup>?</sup> for setting up a common search for gravitational wave bursts of duration of the order of $`1ms`$. This agreement sets the guidelines for the data exchange procedure among the participating groups and the IGEC scientific policy, whose most relevant aspects are:
* each group has responsibility to make available to IGEC its list of candidate gravitational wave events,
* a unanimous agreement of the member groups is required to make public the results based on the IGEC data exchange,
* IGEC is open to new data taking research groups.
In 1999 the first IGEC analysis of the 1997-1998 data was performed and some initial results will be presented in the following sections. At the time of this analysis, not all the 1997-1998 data had been exchanged. Despite the incomplete data set, the simultaneous operation of four gravitational wave detectors was achieved for the first time.
### 1.1 Data exchange protocol
The IGEC data exchange procedure is aimed at searching for coincident excitations at different detectors. For each detector, a list of candidate events, each describing a $`\delta `$-like gravitational wave excitation of the detector, is provided by the corresponding research group. The IGEC protocol requires that in each list the candidate event rate be at most of the order of $`100/day`$, to limit the expected rate of accidental coincidences. Currently, the research groups have been exchanging event lists relative to the past three years, but a future goal is to implement an automatic exchange day by day.
The event lists are then made available to the IGEC collaboration as files under a common protocol,<sup>?</sup> open for future extensions. Under this protocol, it is mandatory for each detector to provide the minimum set of information needed to describe a $`\delta `$-like gw excitation of the detector for each event; namely its Universal Time of arrival, the Fourier component $`H_o`$ of its amplitude in $`Hz^1`$ and the detector noise level at that time. Another mandatory requirement is the declaration of the effective observation time of each detector, so that the IGEC observation time can be calculated. Optional information, such as the time of threshold crossing of the detector output, the duration of the event and its statistical compliance to a $`\delta `$-like gw excitation, can also be exchanged.
### 1.2 IGEC gravitational wave observatory
The relevant parameters of the five resonant bar detectors of the IGEC observatory in the years 1997-1998 are summarised in Table 1. The detectors are sensitive to gw signals in a typical bandwidth of the order of $`1Hz`$ around each one of the two resonances of the detector, which are close to $`700Hz`$ for NIOBE and close to $`900Hz`$ for all other detectors. The relationship between the Fourier amplitude $`H_0`$, averaged on the two resonant frequencies of the detector, and the energy $`E_s`$ deposited by the g.w. burst on the bar, is given by:
$$H_0=\frac{1}{4L_{bar}\nu _{0}^{}{}_{}{}^{2}}\sqrt{E_s/M_{bar}}$$
(1)
where $`L_{bar}`$ is the bar length, $`M_{bar}`$ its mass, $`\nu _0`$ the mean of the detector resonance frequencies.
The typical thresholds used for selecting the events in the 1997-1998 burst search have been in the range $`H_026\times 10^{21}Hz^1`$. The corresponding strain amplitude of the gw can be computed assuming a model for the burst shape: for the conventional $`1ms`$ burst, the Fourier component $`H_0`$ should be multiplied by $`10^3Hz`$ to get the maximum strain amplitude $`h`$.
To maximise the chances of a coincidence detection, the bars have been oriented to be approximately parallel to one another. Neglecting the polarisation effects, the gw amplitude at the detector is $`H_0=H_{gw}sin^2\theta (t)`$, where $`H_{gw}`$ is the incident gw amplitude and $`\theta (t)`$ is the angle between the bar axis and the direction of the source. In this configuration of the observatory, the relative misalignments reported in Tab. 1 for ALLEGRO, AURIGA, EXPLORER and NAUTILUS disperse their $`sin^2\theta (t)`$ responses by at most a few %, while for NIOBE the dispersion is up to a few tenths. Figure 1 shows as an example the resulting amplitude efficiency for the observation of the Galactic Center as the Earth rotates. Since their values of $`sin^2\theta (t)`$ are simultaneously above 0.7 for about 60% of the time, we point out that this configuration of the observatory ensures a rather good and coherent coverage of the central galactic mass during time.
## 2 The Data Set
Each IGEC group independently developed a data acquisition and an optimum filtering procedure for a $`\delta `$-like gw excitation. These procedures differ greatly in the methods, in particular some of them filter for the energy released in the bar (ALLEGRO,<sup>?</sup> NIOBE<sup>?</sup>) while the others filter for the amplitude and phase of the strain excitation of the bar (AURIGA,<sup>?</sup> EXPLORER and NAUTILUS<sup>?</sup>). The output correlation time of the filters ranges from a few tenths up to a few seconds.
A search for maxima on the filter outputs is then used to identify the time and amplitude of the candidate events, which are exchanged only if their amplitude exceeds a selected threshold relative to the noise level. For the current detectors, these thresholds span in the range of signal-to-noise ratio $`SNR35`$ in amplitude. Of the events overcoming the threshold, some are rejected as spurious by different methods. AURIGA implements a $`\chi ^2`$ test on the compliance of the single detected excitations with the expected template of a gw burst.<sup>?</sup> In fact, the filtering procedure implemented is equivalent to a maximum likelihood fit of a signal model to the data and the goodness of the fit is statistically tested for each candidate event; those events not passing the test are rejected. All the other detectors implement spurious signals rejection using sensors of ambient disturbances. In addition, EXPLORER and NAUTILUS reject the events when the local noise is above a certain threshold.
The effective observation times for each detector has been defined by vetoing, a priori, time periods of detector maintenance or malfunctions and times when the detector was excited by the local environment as determined by the local experimentalists. After filtering, AURIGA implements a second level of vetoes a posteriori to reject further periods of unsatisfactory performance due to a lack of self consistency of its data analysis, that is when its noise fails to be compliant with the modeled one used to build its filtering procedure.<sup>?</sup>
The amplitude distributions of all the 1997-1998 exchanged events for each detector are shown in Fig.2. The effective observation time of the exchanged data up to now is summarised in Table 2, together with the mean rate of exchanged events. The net observation time with at least four, three and two detectors simultaneously operating has been respectively 15.5, 90 and 260 days. We expect that the three-way observation time will increase significantly as the exchanged data set will become complete. The ALLEGRO detector has been showing the best duty cycle, close to $`100\%`$ on the exchanged data period, as well as the most stationary noise performance with respect to the other detectors.
## 3 Analysis of Time Coincidences
A search for two, three and four-fold coincidences was carried out on the exchanged data. In the analysis reported here a $`M`$-fold coincidence is observed if the estimated arrival times $`t_i`$ at the $`M`$ detectors are all $`|t_it_j|1s`$. This figure has been chosen as a compromise between the demands for a small accidental background and for a low false dismissal probability. a new approach for multiple time coincidence analysis has been proposed by the Rome group.<sup>?</sup> In fact, the measured uncertainties $`t_w`$ on the estimated arrival times of a burst at each detector are quite similar and are selected to be $`\pm 0.5s`$. This corresponds to a maximum false dismissal of a few % for the exchanged events, even at low $`SNR`$. The maximum separation among coincident events is therefore set to $`2t_w=1s`$.
A preliminary search for three and four-fold coincidences shows none. A detailed analysis is in progress. The two-fold coincidences found for each pair of detectors are shown in Table 3: in all the cases they correspond to the estimated accidental background.
Two methods for estimating the rate of accidental coincidences have been applied for the pairs of detectors: i) performing several time shifts of events times of one detector with respect to the other and then looking for coincidences<sup>?</sup> (the resulting accidental coincidences are reported as $`<n_a>`$ in Tab.3); ii) assuming stationary Poisson distributions of event times and using the mean measured rates of events for each detector during the common observation time. The latter method predicts a number of accidental $`M`$-fold coincidences given by<sup>?</sup>
$$<n_a>_{theory}=M\frac{(2t_w)^{M1}}{T_{obs}^{M1}}\underset{i=1}{\overset{M}{}}n_i,$$
(2)
where $`M`$ is the number of detectors, $`T_{obs}`$ their common observation time, $`t_w`$ the time window describing the timing accuracy of each detector, $`n_i`$ the number of exchanged candidate events of the $`i^{th}`$ detector during $`T_{obs}`$. The agreement of both estimates of the accidental background of coincidences is well within the statistical uncertainties for the detector pairs. In addition, the observed coincidences, $`n_c`$, correspond to both estimates of the accidental coincidence background. This implies that no excess coincidences were observed.
Using Eq.2, we also performed a preliminary analysis of the rates of accidental coincidences for three and four-fold configurations of the IGEC observatory. The most relevant result here is that the statistical significance of three-fold and four-fold time concidences among the current IGEC detectors improves by order of magnitudes.
## 4 Future Plans and Conclusions
This first IGEC joint analysis has shown that, with the current detector performances and the selected coincidence time window, at least three detectors simultaneously operating can perform an autonomous search for gw bursts with a very low false alarm rate even at signal-to-noise ratio as low as $`35`$ in single detectors. Therefore, the continuation of the IGEC international effort is very strongly motivated. Moreover, the joint work of the participating groups has instigated efforts to co-ordinate data analysis techniques between the different groups.
No coincidence above the expected accidental background were found in this preliminary analysis. Work is in progress to complete the analysis to three and four-fold coincidences, as well as to set upper limits on the rate of incoming gw bursts and on the amplitude of single gw bursts associated with selected time windows of astrophysical interest.
Improvements in the detector performances will lead to increase the sensitivity of the IGEC observatory in two respects. On one side the thresholds for gw burst search will be lowered without increasing the level of the accidental background rate. On the other side the effective bandwidths of the detectors will be also widened, thus decreasing the uncertainty in the estimated arrival time of a gw burst. The latter will allow the lowering of the rate of accidental coincidences and, above all, it will give the opportunity to measure the propagation speed and direction of the incoming gw.
In this framework, we think that the participation within the IGEC observatory of the interferometric detectors as they will begin observations would constitute a very important stage towards the establishment of the future worldwide observatory for gravitational waves.
Acknowledgements
The ALLEGRO group was supported by the National Science Foundation and LSO, the NIOBE group by the Australian Research Council. The Italian groups were supported in part by a grant from MURST-COFIN’97. The Rome group thanks F.Campolungo, G.Federici, R.Lenci, G.Martinelli, E.Serrani, R.Simonetti and F.Tabacchioni for their precious technical assistance.
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# Anomalies and symmetries of the regularized action
## I Introduction
Anomalies are one of the most stricking manifestations of the presence of ultraviolet infinities in quantum field theory phenomena. They have important consequences for the quantum consistency (unitarity and renormalizability) of different models, as well as for the study of the theoretical structure of quantum field theory, providing interesting relations with mathematical objects.
The regularization of a theory is a procedure with, by its very definition, produces a violent modification in its large momentum behaviour. Anomalies arise when this modification, whatever the particular regularization method applied, violates a classical symmetry. In the resulting renormalized quantum field theory, the quantum counterpart of that symmetry becomes ‘anomalous’. In particular, continuum symmetries will no longer lead to the conservation of a current, thus modifying the naive Ward identities.
In the functional integral representation, anomalies are usually attributed to the fact that, while the classical action is invariant under a given symmetry transformation, the integration measure is not. When performing a change of variables associated to the classical symmetry, the Jacobian of the transformation is, due to ultraviolet divergences, not well defined. When carefully evaluated, by introducing a proper regularization, this Jacobian gives rise to an anomalous term in the effective action , or, depending on the context, to anomalous Ward identities.
In this approach, the symmetry transformations are implemented on the unregularized action, which is formally invariant, and the regularization is only implemented afterwards, when computing the Jacobian The regularization of this Jacobian is of course related to the definition of the integration measure. of the transformation. It is worth remarking that the unregularized action contains field components at all the momentum scales, in particular the ones that are above the cutoff. To avoid this unpleasant feature, we shall start by considering the regularized action, and restrict our study to transformations that leave this functional invariant. Of course, we shall impose the constraint that the resulting transformations should tend to the usual ones when the cutoff is removed. Also, due to the fact that they avoid the large momentum modes, they will necessarily have to be non local, although in a scale of the order of the inverse of the cutoff.
As a particularly convenient form of the regularized action, which we will use in this note, we shall introduce a set of Pauli-Villars regulator fields, with carefully tunned masses $`M_i`$ that tend to infinity at the end of the calculation. Due to the presence of the regulator masses, the regularized action will be no longer invariant under some the classical symmetries, namely, those that rely on the masslessness of the fields in the action.
In the present letter, we will point out that this regularized action does have a generalized, non local symmetry, that reproduces the classical symmetry in the limit $`M_i\mathrm{}`$. Moreover, when performing a change of variables in the path integral based on this symmetry, the associated Jacobian is finite, reproducing the anomalous term in the effective action, when the cutoff tends to infinity.
In section II we discuss this procedure in detail for the gravitational conformal anomaly of a massless real scalar field in $`1+1`$ dimensions. We exhibit the explicit form of the non local symmetry of the regularized action and compute the conformal anomaly, which of course agrees with the Liouville action , . In section III we consider massless fermions in $`1+1`$ dimensions coupled to the electromagnetic field. Again we present a generalized chiral symmetry of the regularized action and compute the chiral anomaly. Section IV contains our conclusions.
## II The conformal anomaly in $`1+1`$ dimensions
In this section we shall consider the partition function $`𝒵[g_{\mu \nu }]`$, corresponding to a massless scalar field in the presence of an external gravitational background in $`1+1`$ dimensions:
$$𝒵[g_{\mu \nu }]=[𝒟\phi ]_ge^{iS[\phi ,g_{\mu \nu }]}$$
(1)
where
$$S[\phi ,g_{\mu \nu }]=\frac{1}{2}d^2x\sqrt{g}g^{\mu \nu }_\mu \phi _\nu \phi $$
(2)
and $`g=det(g_{\mu \nu })`$. In Eq. (1) we have made it explicit the fact that the definition of the scalar field integration measure depends upon the background metric $`g_{\mu \nu }`$.
The classical action $`S[\phi ,g_{\mu \nu }]`$ is invariant under Weyl transformations of the metric:
$$g_{\mu \nu }g_{\mu \nu }^\omega (x)=\omega (x)g_{\mu \nu }(x)$$
(3)
$$S[\phi ,g_{\mu \nu }^\omega ]=S[\phi ,g_{\mu \nu }]$$
(4)
for any (strictly positive) $`\omega (x)`$. In adequate coordinates, any metric in two dimensions is conformally flat:
$$g_{\mu \nu }=e^{\sigma (x)}\eta _{\mu \nu }$$
(5)
where $`\eta _{\mu \nu }=\mathrm{diag}(1,1)`$ denotes the flat Minkowski metric. The classical dynamics of a scalar field in a gravitational background in $`1+1`$ dimensions is therefore trivial:
$$S[\phi ,g_{\mu \nu }]=S[\phi ,\eta _{\mu \nu }]=\frac{1}{2}d^2x\eta ^{\mu \nu }_\mu \phi _\nu \phi .$$
(6)
This conclusion does not hold true, of course, for the quantum dynamics, the reason being the existence of the conformal or trace anomaly, which spoils the symmetry under the transformations (3) and produces a non vanishing trace in the energy momentum tensor .
In the usual setting, one derives the quantum effects by dealing with the vacuum functional (1). Under a Weyl transformation (3), and using the property (4), we see that
$$𝒵[g_{\mu \nu }^\omega ]=[𝒟\phi ]_{g^\omega }e^{iS[\phi ,g_{\mu \nu }]}.$$
(7)
Namely, any possible quantum effect must come from the non-invariance of the integration measure under Weyl transformations of the metric.
As shown by Polyakov \- , there is, indeed, an anomalous Jacobian, the exponential of the Liouville action
$$[𝒟\phi ]_{g^\omega }=[𝒟\phi ]_gJ[\sigma ]=[𝒟\phi ]_g\mathrm{exp}i𝒜[g_{\mu \nu },w]$$
(8)
where $`𝒜[g_{\mu \nu },w]`$ is the anomaly, a local functional of $`g_{\mu \nu }`$.
By using the decomposition (5), and normalizing $`𝒵`$ such that $`𝒵[\eta _{\mu \nu }]=1`$, one also sees that
$$𝒵[g_{\mu \nu }]=J[\sigma ]=\mathrm{exp}i𝒜[w],$$
(9)
where $`𝒜[w]𝒜[\eta _{\mu \nu },w]`$ is given by
$$𝒜[w]=\frac{1}{96\pi }d^2x\frac{1}{w^2}_\mu w^\mu w.$$
(10)
As can be easily proved from the perturbative derivation of the anomaly given in Ref. , when the system consists of $`d`$ bosonic scalar fields and $`\overline{d}`$ Grassmann scalar fields the Jacobian becomes $`J[\sigma ]=\mathrm{exp}i(d\overline{d})𝒜[g,w]`$.
In what follows we will show that, if one works with the regularized action rather than the classical action $`S[\phi ,g]`$, there is a non local symmetry, and one can compute the anomaly as the Jacobian associated to a non local redefinition of the fields. The Pauli-Villars method, first applied to this system in , amounts to the introduction of three ‘regulator fields’ $`\overline{\chi },\chi ,\eta `$ and a cutoff $`M`$ so that
$$S^{reg}[\phi ,\overline{\chi },\chi ,\eta ,g_{\mu \nu },M]=\frac{1}{2}d^2x\sqrt{g}(g^{\mu \nu }_\mu \phi _\nu \phi +g^{\mu \nu }_\mu \overline{\chi }_\nu \chi M^2\overline{\chi }\chi $$
$$+g^{\mu \nu }_\mu \eta _\nu \eta 2M^2\eta ^2)$$
(11)
where $`\overline{\chi },\chi `$ are complex Grassmann fields, while $`\eta `$ is a real scalar field. The number and massess of the regulators are just right to make the vacuum energy in a non-trivial background finite. For the sake of convenience, we rewrite (11) in a more compact form as follows:
$$S^{reg}[\phi ,\overline{\chi },\chi ,\eta ,g_{\mu \nu },M]=\frac{1}{2}\left[\phi |\mathrm{}^{}|\phi +\overline{\chi }|(\mathrm{}^{}+M^2\sqrt{g})|\chi +\eta |(\mathrm{}^{}+2M^2\sqrt{g})|\eta \right]$$
(12)
where
$$\mathrm{}^{}=\sqrt{g}\mathrm{}=_\mu \sqrt{g}g^{\mu \nu }_\nu $$
(13)
and $`\mathrm{}`$ is the usual curved space Laplacian operator $`\mathrm{}=\frac{1}{\sqrt{g}}_\mu \sqrt{g}g^{\mu \nu }_\nu `$. We used a Dirac braket like notation for the scalar product in flat two dimensional spacetime:
$$\varphi _1|\varphi _2=d^2x\varphi _1(x)\varphi _2(x),$$
(14)
which is convenient, since we have absorbed the factor $`\sqrt{g}`$ of the measure in the operators. It is worth noting that both $`\mathrm{}^{}`$ and $`\sqrt{g}`$ are symmetric (real Hermitian) operators for the scalar product (14). Also, $`\mathrm{}^{}`$ is explicitly invariant under Weyl transformations.
Due to the presence of the mass terms for the regulator fields, it is evident that the regularized action is not invariant under Weyl transformations of the metric tensor:
$$S^{reg}[\phi ,\overline{\chi },\chi ,\eta ,g_{\mu \nu }^\omega ,M]=\frac{1}{2}\left[\phi |\mathrm{}^{}|\phi +\overline{\chi }|(\mathrm{}^{}+M^2\omega \sqrt{g})|\chi +\eta |(\mathrm{}^{}+2M^2\omega \sqrt{g})|\eta \right]$$
$$S^{reg}[\phi ,\overline{\chi },\chi ,\eta ,g_{\mu \nu },M].$$
(15)
We may, however, compensate the non invariance of $`S^{reg}`$ by a transformation of the regulator fields:
$`|\phi ^w>`$ $`=`$ $`|\phi >`$ (16)
$`|\eta ^w>`$ $`=`$ $`(\mathrm{}^{}+2M^2w\sqrt{g})^{1/2}(\mathrm{}^{}+2M^2\sqrt{g})^{1/2}|\eta >`$ (17)
$`|\chi ^w>`$ $`=`$ $`(\mathrm{}^{}+M^2w\sqrt{g})^{1/2}(\mathrm{}^{}+M^2\sqrt{g})^{1/2}|\chi >`$ (18)
after which we see that
$$S_{reg}[\phi ^w,\eta ^w,\overline{\chi }^w,\chi ^w,g^w,M]=S_{reg}[\phi ,\eta ,\overline{\chi },\chi ,g,M].$$
(19)
Eqs. (18 -19) define the non local symmetry of the regularized action. This symmetry must be studied of course at the quantum level, by considering the regularized vacuum functional
$$Z_{reg}[g_{\mu \nu }]=[𝒟\phi 𝒟\eta 𝒟\overline{\chi }𝒟\chi ]_g\mathrm{exp}(iS_{reg}[\phi ,\eta ,\overline{\chi },\chi ,g,M]).$$
(20)
Performing the Weyl transformation (3) for $`g_{\mu \nu }`$ in (20) followed by the change of variables (18) in the functional integral, we see that:
$`Z_{reg}[g_{\mu \nu }^w]`$ $`=`$ $`{\displaystyle [𝒟\phi ^w𝒟\eta ^w𝒟\overline{\chi }^w𝒟\chi ^w]_{g^w}\mathrm{exp}(iS_{reg}[\phi ,\eta ,\overline{\chi },\chi ,g,M])}`$ (21)
$`=`$ $`J_{reg}[g,w,M]{\displaystyle [𝒟\phi 𝒟\eta 𝒟\overline{\chi }𝒟\chi ]_g\mathrm{exp}(iS_{reg}[\phi ,\eta ,\overline{\chi },\chi ,g,M])}`$ (22)
where $`J_{reg}`$ denotes the Jacobian for the transformation (18):
$$J_{reg}[g,w,M]=det[(\mathrm{}^{}+M^2\sqrt{g})^1(\mathrm{}^{}+M^2w\sqrt{g})]$$
$$\times det[(\mathrm{}^{}+2M^2\sqrt{g})^1(\mathrm{}^{}+2M^2w\sqrt{g})]^{1/2},$$
(23)
and the suffix ‘reg’ is used because, as we will see now, this Jacobian is finite. It is worth noting that, although the integration measure for each field $`[𝒟\varphi _i]`$ depends non trivially on the background metric, as in the unregularized case, the product of the integration measures for the four fields does not depend on $`w`$, due to the cancellation between the anomalous factors corresponding to bosonic and Grassmann fields. These Jacobian factors are independent of the masses of the fields.
If we define the finite quantities $`Z`$ and $`J`$ as the limit of $`Z_{reg}`$ and $`J_{reg}`$ for $`M\mathrm{}`$ equation (22) implies that $`Z[g^w]=J[g,w]Z[g]`$. Thus we have to evaluate Eq.(23) in the limit $`M\mathrm{}`$.
Let us now calculate the regulated Jacobian (23). As the metric is assumed to be conformally flat, and $`\mathrm{}^{}`$ is invariant under (3), it is obvious that we may replace $`\mathrm{}^{}`$ by $`\mathrm{}`$, and that $`\sqrt{g}=1`$. We then rewrite $`J_{reg}`$ in the form:
$$J_{reg}[\eta ,w,M]=J_{reg}^{(1)}[w,M]\times J_{reg}^{(2)}[w,M]$$
(24)
where
$`J_{reg}^{(1)}[w,M]`$ $`=`$ $`det\left[(\mathrm{}+M^2)^1(\mathrm{}+M^2w)\right]`$ (25)
$`J_{reg}^{(2)}[w,M]`$ $`=`$ $`\left[J_{reg}^{(1)}[w,\sqrt{2}M]\right]^{1/2}.`$ (26)
We shall now, for calculational purposes, consider $`J_{reg}^{(1)}`$ alone, since the factor $`J_{reg}^{(2)}`$ can be obtained from it by some simple substitutions. However, the two factors have to be taken together for the cancellation between UV divergences to happen.
To evaluate $`J_{reg}^{(1)}[w,M]`$, we take into account the fact that we will, in the end, be interested in the $`M\mathrm{}`$ limit. This justifies the use of some form of expansion in powers of $`\frac{1}{M}`$. A small dimensionless parameter has then to be built using $`M`$, and the only other dimensionful object: derivatives of $`\omega `$ ($`\omega `$ itself is dimensionless). We then follow the derivative expansion technique to split the field $`\omega `$ into a slowly varying part $`\stackrel{~}{\omega }`$ plus a fluctuating piece $`\alpha `$
$$\omega (x)=\stackrel{~}{\omega }(x)+\alpha (x),$$
(27)
where $`\stackrel{~}{\omega }(x)`$ is to be regarded as a constant when acted by the derivative operator. We then rotate to Euclidean spacetime, and expand $`\mathrm{ln}J_{reg}^{(1)}[w,M]`$ in powers of $`\alpha `$, starting from
$$\mathrm{ln}J_{reg}^{(1)}[w,M]=\mathrm{Tr}\mathrm{ln}(\frac{^2+M^2\stackrel{~}{\omega }}{^2+M^2})+\mathrm{Tr}\mathrm{ln}\left(1+\frac{M^2\alpha }{^2+M^2\stackrel{~}{\omega }}\right),$$
(28)
where $`^2=\delta _{\mu \nu }_\mu _\nu `$ is the flat, Euclidean spacetime Laplace operator. Taking into account that the linear term vanishes, and terms with more than two $`\alpha `$’s are suppressed by negative powers of $`M`$, it is sufficient to use the expansion:
$$\mathrm{ln}J_{reg}^{(1)}[w,M]=\mathrm{Tr}\mathrm{ln}(\frac{^2+M^2\stackrel{~}{\omega }}{^2+M^2})$$
$$\frac{1}{2}\mathrm{Tr}\left(\frac{M^2\alpha }{^2+M^2\stackrel{~}{\omega }}\frac{M^2\alpha }{^2+M^2\stackrel{~}{\omega }}\right)+𝒪(\frac{1}{M^2}).$$
(29)
The first, zero derivative term, is divergent (even for finite $`M`$). However, remembering that we are using a Pauli-Villars scheme, we have to evaluate the momentum integral that results from combining it with the corresponding contribution from $`\mathrm{ln}J^{(2)}`$. This produces a finite answer:
$$\mathrm{Tr}\mathrm{ln}(\frac{^2+M^2\stackrel{~}{\omega }}{^2+M^2})\frac{1}{2}\mathrm{Tr}\mathrm{ln}(\frac{^2+2M^2\stackrel{~}{\omega }}{^2+2M^2})=\frac{M^2\mathrm{ln}2}{4\pi }(\stackrel{~}{\omega }1),$$
(30)
proportional to $`M^2`$. This is finite (for finite $`M`$), and this shows that the Jacobian for the non local symmetry transformations is indeed finite. Of course, when $`M`$ tends to infinite, this contribution diverges. The term proportional to $`\stackrel{~}{\omega }`$ requires the introduction of a counterterm of the cosmological constant type. The $`\stackrel{~}{\omega }`$-independent divergence in Eq. 30 can be absorbed into the normalization factor of $`𝒵`$.
For the second order term (which is finite when $`M\mathrm{}`$), a standard calculation yields, for $`M\mathrm{}`$
$$\frac{1}{2}\mathrm{Tr}[\frac{M^2}{^2+M^2\stackrel{~}{\omega }}\alpha \frac{M^2}{^2+M^2\stackrel{~}{\omega }}\alpha ]=\frac{1}{48\pi }d^2x\frac{1}{\stackrel{~}{w}^2}_\mu \alpha ^\mu \alpha .$$
(31)
The derivative expansion technique implies , on the other hand, that $`\stackrel{~}{\omega }`$ may be replaced by $`\omega `$ in a second order term, and that derivatives of $`\alpha `$ are tantamount to derivatives of $`\omega `$. Then,
$$\frac{1}{2}\mathrm{Tr}[\frac{M^2}{^2+M^2\stackrel{~}{\omega }}\alpha \frac{M^2}{^2+M^2\stackrel{~}{\omega }}\alpha ]=\frac{1}{48\pi }d^2x\frac{1}{\omega ^2}_\mu \omega ^\mu \omega .$$
(32)
This contribution has to be combined with the second order term coming from $`J_{reg}^{(2)}`$, which only differs in a $`\frac{1}{2}`$ global factor. Then,
$$\underset{M\mathrm{}}{lim}J_{reg}[\eta ,w,M]=\frac{1}{96\pi }d^2x\frac{1}{\omega ^2}_\mu \omega ^\mu \omega .$$
(33)
which is the (Euclidean) Liouville action. Rotating back to Minkowski spacetime we obtain the result given in Eq. (9).
## III The chiral anomaly
This example shares many properties with the previous one of the conformal anomaly, and helps to understand the general nature of the procedure we have applied in section II.
We shall consider here $`𝒵[A]`$, the vacuum functional for a massless fermion in $`1+1`$ dimensions,
$$𝒵[A]=[𝒟\overline{\psi }𝒟\psi ]_Ae^{iS_F[\overline{\psi },\psi ;A]}$$
(34)
with
$$S_F[\overline{\psi },\psi ;A]=d^2x\overline{\psi }(i\overline{)}e\overline{)}A)\psi .$$
(35)
Again, the integration measure depends on the background field configuration. The background gauge field $`A`$ may be decomposed as follows:
$$eA_\mu =_\mu \phi +ϵ_{\mu \nu }_\nu \rho ,$$
(36)
where $`\phi `$ and $`\rho `$ are scalar and pseudoscalar fields, respectively. Then one easily sees that the fermionic action may be rewritten as
$$S_F[\overline{\psi },\psi ;A]=d^2x\overline{\psi }(i\overline{)}\overline{)}(\phi +\gamma _5\rho )\psi ,$$
(37)
what means that the gauge field may actually be erased by a gauge plus chiral transformation of the fermions. This, as for the previous example, implies that the classical dynamics of the system is trivial. However, the anomaly under chiral fermionic transformations introduces a non trivial quantum dynamics.
In the usual derivation, one performs chiral transformations of the fermions,
$$\psi (x)=e^{i\alpha (x)\gamma _5}\psi (x),\overline{\psi }(x)=\overline{\psi }(x)e^{i\alpha (x)\gamma _5},$$
(38)
and the chiral anomaly appears from the non-invariance of the fermionic measure.
To consider an alternative derivation, we note that the Pauli-Villars regularized action in this case requires the addition of just one (massive) bosonic spinor field $`\varphi `$, such that
$$S_F^{reg}[\overline{\psi },\psi ,\overline{\varphi },\varphi ;A]=d^2x\left[\overline{\psi }(i\overline{)}e\overline{)}A)\psi +\overline{\varphi }(i\overline{)}e\overline{)}AM)\varphi \right].$$
(39)
When $`\alpha `$ is a constant, the non-local infinitesimal chiral symmetry transformations of (39) are,
$`\delta \psi `$ $`=`$ $`i\alpha \gamma _5\psi `$ (40)
$`\delta \overline{\psi }`$ $`=`$ $`i\alpha \overline{\psi }\gamma _5`$ (41)
$`\delta \varphi `$ $`=`$ $`i\alpha \gamma _5{\displaystyle \frac{\overline{)}D}{\overline{)}DiM}}\varphi `$ (42)
$`\delta \overline{\varphi }`$ $`=`$ $`i\alpha \overline{\varphi }{\displaystyle \frac{\overline{)}D}{\overline{)}DiM}}\gamma _5,`$ (43)
where
$$\overline{)}D=\overline{)}+ie\overline{)}A.$$
(44)
The action is invariant under these transformations, while the Jacobian becomes
$$J=det\left[1+i\alpha \gamma _5(1\frac{\overline{)}D}{\overline{)}DiM})\right]^2$$
(45)
which may be rewritten as
$$J=\mathrm{exp}\left[2i\alpha \mathrm{Tr}\gamma _5(\frac{1}{1+\frac{\overline{)}D^2}{M^2}})\right].$$
(46)
The funtional trace is finite, and reproduces the proper result when $`M\mathrm{}`$:
$$J=\mathrm{exp}[i\frac{e}{2\pi }\alpha d^2xϵ^{\mu \nu }_\mu A_\nu ].$$
(47)
When $`\alpha `$ is spacetime dependent, the regularized action is no longer invariant. However, we may use that kind of transformation to get rid of the dependence in $`A_\mu `$. Those transformations are defined by
$$\delta \psi =i\alpha (x)\gamma _5\psi ,\delta \overline{\psi }=i\overline{\psi }\gamma _5\alpha (x)$$
$$\delta \varphi =i\gamma _5\frac{\overline{)}D}{\overline{)}DiM}\alpha (x)\varphi ,\delta \overline{\varphi }=i\overline{\varphi }\alpha (x)\frac{\overline{)}D}{\overline{)}DiM}\gamma _5,$$
(48)
and the corresponding variation of the action is
$$\delta S_F^{reg}=d^2x[\overline{\psi }\overline{)}(\gamma _5\alpha )\psi +\overline{\varphi }\overline{)}(\gamma _5\alpha )\varphi ].$$
(49)
The Jacobian is easily shown to be
$$J=\mathrm{exp}[i\frac{e}{2\pi }d^2x\alpha (x)ϵ^{\mu \nu }_\mu A_\nu ].$$
(50)
## IV Conclusions
We have presented two concrete examples of systems where the regularized action has a non local symmetry which is the natural extension of the standard symmetry of the unregularized action. Moreover, the application of those transformations in the functional integral framework yields regularized Jacobians which properly reproduce the anomalies when the cutoff tends to infinity. This makes the connection between the regularization of the diagrams of a model, and the regularization of its Jacobian more transparent than in the usual setting.
We remark that our method differs also from performing the usual (local) transformations to the regularized action. This procedure would give no Jacobian, due to the cancellation between bare fields and regulators, while the regularized action would be non-invariant. The anomaly would appear in this case from the non-invariance of the regularized action.
Finally, we want to remark that the phenomenon we have described, namely, the existence of a remnant of the classical symmetry for the regularized action is not new. It has been recently emphasized that massless fermions on the lattice, even thought the regulatization breaks the naive chiral symmetry, may have a lattice equivalence of that symmetry, if the lattice Dirac operator satisfies the Ginsparg-Wilson relation . Indeed, this relation can be used to derive the chiral anomaly and the related index theorems on the lattice .
###### Acknowledgements.
This work was supported by Universidad de Buenos Aires, CONICET (Argentina), Fundacion Antorchas, ANPCyT, CNEA and the Abdus Salam International Centre for Theoretical Physics. The authors would like to thank the Abdus Salam International Centre for Theoretical Physics for hospitality during completion of this work.
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# ELECTROWEAK PHYSICS
## 1 INTRODUCTION <sup>1</sup><sup>1</sup>1Section coordinators: W. Hollik, Z. Kunszt.
### 1.1 Electroweak parameters
At the LHC, substantial improvement in the precise determination of electroweak parameters, such as the $`W`$ boson mass, the top-quark mass and the electroweak mixing angle, will become feasible, as well as an accurate measurement of the vector-boson self couplings and of the mass of the Higgs boson. This opens promising perspectives towards very comprehensive and challenging tests of the electroweak theory.
Electroweak precision observables provide the basis for important consistency tests of the Standard Model (SM) or its extensions, in particular the Minimal Supersymmetric Standard Model (MSSM). By comparing precision data with the predictions of specific models, it is possible to derive indirect constraints on the parameters of the model. In the case of the top-quark mass, $`m_t`$, the indirect determination from the precision observables in the framework of the SM turned out to be in remarkable agreement with the direct experimental measurement of $`m_t`$. Since the Higgs boson mass, $`M_H`$, enters the predictions for the precision observables only logarithmically in leading order, the indirect determination of $`M_H`$ requires very accurate experimental data as well as high precision of the theoretical predictions. The uncertainties of the predictions arise from the following sources: a) the unknown higher-order corrections - since the perturbative evaluation is truncated at a certain order, and b) the parametric uncertainties induced by the experimental errors of the input parameters.
The most important universal top-quark contribution to the electroweak precision observables enters via the $`\rho `$ parameter, which deviates from unity by a loop contribution $`\mathrm{\Delta }\rho `$. At the one-loop level, the $`(t,b)`$ doublet yields a term proportional to $`m_t^2`$ , namely $`\mathrm{\Delta }\rho =3G_\mu m_t^2/(8\pi ^2\sqrt{2})`$ in the limit $`m_b0`$. Therefore, it is to be expected that the precision measurement of the top-quark mass at the LHC (see Section 3.1) will significantly improve the theoretical prediction of the $`W`$ mass, $`M_W`$ – at present, the experimental error on $`m_t`$ is a limiting factor for the accuracy in the theoretical predictions of the precision observables. $`M_W`$ itself will be measured at the LHC with a sizably improved accuracy.
The theoretical prediction for $`M_W`$ is obtained from the relation between the vector-boson masses $`M_{W,Z}`$ and the Fermi constant $`G_\mu `$, which is conventionally written in the form
$$M_W^2\left(1\frac{M_W^2}{M_Z^2}\right)=\frac{\pi \alpha }{\sqrt{2}G_\mu }\frac{1}{1\mathrm{\Delta }r}.$$
(1)
The quantity $`\mathrm{\Delta }r=\mathrm{\Delta }r(\alpha ,M_Z,M_W,m_t,M_H)`$, first derived in in one-loop order, summarises the quantum corrections to the vector-boson mass correlation; it is obtained from the calculation of the muon lifetime in the SM beyond the tree-level approximation. At one-loop order, $`\mathrm{\Delta }r`$ can be written as
$$\mathrm{\Delta }r=\mathrm{\Delta }\alpha \frac{c_W^2}{s_W^2}\mathrm{\Delta }\rho +(\mathrm{\Delta }r)_{\mathrm{rem}}.$$
(2)
$`\mathrm{\Delta }\alpha `$ contains the large logarithmic contributions from the light fermions, and $`\mathrm{\Delta }\rho `$ the $`m_t^2`$ dependence; the non-leading terms are collected in $`(\mathrm{\Delta }r)_{\mathrm{rem}}`$ where also the dependence on $`M_H`$ enters. In Equation 1, $`\mathrm{\Delta }r`$ is a quantity that accounts also for terms of higher order than just one-loop. Moreover, a partial resummation of large contributions from light fermions and from the $`\rho `$ parameter is contained in the expression. For a discussion see for example the section on the Electroweak Working Group Report in . Results for $`M_W`$ that were not yet available at the time of the report are the next-to-leading two-loop terms of $`𝒪(G_\mu ^2m_t^2M_Z^2)`$ in an expansion for asymptotically large $`m_t`$ and the result for the Higgs mass dependence of the fermionic two-loop contributions . Recently, the complete result for the fermionic two-loop contributions has been obtained . Furthermore, the QCD corrections to $`\mathrm{\Delta }r`$ of $`𝒪(\alpha \alpha _s^2)`$ have been derived .
The most recent theoretical prediction for $`M_W`$ within the SM is displayed in Figure 1 as a function of $`M_H`$. To illustrate the comparison between theory and experiment, the experimental result is included in the figure for the current uncertainty $`\delta M_W=\pm 0.042`$ GeV and the estimated LHC uncertainty $`\delta M_W=\pm 0.015`$ GeV (see Section 3.1) (assuming the same central value). The uncertainty for the current status and for the case where the LHC will have measured the top-quark mass with much higher accuracy is also displayed, in combination with the theoretical uncertainty from unknown higher-order corrections. It is clear that both improvements, in $`M_W`$ and in $`m_t`$, will lead to a substantial increase in the significance of Standard Model tests, with stringent bounds on the Higgs boson mass to be confronted with the directly measured value of $`M_H`$.
Besides the $`W`$ boson mass, the improvement in $`m_t`$ will also have an effect on the predictions of the $`Z`$ pole observables. They are conveniently described in terms of effective couplings
$$g_V^f=\sqrt{\rho _f}(I_3^f2Q_f\mathrm{sin}^2\theta _{\mathrm{eff}}^f),g_A^f=\sqrt{\rho _f}I_3^f$$
(3)
in the neutral-current vertex at the $`Z`$ resonance for a given fermion species $`f`$, normalised according to $`J_\mu ^{\mathrm{NC}}=(\sqrt{2}G_\mu M_Z^2)^{1/2}(g_V^f\gamma _\mu g_A^f\gamma _\mu \gamma _5)`$. Besides the overall normalisation factor $`\rho _f=1+\mathrm{\Delta }\rho +\mathrm{}`$, we mention in particular the effective mixing angle, which is usually chosen as the on-resonance mixing angle for the leptons $`f=e,\mu ,\tau `$ in Equation 3 and denoted as $`\mathrm{sin}^2\theta _{\mathrm{eff}}^{\mathrm{lept}}`$. This quantity also depends sensitively on the top-quark mass, mainly through $`\mathrm{\Delta }\rho `$. The theoretical prediction of $`\mathrm{sin}^2\theta _{\mathrm{eff}}^{\mathrm{lept}}`$ will definitely be sharpened by the precise measurement of the top-quark mass; a sizable improvement concerning the internal consistency test can be anticipated. The on-resonance mixing angle for the light quarks $`b`$ is numerically very close to the leptonic one. $`\mathrm{sin}^2\theta _{\mathrm{eff}}^{\mathrm{lept}}`$ can therefore be measured at the LHC in the Drell-Yan production of charged-lepton pairs around the $`Z`$ resonance, via $`q\overline{q}l^+l^{}`$, where an accuracy of $`1.4\times 10^4`$ on $`\mathrm{sin}^2\theta _{\mathrm{eff}}^{\mathrm{lept}}`$ may be feasible (see Section 3.2).
Besides these internal consistency checks of the SM, the electroweak precision observables may be useful to distinguish between different models as candidates for the electroweak theory. In Figure 2, the SM prediction of $`M_W`$ as a function of $`m_t`$ is compared with the prediction within the MSSM, where the MSSM prediction is based on results up to $`𝒪(\alpha \alpha _s)`$ . The SM uncertainty arises from the only unknown parameter, the Higgs boson mass. On the other hand, within the MSSM, the Higgs boson mass is not a free parameter , and the uncertainty originates from the unknown SUSY mass scales. In the small overlap region, the MSSM behaves like the SM, i.e. all SUSY particles are heavy and decouple from the precision observables, and the $`M_H`$ value of the SM stays below 130 GeV, the upper bound on the lightest MSSM Higgs boson mass for $`m_t=175`$ GeV (see and references therein). Figure 2 shows the clear improvement from the current status to the LHC era, where eventually, besides direct experimental evidence, a distinction between SM and MSSM might become feasible.
### 1.2 Vector-boson pair production and scattering
At the LHC, the precise measurement of the production of $`W^+W^{}`$, $`W^\pm Z`$, $`ZZ`$, $`W^\pm \gamma `$ or $`Z\gamma `$ pairs is also an important physics goal. In the simplest studies, the gauge-bosons will be detected via their leptonic decays. Already a couple events have been obtained by CDF and D0 for $`WW`$ and $`WZ`$ production and D0 has seen about 100 $`W\gamma `$ and 30 $`Z\gamma `$ events. The data set at Run II will be about 20 times larger and about 1000 times larger at the LHC. For a summary of the experimental situation see .
The production of gauge-boson pairs provide us with the best test of the non-Abelian gauge-symmetry of the Standard Model (SM). Deviation from the SM predictions may come either from the presence of anomalous couplings or the production of new heavy particles and their decays into vector-boson pairs. If the particle spectrum of the SM has to be enlarged with new particles (as in the Minimal Supersymmetric Standard Model (MSSM)) with mass values of $`0.51\mathrm{TeV}`$, small anomalous couplings are generated at low energy. If the Higgs boson is very heavy, it will decay mainly into $`W^+W^{}`$ and $`ZZ`$ pairs. If the symmetry breaking mechanism is dynamical (technicolor models, BESS models), large anomalous couplings might be generated or new heavy particles may be produced. In both of these cases, vector-boson pair production will show deviations from the Standard Model predictions. At the same time, vector-boson pair production gives the most important background for a number of new physics signals. For example, one of the most important physics signal for supersymmetry at hadron colliders is the production of three charged leptons and missing transverse momentum. The dominant background for this process is the production of $`W`$ plus a $`Z`$ (real or virtual) or $`\gamma `$.
The leading order production mechanism of gauge-boson pair production is $`q\overline{q}`$ annihilation. The precise calculation of the cross sections in the QCD improved parton model have received recently a lot of attention. The cross sections of the gauge-boson pair production and its decay into lepton pairs have been calculated in next-to-leading order (NLO) accuracy retaining the full spin correlations of the leptonic decay products. A significant achievement was that the theoretical results in NLO QCD for the production of $`W^+W^{}`$, $`W^\pm Z`$, $`ZZ`$, $`W^\pm \gamma `$ or $`Z\gamma `$ pairs could be documented in short analytic formulae allowing for independent numerical implementations. Subsequently, several so called NLO numerical Monte Carlo programs have been developed and the complete one loop corrections became available for the first time for $`W^+W^{}`$, $`W^\pm Z`$, $`ZZ`$ in , and for $`W^\pm \gamma ,`$ or $`Z\gamma `$ pairs in . These new results have superseded and confirmed previous NLO results on spin averaged production gauge-boson pair production , as well the approximate results where spin correlation have been neglected in the virtual corrections . The agreement between the well documented results in and in is within the precise integration error and the agreement between the results of and the recent programs of is about 3%. Therefore, previous experimental simulation studies based on these programs (see Section 6.5) should not be repeated.
Simple analytic NLO results exist also for the anomalous coupling contributions at NLO accuracy in . Again, the agreement with previous approximate NLO results is also good (see Section 5.5). Future anomalous coupling studies may like to use the more accurate packages. At the LHC, contrary to LEP, the phenomenological studies of anomalous triple gauge-boson coupling constants cannot be treated as constant couplings since they lead to violation of $`SU(2)`$ gauge-symmetry and unitarity. The difficulty comes from truncation of the contribution of an infinite series of higher dimensional non-renormalisable gauge-invariant operators. In the case of $`q\overline{q}`$ annihilation to gauge-boson pairs, a suitable phenomenological approach is the introduction of form factors for the anomalous couplings (which in principle are calculable in the true underlying theory). As long as we do not obtain deviations from the Standard Model, for practical purposes, simple dipole form factors with various cut-off parameters can be used. With better data, one can put limits on the form factor values in small $`\sqrt{\widehat{s}}`$ intervals, assuming constant couplings for each interval. In the case of positive signals, such a form factor measurement will provide us with important information on the underlying theory (see Sections 58 and 10).
At higher energies, the higher order production processes of $`WW`$ and $`ZZ`$ scattering (the weak boson are emitted from the incoming quarks) will become more and more important. These interactions are the most sensitive to the mechanism of the electroweak symmetry breaking. In particular, if the breaking of the electroweak symmetry is due to new particles with strong interactions at the TeV scale, enhanced production of longitudinal gauge-boson pairs will be the most typical signal . The minimal model to describe this alternative is obtained by assuming that the new particles are too heavy to be produced at LHC and the linear $`\sigma `$-model Higgs-sector of the Standard Model is replaced by the non-renormalisable non-linear $`\sigma `$-model which can also be considered as an effective chiral vector-boson Lagrangian with non-linear realisation of the gauge-symmetry . The question is whether this more phenomenological approach is consistent with the precision data. In a recent analysis, a positive answer was obtained . It has been found that due to the screening of the symmetry breaking sector , this alternative still has enough flexibility to be in perfect agreement with the precision data up to a cut-off scale of $`3\mathrm{TeV}`$ (see Sections 10 and 19). In the chiral approach, the gauge-boson observables are obtained as truncated series in powers of the external momenta $`p^n/(4\pi v)^n`$ with $`M_W^2gv^2/8`$. The approximation is valid up to energy scales of $`E=4\pi v3\mathrm{TeV}`$. At the LHC, the partonic centre of mass energy can be higher and the phenomenological implementation is confronted with the problem of unitarisation . Although unitarisation is not unique, the use of the K-matrix formalism or the $`𝒪(p^4)`$ Inverse Amplitude Method appear to give reasonable model independent framework to explore the various possibilities. When extrapolating to higher energies in particular, the masses of resonances are rather sensitive to the actual value of additional chiral parameters. An alternative approach for the phenomenological formulation of the dynamical symmetry breaking consistent with the precision data is offered by the BESS model with an extended strongly interacting gauge-sector with enhanced global symmetries and with important decoupling properties at low energies. The phenomenologically acceptable technicolor models also require an enhanced global symmetry in the spectrum of the theory. In the most pessimistic parameter ranges, it is rather difficult to detect the signals of the strong $`WW`$ and $`WZ`$ scattering; therefore, one has to push the LHC analysis to its limits. In the future, further clever strategies have to be pursued for this case (see Section 19).
## 2 ELECTROWEAK CORRECTIONS TO DRELL-YAN PROCESSES <sup>3</sup><sup>3</sup>3Section coordinator: W. Hollik.
The basic parton processes for single vector-boson production are $`q\overline{q^{}}Wl\nu _l`$ and $`q\overline{q}Zl^+l^{}`$, with charged leptons $`l`$ in the final state. Investigations around the $`W`$ and $`Z`$ resonance allow a precise measurement of the $`W`$ mass and of the electroweak mixing angle from the forward-backward asymmetry. At high invariant masses of the $`l^+l^{}`$ pair, deviations from the standard cross section and $`A_{\mathrm{FB}}`$ could indicate scales of new physics, e.g. associated with an extra heavy $`Z^{}`$ or extra space dimensions. For the envisaged precision, a discussion of the electroweak higher-order contributions is necessary, on top of the QCD corrections. The electroweak corrections consist of the set of electroweak loop contributions, including virtual photons, and of the emission of real photons.
With respect to QCD, the cross sections in this section are all of lowest order, evaluated with parton distribution functions at factorisation scales $`M_W`$ (for $`W`$ production) and $`M_Z`$ (for $`Z`$ production). Hence, the numerical values are not yet directly the physical ones. They are given here to point out the structure and the size of the higher-order electroweak contributions. The QCD corrections are considered in the QCD chapter of this report, where a QCD-related uncertainty of $``$5% is estimated. For illustration, we give the values (in nb) for $`[\sigma (ppW^+)+\sigma (ppW^{})]BR(We\nu )`$ and $`\sigma (ppZ)BR(Ze^+e^{})`$ in the purely electroweak calculation (EW) and with NNLO QCD :
$`W:17.9(\mathrm{EW})`$ and $`20.3\pm 1.0(\mathrm{NNLO}),`$
$`Z:1.71(\mathrm{EW})`$ and $`1.87\pm 0.09(\mathrm{NNLO})`$.
### 2.1 Universal initial-state QED corrections
QED corrections related to the emission of (real or virtual) photons from quarks contain mass singularities which factorise and therefore can be absorbed by a redefinition (renormalisation) of parton distribution functions . This redefinition is well-known in the calculation of QCD radiative corrections where in complete analogy to photon radiation, the emission of gluons leads to mass singularities as well. By the redefinition, the mass singularities disappear from the observable cross section and the renormalised distribution functions become dependent on the factorisation scale $`\mu `$ which is controlled by the well-known Gribov-Lipatov-Altarelli-Parisi (GLAP) equations . The factorisation scale should be identified with a typical scale of the process, i.e. a large transverse momentum, or the mass of a produced particle.
Since mass singularities are universal, i.e. independent of the process under consideration, the definition of renormalised parton distributions is also universal. Therefore it is possible to discuss the bulk of initial-state QED radiative corrections in terms of parton distribution functions. This will be true if there is only one large scale in the process.
The treatment of mass singularities due to gluonic or photonic radiation is identical. Photonic corrections can therefore be taken into account by a straightforward modification of the standard GLAP equations which describe gluonic corrections only. The modification corresponds to the addition of a term of the order of the electromagnetic fine-structure constant, $`\alpha `$. Apart from small non-singular contributions, the resulting modified scale dependence of parton distribution functions is the only observable effect of initial-state QED corrections in high-energy scattering of hadrons.
The modified evolution equation for the charged parton distribution functions, $`q_f(x,\mu ^2)`$ for quarks with flavour $`f`$, can be written as:
$$\begin{array}{cc}\frac{d}{dt}q_f(x,t)=\hfill & \frac{\alpha _s(t)}{2\pi }_x^1\frac{dz}{z}\left[P_{q/q}(z,t)q_f(x/z,t)+P_{q/g}(z,t)g(x/z,t)\right]\hfill \\ & +\frac{\alpha (t)}{2\pi }_x^1\frac{dz}{z}P_{q/q}^\gamma (z,t)q_f(x/z,t)\hfill \end{array}$$
(4)
In the leading logarithmic approximation, the splitting functions $`P_{i/j}`$ are independent of the scale $`t=\mathrm{ln}\mu ^2/\mathrm{\Lambda }^2`$, and the QED splitting function is given by
$$P_{q/q}^\gamma (z)=Q_f^2\left[\frac{1+z^2}{(1z)_+}+\frac{3}{2}\delta (1z)\right]=\frac{Q_f^2}{C_F}P_{q/q}.$$
(5)
Since quarks are coupled through the splitting function $`P_{q/g}(z)=\frac{1}{2}\left[z^2+(1z)^2\right]`$ to gluons, the gluon distribution $`g(x,\mu ^2)`$ is affected by QED corrections as well, although only indirectly, by terms of the order of $`𝒪(\alpha \alpha _s)`$. $`\alpha (t)`$ is the running electromagnetic fine-structure constant and $`Q_f`$ are the fermion charges in units of the positron charge.
The proper treatment of the mass-singular initial-state QED corrections would require not only the solution of the evolution equations with the QED term, but also to correct all data that are used to fit the parton distributions for those QED effects. Apart from a few exceptions, experimental data have not been corrected for photon emission from quarks. However, one can illustrate the QED radiative corrections by comparing the modification of the parton distributions relative to the distribution functions obtained from the evolution equations without the QED terms, which are used as an input.
The solution of the evolution equations corresponds to the resummation of terms containing factors $`\alpha (\alpha _s\mathrm{ln}\mu ^2)^n`$ with arbitrary power $`n`$. In Figures 3a and 3b, we show numerical results for the corrections $`\mathrm{\Delta }_{QED}`$ to the distribution functions $`U(x,\mu ^2)`$ ($`D(x,\mu ^2)`$) for the sum of all up-(down)-type quarks, and the gluon distribution $`g(x,\mu ^2)`$. The figures show the QED corrections in per cent relative to the distribution functions obtained from the GLAP equations without the QED term. The input distributions were taken from . One finds small, negative corrections at the per-mille level for all values of $`x`$ and $`\mu ^2`$ relevant in the LHC experiments. Only at large $`x\stackrel{>}{}0.5`$ and large $`\mu ^2\stackrel{>}{}10^3`$ GeV<sup>2</sup> do the corrections reach the magnitude of one per cent. The increase of corrections for $`x1`$ is due to the $`\mathrm{ln}(1x)`$ terms appearing in the evaluation of the “$`+`$” distributions.
The largest corrections are obtained for up-type quarks due to the larger charge factor $`4/9`$ as compared to $`1/9`$ for down-type quarks. The gluon distribution, being of order $`𝒪(\alpha \alpha _s)`$, is corrected by less than $`0.1\%`$ up to values of $`x`$ of about 0.2.
The corrections vanish for $`\mu ^2\mu _0^2`$ since it was assumed that the input distributions $`q_f(x,\mu _0^2)`$ and $`g(x,\mu _0^2)`$ have been extracted from experiment at the reference scale $`\mu _0^2`$ without subtracting quarkonic QED corrections.
The asymptotic behaviour for $`x0`$ can be checked analytically. The singular behaviour of distributions $`x^\eta `$ for $`x0`$ remains unchanged by the GLAP equations if $`\eta >1`$. Thus the $`𝒪(\alpha )`$ corrected distributions have the same power behaviour as the uncorrected ones, the ratio consequently reaching a constant value for $`x0`$. The valence parts of $`U(x)`$ and $`D(x)`$, however, which vanish at $`x=0`$, receive positive corrections at small $`x`$, thus producing the well-known physical picture: radiation of gluons as well as of photons leads to a depletion at large $`x`$ and an enhancement at small $`x`$, i.e. partons are shifted to smaller $`x`$.
Other input distribution functions lead to differences of QED corrections at the per-mille level which are again irrelevant when compared with the expected experimental precision of structure-function measurements.
### 2.2 Electroweak corrections to $`𝑾`$ production
#### 2.2.1 Physical goals of single $`W`$ production
The Drell-Yan-like production of $`W`$ bosons represents one of the cleanest processes with a large cross section at the LHC. This reaction is not only well suited for a precise determination of the $`W`$ boson mass $`M_W`$, it also yields valuable information on the parton structure of the proton. Specifically, the target accuracy of the order of $`15\mathrm{MeV}`$ in the $`M_W`$ measurement exceeds the precision of roughly $`30\mathrm{MeV}`$ achieved at LEP2 and Tevatron Run II , and thus competes with the one of a future $`e^+e^{}`$ collider . Concerning quark distributions, precise measurements of rapidity distributions provide information over a wide range in $`x`$ ; a measurement of the $`d/u`$ ratio would, in particular, be complementary to HERA results. The more direct determination of parton-parton luminosities instead of single parton distributions is even more precise ; extracting the corresponding luminosities from Drell-Yan-like processes allows us to predict related $`q\overline{q}`$ processes at the per-cent level.
Owing to the high experimental precision outlined above, the predictions for the processes $`ppWl\nu _l`$ should attain per-cent accuracy. To this end, radiative corrections have to be included. In the following some basic features of this processes and recent progress on electroweak corrections are summarised; a discussion of QCD corrections can be found in the QCD chapter of this report.
#### 2.2.2 Lowest-order cross section and preliminaries
We consider the parton process $`u\overline{d}\nu _ll^+(+\gamma )`$, where $`u`$ and $`d`$ generically denote the light up- and down-type quarks, $`u=u,c`$ and $`d=d,s`$. The lepton $`l`$ represents $`l=e,\mu ,\tau `$. In lowest order, only the Feynman diagram shown in Figure 4 contributes to the scattering amplitude,
and the Born amplitude reads
$$_0=\frac{e^2V_{ud}^{}}{2s_\mathrm{W}^2}\left[\overline{v}_d\gamma ^\mu \omega _{}u_u\right]\frac{1}{\widehat{s}M_W^2+iM_W\mathrm{\Gamma }_W(\widehat{s})}\left[\overline{u}_{\nu _l}\gamma _\mu \omega _{}v_l\right],$$
(6)
with $`\widehat{s}`$ being the squared centre-of-mass (CM) energy of the parton system. The notation for the Dirac spinors $`\overline{v}_d`$, etc., is obvious, and $`\omega _{}=\frac{1}{2}(1\gamma _5)`$ is the left-handed chirality projector. The electric unit charge is denoted by $`e`$, the weak mixing angle is fixed by the ratio $`c_\mathrm{W}^2=1s_\mathrm{W}^2=M_W^2/M_Z^2`$ of the $`W`$ and $`Z`$ boson masses $`M_W`$ and $`M_Z`$, and $`V_{ud}`$ is the CKM matrix element for the $`ud`$ transition.
Strictly speaking, Equation (6) already goes beyond lowest order, since the $`W`$ boson width $`\mathrm{\Gamma }_W(\widehat{s})`$ results from the Dyson summation of all insertions of the (imaginary parts of the) $`W`$ self-energy. Defining the mass $`M_W`$ and the width $`\mathrm{\Gamma }_W`$ of the $`W`$ boson in the on-shell scheme (see e.g. ), the Dyson summation directly leads to a running width, i.e. $`\mathrm{\Gamma }_W(\widehat{s})|_{\mathrm{run}}=\mathrm{\Gamma }_W\times (\widehat{s}/M_W^2).`$ On the other hand, a description of the resonance by an expansion about the complex pole in the complex $`\widehat{s}`$ plane corresponds to a constant width, i.e. $`\mathrm{\Gamma }_W(\widehat{s})|_{\mathrm{const}}=\mathrm{\Gamma }_W.`$ In lowest order these two parametrisations of the resonance region are fully equivalent, but the corresponding values of the line-shape parameters $`M_W`$ and $`\mathrm{\Gamma }_W`$ differ in higher orders . The numerical difference is given by $`M_W|_{\mathrm{run}}M_W|_{\mathrm{const}}26\mathrm{MeV}`$ so that it is necessary to state explicitly which parametrisation is used in a precision determination of the $`W`$ boson mass from the $`W`$ line shape.
The differential lowest-order cross section is easily obtained by squaring the lowest-order matrix element $`_0`$ of (6),
$$\left(\frac{\mathrm{d}\widehat{\sigma }_0}{\mathrm{d}\widehat{\mathrm{\Omega }}}\right)=\frac{1}{12}\frac{1}{64\pi ^2\widehat{s}}|_0|^2=\frac{\alpha ^2|V_{ud}|^2}{192s_\mathrm{W}^4\widehat{s}}\frac{\widehat{u}^2}{|\widehat{s}M_W^2+iM_W\mathrm{\Gamma }_W(\widehat{s})|^2},$$
(7)
where $`\widehat{u}=(p_up_l)^2`$ is the squared momentum difference between the up-type quark and the lepton. The explicit factor $`1/12`$ results from the average over the quark spins and colours, and $`\widehat{\mathrm{\Omega }}`$ is the solid angle of the outgoing $`l^+`$ in the parton CM frame. The electromagnetic coupling $`\alpha =e^2/(4\pi )`$ can be set to different values according to different renormalisation schemes. It can be directly identified with the fine-structure constant $`\alpha (0)`$ or the running electromagnetic coupling $`\alpha (Q^2)`$ at a high energy scale $`Q`$. For instance, it is possible to make use of the value of $`\alpha (M_Z^2)`$ that is obtained by analysing the experimental ratio $`R=\sigma (e^+e^{}\text{hadrons})/(e^+e^{}\mu ^+\mu ^{})`$. These choices are called $`\alpha (0)`$-scheme and $`\alpha (M_Z^2)`$-scheme, respectively, in the following. Another value for $`\alpha `$ can be deduced from the Fermi constant $`G_\mu `$, yielding $`\alpha _{G_\mu }=\sqrt{2}G_\mu M_W^2s_\mathrm{W}^2/\pi `$; this choice is referred to as $`G_\mu `$-scheme.
#### 2.2.3 Electroweak corrections
The electroweak $`𝒪(\alpha )`$ corrections consist of virtual one-loop corrections and real-photonic bremsstrahlung. The corrections to resonant $`W`$ production have already been studied in ; detailed discussions of the full calculation, including non-resonant corrections, can be found in . Since in $`𝒪(\alpha ^2)`$ only two-photon bremsstrahlung has been studied so far, the following discussion is restricted to $`𝒪(\alpha )`$ corrections.
The algebraic structure of the virtual corrections allows for a factorisation of the one-loop amplitude $`_1`$ into the Born amplitude $`_0`$ and a relative correction factor $`\delta ^{\mathrm{virt}}`$. Thus, in $`𝒪(\alpha )`$ the correction to the squared amplitude reads
$$|_0+_1|^2=(1+2Re\{\delta ^{\mathrm{virt}}\})|_0|^2+\mathrm{}.$$
(8)
Since only the real part of $`\delta ^{\mathrm{virt}}`$ appears, there is no double-counting of the $`𝒪(\alpha )`$ correction that is already included in $`_0`$ by the $`iM_W\mathrm{\Gamma }_\mathrm{W}`$ term. Moreover, the factorisation trivially avoids potential problems with gauge-invariance after the introduction of the $`W`$ decay width in the resonant terms. Besides the Breit-Wigner factor in $`|_0|^2`$, there are logarithmic terms $`\mathrm{ln}(\widehat{s}M_W^2)`$ in $`\delta ^{\mathrm{virt}}`$ which are singular on resonance. The consistent replacement $`\mathrm{ln}(\widehat{s}M_W^2)\mathrm{ln}(\widehat{s}M_W^2+i\mathrm{\Gamma }_\mathrm{W}M_W)`$ accounts for a Dyson summation of resonant $`W`$ propagators in loop diagrams, without introducing problems with gauge-invariance.
The real corrections are included by adding the lowest-order cross section for the process $`u\overline{d}\nu _ll^++\gamma `$. The only non-trivial condition induced by gauge-invariance is the Ward identity for the external photon, i.e. electromagnetic current conservation. If the $`W`$ width is zero, this identity is trivially fulfilled. This remains true even for a constant width, since the $`W`$ boson mass appears only in the $`W`$ propagator denominators, i.e. the substitution $`M_W^2M_W^2iM_W\mathrm{\Gamma }_W`$ is a consistent reparametrisation of the amplitude in this case. However, if a running $`W`$ width is introduced naively, i.e. in the $`W`$ propagators only, the Ward identity is violated. The identity can be restored by taking into account those part of the fermion-loop correction to the $`\gamma WW`$ vertex that corresponds to the fermion loops in the $`W`$ self-energy leading to the width in the propagator . For an external photon, this modification simply amounts to the multiplication of the $`\gamma WW`$ vertex by the factor $`f_{\gamma WW}|_{\mathrm{run}}=1+i\mathrm{\Gamma }_W/M_W`$.
Adding virtual and real corrections, all IR divergences cancel. Mass singularities of the form $`\alpha \mathrm{ln}m_l`$ related to a final-state lepton drop out for all observables in which photons within a collinear cone around the lepton are treated inclusively, in accordance with the KLN theorem. As already discussed in Section 2.1 (see also ), mass singularities to the initial-state quarks are absorbed into renormalised quark distribution functions.
As long as one is interested in observables that are dominated by resonant $`W`$ boson production, the radiative corrections can be approximated by the corrections to the production and decay subprocesses to resonant $`W`$ bosons. Formally such an approximation can be carried out by a systematic expansion of all amplitudes about the resonance pole and is, therefore, called pole approximation (PA). In PA, the virtual correction consists of two parts. The first contribution is provided by the (constant) correction factors to the $`Wf\overline{f}^{}`$ vertex for stable (on-shell) $`W`$ bosons and is called factorisable. The second contribution, which is called non-factorisable, comprises all remaining resonant corrections. It is entirely due to photonic effects and includes, in particular, the $`\mathrm{ln}(\widehat{s}M_W^2+i\mathrm{\Gamma }_\mathrm{W}M_W)`$ terms. The difference between PA and the exact result can be estimated by $`\delta _{\mathrm{PA}}^{\mathrm{virt}}\delta ^{\mathrm{virt}}(\alpha /\pi )\mathrm{ln}(\widehat{s}/M_W^2)\mathrm{ln}(\mathrm{})`$, where $`\mathrm{ln}(\mathrm{})`$ indicates any logarithmic enhancements. In principle, also the real corrections can be treated in PA. However, since a reliable error estimate is not obvious, they are usually calculated exactly. More details about PA can be found in .
#### 2.2.4 Numerical results
The following numerical results have been obtained with the input parameters of and a constant $`W`$ width; in particular, we have $`M_W=80.35\mathrm{GeV}`$ and $`\mathrm{\Gamma }_W=2.08`$ GeV. The QED factorisation is performed in the $`\overline{\text{MS}}`$ scheme with $`M_W`$ being the factorisation scale, and the CTEQ4L quark distributions are used in the evaluation of the $`pp`$ cross section. For the partonic cross section, the CKM matrix element $`V_{ud}`$ is set to 1; for the $`pp`$ cross section a non-trivial CKM matrix is included in the parton luminosities (see ).
Figure 5 shows the total partonic cross section $`\widehat{\sigma }`$ and the corresponding relative correction $`\delta `$ for intermediate energies. Note that the total cross section and its correction is the same for all final-state leptons $`l=e,\mu ,\tau `$ in the limit of vanishing lepton masses. As expected, the $`G_\mu `$ parametrisation of the Born cross section minimises the correction at low energies, since the universal corrections induced by the running of $`\alpha `$ and by the $`\rho `$ parameter are absorbed in the lowest-order cross section. Moreover, the naive error estimate for the PA taken from above turns out to be realistic. The PA describes the correction in the resonance region within a few per mille. Table 1 contains some results on the partonic cross section and its correction up to energies in the TeV range. Far above resonance, the PA cannot follow the exact correction anymore, since non-resonant corrections become more and more important. The leading corrections are due to Sudakov logarithms of the form $`\alpha \mathrm{ln}^2(\widehat{s}/M_W^2)`$.
Figure 6 shows the transverse-momentum distribution for the lepton $`l^+`$ produced in $`ppW^+\nu _ll^+(+\gamma )`$ for the $`pp`$ CM energy $`\sqrt{s}=14\mathrm{TeV}`$ of the LHC. The transverse momenta $`p_\mathrm{T}`$ and the lepton pseudorapidity $`\eta _l`$ are restricted by $`p_{\mathrm{T},l},/p_\mathrm{T}>25\mathrm{GeV}`$ and $`|\eta _l|<1.2`$. Since we do not recombine collinear photons and leptons, the corrections for different leptons do not coincide, but differ by corrections of the form $`\mathrm{ln}(m_l/M_W)`$. In the total cross section without any cuts these logarithms cancel, and the correction is again universal for all leptons in the massless limit. Since the $`\mathrm{ln}m_l`$ corrections are strongest for electrons, and since collinear photon emission reduces the momentum of the produced lepton, the correction $`\delta `$ for electrons is more negative (positive) for large (small) momenta than in the case of the muon. In particular, Figure 6 demonstrates the reliability of the PA for transverse lepton momenta $`p_{\mathrm{T},l}\stackrel{<}{}M_W/2`$, where resonant $`W`$ bosons dominate. However, high $`p_{\mathrm{T},l}`$ values may also be interesting in searches for new physics. Table 2 shows the contributions to the total cross section divided by different ranges in $`p_{\mathrm{T},l}`$. From the above discussion of the parton cross section it is clear that the PA is not applicable for very large $`p_{\mathrm{T},l}`$, where the $`W`$ boson is far off shell.
The above results underline the importance of electroweak radiative corrections in a precise description for the $`W`$ boson cross section at the LHC. Although the corrections of $`𝒪(\alpha )`$ are well under control now, there are still some topics to be studied, such as the impact of realistic detector cuts and photon recombination procedures or the inclusion of higher-order effects.
The impact of final state photon radiation on $`W`$ observables strongly depends on the lepton identification requirements imposed by the experiment. In addition to the lepton $`p_T`$, $`p\text{/}_T`$ and pseudorapidity cuts, one usually imposes requirements on the separation of the charged lepton and the photon. For muons, the energy of the photon is required to be less than a critical value, $`E_c^\gamma `$, in a cone of radius $`R_c^\mu `$ around the muon. For electrons, the finite resolution of the electromagnetic calorimeter makes it difficult to separate electrons and photons for small opening angles between the particles. Their four momentum vectors are therefore recombined if their separation is smaller than a critical value $`R_c^e`$. Finally, uncertainties in the energy and momentum measurements of the charged lepton and the missing transverse energy need to be taken into account. They can be simulated by Gaussian smearing of the particle four-momentum vectors with standard deviation $`\sigma `$ which depends on the particle type and the detector.
To illustrate how finite detector resolution affects the size of the electroweak corrections, we show in Figure 7 the ratio of the NLO and lowest-order cross sections as a function of the $`p_T`$ of the electron in $`pp\nu _ee^+(\gamma )`$ obtained with the Monte Carlo generator WGRAD . The solid histogram shows the cross section ratio taking only transverse-momentum and pseudorapidity cuts into account. The dashed histogram displays the result obtained when, in addition, the four-momentum vectors are smeared according to the ATLAS specifications , and electron and photon momenta are combined if $`\mathrm{\Delta }R(e,\gamma )<0.07`$ . Recombining the electron and photon four-momentum vectors eliminates the mass-singular logarithmic terms of the form $`\alpha \mathrm{ln}m_e`$, and strongly reduces the size of the electroweak corrections.
### 2.3 Electroweak corrections to $`𝒁`$ production and continuum neutral-current processes
#### 2.3.1 QED corrections
The mass-singular universal QED corrections from initial-state radiation from quarks have already been discussed in Section 2.1. They are part of the quark distribution functions. The residual QED initial-state corrections, together with final-state corrections and interference of initial-final radiation are treated separately by an explicit diagrammatic computation.
A complete calculation of the QED $`𝒪(\alpha )`$ radiative corrections to $`ppZ,\gamma l^+l^{}`$ has been carried out in . The calculation is based on an explicit diagrammatic approach. The collinear singularities associated with initial-state photon radiation are factorised into the parton distribution functions (see Section 2.1). Absorbing the initial-state mass singularities into the pdf’s introduces a QED factorisation-scale dependence. The results presented here are obtained within the QED DIS scheme which is defined analogously to the QCD DIS factorisation scheme. The MRS(A) parton distributions are used, with a factorisation scale $`M_Z`$. Due to mass-singular logarithmic terms associated with photons emitted collinear with one of the final-state leptons, QED radiative corrections strongly affect the shape of the di-lepton invariant mass distribution, the lepton transverse momentum spectrum, and the forward-backward asymmetry, $`A_{\mathrm{FB}}`$.
The effect of the QED corrections on the di-muon invariant mass distribution in the region $`45`$ GeV $`<m(\mu ^+\mu ^{})<105`$ GeV is shown in Figure 8a where we plot the ratio of the $`𝒪(\alpha ^3)`$ and lowest-order differential cross sections as a function of $`m(\mu ^+\mu ^{})`$. The lowest-order cross section has been evaluated in the effective Born approximation (EBA) which already takes into account those higher-order corrections which can be absorbed into a redefinition of the coupling constants and the effective weak mixing angle. More details on the EBA can be found in Section 2.3.2. In the region shown in the figure, the cross-section ratio is seen to vary rapidly. Below the $`Z`$ peak, QED corrections significantly enhance the cross section. At the $`Z`$ pole, the differential cross section is reduced by about 20%. Photon radiation from one of the leptons lowers the di-lepton invariant mass. Therefore, events from the $`Z`$ peak region are shifted towards smaller values of $`m(\mu ^+\mu ^{})`$, thus reducing the cross section in and above the peak region, and increasing the rate below the $`Z`$ pole. Final-state radiative corrections completely dominate over the entire mass range considered. They are responsible for the strong modification of the di-lepton invariant mass distribution. In contrast, initial-state corrections are uniform and small ($`+0.4\%`$ in the QED DIS scheme).
As pointed out earlier, at the LHC a precise measurement of the effective mixing angle $`\mathrm{sin}^2\theta _{\mathrm{eff}}^{\mathrm{lept}}`$ using the forward-backward asymmetry may be possible. In Figure 8b, the forward-backward asymmetry is shown in the EBA (dashed line), and including QED corrections (solid line) for $`pp\mu ^+\mu ^{}(\gamma )`$ in the di-muon invariant mass range from 45 GeV to 105 GeV. Here, $`A_{\mathrm{FB}}`$ is defined by
$$A_{\mathrm{FB}}=\frac{FB}{F+B}$$
(9)
where
$$F=_0^1\frac{\mathrm{d}\sigma }{\mathrm{d}\mathrm{cos}\theta ^{}}\mathrm{d}\mathrm{cos}\theta ^{},B=_1^0\frac{\mathrm{d}\sigma }{\mathrm{d}\mathrm{cos}\theta ^{}}\mathrm{d}\mathrm{cos}\theta ^{}.$$
(10)
$`\mathrm{cos}\theta ^{}`$ is given by
$$\mathrm{cos}\theta ^{}=\frac{|p_z(\mu ^+\mu ^{})|}{p_z(\mu ^+\mu ^{})}\frac{2}{m(\mu ^+\mu ^{})\sqrt{m^2(\mu ^+\mu ^{})+p_T^2(\mu ^+\mu ^{})}}\left[p^+(\mu ^{})p^{}(\mu ^+)p^{}(\mu ^{})p^+(\mu ^+)\right]$$
(11)
in the Collins-Soper frame , with
$$p^\pm =\frac{1}{\sqrt{2}}\left(E\pm p_z\right),$$
(12)
where $`E`$ is the energy and $`p_z`$ is the longitudinal component of the momentum vector. As expected, the $`𝒪(\alpha )`$ QED corrections to $`A_{\mathrm{FB}}`$ are large in the region below the $`Z`$ peak. Since events from the $`Z`$ peak, where $`A_{\mathrm{FB}}`$ is positive and small, are shifted towards smaller values of $`m(\mu ^+\mu ^{})`$ by photon radiation, the forward-backward asymmetry is significantly reduced in magnitude by radiative corrections for $`55\mathrm{GeV}<m(\mu ^+\mu ^{})<90`$ GeV. It should be noted that the forward-backward asymmetry is rather sensitive to the rapidity cuts imposed on the leptons. More details on $`A_{\mathrm{FB}}`$ and the measurement of the effective weak mixing angle can be found in Section 3.2.4.
The mass singular terms arising from final-state photon radiation are proportional to $`\alpha \mathrm{log}(\widehat{s}/m_l^2)`$, where $`m_l`$ is the lepton mass. Thus, the corrections to the $`Z`$ line shape and $`A_{\mathrm{FB}}`$ for electrons in the final state are considerably larger than those in the muon case .
To simulate detector acceptances, we have imposed a $`p_T(\mu )>20`$ GeV and a $`|\eta (\mu )|<3.2`$ cut in Figure 8. Except for the threshold region, the effects of the lepton acceptance cuts approximately cancel in the cross section ratio. In a more realistic simulation of how QED corrections affect observables in Drell-Yan production, lepton and photon identification requirements need to be taken into account in addition to the lepton acceptance cuts. Muons are identified in a hadron collider detector by hits in the muon chambers. In addition to a hit in the muon chambers, one requires that the associated track is consistent with a minimum ionising particle. This limits the energy of a photon which traverses the same calorimeter cell as the muon to be smaller than a critical value $`E_c^\gamma `$. For electrons, the finite resolution of the electromagnetic calorimeter makes it difficult to separate electrons and photons for small opening angles between their momentum vectors. Therefore, electron and photon four-momentum vectors are recombined if their separation in the azimuthal angle–pseudorapidity plane is smaller than a critical value, $`R_c`$. This eliminates the mass-singular terms associated with final-state photon radiation (KLN theorem) and thus may reduce significantly the effect QED corrections have on physical observables in $`ppe^+e^{}(\gamma )`$ . Specific results sensitively depend on the value of $`R_c`$, which is detector dependent.
#### 2.3.2 Non-QED corrections and effective Born description
The amplitude for the parton process $`q(p)+\overline{q}(\overline{p})l^+(k_+)+l^{}(k_{})`$ of quark-antiquark annihilation into charged-lepton pairs is in lowest order described by photon and $`Z`$ boson exchange. In the kinematical variables for the parton system
$$\widehat{s}=(k_++k_{})^2,t=(pk_{})^2,u=(pk_+)^2$$
(13)
the differential parton cross section can be written as follows ($`\theta `$ denotes the scattering angle in the parton CMS):
$$64\pi ^2\widehat{s}\frac{\mathrm{d}\widehat{\sigma }}{\mathrm{d}\mathrm{\Omega }}=\mathrm{\hspace{0.17em}2}𝒜_0\frac{u^2+t^2}{\widehat{s}^2}+𝒜_1\frac{u^2t^2}{\widehat{s}^2}=𝒜_0(1+\mathrm{cos}^2\theta )+𝒜_1\mathrm{cos}\theta $$
(14)
with
$`𝒜_0`$ $`=`$ $`Q_q^2Q_l^2e(\widehat{s})^4+\mathrm{\hspace{0.17em}2}v_qv_lQ_qQ_le(\widehat{s})^2\mathrm{Re}\chi (\widehat{s})+(v_q^2+a_q^2)(v_l^2+a_l^2)|\chi (\widehat{s})|^2,`$
$`𝒜_1`$ $`=`$ $`4Q_qQ_la_qa_le(\widehat{s})^2\mathrm{Re}\chi (\widehat{s})+\mathrm{\hspace{0.17em}8}v_qa_qv_la_l|\chi (\widehat{s})|^2.`$ (15)
This expression is an effective Born approximation, which incorporates several entries from higher-order calculations: the effective (running) electromagnetic charge containing the photon vacuum polarisation (real part)
$$e(\widehat{s})^2=\frac{4\pi \alpha }{1\mathrm{\Delta }\alpha (\widehat{s})};$$
(16)
the $`Z`$ propagator, together with the overall normalisation factor of the neutral-current couplings in terms of the Fermi constant $`G_\mu `$,
$$\chi (\widehat{s})=(G_\mu M_Z^2\sqrt{2})^2\frac{\widehat{s}}{\widehat{s}M_Z^2+\mathrm{i}\widehat{s}\mathrm{\Gamma }_Z/M_Z},$$
(17)
containing the $`Z`$ width as measured from the $`Z`$ resonance at LEP; the vector and axial-vector coupling constants for $`f=l,q`$
$$v_f=I_3^f2Q_f\mathrm{sin}^2\theta _{\mathrm{eff}},a_f=I_3^f,$$
(18)
which contain the effective (leptonic) mixing angle at the $`Z`$ peak, which is measured at LEP and SLC. Taking $`\mathrm{\Gamma }_Z`$ and $`\mathrm{sin}^2\theta _{\mathrm{eff}}`$ from higher-order calculations, the formulae above yield a good description in the region around the $`Z`$ resonance.
From the cross section (14) a forward-backward asymmetry for the produced $`l^+l^{}`$ system can be derived, which at the parton level is given by
$$\widehat{A}_{\mathrm{FB}}=\frac{\widehat{\sigma }_\mathrm{F}\widehat{\sigma }_\mathrm{B}}{\widehat{\sigma }_\mathrm{F}+\widehat{\sigma }_\mathrm{B}}=\frac{3}{8}\frac{𝒜_1}{𝒜_0}.$$
(19)
Around the $`Z`$ peak, this quantity depends sensitively on $`\mathrm{sin}^2\theta _{\mathrm{eff}}`$. Using a parametrisation of the Born-like expressions in Equation 2.3.2, a measurement of $`\widehat{A}_{\mathrm{FB}}`$ allows a determination of the mixing angle (see Section 5). Below we give a quantitative evaluation of the higher-order electroweak effects in the integrated cross section and in $`\widehat{A}_{\mathrm{FB}}`$ to demonstrate the quality of the approximation around the $`Z`$ pole and to point out deviations at higher invariant masses of the lepton pairs.
Besides the universal and non-universal QED corrections, the following IR-finite next-order electroweak terms contribute, which are schematically depicted in Figure 9: self-energy contributions to the photon and $`Z`$ propagators, vertex corrections to the $`\gamma /Z`$-$`ll`$ and $`\gamma /Z`$-$`q\overline{q}`$ 3-point couplings, and box diagrams with two massive boson exchanges. Details of the treatment of the resonance region at higher order is equivalent to that in $`e^+e^{}`$ annihilation in fermion pairs and can be found e.g. in . Around the $`Z`$ pole, the box graphs are negligible, but they increase strongly with the energy and hence contribute sizeably at high invariant masses of the lepton pair. A description in terms of an effective-Born cross section far away from the $`Z`$ pole becomes insufficient for two reasons: the effective couplings (based on self-energies and vertex corrections only) are not static but grow as functions of $`\widehat{s}`$, and the presence of the box contributions, which cannot be absorbed in effective vector and axial-vector couplings in a Born-like structure.
In Figures 10 and 11 we compare the integrated cross section $`\widehat{\sigma }`$ and the asymmetry $`A_{\mathrm{FB}}`$ at the parton level in the approximation corresponding to Equations 14 and 2.3.2 with results obtained by a complete one-loop calculation with proper treatment of higher-order terms around the $`Z`$ resonance. For demonstrational purpose, the effect of the box diagrams is displayed separately. As one can see, the region where the effective Born description starts to become unsatisfactory is at rather high values of the parton energy.
In order to give an idea of the effects remaining in the hadronic cross section after convolution with the quark distribution functions, Table 3 contains the relative deviations of the cross section based on the higher-order parton results from those based on the Born approximation Equation 2.3.2. Also listed are the estimated experimental accuracies with which the cross section in the various bins can be measured. The comparison shows that at high invariant masses the radiative corrections remain sizeable and should be taken into account for studies at high $`\widehat{s}`$, for example in the search for new physics effects originating from a heavy extra gauge-boson $`Z^{}`$.
#### 2.3.3 The full electroweak $`𝒪(\alpha )`$ corrections: Monte Carlo simulations with ZGRAD2
The QED corrections described in Section 2.3.1 have been combined with the weak corrections summarised in the previous section in a new Monte Carlo program called ZGRAD2 . In Figure 12a we show the ratio of the full $`𝒪(\alpha ^3)`$ electroweak and the $`𝒪(\alpha ^3)`$ QED differential cross sections for $`pp\mu ^+\mu ^{}(\gamma )`$ obtained with ZGRAD2 as a function of the $`\mu ^+\mu ^{}`$ invariant mass. As in Section 2.3.1, we have imposed a $`p_T(\mu )>20`$ GeV and a $`|\eta (\mu )|<3.2`$ cut, and used the EBA to evaluate the lowest-order contribution to the $`𝒪(\alpha ^3)`$ QED cross section. Thus, the ratio directly displays the effect of the weak box-diagrams and the energy dependence of the weak coupling form factors. While the additional weak contributions only change the differential cross section by 0.6% at most, they do modify the shape of the $`Z`$ resonance curve.
Figure 12b compares the effect of the $`𝒪(\alpha ^3)`$ QED corrections and the full $`𝒪(\alpha ^3)`$ electroweak corrections on the di-muon invariant mass distribution for $`m(\mu ^+\mu ^{})`$ values between 200 GeV and 2 TeV. Due to the presence of logarithms of the form $`\mathrm{log}(\widehat{s}/M_Z^2)`$, the weak corrections become significantly larger than the QED corrections at large values of $`m(\mu ^+\mu ^{})`$, and, eventually, may have to be resummed . For $`m(\mu ^+\mu ^{})=2`$ TeV, the full $`𝒪(\alpha ^3)`$ electroweak corrections are found to reduce the differential cross section by more than 20%.
Finally, in Figure 13 we show how the $`𝒪(\alpha ^3)`$ corrections affect the forward-backward asymmetry (see Equations 9 to 11). Both QED and weak corrections reduce $`A_{\mathrm{FB}}`$, and their size increases with growing di-muon masses. For $`m(\mu ^+\mu ^{})=2`$ TeV, the weak corrections are about twice as large as the QED corrections. Note that the electroweak corrections affect $`A_{\mathrm{FB}}`$ much less than the lepton pair invariant mass distribution. In the $`Z`$ pole region, $`75\mathrm{GeV}<m(\mu ^+\mu ^{})<105`$ GeV, the weak corrections change the forward-backward asymmetry by at most $`5\times 10^4`$. Results qualitatively similar to those shown in Figures 12 and 13 are obtained for $`ppe^+e^{}(\gamma )`$.
ZGRAD2 includes the complete weak one-loop corrections and the full non-universal QED $`𝒪(\alpha )`$ corrections. The collinear singularities associated with initial-state photon radiation are factorised into the parton distribution functions. However, QED corrections to the evolution of the parton distribution functions (see Section 2.1) are not included in ZGRAD2. These corrections should be part of a complete global fit of the pdf’s including all QED effect - this is beyond the scope of the calculation presented here. None of the current fits to the pdf’s include QED corrections.
### 2.4 $`𝒁^{\mathbf{}}`$ indication from new APV data in cesium and searches at LHC
The weak charge $`Q_W`$ for a heavy atom is defined in terms of the number of $`u,d`$ quarks $`N_u=2Z+N`$, $`N_d=2N+Z`$ in the nucleus $`(Z,N)`$ and the coefficients $`C_{1u,d}`$ in the parity-violating part of the electron-quark Hamiltonian,
$$_{PV}=\frac{G_F}{\sqrt{2}}\overline{e}\gamma _\mu \gamma _5e\left(C_{1u}\overline{u}\gamma ^\mu u+C_{1d}\overline{d}\gamma ^\mu d\right),$$
(20)
via the relation
$$Q_W=2(N_uC_{1u}+N_dC_{1d}).$$
(21)
In the SM: $`C_{1q}=I_3^q2Q_q\mathrm{sin}^2\theta _W`$.
In a recent paper a new determination of the weak charge of atomic cesium has been reported. The most precise atomic parity violating (APV) experiment compares the mixing among $`S`$ and $`P`$ states due to neutral weak interactions to an induced Stark mixing . The 1.2% uncertainty on the previous measurement of the weak charge $`Q_W`$ was dominated by the theoretical calculations on the amount of Stark mixing and on the electronic parity violating matrix elements. In the Stark mixing was measured and, incorporating new experimental data, the uncertainty in the electronic parity violating matrix elements was reduced. The new result $`Q_W(_{55}^{133}\mathrm{Cs})=72.06\pm (0.28)_{\mathrm{expt}}\pm (0.34)_{\mathrm{theor}}`$ represents a considerable improvement with respect to the previous determination . The discrepancy between the standard model (SM) and the experimental data is now given by $`Q_W^{\mathrm{expt}}Q_W^{SM}=1.18(1.28)\pm 0.46`$ (for $`m_t=175`$ GeV and $`M_H=100(300)`$ GeV). This corresponds to 2.6(2.8) standard deviations , excluding the SM at 99% CL and, a fortiori, all the models leading to negative additional contributions to $`Q_W`$, as for example models with a sequential $`Z^{}`$ . This deviation could be explained by assuming the existence of an extra $`Z^{}`$ from $`E_6`$ or $`O(10)`$ or from $`Z_{LR}^{}`$ of left-right (LR) models . The high-energy data at the $`Z`$ resonance strongly bound the $`ZZ^{}`$ mixing ; for this reason we will assume zero mixing. In this case, the new physics contribution to $`Q_W`$ is due to the direct exchange of the $`Z^{}`$ and is completely fixed by the $`Z^{}`$ parameters, $`\delta _NQ_W=16a_e^{}[(2Z+N)v_u^{}+(Z+2N)v_d^{}]M_Z^2/M_Z^{}^2`$, where $`a_f^{},v_f^{}`$ are the couplings $`Z^{}`$ to fermions and, for $`{}_{55}{}^{133}\mathrm{Cs}`$, $`Z=55`$ and $`N=78`$. The relevant couplings of the $`Z^{}`$ to the electron and to the up and down quarks are given in the Table 1 of .
In the case of the LR model considered in , the extra contribution to the weak charge is $`\delta _NQ_W=M_Z^2/M_Z^{}^2Q_W^{SM}`$. For this model one has a 95% CL lower bound on $`M_{Z_{LR}^{}}`$ from the Tevatron given by $`M_{Z_{LR}^{}}630`$ GeV. An LR model could then explain the APV data allowing for a mass of the $`Z_{LR}^{}`$ varying between the intersection from the 95% CL bounds $`540M_{Z_{LR}^{}}`$ (GeV) $`1470`$ deriving from $`Q_W`$ and the lower bound of $`630`$ GeV. In the case of the extra-U(1) models, the CDF experimental lower bounds for the masses vary according to the values of the angle $`\theta _6`$ which parameterises different extra-U(1) models, but in general they are about 600 GeV at 95% CL . For the particular models $`\eta `$, $`\psi `$, $`\chi `$, corresponding to $`\theta _6=\mathrm{arctan}(\sqrt{5/3}),\pi /2,0`$, the 95% CL lower bounds are $`M_{Z_\eta ^{}}620`$ GeV, $`M_{Z_\psi ^{}}590`$ GeV, $`M_{Z_\chi ^{}}595`$ GeV. In Figure 14, the 95% CL bounds on $`M_Z^{}`$ from APV are plotted versus $`\theta _6`$ (the direct lower bounds from the Tevatron are about $`600`$ GeV). We see that an extra $`Z^{}`$ can explain the discrepancy with the SM prediction for the $`Q_W`$ for a wide range of $`\theta _6`$ angle. In particular, the models $`\eta `$ and $`\psi `$ are excluded, whereas the $`\chi `$ model is allowed for $`M_{Z_\chi ^{}}`$ less than about 1.2 TeV.
In the near future, the Tevatron upgrade and LHC can confirm or disprove this indication coming from $`Q_W`$. The existing bounds for $`E_6`$ models from direct searches at the Tevatron will be upgraded by the future run with $`\sqrt{s}=2`$ TeV and 1 fb<sup>-1</sup> to $`M_Z^{}800900`$ GeV and pushed to $`1`$ TeV for 10 fb<sup>-1</sup>. The bounds are based on 10 events in the $`e^+e^{}+\mu ^+\mu ^{}`$ channels and decays to SM final-states only are assumed . At the LHC with an integrated luminosity of $`100`$ fb<sup>-1</sup>, one can explore a mass range up to $`44.5`$ TeV depending on the $`\theta _6`$ value. Concerning LR models, the 95% CL lower limits from the Tevatron run with $`\sqrt{s}=2`$ TeV and 1(10) fb<sup>-1</sup> are $`900(1000)`$ GeV and extend to $`4.5`$ TeV at LHC . Ratios of $`Z^{}`$ couplings to fermions can be probed at LHC, by considering the forward-backward asymmetries, ratios of cross sections in different rapidity bins and other observables. For example, for $`M_Z^{}=1`$ TeV, the LHC can determine the magnitude of normalised $`Z^{}`$ quark and lepton couplings to around $`1020\%`$ . Therefore, if the deviation for the weak charge $`Q_W`$ with respect to the SM prediction is not due to a statistical fluctuation, the new physics described by an extra gauge-boson model like $`Z_\chi ^{}`$ can explain the discrepancy and the LHC will be able to verify this possible evidence.
## 3 PRECISION MEASUREMENTS <sup>5</sup><sup>5</sup>5Section coordinator: S. Haywood
### 3.1 Measurement of the $`𝑾`$ mass
At the time of the LHC start-up, the $`W`$ mass will be known with a precision of about 30 MeV from measurements at LEP2 and Tevatron . The motivation to improve on such a precision is discussed briefly below. The $`W`$ mass, which is one of the fundamental parameters of the Standard Model, is related to other parameters of the theory, i.e. the QED fine structure constant $`\alpha `$, the Fermi constant $`G_F`$ and the Weinberg angle $`\mathrm{sin}\theta _W`$, through the relation
$$M_W=\sqrt{\frac{\pi \alpha }{G_F\sqrt{2}}}\frac{1}{\mathrm{sin}\theta _W\sqrt{1\mathrm{\Delta }r}}$$
(22)
where $`\mathrm{\Delta }r`$ accounts for the radiative corrections which amount to about 4%. The radiative corrections depend on the top mass as $`m_t^2`$ and on the Higgs mass as $`\mathrm{log}M_H`$. Therefore, precise measurements of both the $`W`$ mass and the top mass constrain the mass of the Standard Model Higgs boson or of the $`h`$ boson of the MSSM. This constraint is relatively weak because of the logarithmic dependence of the radiative corrections on the Higgs mass. When it comes to making a comparison of the measurements of $`(M_W,m_t)`$ with the SM predictions, it is not very useful if one measurement is much more restrictive than the other. To ensure that the two mass determinations have equal weight in a $`\chi ^2`$ test, the precision on the top mass and on the $`W`$ mass should be related by the expression
$$\mathrm{\Delta }M_W0.7\times 10^2\mathrm{\Delta }m_t$$
(23)
Since the top mass will be measured with an accuracy of about 2 GeV at the LHC , the $`W`$ mass should be known with a precision of about 15 MeV, so that it does not become the dominant error in the test of the radiative corrections and in the estimation of the Higgs mass. Such a precision is beyond the sensitivity of Tevatron and LEP2.
A study was performed to assess whether the LHC will be able to measure the $`W`$ mass to about 15 MeV . The ATLAS experiment was taken as an example, but similar conclusions hold also for CMS. Such a precise measurement, which will be performed already in the initial phase at low luminosity as will the top mass measurement, would constrain the mass of the Higgs boson to better than 30%. When and if the Higgs boson will be found, such constraints would provide an important consistency check of the theory, and in particular of its scalar sector. Distinguishing between the Standard Model and the MSSM might be possible, since the radiative corrections to the $`W`$ mass are expected to be a few percent larger in the latter case.
The measurement of the $`W`$ mass at hadron colliders is sensitive to many subtle effects which are difficult to predict before the experiments start. However, based on the present knowledge of the LHC detector performance and on the experience from the Tevatron, it is possible to make a reasonable estimate of the total uncertainty and of the main contributions to be expected. In turn, this will lead to requirements for the detector performance and the theoretical inputs which are needed to achieve the desired precision. This is the aim of the study which is described in the next sections.
#### 3.1.1 The method
The measurement of the $`W`$ mass at hadron colliders is performed in the leptonic channels. Since the longitudinal momentum of the neutrino cannot be measured, the transverse mass $`m_T^W`$ is used. This is calculated using the transverse momenta of the neutrino and of the charged lepton, ignoring the longitudinal momenta:
$$m_T^W=\sqrt{2p_T^lp_T^\nu (1\mathrm{cos}\mathrm{\Delta }\varphi )}$$
(24)
where $`l=e,\mu `$. The lepton transverse momentum $`p_T^l`$ is measured, whereas the transverse momentum of the neutrino $`p_T^\nu `$ is obtained from the transverse momentum of the lepton and the momentum $`\stackrel{}{u}`$ of the system recoiling against the $`W`$ in the transverse plane (hereafter called “the recoil”):
$$p_T^\nu =|\stackrel{}{p_T}^l+\stackrel{}{u}|$$
(25)
The angle between the lepton and the neutrino in the transverse plane is denoted by $`\mathrm{\Delta }\varphi `$. The distribution of $`m_T^W`$, and in particular the trailing edge of the spectrum, is sensitive to the $`W`$ mass. Therefore, by fitting the experimental distribution of the transverse mass with Monte Carlo samples generated with different values of $`M_W`$, it is possible to obtain the mass which best fits the data. The trailing edge is smeared by several effects, such as the $`W`$ intrinsic width and the detector resolution. This is illustrated in Figure 15, which shows the distribution of the $`W`$ transverse mass as obtained at particle level (no detector resolution) and by including the energy and momentum resolution as implemented in a fast particle-level simulation and reconstruction of the ATLAS detector (ATLFAST, ). The smearing due to the finite resolution reduces the sharpness of the end-point and therefore the sensitivity to $`M_W`$.
When running at high luminosity, the pile-up will smear significantly the transverse mass distribution, therefore the use of the transverse-mass method will probably be limited to the initial phase at low luminosity. Alternative methods are mentioned in Section 3.1.4.
#### 3.1.2 $`W`$ production and selection
At the LHC, the cross-section for the process $`ppW+X`$ with $`Wl\nu `$ and $`l=e,\mu `$ is 30 nb. Therefore, about 300 million events are expected to be produced in each experiment in one year of operation at low luminosity (integrated luminosity 10 fb<sup>-1</sup>). Such a cross-section is a factor of ten larger than at the Tevatron ($`\sqrt{s}=`$ 1.8 TeV).
To extract a clean $`W`$ signal, one should require:
* An isolated charged lepton ($`e`$ or $`\mu `$) with $`p_T>`$ 25 GeV inside the pseudorapidity region devoted to precision physics $`|\eta |<2.4`$.
* Missing transverse energy $`E_T^{\mathrm{miss}}>`$ 25 GeV.
* No jets in the event with $`p_T>`$ 30 GeV.
* The recoil should satisfy $`|\stackrel{}{u}|<`$ 20 GeV.
The last two cuts are applied to reject $`W`$’s produced with high $`p_T`$, since for large $`p_T^W`$ the transverse mass resolution deteriorates and the QCD background increases. The acceptance of the above cuts is about 25%. By assuming a lepton reconstruction efficiency of 90% and an identification efficiency of 80% , a total selection efficiency of about 20% should be achieved. Therefore, after all cuts about 60 million $`W`$’s are expected in one year of data taking at low luminosity in each experiment, which is a factor of about 50 larger than the statistics expected from the Tevatron Run 2.
#### 3.1.3 Expected uncertainties
Due to the large event sample, the statistical uncertainty on the $`W`$ mass should be smaller than 2 MeV for an integrated luminosity of 10 fb<sup>-1</sup>.
Since the $`W`$ mass is obtained by fitting the experimental distribution of the transverse mass with Monte Carlo samples, the systematic uncertainty will come mainly from the Monte Carlo modelling of the data, i.e. the physics and the detector performance. Uncertainties related to the physics include the knowledge of: the $`W`$ $`p_T`$ spectrum and angular distribution, the parton distribution functions, the $`W`$ width, the radiative decays and the background. Uncertainties related to the detector include the knowledge of: the lepton energy and momentum scale, the energy and momentum resolution, the detector response to the recoil and the effect of the lepton identification cuts. At the LHC, as now at the Tevatron, most of these uncertainties will be constrained in situ by using data samples such as $`Zll`$ decays. The latter will be used to determine the lepton energy scale, to measure the detector resolution, to model the detector response to the $`W`$ recoil and the $`p_T`$ spectrum of the $`W`$, etc..
The advantages of the LHC with respect to the Tevatron experiments are:
* The large number of $`W`$ events mentioned above.
* The large size of the ‘control samples’. About six million $`Zll`$ decays, where $`l=e,\mu `$, are expected in each experiment in one year of data taking at low luminosity after all selection cuts. This is a factor of about 50 larger than the event sample from the Tevatron Run 2.
* ATLAS and CMS are in general more powerful than CDF and D0 are, in terms of energy resolution, particle identification capability, geometrical acceptance and granularity. What may be more important for this measurement is the fact that ATLAS and CMS will benefit from extensive and detailed simulations and test-beam studies of the detector performance, undertaken even before the start of data-taking
Nevertheless, the LHC experiments have complex detectors, which will require a great deal of study before their behaviour is well understood.
To evaluate the systematic uncertainty on the $`W`$ mass to be expected in ATLAS, $`Wl\nu `$ decays were generated with PYTHIA 5.7 and processed with ATLFAST. After applying the selection cuts discussed above, a transverse mass spectrum was produced for a reference mass value (80.300 GeV). All sources of systematic uncertainty affecting the measurement of the $`W`$ mass from CDF Run 1 were then considered as an example<sup>7</sup><sup>7</sup>7Similar results have been obtained by the D0 experiment .. Their magnitude was evaluated in most cases by extrapolating from the Tevatron results, on the basis of the expected ATLAS detector performance. The resulting error on the $`W`$ mass was determined by generating new $`W`$ samples, each one including one source of uncertainty, and by comparing the resulting transverse mass distributions with the one obtained for the reference mass. A Kolmogorov test was used to evaluate the compatibility between distributions.
Since the goal is a total error of $`20`$ MeV per experiment, the individual contributions should be much smaller than 10 MeV. A large number of events was needed to achieve such a sensitivity. With three million events after all cuts, corresponding to twelve million events at the generation level, a sensitivity at the level of 8 MeV was obtained.
The main sources of uncertainty and their impact on the $`W`$ mass measurement are discussed one by one in the remainder of this section. The total error and some concluding remarks are presented in Section 3.1.4.
##### Lepton energy and momentum scale
This is the dominant source of uncertainty on the measurement of the $`W`$ mass from Tevatron Run 1, where the absolute lepton scale is known with a precision of about 0.1% . Most likely, this will be the dominant error also at the LHC. In order to measure the $`W`$ mass with a precision of better than 20 MeV, the lepton scale has to be known to 0.02%. The latter is the most stringent requirement on the energy and momentum scale from LHC physics. It should be noted that a very high precision (0.04%) must be achieved also by the Tevatron experiments in Run 2, in order to measure the $`W`$ mass to 40 MeV . If such a precision will indeed be demonstrated at the Tevatron, it would represent a good benchmark for the LHC experiments.
The lepton energy and momentum scale will be calibrated in situ at the LHC by using physics samples, which will complement the information coming from the hardware calibration, from the magnetic field mapping of solenoids and toroids, and from test-beam measurements. The muon scale will be calibrated by using mainly $`Z\mu \mu `$ events, and the electromagnetic calorimeter scale will be calibrated by using mainly $`Zee`$ events or $`E/p`$ measurements for isolated electrons, where $`E`$ and $`p`$ are the electron energy and momentum as measured in the electromagnetic calorimeter and in the inner detector respectively. Leptonic decays of other resonances ($`\mathrm{{\rm Y}}`$, $`J/\psi `$) should provide additional constraints which minimise the extrapolation error to lower masses than the $`Z`$ boson mass.
Similar methods are used today at the Tevatron, where the uncertainty on the absolute lepton scale is dominated by the statistical error due to the limited $`Z`$ data sample. The main advantage of the LHC compared to the Tevatron is the large sample of $`Zll`$ decays. The $`Z`$ boson is close in mass to the $`W`$ boson, therefore the extrapolation error from the point where the scale is determined to the point where the measurement is performed is small.
A preliminary study of the error on the absolute electron scale to be expected in ATLAS was performed by using a sample of 500000 $`Zee`$ decays processed through a full GEANT-based simulation of the ATLAS detector . Several possible sources of uncertainties were considered: the knowledge of the amount of material in the inner detector, which affects the electromagnetic calorimeter scale because of photon bremsstrahlung; radiative $`Z`$ decays, which distort the reconstructed mass spectrum; the modelling of the underlying event and of the pile-up at low and high luminosity. Table 4 shows that the impact of these uncertainties on the electron scale in the calorimeter can most likely be kept below 0.02%. The most stringent requirement to achieve this goal is the knowledge of the material in the inner detector to 1%, which will require scrutiny during construction plus in situ measurements with photon conversions and $`E/p`$ for isolated electrons. More details can be found in .
Several experimental constraints will be needed to achieve a 0.02% uncertainty on the inner detector muon scale: the solenoidal magnetic field in the inner cavity must be known locally to better than 0.1%, the alignment must be understood locally to $``$ 1 $`\mu `$m in the bending plane, etc.. A detailed discussion on how to meet these goals can be found in .
The scale calibration of the external muon spectrometer depends on the knowledge of the magnetic field, on the chamber alignment and on the knowledge of the muon energy losses in the calorimeters. The latter must be understood to a precision of 0.25% in order to achieve the goal uncertainty of 0.02% on the absolute scale. A preliminary study based on a full GEANT simulation of the ATLAS detector demonstrated that with a sample of only 10000 $`Z\mu \mu `$ decays a scale uncertainty of 0.1% should be attained in the muon spectrometer. More details can be found in .
In conclusion, to achieve the needed precision on the lepton scale, several experimental constraints will have to be satisfied. In addition, cross-checks and combined fits between different sub-detectors (inner detector and electromagnetic calorimeter for the electron scale, inner detector and muon system for the muon scale) will be needed. Indeed, only in an over-constrained situation will it be possible to disentangle the various contributions to the detector response, and therefore to derive a reliable systematic error.
##### Lepton energy and momentum resolution
To keep the uncertainty on the $`W`$ mass from the lepton resolution to less than 10 MeV, the energy resolution of the electromagnetic calorimeter and the momentum resolution of the inner detector and muon system have to be known with a precision of better than 1.5%.
The lepton energy and momentum resolutions will be determined at the LHC by using information from test-beam data and from Monte Carlo simulations of the detector, as well as in situ measurements of the $`Z`$ width in $`Zll`$ final states. The $`E/p`$ distribution for electrons from $`W`$ decays provides an additional tool. These methods are used presently at the Tevatron. As an example, the statistical error on the momentum resolution obtained by CDF in Run 1A is 10%, whereas the systematic error is only 1% and is dominated by the uncertainty on the radiative decays of the $`Z`$ . Since the ATLAS performance in terms of momentum resolution is expected to be similar to that of CDF in the momentum range relevant to $`W`$ production and decays, and since the statistical error at the LHC will be negligible, a total error of much less than 1.5% should be achieved. This uncertainty might further be decreased if improved theoretical calculations of radiative $`Z`$ decays will become available.
##### Recoil modelling
The transverse momentum of the system recoiling against the $`W`$, together with the lepton transverse momentum, is used to determine the $`p_T`$ of the neutrino (see Equation 25). The recoil is mainly composed of soft hadrons from the underlying event, for which neither the physics nor the detector response are known with enough accuracy. Therefore, in order to get a reliable recoil distribution in the Monte Carlo, information from data is used at the Tevatron. By exploiting the similar production mechanisms of $`W`$ and $`Z`$ bosons, in each Monte Carlo event with a given $`p_T^W`$ (determined from the truth information) the recoil is replaced by the recoil measured in the data for $`Z`$ events characterised by a $`p_T^Z`$ (measured by the leptons) similar to $`p_T^W`$. The resulting error on the $`W`$ mass from CDF Run 1A is 60 MeV per channel, and is dominated by the limited statistics of $`Z`$ data. The result obtained from Run 1B (about 30 MeV) shows that this uncertainty scales with $`\sqrt{N}`$, where $`N`$ is the number of events. Extrapolating to the LHC data sample, an error of smaller than 10 MeV per channel should be achieved. It should be noted that the recoil includes the contribution of the pile-up expected at low luminosity (two minimum-bias events per bunch crossing on average).
##### $`𝑾`$ $`𝒑_𝑻`$ spectrum
The modelling of $`p_T^W`$ in the Monte Carlo is affected by both theoretical and experimental uncertainties. Theoretical uncertainties arise from the difficulty in predicting the non-perturbative regime of soft-gluon emission, as well as from missing higher-order QCD corrections. Experimental uncertainties are mainly related to the difficulty of simulating the detector response to low-energy particles.
Therefore, the method used at the Tevatron to obtain a reliable estimate of $`p_T^W`$ consists of measuring the $`p_T`$ distribution of the $`Z`$ boson from $`Zll`$ events in the data, exploiting the fact that both gauge-bosons have similar $`p_T`$ distributions, and using the theoretical prediction for the ratio $`p_T^W/p_T^Z`$ (in this ratio several uncertainties cancel) to convert the measured $`p_T^Z`$ into $`p_T^W`$. The resulting error on the $`W`$ mass obtained by CDF is 20 MeV, dominated by the limited $`Z`$ statistics.
At the LHC, the average transverse momentum of the $`W`$ ($`Z`$) is 12 GeV (14 GeV), as given by PYTHIA 5.7. Over the range $`p_T`$ ($`W`$,$`Z`$) $`<`$ 20 GeV, both gauge-bosons have $`p_T`$ spectra which agree to within $`\pm `$10%. By assuming a negligible statistical error on the knowledge of $`p_T^Z`$, which will be measured with high-statistics data samples, and by using the $`p_T^Z`$ spectrum instead of the $`p_T^W`$ distribution, an error on the $`W`$ mass of about 10 MeV per channel was obtained without any further tuning. Although the leading-order parton shower approach of PYTHIA is only an approximation to reality, this result is encouraging. Furthermore, improved theoretical calculations for the ratio of the $`W`$ and $`Z`$ $`p_T`$ distributions should become available at the time of the LHC, so that the final uncertainty will most likely be smaller than 10 MeV.
##### Parton distribution functions
Parton momentum distributions inside the protons determine the $`W`$ longitudinal momentum, and therefore affect the transverse mass distribution through lepton acceptance effects. At the Tevatron, parton distribution functions (pdf), in particular the $`u/d`$ ratio, are constrained by measuring the forward-backward charge asymmetry of the $`W`$ rapidity distribution. Such an asymmetry, which is typical of $`p\overline{p}`$ collisions, is not present in $`pp`$ collisions and therefore cannot be used at the LHC. However, it has been shown that pdf can be constrained to a few percent at the LHC by using mainly the pseudorapidity distributions of leptons produced in $`W`$ and $`Z`$ decays. The resulting uncertainty on the $`W`$ mass should be smaller than 10 MeV.
##### $`𝑾`$ width
At hadron colliders, the $`W`$ width can be obtained from the measurement of $`R`$, the ratio between the rate of leptonically decaying $`W`$’s and leptonically decaying $`Z`$’s:
$$R=\frac{\sigma _W}{\sigma _Z}\times \frac{BR(Wl\nu )}{BR(Zll)}$$
(26)
where the $`Z`$ branching ratio ($`BR`$) is obtained from LEP measurements, and the ratio between the $`W`$ and the $`Z`$ cross-sections is obtained from theory. By measuring $`R`$, the leptonic branching ratio of the $`W`$ can be extracted from the above formula, and therefore $`\mathrm{\Gamma }_W`$ can be deduced assuming Standard Model couplings for $`Wl\nu `$. The precision achievable with this method is limited by the theoretical knowledge of the ratio of the $`W`$ to the $`Z`$ cross-sections. Another method consists of fitting the high-mass tails of the transverse mass distribution, which are sensitive to the $`W`$ width.
By using these methods, the $`W`$ width was measured with a precision of about 60 MeV by CDF in Run 1, which translates into an error of 10 MeV per channel on the $`W`$ mass measurement.
In Run 2, the $`W`$ width should be measured with a precision of 30 MeV , which contributes an error of 7 MeV per channel on the $`W`$ mass. This is however a conservative estimate for the LHC, where the $`W`$ width should be measured with higher precision than at Tevatron by using the high-mass tails of the transverse mass distribution. The measurement of $`R`$, on the other hand, in addition to being model-dependent would require very precise theoretical inputs. It should be noted that one could also use the value of the $`W`$ width predicted by the Standard Model.
##### Radiative decays
Radiative $`Wl\nu \gamma `$ decays produce a shift in the reconstructed transverse mass, which must be precisely modelled in the Monte Carlo. Uncertainties arise from missing higher-order corrections, which translate into an error of 20 MeV on the $`W`$ mass as measured by CDF in Run 1. Improved theoretical calculations have become recently available . Furthermore, the excellent granularity of the ATLAS electromagnetic calorimeter, and the large statistics of radiative $`Z`$ decays, should provide useful additional information. Therefore, a $`W`$ mass error of 10 MeV per channel was assumed in this study. This is a conservative estimate, since the D0 error from Run 1 is smaller than 10 MeV .
##### Background
Backgrounds distort the $`W`$ transverse mass distribution, contributing mainly to the low-mass region. Therefore, uncertainties on the background normalisation and shape translate into an error on the $`W`$ mass. This error is at the level of 5 MeV (25 MeV) in the electron (muon) channel for the measurement performed by CDF in Run 1, where the background is known with a precision of about 10%.
A study was made of the main backgrounds to $`Wl\nu `$ final states to be expected in ATLAS. The contribution from $`W\tau \nu `$ decays should be of order 1.3% in both the electron and the muon channel. The background from $`Zee`$ decays to the $`We\nu `$ channel is expected to be negligible, whereas the contribution of $`Z\mu \mu `$ decays to the $`W\mu \nu `$ channel should amount to 4%. The difference between these two channels is due to the fact that the calorimetry coverage extends up to $`|\eta |5`$, whereas the coverage of the muon spectrometer is limited to $`|\eta |<2.7`$. Therefore, muons from $`Z`$ decays which are produced with $`|\eta |>2.7`$ escape detection and thus give rise to a relatively large missing transverse momentum. On the other hand, electrons from $`Z`$ decays produced with $`|\eta |>2.4`$ are not efficiently identified, because of the absence of tracking devices and of fine-grained calorimetry, however their energy can be measured up to $`|\eta |5`$. Therefore these events do not pass the $`E_T^{miss}`$ cut described in Section 3.1.2. Finally, $`t\overline{t}`$ production and QCD processes are expected to give negligible contributions.
In order to limit the error on the $`W`$ mass to less than 10 MeV, the background to the electron channel should be known with a precision of 30%, which is easily achievable, and the background to the muon channel should be known with a precision of 7%. The latter could be monitored by using $`Zee`$ decays.
#### 3.1.4 Results
The expected contributions to the uncertainty on the $`W`$ mass measurement, of which some are discussed in the previous sections, are summarised in Table 5. For comparison, the errors obtained by CDF in Run 1A (integrated luminosity $``$ 20 pb<sup>-1</sup>) and Run 1B (integrated luminosity $``$ 90 pb<sup>-1</sup>) are also shown separately. The evolution of the uncertainty between Run 1A and Run 1B shows the effect of the increased statistics and of the improved knowledge of the detector performance and of the physics, and provides a solid basis for the LHC results presented here.
With an integrated luminosity of 10 fb<sup>-1</sup>, which should be collected in one year of LHC operation, and by considering only one lepton species ($`e`$ or $`\mu `$), a total uncertainty of smaller than 25 MeV should be achieved by each LHC experiment. By combining both lepton channels, which should also provide useful cross-checks since some of the systematic uncertainties are different for the electron and the muon sample, and taking into account common uncertainties, the total error should decrease to less than 20 MeV per experiment. Finally, the total LHC uncertainty could be reduced to about 15 MeV by combining ATLAS and CMS together. Such a precision would allow the LHC to compete with the expected precision at a Next Linear Collider .
The most serious experimental challenge in this measurement is the determination of the lepton absolute energy and momentum scale to 0.02%. All other uncertainties are expected to be of the order of (or smaller than) 10 MeV. However, to achieve such a goal, improved theoretical calculations of radiative decays, of the $`W`$ and $`Z`$ $`p_T`$ spectra, and of higher-order QCD corrections will be needed.
The results presented here have to be considered as preliminary and far from being complete. It may be possible that, by applying stronger selection cuts, for instance on the maximum transverse momentum of the $`W`$, the systematic uncertainties may be reduced further. Moreover, two alternative methods to measure the $`W`$ mass can be envisaged. The first one uses the $`p_T`$ distribution of the charged lepton in the final state. Such a distribution features a Jacobian peak at $`p_T^lM_W/2`$ and has the advantage of being affected very little by the pile-up, therefore it could be used at high luminosity. However, the lepton momentum is very sensitive to the $`p_T`$ of the $`W`$ boson, whereas the transverse mass is not, and hence a very precise theoretical knowledge of the $`W`$ $`p_T`$ spectrum would be needed to use this method. Another possibility is to use the ratio of the transverse masses of the $`W`$ and $`Z`$ bosons . The $`Z`$ transverse mass can be reconstructed by using the $`p_T`$ of one of the charged leptons, while the second lepton is treated like a neutrino whose $`p_T`$ is measured by the first lepton and the recoil. By shifting the $`m_T^Z`$ distribution until it fits the $`m_T^W`$ distribution, it is possible to obtain a scaling factor between the $`W`$ and the $`Z`$ masses and therefore the $`W`$ mass. The advantage of this method is that common systematic uncertainties cancel in the ratio. The main disadvantage is the loss of a factor of ten in statistics, since the $`Zll`$ sample is a factor of ten smaller than the $`Wl\nu `$ sample (and only events near to the Jacobian peak contribute significantly to the mass determination). Furthermore, differences in the production mechanism between the $`W`$ and the $`Z`$ ($`p_T`$, angular distribution, etc.), and possible biases coming from the $`Z`$ selection cuts, will give rise to a non-negligible systematic error.
The final measurement will require using all the methods discussed above, in order to cross-check the systematic uncertainties and to achieve the highest precision.
#### 3.1.5 Conclusions
Preliminary studies indicate that measuring the $`W`$ mass at the LHC with a precision of about 15 MeV should be possible, although very challenging. The biggest single advantage of the LHC is the large statistics, which will result in small statistical errors and good control of the systematics. To achieve such unprecedented precision, improved theoretical calculations in many areas will be needed (e.g. radiative decays, pdf’s, $`p_T^W`$), and many stringent experimental requirements will have to be satisfied.
### 3.2 Drell-Yan production of lepton pairs
#### 3.2.1 Introduction
##### Parton level:
In the Standard Model (SM), the production of lepton pairs in hadron-hadron collisions (the Drell-Yan process) is described by $`s`$-channel exchange of photons or $`Z`$ bosons. The parton cross section in the centre-of-mass system has the form:
$$\frac{d\widehat{\sigma }}{d\mathrm{\Omega }}=\frac{\alpha ^2}{4s}[A_0(1+\mathrm{cos}^2\theta )+A_1\mathrm{cos}\theta ]$$
(27)
where $`\widehat{\sigma }=\frac{4\pi \alpha ^2}{3s}A_0`$ and $`A_{FB}=\frac{3}{8}\frac{A_1}{A_0}`$ give the total cross section and the forward-backward asymmetry, respectively. The terms $`A_0`$ and $`A_1`$ are fully determined by the electroweak couplings of the initial- and final-state fermions. At the $`Z`$ peak, the $`Z`$ exchange dominates while the interference term is vanishing. At higher energies, both photon and $`Z`$ exchange contribute and the large value of the forward-backward asymmetry arises from the interference between the neutral currents.
Fermion-pair production above the $`Z`$ pole is a rich search field for new phenomena at present and future high-energy colliders . The differential cross section is given by
$$\frac{d\widehat{\sigma }}{d\mathrm{\Omega }}|\gamma _s+Z_s+\mathrm{New}\mathrm{Physics}\mathrm{?}!|^2$$
(28)
where many proposed types of new physics can lead to observable effects by adding new amplitudes or through their interference with the neutral currents of the SM.
##### At hadron colliders:
The parton cross sections are folded with the parton distribution functions (pdf’s):
$$\frac{\text{d}^2\sigma }{\text{d}M_{ll}\text{d}y}(ppl_1l_2)\underset{ij}{}(f_{i/p}(x_1)f_{j/p}(x_2)+(ij))\widehat{\sigma }$$
(29)
where $`\widehat{\sigma }`$ is the cross section for the partonic subprocess $`ijl_1l_2`$, $`M_{ll}=\sqrt{\tau s}=\sqrt{\widehat{s}}`$ and $`y`$ are the invariant mass and rapidity of the lepton pair, $`x_1=\sqrt{\tau }e^y`$ and $`x_2=\sqrt{\tau }e^y`$ are the parton momentum fractions, and $`f_{i/p(\overline{p})}(x_i)`$ is the probability to find a parton $`i`$ with momentum fraction $`x_i`$ in the (anti)proton.
$`\sigma _{F\pm B}(y,M)`$ $`=`$ $`[{\displaystyle _0^1}\pm {\displaystyle _1^0}]\sigma _{ll}d(\mathrm{cos}\theta ^{})`$ (30)
$`A_{FB}(y,M)`$ $`=`$ $`{\displaystyle \frac{\sigma _{FB}(y,M)}{\sigma _{F+B}(y,M)}}`$ (31)
The total cross section and the forward-backward asymmetry are functions of observables which are well measured experimentally: the invariant mass and the rapidity of the final state lepton-pair. For a pair of partons ($`x_1x_2`$), there are four combinations of quarks initiating Drell-Yan production: $`u\overline{u},\overline{u}u,d\overline{d},\overline{d}d`$. In $`pp`$ collisions, the antiquarks come always from the sea while the quarks can have valence or sea origin. The $`x`$-range probed depends on the mass and rapidity of the lepton-pair as shown in Table 6. Going to higher rapidities increases the difference between $`x_1`$ and $`x_2`$ and hence the probability that the first quark is a valence one.
#### 3.2.2 Event rates
The expected numbers of events for the Tevatron Run 2 (TEV2) and the LHC are shown in Table 7 and Figure 16. The estimation is based on simulations with PYTHIA 5.7 by applying the following cuts:
1. For LHC: both leptons $`|\eta |<2.5`$; for TEV2: one lepton $`|\eta |<1`$, the other $`|\eta |<2.5`$.
2. For both leptons, $`p_T>20`$ GeV.
The data sample can be divided into three classes:
Events near the $`Z`$ pole:
* There will be a huge sample of $`Z`$ events at the LHC. These will allow study of the interplay between $`\mathrm{sin}^2\theta _{\mathrm{eff}}^{\mathrm{lept}}(M_Z^2)`$ and the pdf’s.
High mass pairs (110-400 GeV):
* LEP2 will study this region up to 200 GeV.
* TEV2 will collect a sizeable sample of events in this region.
* LHC will be able to do precision studies.
Very high mass pairs (400-4000 GeV):
* TEV2 will have a first glance.
* LHC will collect a sizeable sample for tests of the SM at the highest momentum transfers ($`Q^2`$) and for searches of new phenomena at the TeV scale.
#### 3.2.3 Measurements of $`\sigma `$ and $`A_{FB}`$
The experimental signature for Drell-Yan events is distinctive: a pair of well isolated leptons with opposite charge. This should be straight forward for the ATLAS and CMS detectors to identify. The backgrounds are low: $`W^+W^{}`$, $`\tau ^+\tau ^{}`$, $`c\overline{c}`$, $`b\overline{b}`$, $`t\overline{t}`$; fakes, cosmics etc.. If the need arises, they can be further suppressed by acoplanarity and isolation cuts. The selection cuts used in this study have already been described in the section on simulations.
An important ingredient in the cross section measurement is the precise determination of the luminosity. A promising possibility is to go directly to the parton luminosity by using the $`W^\pm `$ ($`Z`$) production of single (pair) leptons:
* Constrain the pdf’s.
* Measure directly the parton-parton luminosity.
In this way, the systematic error on $`\sigma _{\mathrm{DY}}^{\mathrm{high}Q^2}`$ relative to $`\sigma _Z`$ can be reduced to $``$ 1%.
In order to measure the forward-backward asymmetry, it is necessary to tag the directions of the incoming quark and antiquark. At the Tevatron, the $`p\overline{p}`$ collisions provide a natural label for the valence (anti)quark. In contrast at the LHC, the $`pp`$ initial state is symmetric. But in the reaction $`q\overline{q}l^+l^{}`$ only $`q`$ can be a valence quark, carrying on average a higher momentum compared to the sea antiquarks. Therefore at the LHC, $`A_{FB}`$ will be signed according to the sign of the rapidity of the lepton pair $`y(ll)`$. Consequently, $`A_{FB}`$ increases as a function of $`y(ll)`$ (see Figure 18).
A precise measurement of $`\sigma `$ and $`A_{FB}`$ at large $`\widehat{s}`$ requires good knowledge of the different types of electroweak radiative corrections to the DY process: vertex, propagator, EW boxes. A complete one-loop parton cross section calculation has been performed . The size of these corrections after folding with the pdf’s and the expected experimental precision on the cross section measurement are compared in Figure 17. The LHC experiments can probe these corrections up to $``$ 2 TeV.
#### 3.2.4 Determination of $`\mathrm{sin}^2\theta _{\mathrm{eff}}^{\mathrm{lept}}(M_Z^2)`$
A very precise determination of $`\mathrm{sin}^2\theta _{\mathrm{eff}}^{\mathrm{lept}}(M_Z^2)`$ will constrain the Higgs mass or, if the Higgs boson is discovered, will check the consistency of the SM . The latest result of the LEP Electroweak Working Group from the summer of 1999 is:
$$\mathrm{sin}^2\theta _{\mathrm{eff}}^{\mathrm{lept}}(\mathrm{M}_\mathrm{Z}^2)=0.23151\pm 0.00017$$
(32)
##### Event selection
A careful study of the precision which can be obtained from the $`Zee`$ decay by ATLAS and CMS has been made using PYTHIA 5.7 and JETSET 7.2. Background processes from $`pp2`$ jets and $`ppt\overline{t}e^+e^{}`$ have been included. In the regions of precision measurements ($`|\eta |2.5`$), the precision which can be obtained from $`Z\mu \mu `$ decays should be comparable to that from the electron channel. In addition, the detectors have calorimetry extending to $`|\eta |5`$ and hence, if it is possible to tag very forward electrons, albeit with significantly lower quality, it may be possible to improve dramatically the measurement of $`\mathrm{sin}^2\theta _{\mathrm{eff}}^{\mathrm{lept}}(M_Z^2)`$.
The following cuts were made:
1. $`p_T^{\mathrm{electron}}>20`$ GeV/$`c`$
2. $`85.2`$ GeV/$`c^2<M(e^+e^{})<97.2`$ GeV/$`c^2`$
In all cases, one electron was required in the precision calorimetry $`|\eta |2.5`$. Efficiencies after typical electron identification cuts were taken from detailed studies reported in . These are typically around 70%, with corresponding jet rejections of $`>10^4`$ (there was no advantage for this measurement of larger rejection factors). For the second electron, the possibility for it to be identified in the forward calorimetry $`2.5<|\eta |4.9`$ was considered. In this region, there is no magnetic tracking. An electron identification efficiency of 50% was assumed with a corresponding jet rejection of $`\rho `$. Extending the pseudorapidity coverage for the second electron increases the range of lepton pair rapidity from $`|y(e^+e^{})|2`$ to $`|y(e^+e^{})|3`$. Figure 18 shows how the asymmetry varies as a function of $`|y(e^+e^{})|`$.
##### Statistical sensitivity
The sensitivity of $`A_{FB}`$ to $`\mathrm{sin}^2\theta _{\mathrm{eff}}^{\mathrm{lept}}(M_Z^2)`$ can be parametrised as follows:
$`A_{FB}`$ $`=`$ $`b(a\mathrm{sin}^2\theta _{\mathrm{eff}}^{\mathrm{lept}}(M_Z^2))`$ (33)
$`a^{O(\alpha ^3)}`$ $`=`$ $`a^{Born}+\mathrm{\Delta }a^{QED}+\mathrm{\Delta }a^{QCD}`$
$`b^{O(\alpha ^3)}`$ $`=`$ $`b^{Born}+\mathrm{\Delta }b^{QED}+\mathrm{\Delta }b^{QCD}`$
Values of $`a`$ and $`b`$ were calculated in and have been re-evaluated by Baur corresponding to the above cuts - see Table 8.
A summary of the statistical errors which can be obtained with 100 fb<sup>-1</sup> are indicated in Table 9. With the best rejection factors shown in the table, the effect of the background is negligible. If no jet rejection is possible in the forward calorimetry, the statistical precisions which can be obtained on $`\mathrm{sin}^2\theta _{\mathrm{eff}}^{\mathrm{lept}}(M_Z^2)`$ are $`3.4\times 10^4`$ and $`4.1\times 10^4`$ for no $`y`$ cut and $`|y(e^+e^{})|>1.0`$ respectively. While the sensitivity to $`\mathrm{sin}^2\theta _{\mathrm{eff}}^{\mathrm{lept}}(M_Z^2)`$ is increased by cutting on $`|y(e^+e^{})|`$ (see Table 8), the gain is reduced by the loss of acceptance and increased significance of the background when the forward calorimetry is used. It is probable that greater sensitivity could be obtained by fitting $`A_{FB}`$ as a function of $`|y(e^+e^{})|`$.
From Table 9, it can be seen that for a single lepton species from one LHC experiment, using leptons measured in $`|\eta |<2.5`$, a statistical precision of $`4.0\times 10^4`$ on $`\mathrm{sin}^2\theta _{\mathrm{eff}}^{\mathrm{lept}}(M_Z^2)`$ could be obtained. With the combination of electrons and muons in two experiments, $`2.0\times 10^4`$ could be obtained.
The table shows that for moderate jet rejection ($`10^2`$) in the forward calorimetry, a statistical precision of $`1.4\times 10^4`$ could be reached by a single experiment using just the electron channel (cannot include the muons). Even a poor rejection $`10`$, would provide a useful measurement. While no studies with full detector simulation have been done, its seems likely that both the ATLAS and CMS forward calorimetry will be able to provide useful electron identification because of moderate longitudinal and transverse segmentation. Combining both experiments will permit a further $`\sqrt{2}`$ reduction in the statistical uncertainty.
##### Systematic uncertainties
In order to be able to exploit the possibility of measuring $`\mathrm{sin}^2\theta _{\mathrm{eff}}^{\mathrm{lept}}(M_Z^2)`$ with such high precision, the systematic errors have to be comparably small. Quick estimates indicate that the following factors are the most important ones:
1. pdf’s: affect both the lepton acceptance as well as the results of radiative correction calculations.
2. Lepton acceptance and reconstruction efficiency as a function of lepton rapidity: while there is some cancellation in the determination of the asymmetry, the product will need to be known to better than 0.1%. CDF has shown that it is possible to achieve a precision of about 1%, with the largest contribution being due to the uncertainty in the pdf’s.
3. Effects of higher order QCD (and electroweak) corrections: can be estimated by varying the errors on the parameters $`a`$ and $`b`$.
4. Mass Scale: $`A_{FB}`$ varies as a function of the invariant mass of the lepton pair. Since the measured asymmetry corresponds to an integration over the $`Z`$ resonance, it is important to understand the mass scale. It is expected that this will be known to $`0.02`$% (see 3.1.3) by direct comparison of the $`Z`$ peak with the measured LEP parameters.
The most important systematic contribution is that coming from the uncertainties in the pdf’s. A study using several “modern” pdf’s (MRST, CTEQ3 and CTEQ4) gave agreement between the resulting values of $`A_{FB}`$ within the 1% statistical errors of the study ($`5\times 10^5`$ events were generated for each pdf set). This uncertainty must be reduced by a factor of 10 if it is to be smaller than the expected statistical precision on $`A_{FB}`$ shown in Table 9. It remains to be seen whether (a) the differences arising from the various pdf’s will shrink with increased statistical sensitivity of the study and (b) whether the current pdf’s actually describe the measured data sufficiently well (since the pfd’s are fitted to common data, variations are not necessarily indicative of the actual uncertainties). New measurements from the Tevatron (and ultimately the LHC itself) will improve the understanding of the pdf’s, but it is unclear at this stage whether this will be sufficient. It may be possible to fit simultaneously $`\mathrm{sin}^2\theta _{\mathrm{eff}}^{\mathrm{lept}}(M_Z^2)`$ and the pdf’s, or alternatively, it may be necessary to reverse the strategy and use the measurement of $`A_{FB}`$ combined with existing measurements of $`\mathrm{sin}^2\theta _{\mathrm{eff}}^{\mathrm{lept}}(M_Z^2)`$ to constrain the pdf’s.
#### 3.2.5 Search for new phenomena
##### Contact interactions
Contact interactions offer a general framework for a new interaction with coupling $`g`$ and typical energy scale $`\mathrm{\Lambda }\sqrt{\widehat{s}}`$. At LEP2, the current limits for quark-lepton compositeness at 95% CL vary between 3 and 8 TeV, depending on the model. At the LHC scales up to 25-30 TeV are reachable, as illustrated in Figure 19.
##### Search for resonances
The other extreme is the search for resonances like $`Z^{}`$ or $`\stackrel{}{\nu }`$, which produce peaks in the mass distributions. A neutral heavy gauge-boson $`Z^{}`$ is characterised by its mass $`m_Z^{}`$, by its couplings and by its mixing angle $`\theta _M`$ with the standard $`Z`$ boson. If $`\theta _M=0`$ and the $`Z^{}`$ has SM couplings, the current limit is $`m_Z^{}>1050`$ GeV . For other coupling scenarios the lower limits are model dependent and typically of the order of several hundred GeV. Resonances with masses up to $``$ 4-5 TeV can be probed at LHC, as shown in Figure 19.
##### R-parity violation
In SUSY theories with R-parity violation, it is possible to couple sleptons to pairs of SM leptons or quarks through new independent Yukawa couplings (9 $`\lambda `$ couplings for the slepton-lepton sector and 27 $`\lambda ^{^{}}`$ couplings for the slepton-quark sector). This makes the resonance formation of single scalar neutrino $`\stackrel{}{\nu }`$ in $`d\overline{d}`$ scattering possible. It can be observed through the decay of the $`\stackrel{}{\nu }`$ to lepton pairs, if a suitable combination of two couplings (e.g. $`\lambda _{311}^{}\lambda _{131}`$) is present . The K-factor for slepton production is not calculated yet, leading to an uncertainty $``$ 10% in the estimate of the $`\lambda \lambda ^{}`$ sensitivity.
##### Low-scale gravity
An exciting possibility is the search for low-scale gravity effects in theories with extra spatial dimensions, leading to virtual graviton exchange. The best limits at LEP2 come from combined analysis of Bhabha scattering :
$`\mathrm{\Lambda }_T=1.412(1.077)`$ TeV for $`\lambda =+1(1)`$ at 95% CL
In the Drell-Yan process there is an unique contribution from $`s`$-channel graviton exchange , which modifies the form of the differential cross section and gives a distinct signature:
$$ggl^+l^{}$$
(34)
$$\frac{d\sigma }{d\mathrm{cos}\theta }=\frac{\lambda ^2\widehat{s}^3}{64\pi M_s^8}(1\mathrm{cos}^4\theta )$$
(35)
The large parton luminosity for gluons at LHC may also compensate the $`M_s^8`$ suppression. Scales up to $``$ 5 TeV can be probed with luminosity 100 fb<sup>-1</sup>.
#### 3.2.6 Summary
The main results of this study are:
* A competitive measurement of $`\mathrm{sin}^2\theta _{\mathrm{eff}}^{\mathrm{lept}}(M_Z^2)`$ is hard due to the central acceptance of the experiments and the difficulty of controlling the pdf’s (parton distribution functions) with the required precision. However, a detector with extended forward acceptance for one of the leptons offers the possibility to measure $`\mathrm{sin}^2\theta _{\mathrm{eff}}^{\mathrm{lept}}(M_Z^2)`$ with a statistical precision of $`1.4\times 10^4`$.
* The total cross-section can be measured with systematic error $`\frac{\sigma _{\mathrm{DY}}^{\mathrm{high}Q^2}}{\sigma _Z}`$ $`<`$ 1%.
* A non-zero forward-backward asymmetry $`A_{FB}`$ can be measured up to 2 TeV with statistical precision $`>3\sigma `$.
* The Drell-Yan process can probe electroweak radiative corrections up to 1.5 TeV with statistical precision at the $`2\sigma `$ level as a function of $`Q^2`$.
* The high energy and luminosity of LHC offers a rich search field at the TeV scale in the Drell-Yan channel: contact interactions, resonance formation ($`Z^{}`$, scalar neutrinos), low-scale gravity, etc..
Further studies will refine the following points:
* The effect of higher order QED corrections (initial- and final-state radiation and their interference).
* The effect of experimental cuts on the electroweak corrections.
* The careful separation of the $`\sigma _{u\overline{u}}`$ and $`\sigma _{d\overline{d}}`$ contributions.
### 3.3 Tau physics
The $`\tau `$ lepton is a member of the third generation which decays into particles belonging to the first and second ones. Thus, $`\tau `$ physics could provide some clues to the puzzle of the recurring families of leptons and quarks. One naïvely expects the heavier fermions to be more sensitive to whatever dynamics is responsible for the fermion-mass generation. The pure leptonic or semileptonic character of $`\tau `$ decays provides a clean laboratory to test the structure of the weak currents and the universality of their couplings to the gauge-bosons.
The last few years have witnessed a substantial change in our knowledge of the $`\tau `$ properties . The large (and clean) data samples collected by the most recent experiments have improved considerably the statistical accuracy and, moreover, have brought a new level of systematic understanding.
A high-energy hadron collider does not provide a very good environment to perform precision $`\tau `$ physics. Nevertheless, there are a few topics where LHC could contribute in a relevant and unique way. Moreover, since the $`\tau `$ is the heaviest known lepton, it can play a very important role in searches for new particles (for example, as in Section 6.1).
#### 3.3.1 Charged-current universality
Table 10 shows the present experimental tests on the universality of the leptonic charged-current couplings. The leptonic $`\tau `$ branching ratios are already known with a quite impressive precision of $`0.3\%`$; this translates into a test of $`g_\mu /g_e`$ universality at the 0.22% level. However, in order to test the ratios $`g_\tau /g_\mu `$ and $`g_\tau /g_e`$, one needs precise measurements of the $`\tau `$ mass and lifetime, in addition. At present, these quantities are known with a precision of 0.016% ($`m_\tau =1777.05_{0.26}^{+0.29}`$ MeV) and 0.34% ($`\tau _\tau =290.77\pm 0.99`$ fs), respectively , which leads to a sensitivity of 0.23% for the three $`g_l/g_l^{}`$ ratios.
Future high-luminosity $`e^+e^{}`$ colliders running near the $`\tau ^+\tau ^{}`$ production threshold could perform more precise measurements of the leptonic $`\tau `$ branching fractions and the $`\tau `$ mass. However, one needs a high-energy machine to measure the $`\tau `$ lifetime. Clearly, the future tests of lepton universality will be limited by the $`\tau _\tau `$ accuracy. It is not clear whether the $`B`$-factories would be able to improve the present $`\tau _\tau `$ measurement in a significant way. Thus, it is important to know how well $`\tau _\tau `$ can be determined at LHC.
A less precise but more direct test on the universality of the leptonic $`W^\pm `$ couplings is provided by the comparison of the different $`W^+l^+\nu _l`$ branching fractions. LEP2 has already achieved a better sensitivity than the Tevatron collider, and a further improvement is expected when the full LEP2 statistics will be available. It is an open question whether LHC could be competitive at this level ($`1\%`$) of precision.
#### 3.3.2 Tau lifetime
The current world average for the $`\tau `$ lifetime is $`290.8\pm 1.0`$ fs ($`c\tau =87\mu `$m) . Improvements in this measurement would be welcome in order to give better tests of the Standard Model, in particular lepton universality and electroweak calculations. In this section, the results of a preliminary study to examine the LHC potential are given.
In LEP experiments, $`\tau `$ pairs are produced back-to-back with well defined momenta - this will not be the case at the LHC. The first feature allows valuable correlations to be made between the two $`\tau `$ decays, while the second provides the boost required to obtain proper lifetime estimates. At the LHC, $`Z\tau \tau `$ events will be triggered by requiring one $`\tau `$ to decay to an electron or muon, while the lifetime is estimated from the other $`\tau `$ which is required to decay to three charged particles.
##### Tau reconstruction
A study was made using fully simulated events in the ATLAS detector (see for more details). When the $`Z`$ has some transverse momentum, the momenta of the $`\tau `$’s can be deduced by projecting the recoil momentum vector measured by the calorimetry along the lines of flight of the two $`\tau `$’s (determined from the direction of the lepton and the hadronic jet, respectively). Due to resolution effects, this procedure works best when the $`\tau `$’s are not back-to-back. The following cuts were made:
* The lepton should have $`p_T>24`$ GeV, $`|\eta |<2.5`$.
* The identified hadronic jet should contain three charged tracks and satisfy $`E_T>30`$ GeV, $`|\eta |<2.5`$.
* Transverse mass of lepton and missing energy should be $`<50`$ GeV.
* The angle $`\mathrm{\Delta }\varphi `$ between the $`\tau `$’s in the transverse plane should satisfy: $`1.8<\mathrm{\Delta }\varphi <2.7`$ or $`3.6<\mathrm{\Delta }\varphi <4.5`$.
* The invariant mass of the $`\tau `$ pair should satisfy: $`60<m_{\tau \tau }<120`$ GeV.
These cuts result in an efficiency of 1.5%. For these events, the $`\tau `$ momenta could be estimated with a resolution of 15%.
A vertex was formed from the charged tracks in the hadronic jet. It was required that the vertex should be within 2 cm of the interaction point and the invariant mass of the particles should be between 0.4 and 1.78 GeV. The efficiency for this was 70% and resulted in a resolution on the vertex position in the transverse plane of $`490\mu `$m, corresponding to a resolution on the proper decay length of $`17\mu `$m.
##### Lifetime estimate
The statistical resolution on the proper decay length from the combination of the vertexing and the estimate of the tau momentum is of the order of $`21\mu `$m (corresponding to 55 fs). A simple Monte Carlo study was made to estimate the statistical uncertainty on the $`\tau `$ lifetime ($`\tau _\tau `$) which could be achieved with $`N`$ hadronic $`\tau `$ decays. Since the resolution of the lifetime for a single event (55 fs) is a fair bit smaller than the $`\tau `$ lifetime (291 fs), the statistical error which can be obtained is dominated by the number of events: $`\sigma (\tau _\tau )\tau _\tau /\sqrt{N}`$.
At the LHC, the cross section for $`Z\tau \tau `$ will be 1.5 nb, with a branching ratio of 11% for a lepton and a three-prong hadronic decay. The reconstruction and selection described above results in an efficiency of 0.54%. If 30 fb<sup>-1</sup> were collected in a low luminosity run, then 26,000 reconstructed $`\tau `$’s could be used, leading to a statistical error on the lifetime of 1.8 fs. To make this competitive would require increased efficiency for selecting the $`\tau `$ decays - this is probably a low luminosity measurement and so cannot benefit from the statistics of a high luminosity run.
Increasing the efficiency may not be simple, since the cuts were designed to control the background. $`W+`$jet events will be removed by the mass cuts, and apart from a small amount of gluon splitting to heavy flavour, the jets should not contain significant lifetime information, hence this background should not be a problem. The $`B`$ lifetime is a factor of five larger than that of the $`\tau `$, and hence more care will be required with $`b\overline{b}`$ events.
Concerning systematic errors coming from the determination of the decay length in the silicon tracking, the average radial position of the detectors in the vertexing layer will need to be understood to better than $`10\mu `$m. This will be challenging but studies suggest this may be feasible . It should be possible to control the systematics on the measurement of recoil momentum of the $`Z`$ by comparison with $`Zee`$ or $`Z\mu \mu `$ events, where the recoil can be measured accurately by the leptons.
###### The use of $`𝑾\mathbf{}𝝉𝝂`$
It may be possible to use the decays $`W\tau \nu `$ which have a higher cross section than $`Z\tau \tau `$. In ATLAS, such events could be triggered by a special $`\tau `$-jet and missing $`E_T`$ trigger . Information about the $`\tau `$ momentum can be deduced by comparing the energy and direction of the hadronic jet with the direction of the $`\tau `$ and using the $`\tau `$ mass constraint, where the $`\tau `$ direction can be determined from the reconstructed decay vertex. In principle, it is possible to solve completely for the $`\tau `$ momentum, although resolution effects on the vertex position and complications arising from $`\pi ^0`$’s in the hadronic jet mean that sometimes solutions are not physical. Alternatively, an approximate estimator can be formed which does not employ the mass constraint . This uses the $`\tau `$-jet energy, mass and $`p_T`$ relative to the $`\tau `$ direction - all three quantities being determined from the charged tracks alone. This is more robust but its behaviour is sensitive to the selection cuts. It is yet to be proved that a $`W\tau \nu `$ signal can be identified with sufficient efficiency above the huge QCD (and in particular, $`b\overline{b}`$) background.
#### 3.3.3 Rare decays
Owing to the huge backgrounds, it will not be possible to make a general search for rare decay modes of the $`\tau `$. However, the lepton-number violating decay $`\tau ^{}\mu ^+\mu ^{}\mu ^{}`$ has a clean signature, which is well suited for the LHC detectors. The present experimental bound is
$$BR(\tau ^{}\mu ^+\mu ^{}\mu ^{})<1.9\times 10^6(90\%\text{CL})$$
This limit reflects the size of the existing $`\tau `$ data samples. LHC will produce a huge statistics, several orders of magnitude larger than the present one. The achievable limit will then be set by systematics and backgrounds, which need to be properly estimated. A sensitivity at the level of $`10^8`$ does not seem out of reach. This could open a very interesting window into new physics phenomena, since many extensions of the Standard Model framework can lead to signals in the $`10^6`$ to $`10^8`$ range.
Although more difficult to detect, other lepton-number violating decays such as $`\tau \mu \mu e`$,$`\mu ee`$,$`eee`$, $`\mu \gamma `$ are worth studying.
## 4 VECTOR-BOSON PAIR PRODUCTION <sup>8</sup><sup>8</sup>8Section coordinator: Z. Kunszt
### 4.1 $`𝑾^\mathbf{+}𝑾^{\mathbf{}}\mathbf{,}𝑾^\mathbf{\pm }𝒁\mathbf{,}𝒁𝒁`$ production
#### 4.1.1 Recent numerical implementations
As already is noted in the introduction, for the description of $`W^+W^{},W^\pm Z,ZZ`$ production with their subsequent decays into lepton pairs two new numerical parton-level Monte Carlo programs have recently become available (MCFM), (DKS). These packages consider the production of four leptons in the double resonance approximation with complete $`𝒪(\alpha _s)`$ corrections. They can be used to compute any infra-red safe quantity with arbitrary experimental cuts on the leptonic decay products. These packages have already been used for updating and cross-checking previous results. The DKS program is available in fortran90 and fortran77 versions and includes anomalous triple gauge-boson couplings. The MCFM program is more complete in the sense that single resonance background diagrams are also added and finite width effects are included in some approximation which respects gauge-invariance. However, it does not include anomalous triple gauge-boson couplings. The results of the MCFM and DKS programs agree with each other within the integration error of $`0.5\%`$. Similar agreement is found with the spin averaged cross section indicated in . In the past years the majority of the experimental studies used the programs described in (BHO). A recent comparison between the DKS and BHO programs finds agreement at the level of 1% for $`WZ`$ production and 3-4% for $`WW`$ production (further details see Section 5.5). This confirms the assumptions of that the spin correlations effects coming from virtual corrections are small. Note that recently a new $`𝒪(\alpha _s`$) package has been written also for $`W\gamma `$ and $`Z\gamma `$ production with anomalous couplings and for the first time the complete one loop QCD corrections are available also for these processes (see Section 4.2).
#### 4.1.2 Input parameters and bench mark cross sections
In using these packages, one should be careful with input parameters. The QCD input is standard: the latest next-to-leading order parton number densities have to be used with the corresponding running coupling constant at some physical scale defined in terms of the kinematics of the outgoing particles.
The helicity amplitudes coded into these programs are calculated in $`𝒪(\alpha _s)`$ but they are leading order in the electroweak theory. However, the one loop electroweak radiative corrections are not completely negligible. The dominant corrections are given by light fermion loops and large custodial symmetry violating contributions of the top quark. Fortunately, they are universal and can be taken into account in the spirit of the “improved Born approximation” . Universality means that their contributions can modify only the leading order relation between $`M_Z`$, $`M_W`$ and $`\mathrm{sin}^2\theta _W`$ which can be taken into account with the use of the effective coupling
$$\mathrm{sin}^2\theta _W\frac{\pi \alpha (M_Z)}{\sqrt{2}G_FM_W^2},$$
(36)
where $`G_F=1.16639\times 10^5`$ GeV<sup>-2</sup> is the Fermi constant and $`\alpha (\mu )`$ is the running QED coupling. With the values of the gauge-bosons masses of $`M_Z=91.187`$ GeV and $`M_W=80.41`$ GeV, one obtains $`\alpha =\alpha (M_Z)=1/128`$ and $`\mathrm{sin}^2\theta _W=0.230`$. Ignoring this correlation may lead to about 5-6% discrepancy in the cross section values. The remaining electroweak corrections are estimated to be less than 2% as long the parton sub-energy is below $`0.51\mathrm{TeV}`$. However, above the $`1\mathrm{TeV}`$ scale double logarithmic corrections of $`𝒪(\alpha _W\mathrm{log}^2\widehat{s}/M_W^2)`$ become non-negligible. The origin of these large contributions is the incomplete cancellation of the soft singularities of massless gauge-boson emission (the Bloch-Nordsieck theorem is not valid for non-Abelian theories ). Since the physical cross section decreases strongly with the increase of the invariant mass of the gauge-boson pairs, these corrections are not important at the LHC. The validity of the improved Born approximation and the presence of the double logarithmic corrections has been tested for $`W`$ pair production at LEP2 where the full next-to-leading order corrections are available .
Additional electroweak input parameters are the matrix elements of the CKM mixing matrix. In the light quark sector, one should use the best experimental values . In the case of the heavy quark contributions, the calculation is approximate since the $`𝒪(\alpha _s)`$ helicity amplitudes have been calculated assuming massless quarks . This assumption is clearly not valid for the top contributions. $`WW`$ pair production receives contributions from diagrams with the $`t`$-channel exchange of the top quark (with $`|V_{td}|=|V_{ts}|=0`$ and $`|V_{tb}|=1`$). However, it is suppressed due to the large top mass and small $`b`$-quark parton densities; therefore, it is reasonable to use $`|V_{tb}|=0`$. The contribution of the subprocess $`b\overline{b}W^+W^{}`$ (treating the top as massless) is of the order of 2% for the LHC giving an upper limit on the theoretical ambiguity coming from this source. In the case of $`W^\pm Z`$ production, one can neglect the subprocess $`bgW^{}Zt`$. It is present at next-to-leading order but again it is strongly suppressed by the large top quark mass, as well as the small $`b`$-quark distribution function. For the numerical results presented here, values $`|V_{ud}|=|V_{cs}|=0.975;|V_{us}|=|V_{cd}|=0.222`$ and $`|V_{ub}|=|V_{cb}|=|V_{td}|=|V_{ts}|=|V_{tb}|=0`$ are used. We present cross-section values without including the branching ratios. To get event signals, they have to be multiplied with the leptonic branching ratios of the vector-bosons. We use
$$BR(Ze^+e^{}\mathrm{or}\mu ^+\mu ^{})=3.37\%BR(Z\underset{i=e,\mu ,\tau }{}\nu _i\overline{\nu _i})=20.1\%$$
$$BR(W^+e^+\nu _e\mathrm{or}\mu ^+\nu _\mu ^{})=10.8\%$$
These ratios implicitly incorporate QCD corrections to the hadronic decay widths of the $`W`$ and $`Z`$.
Most of the results are obtained with some “standard cuts” defined as follows: a transverse momentum cut of $`p_T>20`$ GeV and pseudorapidity cut of $`|\eta |2.5`$ is applied for all charged leptons and $`p_T^{\mathrm{miss}}20\mathrm{GeV}`$ is required for $`WZ`$ production while $`p_T^{\mathrm{miss}}25\mathrm{GeV}`$ for $`W`$ pair production. We use two different parton distributions, MRST with $`M_W=80.41\mathrm{GeV}`$ and CTEQ(4M) with $`M_W=80.33\mathrm{GeV}`$ which we refer to simply as MRST and CTEQ. $`\alpha _s(M_Z)=0.1175`$ is used for MRST and $`\alpha _s(M_Z)=0.116`$ is used for CTEQ. In all computations, we set the renormalisation and factorisation scales equal to each other.
In Table 11, we present the total cross section values for the various processes at the LHC, for the MRST and CTEQ parton distributions. We tabulated the results for $`\sigma ^{\mathrm{tot}}`$ (the cross sections without any cuts applied) as well as $`\sigma ^{\mathrm{cut}}`$ (the cross sections with the standard cuts defined above). The cross section values are given for the scale
$$\mu =(M_{V_1}+M_{V_2})/2,$$
(37)
where $`M_{V_i}`$ are the masses of the two produced vector bosons.
In previous publications a number of phenomenologically interesting questions have been considered. Here we restrict ourselves to recall two interesting and typical features: the scale dependence of the radiative corrections for $`WW`$ production and radiation zeros for $`WZ`$ production.
#### 4.1.3 Scale dependence
The one-loop corrections to the total cross sections are of the order 50% of the leading order term and they can be much larger for the kinematical range of larger transverse momenta or invariant mass of the vector-boson pair. For differential distributions where $`p_T`$ is not integrated out completely, the scale choice
$$\mu ^2=\mu _{\mathrm{st}}^2\frac{1}{2}(p_T^2(V_1)+p_T^2(V_2)+M_{V_1}^2+M_{V_2}^2)$$
(38)
appears to be appropriate. For the total cross section, the difference between the two scale choices expressed in Equations 37 and 38 is very small since it is dominated by low-$`p_T`$ vector-bosons. However, for more exclusive quantities, the differences can be substantial. At the LHC, the huge one-loop corrections in the tails of the distributions are dominated by the bremsstrahlung contributions; therefore it is natural to consider the cross sections with and without the jet veto (that is, with or without the cut $`E_T^{\mathrm{jet}}<40`$ GeV).
In Figure 20, the scale dependence of $`\sigma ^{\mathrm{cut}}`$ is shown for standard cuts, with a jet veto and with stronger cuts on the transverse momenta of the charged leptons. We can see that the corrections are large and increase with the additional cuts applied. The scale dependence at LO is reduced at NLO and it is reduced further when a jet veto is applied. In particular, the size of the correction is strongly reduced when applying the jet veto - an important feature for background studies.
#### 4.1.4 Approximate radiation zeros in $`WZ`$ production
In leading order, the angular distribution of $`WZ`$ production exhibits an approximate radiation zero for $`\mathrm{cos}\theta =(g_1+g_2)/(g_1g_2)`$ where $`g_1,g_2`$ denote the $`Z`$ boson couplings to the left handed up and down quarks, respectively. Since the precise flight direction of the $`W`$ boson is not known (due to the uncertainty in the longitudinal momentum carried by the neutrino) it is convenient to plot a distribution in the (true) rapidity difference between the $`Z`$ boson and the charged lepton coming from the decay of the $`W`$: $`\mathrm{\Delta }y_{Zl}y_Zy_l`$. This quantity is similar to the rapidity difference $`\mathrm{\Delta }y_{WZ}|y_Wy_Z|`$ studied in , but uses only the observable charged-lepton variables. It is the direct analogue of the variable $`y_\gamma y_{l^+}`$ considered in for the case of $`W\gamma `$ production. It is possible to determine $`\mathrm{cos}\theta `$ in the $`W\gamma `$ or $`WZ`$ rest frame, by solving for the neutrino longitudinal momentum using the $`W`$ mass as a constraint, up to a two-fold discrete ambiguity for each event . However, it has been found that the ambiguity degrades the radiation zero - at least if each solution is given a weight of 50% - so that the rapidity difference $`y_\gamma y_{l^+}`$ is more discriminating than $`\mathrm{cos}\theta `$. As one can see from Figure 21, there is a residual dip in the $`\mathrm{\Delta }y_{Zl}`$ distribution, even at order $`\alpha _s`$. This dip can be enhanced easily by requiring a minimal energy for the decay lepton from the $`W`$ and by cutting on the rapidity of the $`Z`$ boson. In Figure 21, we have chosen $`E(l)>100`$ GeV with and without $`y_Z<0`$. Note that the latter two curves are scaled up by a factor of 5. At the LHC, for the first time, we shall have enough statistics to test experimentally for the presence of approximate radiation zeros.
New physics contributions can modify the self-interactions of vector-bosons, in particular the triple gauge-boson vertices. If new physics occurs at an energy scale well above that being probed experimentally, it can be integrated out, and the result expressed as a set of anomalous (non-Standard Model) interaction vertices. (The physics of anomalous coupling will be considered in detail in Section 10 and our standard notation for the anomalous triple gauge-boson couplings is given there.) It is interesting to know what is the effect of the anomalous $`W^+W^{}Z`$ couplings on the approximate radiation zero of $`WZ`$ production . In Figure 22, the $`\mathrm{\Delta }y_{Zl}`$ distribution is plotted for two different sets of anomalous couplings at vanishing $`q^2`$ $`(\mathrm{\Delta }g_1=0.013,\lambda ^Z=0.02,\mathrm{\Delta }\kappa ^Z=0.028)`$ and$`(\mathrm{\Delta }g_1=0.065,\lambda ^Z=0.04,\mathrm{\Delta }\kappa ^Z=0.071)`$. For the $`q^2`$ dependence we assumed dipole form factors of the generic form
$$\widehat{a}(q^2)=\frac{a}{\left(1+\frac{q^2}{\mathrm{\Lambda }^2}\right)^2}$$
(39)
with $`\mathrm{\Lambda }=2\mathrm{TeV}`$. As one can see in Figure 22, the contributions of anomalous couplings have the tendency to make the dip less pronounced.
#### 4.1.5 Future improvements
The present state of art of the description of gauge-boson pair production is not completely satisfactory yet. Of the various issues, there are three which require further theoretical studies. First, the double resonant approximation is expected to be correct only up to a few percent accuracy - it is important to go beyond this approximation. A first attempt has been made by Campbell and Ellis where, as already mentioned above, the singly-resonant diagrams have also been included. These additions are obviously relevant in the off-resonant regions. The inclusion of finite width effect is not completely straightforward because of possible conflict with gauge-invariance. This issue requires further theoretical study. Secondly, we need NLO results also for the semi-leptonic channels when one of the gauge-bosons decays hadronically. This requires the inclusion of the contributions of diagrams describing the gluonic corrections to the final-state quarks. Thirdly, fixed order perturbative QCD description is not applicable for the description of the low-$`p_T`$ behaviour of the gauge-boson pair. The technique for the resummation of the low-$`p_T`$ contributions is well known and it can be applied also to the case of gauge-boson pair production. For example, one calculation for the $`ZZ`$ has been carried out .
#### 4.1.6 Comparison with PYTHIA
In most of the studies carried out so far for the LHC, where the production of vector boson pairs played an important role, the usual Monte Carlo simulation tool has been PYTHIA based on LO matrix elements with parton shower. In particular, it is expected that for some optimisation cuts, where the large corrections provided by NLO diagrams (for example by choosing high-$`p_T(V)`$ or high-$`M_{VV}`$ regions) its predictions are not acceptable. By making comparison between the predictions of PYTHIA and the the DKS parton level NLO Monte Carlo , we investigate here how accurate does PYTHIA simulate the di-boson cross sections at the LHC, especially in some kinematic regions. We relate our analysis to the special case of the CMS detector .
In all results presented in this analysis, we assume that the vector-bosons always decay leptonically. We use the CTEQ(4M) parton distribution in both Monte Carlos and the cross section values are for the scale $`\mu =(M_{V_1}+M_{V_2})/2`$, where $`M_{V_i}`$ are the masses of the two produced vector-bosons. If the DKS Monte Carlo is run at Born-level, we obtain very good agreement with the total cross sections given by PYTHIA.
Figure 24 shows the transverse momentum of the $`WW`$ pairs. The comparison between PYTHIA and DKS indicates the large difference in cross section observables at high-$`p_T^{WW}`$ values. This is related to the fact that at NLO, the sub-processes $`qgV_1V_2q`$ have to be taken into account . This is also reported in Table 12. The leptons are selected following the CMS criteria, where a $`p_T`$ larger than 20 GeV and a pseudorapidity $`|\eta |<`$ 2.5 are required. Jets are selected by: $`p_T>`$ 20 GeV and $`|\eta |<`$ 3. The K-factor increases then from 1.5 for the total cross sections up to values of about 60 if the jets are required to have a $`p_T`$ larger than 150 GeV. The same effect is shown in figure 24 for the $`WZ`$ production, where the $`p_T`$ of the jets is shown (the jet balances the $`p_T^{VV}`$). For this process the K-factors at large $`p_T`$-values are even larger than in the $`WW`$ case (as shown in the table). The transverse momentum of the di-boson system (or of the jet(s)) are not the only variables affected by large NLO corrections. Other variables can show significant differences within their distributions: for example the lepton $`p_T`$, the invariant mass of the lepton pair $`M_{ll}`$, the missing transverse energy $`\overline{)}E_T`$ (as shown in Figure 26), the maximal transverse momentum of the two charged leptons $`p_T^{max}`$, the lepton pseudorapidities $`\eta ^l`$, their difference $`\mathrm{\Delta }\eta ^l=\eta ^l^{}\eta ^{l^+}`$, the angle between leptons $`cos\theta _{ll}`$, the transverse angle between leptons $`cos\varphi _{ll}`$ and so on.
Therefore, it is extremely important to take into account the possible influence of NLO corrections for the vector-boson production at the LHC energy. Every time one is performing an optimisation of signal selection, one should be aware of the possible deviations due to the use of a LO generator like PYTHIA. This is especially true for complicated cuts, where it is difficult to judge whether the effects are large or not. An example is shown for $`WW`$ events in Figure 26, where the smallest angle between one of the $`W`$’s and the jet is shown for events with a high-$`p_T`$ jet. Not only is the cross section clearly smaller in PYTHIA but also the shape of the distribution is quite different, changing the result of a possible cut. Another good example is the Higgs search through the decay channel $`HZZ4l`$ (see Figure 27). The idea of using $`p_T`$-cuts to improve the signal-to-background ratio may not be as effective as one would expect from using only PYTHIA. The figure shows indeed that, if the NLO corrections are included, the $`p_T`$ distribution of the non-resonant background follows much more closely those of the signal, reducing the gain considerably.
### 4.2 $`𝑾𝜸`$ and $`𝒁𝜸`$ production at NLO
In this section, we present order $`\alpha _s`$ results for $`W\gamma `$ and $`Z\gamma `$ production at the LHC, including the full leptonic correlations and anomalous couplings in the narrow-width approximation . Previous analyses included decay correlations only in the bremsstrahlung amplitudes implementing, as an approximation, the finite part of the spin-summed one-loop amplitudes.
To perform the calculation, we use the helicity amplitudes presented in . The amplitudes relevant for the inclusion of anomalous couplings are given in . In order to cancel analytically the soft and collinear singularities coming from the bremsstrahlung and one loop parts, we have used the version of the subtraction method presented in . Therefore, the amplitudes are implemented into a numerical Monte Carlo style program which allows calculation of any infrared-safe physical quantity with arbitrary cuts.
The results presented in this section correspond to $`pp`$ scattering at $`\sqrt{s}=14`$ TeV using the following cuts: a transverse momentum cut of $`p_T^l>25`$ GeV for the charged leptons is imposed and the pseudorapidity is limited to $`|\eta |<2.4`$ for all detected particles. The photon transverse momentum cut is $`p_T^\gamma >50(100)`$ GeV for $`W\gamma `$ ($`Z\gamma `$) production. For the $`W\gamma `$ case, we require a minimum missing transverse momentum carried by the neutrinos $`p_T^{\mathrm{miss}}>50`$ GeV. Additionally, charged leptons and the photons must be separated in the pseudorapidity-azimuthal angle by $`\mathrm{\Delta }R_{l\gamma }=\sqrt{(\eta _\gamma \eta _l)^2+(\varphi _\gamma \varphi _l)^2}>0.7`$. In order to suppress the contribution from the off-resonant diagrams, we require the transverse mass $`M_T>90`$ GeV for $`W\gamma `$ production and the invariant mass of the $`ll\gamma `$ system $`M_{ll\gamma }>100`$ GeV for the $`Z\gamma `$ case.
Finally, in order to suppress the contribution from the fragmentation of partons into photons, computed only to LO accuracy, the photons are required to be isolated from hadrons: the transverse hadronic momentum in a cone of size $`R_0=0.7`$ around the photon should be smaller than a fraction of the transverse momentum of the photon
$$\underset{\mathrm{\Delta }R<R_0}{}p_T^{\mathrm{had}}<0.15p_T^\gamma $$
(40)
This completes the definition of the “standard” cuts.
In the results presented here, the branching ratios of the vector-bosons into leptons are not included. For both the LO and NLO results, we use the latest set of parton distributions of MRST(cor01) and the two loop expression for the strong coupling constant. For the fragmentation component, we use the fragmentation functions from .
The “standard” scale for both the factorisation and renormalisation scales is
$$\mu ^2=\mu _{\mathrm{st}}^2M_V^2+\frac{1}{2}\left[(p_T^V)^2+(p_T^\gamma )^2\right].$$
(41)
The masses of the vector-bosons have been set to $`M_Z=91.187`$ GeV and $`M_W=80.41`$ GeV and the following values have been used for the Cabibbo-Kobayashi-Maskawa (CKM) matrix elements: $`|V_{ud}|=|V_{cs}|=0.975`$ and $`|V_{us}|=|V_{cd}|=0.222`$. We do not include any QED or electroweak corrections but choose the coupling constants $`\alpha `$ and $`\mathrm{sin}^2\theta _W`$ in the spirit of the “improved Born approximation” , with $`\mathrm{sin}^2\theta _W=0.230`$. Notice that the observable is order $`\alpha ^2`$; within the same spirit, we use the running $`\alpha =\alpha (M_Z)=1/128`$ for the coupling between the vector-boson and the quarks (to take into account effectively the EW corrections) whereas we keep $`\alpha =1/137`$ for the photon coupling. It is worth noticing that this modification results already in more than a 6% change in the normalisation of the cross section with respect to the standard approach of using both running coupling constants.
#### 4.2.1 Results at NLO
For future checks, and for an estimate of the number of events to be observed at the LHC, some benchmark total cross section numbers are presented in Table 13. The first ones were obtained by imposing only the cut on the transverse momentum of the photon $`p_T^\gamma >50(100)`$ GeV for $`W\gamma `$ ($`Z\gamma `$) production. The importance of the NLO corrections, as well as the size of the fragmentation contribution before applying the isolation cut prescription, can be seen from the table. Furthermore, we also include the result for the total cross section obtained after the implementation of the standard cuts.
In what follows, we will estimate the theoretical uncertainty of the results by analysing the changes on different distributions when varying the scale by a factor of two in both directions $`\frac{\mu _{\mathrm{st}}}{2}<\mu <2\mu _{\mathrm{st}}`$.
In Figure 28, we show the scale dependence of the $`p_T`$ distribution of the photon in $`W^+\gamma `$ production with the standard cuts (upper curves) and also with the additional requirement of a jet-veto. As can be observed, the scale dependence is still large ($`\pm `$ 10%) but is considerably reduced when the jet-veto is applied. The situation is similar to what has been observed in the case of $`WW`$ production and is caused by the suppression of the contribution from the $`qg`$ initial state appearing for the first time at NLO. Since this initial state dominates the cross section, the NLO result behaves effectively like a LO one, as far as the scale dependence is concerned.
In the inset plot, we present the ratio between the NLO and LO results (with the standard scale), which remains larger than 3 and increases with the photon transverse momentum. This clearly shows that the LO calculation is not even sufficient for an understanding of the shape of the distribution, since the NLO effect goes beyond a simple normalisation. As is well known , the relevance of the NLO corrections for this process is mainly due to the breaking of the radiation amplitude zero appearing at LO and to the large $`qg`$ initial state parton luminosity at the LHC. It is worth mentioning that the scale dependence of the LO result turns out to be very small. This is an artificial effect and illustrates that a small scale dependence is by no means a guarantee for small NLO corrections. Furthermore, we present the ratio of the NLO jet-veto and the LO result. As expected, this ratio is closer to 1, again due to the fact that most of the contributions coming from the new subprocesses appearing at NLO are suppressed by the jet-veto.
In Figure 29, we study the lepton correlation in the azimuthal angle for $`Z\gamma `$ production $`\mathrm{\Delta }\varphi _{ll}=|\varphi _l^{}\varphi _{l^+}|`$. Notice that this observable can be studied at NLO since the spin correlations between the leptons are fully taken into account in the implementation of the one-loop corrections. In this case, we observe that the NLO corrections are rather sizeable and increase the cross section by $`50\%`$ for small $`\mathrm{\Delta }\varphi _{ll}`$. The region $`\mathrm{\Delta }\varphi _{ll}>2`$ (with the standard cuts) is kinematically forbidden unless a jet with a high transverse momentum is produced; therefore, the cross section vanishes at LO and it is strongly suppressed for the NLO calculation with jet-veto. In this region, the full NLO calculation is effectively LO and its scale dependence becomes larger, as expected.
Because there is no radiation amplitude zero appearing at LO for $`Z\gamma `$ production, the NLO corrections are under better control in the kinematical region where the LO cross section does not vanish. Nevertheless, for large transverse momentum, the $`qg`$ initial state again dominates the NLO contribution and the corrections increase considerably.
#### 4.2.2 Anomalous couplings without form factors
The study of triple vector-boson couplings is motivated by the hope that some physics beyond the Standard Model leads to a modification of these couplings which eventually could be detected. In order to quantify the effects of the new physics, an effective Lagrangian is introduced which contains all Lorentz invariant terms, in principle. The new terms spoil the gauge-cancellation in the high energy limit and, therefore, will lead to violation of unitarity for increasing partonic centre of mass energy $`\widehat{s}`$. Usually, in an analysis of anomalous couplings from experimental data in hadronic collisions, this problem is circumvented by supplementing the anomalous couplings $`\alpha _{\mathrm{AC}}`$ with form factors. A common choice for the form factor is
$$\alpha _{\mathrm{AC}}\frac{\alpha _{\mathrm{AC}}}{(1+\frac{\widehat{s}}{\mathrm{\Lambda }^2})^n}$$
(42)
where $`n`$ has to be large enough to ensure unitarity and $`\mathrm{\Lambda }`$ is interpreted as the scale for new physics. Obviously, this procedure is rather ad hoc and introduces some arbitrariness. Therefore, it would be very convenient to avoid it in an analysis of anomalous couplings at hadron colliders. This would bring these analyses more into line with those at $`e^+e^{}`$ colliders. In order to do so, one should analyse the data at fixed values of $`\widehat{s}`$, as it is done at LEP. This results in limits for the anomalous parameters which are a function of $`\widehat{s}`$.
Clearly, it is possible to do such analysis for the production of $`Z\gamma `$ when both leptons are detected , since the partonic centre of mass energy can be reconstructed from the kinematics of the final state particles and therefore the cross section can be measured for different bins of fixed $`\widehat{s}`$.
The situation is more complicated for $`W\gamma `$ production since the neutrino is not observed. Nevertheless, by identifying the transverse momentum of the neutrino with the missing transverse momentum, and assuming the $`W`$ boson to be on shell, it is possible to reconstruct the neutrino kinematics (particularly the longitudinal momentum) with a two-fold ambiguity. In the case of the Tevatron, since it is a $`p\overline{p}`$ collider, it is possible to choose the “correct” neutrino kinematics 73% of the times by selecting the maximum (minimum) of the two reconstructed values for the longitudinal momentum of the neutrino for $`W^+\gamma `$($`W^{}\gamma `$).
This is not true at the LHC where, due to the symmetry of the colliding beams, both reconstructed kinematics have equal chances to be correct. Fortunately, in the case of anomalous couplings, we are interested in a efficient way to reconstruct the $`\widehat{s}`$ rather than the full kinematics. Again there are two possible values of $`\widehat{s}`$. It turns out that there is a simple method to choose the “correct” one 66% of the times at the LHC (73% of the times at Tevatron) by selecting the minimum $`\widehat{s}`$, $`\widehat{s}_{\mathrm{min}}`$, of the two reconstructed values (for both $`W^+\gamma `$ and $`W^{}\gamma `$). Furthermore, we checked that the selected value $`\widehat{s}_{\mathrm{min}}`$ differs in almost 90% of the events by less than 10% from the exact value $`\widehat{s}`$. This is likely to be enough precision, since the data will be collected in sizeable bins of $`\widehat{s}`$ and the anomalous parameters are not expected to change very rapidly with the energy in any case.
To quantify the advantage of the method, we show in Figure 30 the correlations of $`\sqrt{\widehat{s}_{\mathrm{min}}}`$ with $`\sqrt{\widehat{s}}`$. The left plot corresponds to the case of pure Standard Model, whereas the right plot presents results for (already experimentally ruled out) huge values of anomalous couplings $`\mathrm{\Delta }\kappa =0.8`$ and $`\lambda =0.2`$ with an ordinary form factor ($`n=2`$, $`\mathrm{\Lambda }=1`$ TeV).
The cross section drops very rapidly for increasing $`\sqrt{\widehat{s}}\sqrt{\widehat{s}_{\mathrm{min}}}`$. This correlation clearly holds in the particularly interesting large $`\sqrt{\widehat{s}}`$ region and for both Standard Model and anomalous contribution.
As a result of this investigation, we conclude that even in the case of $`W\gamma `$ production, reliable bounds for anomalous couplings as a function of $`\widehat{s}`$ (using $`\widehat{s}_{\mathrm{min}}`$) can be obtained. Such a procedure would certainly allow a comparison of various bounds from different experiments.
## 5 ANOMALOUS VECTOR-BOSON COUPLINGS <sup>10</sup><sup>10</sup>10Section coordinators: P.R. Hobson, W. Hollik
The principle of gauge-invariance is used as the basis for the Standard Model. The non-Abelian gauge-group structure of the theory of electroweak interactions predicts very specific couplings between the electroweak gauge-bosons. Measurements of these triple gauge-boson couplings (TGCs) of the $`W`$, $`Z`$ and $`\gamma `$ gauge-bosons therefore provide powerful tests of the Standard Model.
In the most general Lorentz invariant parametrisation, the three gauge-boson vertices, $`WW\gamma `$ and $`WWZ`$, can be described by fourteen independent couplings , seven for each vertex. The possible four quadruple gauge-boson vertices: $`\gamma \gamma WW`$, $`Z\gamma WW`$, $`ZZWW`$ and $`WWWW`$ require 36, 54, 81 and 81 couplings, respectively for a general description. Assuming electromagnetic gauge-invariance, C- and P-conservation, the set of 14 couplings for the three gauge-boson vertices is reduced to 5: $`g_1^Z`$, $`\kappa _\gamma `$, $`\kappa _Z`$, $`\lambda _\gamma `$ and $`\lambda _Z`$ , where their Standard Model values are equal to $`g_1^Z=\kappa _\gamma =\kappa _Z=1`$ and $`\lambda _\gamma =\lambda _Z=0`$ at tree level.
The TGCs related to the $`WW\gamma `$ vertex determine properties of the $`W`$, such as its magnetic dipole moment $`\mu _W`$ and electric quadrupole moment $`q_W`$:
$$\mu _W=\frac{e}{2M_W}(g_1^Z+\kappa _\gamma +\lambda _\gamma )$$
(43)
$$q_W=\frac{e}{M_W^2}(\kappa _\gamma \lambda _\gamma )$$
(44)
In the following, the anomalous TGCs are denoted by $`\mathrm{\Delta }g_1^Z`$, $`\mathrm{\Delta }\kappa _\gamma `$, $`\mathrm{\Delta }\kappa _Z`$, $`\lambda _\gamma `$ and $`\lambda _Z`$, where the $`\mathrm{\Delta }`$ denotes the deviations of the respective quantity from its Standard Model value.
### 5.1 Introduction
The Standard Model is well established by the experiments at LEP and the Tevatron. Any deviations of the Standard Model can therefore be introduced only with care. Changes to the Standard Model come with different forms of severity. In order to see at what level anomalous vector-boson couplings can be reasonably discussed, one has to consider these cases separately. Changes to the gauge-structure of the theory, that do not violate the renormalisability of the theory, i.e. the introduction of extra fermions or possible extensions of the gauge-group are the least severe. They will typically generate small corrections to vector-boson couplings via loop effects. In this case also, radiative effects will be generated at lower energies. For the LHC, the important thing in this case is not to measure the anomalous couplings precisely, but to look for the extra particles. However, this is beyond the scope of this chapter. In the other case, a more fundamental role is expected for the anomalous couplings, implying strong interactions. In this case, one has to ask oneself whether one should study a model with or without a fundamental Higgs boson.
Simply removing the Higgs boson from the Standard Model is a relatively mild change. The model becomes non-renormalisable, but the radiative effects grow only logarithmically with the cut-off at the one-loop level. The question is whether this scenario is ruled out by the LEP1 precision data. The LEP1 data appear to be in agreement with the Standard Model, preferring a low Higgs mass. One is sensitive to the Higgs mass in three parameters, labelled $`S`$, $`T`$, $`U`$ or $`ϵ_1,ϵ_2,ϵ_3`$. These receive corrections of the form $`g^2(\mathrm{log}(M_H/M_W)+constant)`$, where the constants are of order one. The logarithmic enhancement is universal and would also appear in models without a Higgs as $`\mathrm{log}(\mathrm{\Lambda })`$, where $`\mathrm{\Lambda }`$ is the cut-off at which new interactions should appear. Only when one can determine the three different constants independently, can one say that one has established the Standard Model. At present, the data do not provide sufficient precision to do this.
A much more severe change to the Standard Model is the introduction of vector-boson couplings not of the gauge-interaction type. These new couplings violate renormalisability much more severely than simply removing the Higgs boson. Typically, quadratically and quartically divergent corrections would appear to physical observables. Therefore, it is questionable as to whether one should study models with a fundamental Higgs boson, but with extra anomalous vector-boson couplings. It is hard to imagine a form of dynamics that could do this. If the vector-bosons become strongly interacting, the Higgs probably would exist at most in an “effective” way. Therefore, the most natural way is to study anomalous vector-boson couplings in models without a fundamental Higgs. Actually when one removes the Higgs boson, the Standard Model becomes a gauged non-linear sigma-model. It is well known that the nonlinear sigma-model describes low-energy pion physics. The “pions” correspond to the longitudinal degrees of freedom of the vector-bosons and $`f_\pi `$ corresponds to the vacuum expectation value of the Higgs field. Within this description, the Standard Model corresponds to the lowest-order term quadratic in the momenta, anomalous couplings to higher derivative terms. The systematic expansion in terms of momenta is known as chiral perturbation theory and is extensively used in meson physics.
Writing down the most general non-linear chiral Lagrangian containing up to four derivatives gives rise to a large number of terms, which are too general to be studied effectively. One therefore has to look for dynamical principles that can limit the number of terms. Of particular importance are approximate symmetry principles. In the first place one, expects CP-violation to be small. We limit ourselves therefore to CP-preserving terms. In order to see what this means in practice, it is advantageous to describe the couplings in a manifestly gauge-invariant way, using the Stückelberg formalism . One needs the following definitions:
$$F_{\mu \nu }=\frac{i\tau _i}{2}(_\mu W_\nu ^i_\nu W_\mu ^i+gϵ^{ijk}W_\mu ^jW_\nu ^k)$$
(45)
is the $`SU(2)`$ field strength with the $`SU(2)`$ gauge-coupling $`g`$;
$$D_\mu U=_\mu U+\frac{ig}{2}\tau _iW_\mu ^iU+ig\mathrm{tan}\theta _WU\tau _3B_\mu $$
(46)
is the gauge-covariant derivative of the $`SU(2)`$-valued field $`U`$, which describes the longitudinal degrees of freedom of the vector fields in a gauge-invariant way;
$$B_{\mu \nu }=_\mu B_\nu _\nu B_\mu $$
(47)
is the hypercharge field strength. In addition,
$`V_\mu `$ $`=`$ $`(D_\mu U)U^{}/g,`$ (48)
$`T`$ $`=`$ $`U\tau _3U^{}/g`$ (49)
are auxiliary quantities having simple transformation properties. Excluding CP violation, the non-standard three and four vector-boson couplings are described in this formalism by the following set of operators:
$`_1`$ $`=`$ $`\mathrm{Tr}(F_{\mu \nu }[V_\mu ,V_\nu ])`$ (50)
$`_2`$ $`=`$ $`i{\displaystyle \frac{B_{\mu \nu }}{2}}\mathrm{Tr}(T[V_\mu ,V_\nu ])`$ (51)
$`_3`$ $`=`$ $`\mathrm{Tr}(TF_{\mu \nu })\mathrm{Tr}(T[V_\mu ,V_\nu ])`$ (52)
$`_4`$ $`=`$ $`(\mathrm{Tr}[V_\mu V_\nu ])^2`$ (53)
$`_5`$ $`=`$ $`(\mathrm{Tr}[V_\mu V_\mu ])^2`$ (54)
$`_6`$ $`=`$ $`\mathrm{Tr}(V_\mu V_\nu )\mathrm{Tr}(TV_\mu )\mathrm{Tr}(TV_\nu )`$ (55)
$`_7`$ $`=`$ $`\mathrm{Tr}(V_\mu V_\mu )(\mathrm{Tr}[TV_\nu ])^2`$ (56)
$`_8`$ $`=`$ $`{\displaystyle \frac{1}{2}}[(\mathrm{Tr}[TV_\mu ])(\mathrm{Tr}[TV_\nu ])]^2`$ (57)
In the unitary gauge $`U=1`$, one has (with $`c_W=\mathrm{cos}\theta _W`$, $`s_W=\mathrm{sin}\theta _W`$)
$`_1`$ $`=`$ $`i(c_WZ_{\mu \nu }+s_WF_{\mu \nu })W_\mu ^+W_\nu ^{}+Z_\nu /c_W(W_{\mu \nu }^+W_\mu ^{}W_{\mu \nu }^{}W_\mu ^+)`$
$`+\text{gauge-induced four boson vertices},`$
$`_2`$ $`=`$ $`i(c_WF_{\mu \nu }s_WZ_{\mu \nu })W_\mu ^+W_\nu ^{},`$ (59)
$`_3`$ $`=`$ $`i(c_WZ_{\mu \nu }+s_WF_{\mu \nu })W_\mu ^+W_\nu ^{}.`$ (60)
where $`Z_{\mu \nu }=_\mu Z_\nu _\nu Z_\mu `$ and $`W_{\mu \nu }^{+,}=_\mu W_\nu ^{+,}_\nu W_\mu ^{+,}`$. The Standard Model without a Higgs corresponds to
$$_{EW}=\frac{1}{2}\mathrm{Tr}(F_{\mu \nu }F^{\mu \nu })\frac{1}{4}B_{\mu \nu }B^{\mu \nu }+\frac{g^2v^2}{4}\mathrm{Tr}(V_\mu V^\mu ).$$
(61)
### 5.2 Dynamical constraints
The list given in the previous section contains terms that give rise to vertices with minimally three or four vector-bosons. Already with the present data a number of constraints and/or consistency conditions can be put on the vertices. The most important of these come from the limits on the breaking of the so-called custodial symmetry. If the hypercharge is put to zero, the effective Lagrangian has a larger symmetry than $`SU_L(2)\times U_Y(1)`$, i.e. it has the symmetry $`SU_L(2)\times SU_R(2)`$. The $`SU_R(2)`$ invariance is a global invariance. Within the Standard Model this invariance is an invariance of the Higgs potential, but not of the full Lagrangian. It is ultimately this invariance that is responsible for the fact that the $`\rho `$ parameter, which is the ratio of charged to neutral current strength, is equal to one at the tree level. Some terms in the Lagrangian, i.e. the ones containing the hypercharge field explicitly or the terms with $`T`$, that project out the third isospin component violate this symmetry explicitly. These terms, when inserted in a loop graph, give rise to quartically divergent contributions to the $`\rho `$ parameter. Given the measurements, this means that the coefficients of these terms must be extremely small. It is therefore reasonable to limit oneself to a Lagrangian, where hypercharge appears only indirectly via a minimal coupling, so without explicit $`T`$. This assumption means physically that the ultimate dynamics that is responsible for the strong interactions among the vector-bosons acts in the non-Abelian sector. Indeed one would not normally expect the hypercharge alone to become strong. However, we know that there is a strong violation of the custodial symmetry in the form of the top-quark mass. Actually the top-mass almost saturates the existing corrections to the $`\rho `$ parameter, leaving no room for violations of the custodial symmetry in the anomalous vector-boson couplings. Therefore, we conclude: If there really are strong vector-boson interactions, the mechanism for mass generation is unlikely to be the same for bosons and fermions.
Eliminating the custodial symmetry violating interactions, we are left with the simplified Lagrangian, containing $`_1`$, $`_4`$, $`_5`$. Besides the vertices, there are also propagator corrections, in principle. We take the two-point functions without explicit $`T`$. Specifically, we add to the theory
$$_{hc,tr}=\frac{1}{2\mathrm{\Lambda }_W^2}\mathrm{Tr}[(D_\alpha F_{\mu \nu })(D^\alpha F^{\mu \nu })]+\frac{1}{2\mathrm{\Lambda }_B^2}\mathrm{Tr}[(_\alpha B_{\mu \nu })(^\alpha B^{\mu \nu })]$$
(62)
for the transverse degrees of freedom of the gauge-fields, and
$$_{hc,lg}=\frac{g^2v^2}{4\mathrm{\Lambda }_V^2}\mathrm{Tr}[(D^\alpha V^\mu )(D_\alpha V_\mu )]$$
(63)
for the longitudinal ones, where the $`\mathrm{\Lambda }_X`$ parametrise the quadratic divergences and are expected to represent the scales where new physics comes in. In phenomenological applications, these contributions give rise to form factors in the propagators . Introducing such cut-off dependent propagators in the analysis of the vector-boson pair production is similar to having $`s`$-dependent triple vector-boson couplings, which is the way the data are usually analysed.
This effective Lagrangian is very similar to the one in pion-physics. Indeed, if one takes the limit vacuum expectation value (vev) fixed and gauge-couplings to zero, one finds the standard pion Lagrangian. As it stands, one can use the LEP1 data to put a limit on the terms in the two point vertices. Using a naive analysis one finds $`1/\mathrm{\Lambda }_B^2=0`$. For the other two cut-offs one has:
A. The case $`\mathrm{\Lambda }_V^2>0,\mathrm{\Lambda }_W^2<0`$: $`\mathrm{\Lambda }_V>0.49`$ TeV, $`|\mathrm{\Lambda }_W|>1.3`$ TeV.
B. The case $`\mathrm{\Lambda }_V^2<0,\mathrm{\Lambda }_W^2>0`$: $`|\mathrm{\Lambda }_V|>0.74`$ TeV, $`\mathrm{\Lambda }_W>1.5`$ TeV.
This information is important for further limits at high-energy colliders, as it tells us, how one has to cut off off-shell propagators. We notice that the limits on the form factors are different for the transverse, longitudinal and hypercharge form factors. The precise limits are somewhat qualitative and should be taken as such. The current data show that $`\mathrm{\Lambda }=0.5`$ TeV, which thus has to be considered as a minimal possible value as long as a dipole form factor is used. Further information comes from the direct measurements of the three-point couplings at LEP2, which tell us that they are small. Similar limits at the Tevatron have to be taken with some care, as there is a cut-off dependence. As there is no known model that can give large three-point interactions, we assume for the further analysis of the four-point vertices, that the three-point anomalous couplings are absent. Two more constraints can be put on the remaining two four-point vertices . The first comes from consistency of chiral perturbation theory . Not every effective chiral Lagrangian can be generated from a physical underlying theory.
A second condition comes from the $`\rho `$ parameter. Even the existing violation of the custodial symmetry, though indirect via the minimal coupling to hypercharge, gives a contribution to the $`\rho `$ parameter. It constrains the combination $`5g_4+2g_5`$. The remaining combination $`2_45_5`$ is fully unconstrained by experiment and in principle gives a possibility for very strong interactions to be present. However, this particular combination does not seem to have any natural interpretation from underlying dynamics. Therefore, one can conclude presumably that both couplings $`g_4,g_5`$ are small. There is a loophole to this conclusion, namely when the anomalous couplings are so large that the one-loop approximation, used to arrive at the limits, is not consistent and resummation has to be performed everywhere. This is a somewhat exotic possibility that could lead to very low-lying resonances and which ought to be easy to discover at the LHC .
### 5.3 LHC processes
Given the situation described above, one has to ask oneself, what the LHC can do and in which way the data should be analysed. There are essentially three processes that can be used to study vector-boson vertices: vector-boson pair production, vector-boson scattering, triple vector-boson production. About the first two we have only a few remarks to make. They are discussed more fully in other contributions to the workshop.
#### 5.3.1 Vector-boson pair production
Vector-boson pair production can be studied in a relatively straightforward way. The reason is that here the Higgs boson does not play a role in the Standard Model, as we take the incoming quarks to be massless. Therefore naive violations of unitarity can be compensated by the introduction of smooth form-factors.
One produces two vector-bosons via normal Standard Model processes with an anomalous vertex added. The extra anomalous coupling leads to unitarity-violating cross sections at high energy. As a total energy of 14 TeV is available this is a serious problem, in principle. It is cured by introducing a form factor for the incoming off-shell line connected to the anomalous vertex. Naively this leads to a form-factor dependent limit on the anomalous coupling in question. The LEP1 data gives a lower limit on the cut-off to be used inside the propagator. When one wants an overall limit on the anomalous coupling, one should use this value. This is particularly relevant for the Tevatron. Here one typically takes a cut-off of 2 TeV. This might give too strict a limit, as the LEP1 data indicate that the cut-off can be as low as 500 GeV. For practical purposes the analysis at the Tevatron should give limits on anomalous couplings for different values of the cut-off form factors, including low values of the cut-off. For the analysis at the LHC, one has much larger statistics. This means that one can do better and measure limits on the anomalous couplings as a function of the invariant mass of the produced system. This way one measures the anomalous form factor completely.
#### 5.3.2 Vector-boson scattering
Here the situation is more complicated than in vector-boson pair production. The reason is that within the Standard Model the process cannot be considered without intermediate Higgs contribution. This would violate unitarity. However the incoming vector-bosons are basically on-shell and this allows the use of unitarisation methods, as are commonly used in chiral perturbation theory in pion physics. These methods tend to give rise to resonances in longitudinal vector-boson scattering. The precise details depend on the coupling constants. The unitarisation methods are not unique, but generically give rise to large $`I=J=0`$ and/or $`I=J=1`$ cross section enhancements. The literature is quite extensive: a good introduction is ; a recent review is .
#### 5.3.3 Triple vector-boson production
In this case it is not clear how one should consistently approach an analysis of anomalous vector-boson couplings. Within the Standard Model the presence of the Higgs boson is essential in this channel. Leaving it out, one has to study the unitarisation. This unitarisation has to take place not only on the two-to-two scattering subgraphs, as in vector-boson scattering, but also on the incoming off-shell vector-boson, decaying into three real ones. The analysis here becomes too arbitrary to derive very meaningful results. One cannot calculate confidently anything here without a fully known underlying model of new strong interactions. Also measurable cross sections tend to be small, so that the triple vector-boson production is best used as corroboration of results in vector-boson scattering. Deviations of Standard Model cross sections could be seen, but the vector-boson scattering would be needed for interpretation.
One therefore needs the Standard Model results. The total number of events with three vector-bosons in the final state is given in Table 14. We used an integrated luminosity of 100 fb<sup>-1</sup> and an energy of 14 TeV throughout.
One sees from this table that a large part of the events comes from associated Higgs production, when the Higgs is light. However for the study of anomalous vector-boson couplings, the heavier Higgs results are arguably more relevant. Not all the events can be used for the analysis. If we limit ourselves to events, containing only electrons, muons and neutrinos, assuming just acceptance cuts we find the results shown in Table 15.
We see that very little is left, in particular in the processes with at least two $`Z`$ bosons, where the events can be fully reconstructed. In order to see how sensitive we are to anomalous couplings, we assumed a 4$`Z`$ coupling with a form factor cut-off at 2 TeV. We make here no correction for efficiencies etc.. Using the triple $`Z`$ boson production, assuming no events are seen in 100 fb<sup>-1</sup>, we find a limit $`|g_4+g_5|<0.09`$ at the 95% CL, where $`g_4`$ and $`g_5`$ are the coefficients multiplying the operators $`_4`$ and $`_5`$. This is to be compared with $`0.15<5g_4+2g_5<0.14`$ or $`0.066<(5g_4+2g_5)\mathrm{\Lambda }^2(\mathrm{TeV})<0.026`$ . So the sensitivity is not better than present indirect limits. Better limits exist in vector-boson scattering or at a linear collider .
In the following tables we present numbers for observable cross sections in different decay modes of the vector-bosons. We used the following cuts.
$`|\eta |_{lepton}<3,|\eta |_{jet}<2.5,`$
$`|p_T|_{lepton}>20\mathrm{GeV},|p_T|_{jet}>40\mathrm{GeV},|p_T|_{2\nu }>50\mathrm{GeV},`$
$`\mathrm{\Delta }R_{jet,lepton}>0.3,\mathrm{\Delta }R_{jet,jet}>0.5.`$
States with more than two neutrinos are not very useful because of the background from two vector-boson production. We did not consider final states containing $`\tau `$-leptons.
With the given cuts, the total number of events to be expected is rather small. In particular, this is the case because we did not consider the reduction in events due to experimental inefficiencies, which may be relatively large because of the large number of particles in the final state. For the processes containing jets in the final state, there will be large backgrounds due to QCD processes. A final conclusion on the significance of the triple vector-boson production for constraining the four vector-boson couplings will need more work, involving detector Monte Carlo calculations.
However it is probably fair to say from the above results, that no very strong constraints will be found from this process at the LHC, but it is useful as a cross-check with other processes. It may provide complementary information if non-zero anomalous couplings are found.
### 5.4 Unitarity limits and form factors
Unitarity in the Standard Model depends directly on its gauge-structure. Departure from this structure can violate unitarity at relatively low energies and so protection is provided in the effective Lagrangian for triple gauge-boson vertices by expressing the anomalous couplings as energy dependent form factors. For experimental results at a given subprocess energy $`\widehat{s}`$ (i.e. $`e^+e^{}`$ colliders), the choice of form factor parametrisation is not important since one can unambiguously translate between parametrisations. However, when results are integrated over a range of $`\widehat{s}`$ as they will be at the LHC, no simple translation is possible and results depend crucially on the choice of the form factors. The form factor behaviour of anomalous couplings should not be neglected, particularly in regions of $`\widehat{s}`$ near to unitarity limits. Any measurement of anomalous couplings over integrated energies carries with it assumptions on the parametrisation of the form factor.
This section outlines the considerations which influence the choice of form factor and suggests a method for measuring energy dependent anomalous couplings.
#### 5.4.1 Form factor parametrisation
Triple gauge-boson vertices in di-boson production arise in the $`J=1`$ partial wave amplitude only ($`s`$-channel exchange of a gauge-boson coupled to massless fermions). $`S`$-matrix unitarity implies a constant bound to any partial wave amplitude. This means unitarity is violated at asymptotically high energies if constant anomalous couplings are assumed. Unambiguous and model-independent constant unitarity constraints for $`WV`$ production have been derived<sup>12</sup><sup>12</sup>12 Cancellations may occur if more than one anomalous coupling is allowed non-zero at a time, which weakens the unitarity limits somewhat. .
To conserve unitarity at arbitrary energies, anomalous couplings must be introduced as form factors. Thus, an arbitrary anomalous coupling $`\stackrel{~}{A}=\stackrel{~}{A}_0\times (q_1^2,q_2^2,P^2)`$ vanishes when $`q_1^2,q_2^2,`$ or $`P^2`$ becomes large, where $`q_1^2`$ and $`q_2^2`$ are the invariant masses squared of the production bosons and $`P^2=\widehat{s}`$ is the virtual exchange boson invariant mass squared. We refer to $`\stackrel{~}{A}_0`$ as the “bare coupling” and $`\stackrel{~}{A}`$ as the form factor ($`\stackrel{~}{A}ϵ\lambda ^V,\mathrm{\Delta }\kappa ^V,h_i^V,`$…). For di-boson production, the final state bosons are nearly on-shell $`q_1^2,q_2^2M_V^2`$ even when finite width effects are taken into account, though large virtual exchange boson masses $`\sqrt{\widehat{s}}`$ will be probed at the LHC.
The choice of parametrisation for the form factors is arbitrary provided unitarity is conserved at all energies for a sufficiently small value of anomalous coupling. A step function operating at a cutoff scale $`\mathrm{\Lambda }_{\text{FF}}`$ is sufficient<sup>13</sup><sup>13</sup>13 i.e. assuming a step function form factor operating at 2 TeV, the $`\lambda ^\gamma `$ coupling conserves unitarity for $`\lambda ^\gamma <0.99`$ \[145, Equation 23\]. though discontinuous and thus unphysical. More common in the literature is a generalised dipole form factor which is motivated by the well known nucleon form factors and has further appeal because it enters the Lagrangian in a form similar to that of a propagator of mass $`\mathrm{\Lambda }_{\text{FF}}`$. The parametrisation is
$$\stackrel{~}{A}=\frac{\stackrel{~}{A}_0}{(1+\frac{\widehat{s}}{\mathrm{\Lambda }_{\text{FF}}^2})^n}$$
(64)
where $`n>1/2(n>1)`$ is sufficient for the $`WWV`$ vertex anomalous couplings $`\mathrm{\Delta }\kappa ^V(\lambda ^V,\mathrm{\Delta }g_1^V)`$ which grow like $`\widehat{s}^{1/2},(\widehat{s})`$. For the $`ZV\gamma `$ vertex $`n>3/2(n>5/2)`$ is sufficient for anomalous couplings $`h_{1,3}^V,(h_{2,4}^V)`$ which grow like $`\widehat{s}^{3/2},(\widehat{s}^{5/2})`$. The usual assumptions are $`n=2`$ for $`g_1^V,\lambda ^V,\kappa ^V`$ and $`n=3(n=4)`$ for $`h_{1,3}^V,(h_{2,4}^V)`$ . Unitarity limits for generalised dipole form factors have been enumerated \[147, Equations 22-26\].
The form factor scale $`\mathrm{\Lambda }_{\text{FF}}`$ can be regarded as a regularisation scale. It is related to (but not necessarily identical to) the energy scale at which new physics becomes important in the weak boson sector.
#### 5.4.2 Impact of form factor on $`\widehat{s}`$ dependent distributions
The impact of the form factor parametrisation on $`\widehat{s}`$ dependent distributions is illustrated in Figure 31 where the reconstructed <sup>14</sup><sup>14</sup>14 Reconstructing $`M_{\text{inv}}(WZ)`$ requires knowledge of the neutrino longitudinal momentum which is obtained up to a two-fold ambiguity using the $`W`$ mass constraint. Each solution is given half weight in the $`M_{\text{inv}}(WZ)`$ spectrum. $`M_{\text{inv}}(WZ)`$ and $`p_T(Z)`$ spectra are plotted for LHC $`W^+Z`$ production with leptonic decays at $`O(\alpha _s)`$. The Standard Model expectation is compared to scenarios with a modest $`\lambda _0^Z=0.05`$ coupling for various generalised dipole form factor parametrisations.
For the region of low invariant mass where $`\sqrt{\widehat{s}}\mathrm{\Lambda }_{\text{FF}}`$, the form factors remain essentially constant and distributions with the same bare coupling agree well. As the form factor scale $`\mathrm{\Lambda }_{\text{FF}}`$ is approached, the distributions begin to be pushed back to the SM expectation (visible at about $`M_{\text{inv}}(WZ)=500`$ GeV for the $`\mathrm{\Lambda }_{\text{FF}}`$=2 TeV case). For $`\sqrt{\widehat{s}}>\mathrm{\Lambda }_{\text{FF}}`$ the distribution returns to the SM expectation. The exponent of the form factor $`n`$ dictates how fast the “pushing” occurs as $`\mathrm{\Lambda }_{\text{FF}}`$ is approached. Thus distributions sensitive to the $`ZV\gamma `$ vertex (for which $`n=`$3 or 4 is the usual choice) exhibit a more pronounced form factor behaviour than distributions sensitive to the $`WWV`$ vertex (for which $`n=2`$ is usual).
Since distributions are constrained to the SM expectation at invariant masses above the form factor scale, great care should be taken when fitting to a form factor parametrised model in a region with data where $`\sqrt{\widehat{s}}\mathrm{\Lambda }_{\text{FF}}`$. Effectively, since the anomalous couplings are constrained near zero above $`\mathrm{\Lambda }_{\text{FF}}`$ by the parametrisation model, there are no free parameters for the fit in this $`\widehat{s}`$ region. For the case of observable non-zero anomalous couplings, an analysis assuming a parametrisation of the form factor with fixed $`\mathrm{\Lambda }_{\text{FF}}`$ smaller than that provided by nature but within the $`\widehat{s}`$ accessible by the machine would overestimate the anomalous coupling. This is because large bare coupling fit values are necessary in the $`\sqrt{\widehat{s}}\mathrm{\Lambda }_{\text{FF}}`$ region to counter the (artificially imposed) form factor behaviour.
#### 5.4.3 Impact of form factor scale on sensitivity limits
If triple gauge-coupling (TGC) measurements are consistent with the SM and confidence limits are to be derived, it is impossible to avoid form factor parametrisation assumptions.
The dependence of anomalous coupling limits on the form factor scale $`\mathrm{\Lambda }_{\text{FF}}`$ is illustrated in Figure 32 where the 95% confidence limits for $`WW\gamma `$ vertex anomalous $`\lambda _0^\gamma ,\mathrm{\Delta }\kappa _0^\gamma `$ couplings in $`W\gamma `$ production with $`We\nu _e,\mu \nu _\mu `$ are presented as a function of $`\mathrm{\Lambda }_{\text{FF}}`$ for a dipole form factor with $`n=2`$. The limits are for illustrative purposes only and have been derived at NLO generator level using a binned maximum likelihood fit to the $`p_T(\gamma )`$ distribution. No detector simulation has been applied and the specific choice of cuts are unimportant.
The unitarity limit curve is superimposed. The region above this is non-physical (violates unitarity). The curve is independent of experiment and analysis but depends on the form factor parametrisation. It goes asymptotically to zero for large $`\mathrm{\Lambda }_{\text{FF}}`$ indicating TGC couplings are restricted to SM values at extreme energies.
Simulated experimental limits for the Tevatron (2 TeV $`p\overline{p}`$ collisions, $`=100`$ pb<sup>-1</sup>) and the LHC (14 TeV $`pp`$ collisions, $`=300`$ fb<sup>-1</sup>) are presented. The limits depend on the analysis and machine parameters. The restricted $`\widehat{s}`$ accessible by the machines result in an asymptotic behaviour wherein an optimal limit for anomalous couplings is reached. We refer to the scale at which this occurs as $`\mathrm{\Lambda }_{\text{machine}}`$. A measurement with this scale reflects the maximal discovery potential for anomalous couplings for a given machine (since the full spectra in $`\widehat{s}`$ contributes to the limit). It occurs at about 2 TeV for the Tevatron and about 5-10 TeV for the LHC for $`\lambda ^\gamma ,\mathrm{\Delta }\kappa ^\gamma `$ and lies below the unitarity limit in both cases. The experimental limits are not sensitive to changes in $`\mathrm{\Lambda }_{\text{FF}}`$ for $`\mathrm{\Lambda }_{\text{FF}}\mathrm{\Lambda }_{\text{machine}}`$. Indeed, in this region the distributions behave exactly as if the form factors were constants $`\stackrel{~}{A}\stackrel{~}{A}_0`$. There is no contradiction with unitarity in approximating them as such, provided we consider sufficiently small anomalous couplings so as to remain far from the unitary limit at the energy regimes accessible by the machines. This is consistent with the basic assumption ($`\mathrm{\Lambda }\sqrt{\widehat{s}}`$) which allows for the effective Lagrangian parametrisation of the TGC vertex keeping only the lowest dimensions: it is sufficient to assume the form factor behaviour commences above the observable scale so as to regulate the distributions before the unitarity limit.
There is also a region on the extreme left side of the plots in Figure 32 (although not indicated) which is excluded by direct experimental searches. This is the region where physics is believed to be well described by the SM.
Experimentally it is desirable to report confidence limits as a function of $`\mathrm{\Lambda }_{\text{FF}}`$. A result using $`\mathrm{\Lambda }_{\text{FF}}=\mathrm{\Lambda }_{\text{machine}}`$ should be included (so long as $`\mathrm{\Lambda }_{\text{machine}}`$ lies below the unitarity limit) as it is motivated by machine parameters and provides a reasonable point of reference for comparisons between different experiments. Other scales (particularly those of theoretical interest) should not be neglected<sup>15</sup><sup>15</sup>15 It should be noted that particularly for small choices of $`\mathrm{\Lambda }_{\text{FF}}`$, a change in the analysis strategy may be necessary to increase sensitivity to the relevant regions of $`\widehat{s}`$. .
#### 5.4.4 Measuring form factors
For a machine of sufficient luminosity such as the LHC, it is possible to measure the energy dependence of anomalous couplings<sup>16</sup><sup>16</sup>16 The suggestion of making such a measurement is not new but has received little attention in the literature. by grouping the data into bins of invariant mass and extracting constant anomalous couplings within these restricted domains. Such a measurement does not carry any assumptions about the form factor (until a fit to a given parametrisation is performed). It is a viable method for measuring form factors, but due to the restricted number of events in each bin, will not produce competitive limits. The method is best employed in the case where non-zero anomalous couplings have been observed.
The method is illustrated in Figure 33 for the case of the $`W\gamma `$ channel with $`We\nu _e,\mu \nu _\mu `$ assuming nature provides an anomalous $`\lambda _0^\gamma =0.025`$ coupling described by an $`n=2`$ dipole form factor with $`\mathrm{\Lambda }_{\text{FF}}=`$2 TeV. Three years of high luminosity (300 fb<sup>-1</sup>) LHC events generated at NLO are binned according to the reconstructed $`M_{\text{inv}}(W\gamma )`$. The corresponding points derived using the generated (unobservable) $`M_{\text{inv}}(W\gamma )`$ are superimposed for comparison. Bin widths (denoted by arrows along the x-axis) are chosen so as to ensure sufficient data in each $`M_{\text{inv}}(W\gamma )`$ domain. A measurement of the anomalous coupling (assumed constant) is performed within each domain using a binned maximum likelihood fit to the $`p_T(\gamma )`$ distribution. No detector simulation has been applied and the specific choice of cuts is unimportant for this illustration. The results of the likelihood fits are plotted as a function of $`M_{\text{inv}}(W\gamma )`$ and a fit to an $`n=2`$ dipole form factor is performed. With this simple illustration, the bare coupling and form factor scale are reconstructed as $`\lambda _0^\gamma =0.029`$ and $`\mathrm{\Lambda }_{\text{FF}}=1.67`$ TeV. Sensitivity to the anomalous coupling increases in the larger invariant mass domains, reflecting the $`\widehat{s}`$ growth of the $`\lambda _0^\gamma `$ coupling (indeed the measurement in the first bin is consistent with zero). Systematic effects related to the fit method (such as the non-uniform distribution of events within the bins) have not been accounted for in this illustration.
### 5.5 Partonic simulation tools for di-boson production
Several Monte Carlo programs for hadronic di-boson event simulation are in common use. General purpose programs such as PYTHIA evaluate the matrix element at leading order (LO) with no spin correlations for boson decay products. Limited or no anomalous couplings are included. In the past decade, programs have been implemented to calculate di-boson production with leptonic decays to next-to-leading order (NLO) in QCD. The diagrams contributing to $`O(\alpha _s)`$ are: the squared Born (LO) graphs, the interference of the Born with the virtual one-loop graphs, and the squared real emission graphs.
The NLO generators by Baur, Han, and Ohnemus (BHO) have been available for several years. They employ the phase space slicing method and the calculation is performed in the narrow width approximation for the leptonically decaying gauge-bosons. Non-standard TGC couplings are included. Spin correlations in the leptonic decays are included everywhere except in the virtual contribution. The authors expect a negligible overall effect from neglecting the spin correlations in the virtual corrections as compared to the uncertainty from parton distribution functions and the choice of factorisation scale. More recently Dixon, Kunszt, and Signer (DKS) have implemented a program with full lepton decay spin correlations (helicity amplitudes are presented in ). The subtraction method is employed in the narrow width approximation including non-standard TGC couplings. A third Monte Carlo program, `MCFM`, by Campbell and Ellis exists. It does not assume the narrow width approximation and includes singly resonant diagrams but does not allow for non-standard TGC couplings. The effects of these improvements in `MCFM` are largest in off-resonant regions - such as near di-boson production thresholds. The regions are of importance to studies of SM backgrounds to new physics but contribute negligibly to the cross section in TGC studies for typical choices of kinematic cuts .
A common feature of the NLO generators is the inability to produce unweighted events. Both the phase space slicing and subtraction methods produce events for which the weight may be either positive or negative - thus it is only the integrated cross section over a region of phase space (i.e. histogram bin) which is physical. This makes traditional Monte Carlo techniques for unweighting events (such as hit-and-miss) difficult to apply, and we are aware of no universally satisfactory technique for producing unweighted events using the NLO generators<sup>17</sup><sup>17</sup>17 One method involves reweighting events from a LO generator using a “look-up table” constructed at NLO. . Computationally this can render analyses very slow, since a large fraction of CPU time can be spent processing events with near-vanishing cross sections.
#### 5.5.1 Comparison of NLO particle level generators
In this section, we present a comparison of the predictions from the BHO and DKS generators, for which no published consistency check exists, restricting ourselves to $`W^+Z`$ and $`WW`$ production for simplicity. The DKS and `MCFM` packages have been found to be in good agreement .
The comparison is performed at LHC energy ($`14`$ TeV $`pp`$ collisions) using CTEQ4M structure functions<sup>18</sup><sup>18</sup>18 The choice of parton distribution function has an $`𝒪(5\%)`$ effect on the cross section.. Input parameters are taken as $`\alpha _{EM}=\frac{1}{128}`$, $`\mathrm{sin}^2\theta _W=0.23`$, $`\alpha _s(M_Z)=0.116`$, $`M_W=80.396`$ GeV, $`M_Z=91.187`$ GeV, factorisation scale $`Q^2=M_W^2`$, and Cabibbo angle $`\mathrm{cos}\theta _\mathrm{C}=0.975`$ with no 3rd generation mixing. Branching ratios are taken as $`BR(Zl^+l^{})`$ = 3.36%, $`BR(W^\pm l^\pm \nu )`$ = 10.8%. The $`b`$ quark contribution to parton distributions has been taken as zero ($`b\overline{b}W^+W^{}`$ contributes $`𝒪(2\%)`$ at LHC ). Kinematic cuts motivated by TGC analyses are chosen. The transverse momentum of all leptons must exceed 25 GeV and the rapidity of all leptons must be less than 3. Missing transverse momentum must be greater than 25 GeV. A jet is defined when the transverse momentum of a parton exceeds 30 GeV in the pseudorapidity interval $`|\eta |<3`$.
For $`W^+Z`$ production, the transverse momentum distribution of the $`Z`$ boson $`p_T(Z)`$, the distribution of rapidity separation between the $`W^+`$ decay lepton and the $`Z`$ boson $`y(l)y(Z)`$, and total cross section are compared at LO, inclusive NLO, and NLO with a jet veto. Branching ratios to $`e,\mu `$-type leptons are applied. For $`WW`$ production, the transverse momentum distribution of the lepton pair from the $`W^\pm `$ decays $`|\stackrel{}{p_T}(e^{})+\stackrel{}{p_T}(e^+)|`$, the distribution of rapidity separation between the $`W`$ decay leptons $`y(e^{})y(e^+)`$, the angle between the $`W`$ decay leptons in the transverse plane $`\mathrm{cos}\mathrm{\Phi }(e^{},e^+)`$, and the total cross section are compared at LO, inclusive NLO, and NLO with a jet veto. Branching ratios to one lepton flavour are applied.
The cross section results are presented in Table 21 and the distributions in Figure 34. Consistency between generators is at the 1% level for $`WZ`$ production and 3-4% level for $`WW`$ production. Qualitative agreement is observed in the distribution shapes.
#### 5.5.2 Effects of NLO corrections
NLO corrections in hadronic di-boson production are large at LHC energies, particularly in the region of high transverse momentum and small rapidity separation (see Figure 34) which is the same region of maximum sensitivity to anomalous TGCs. The corrections can amount to more than an order of magnitude. The high quark-gluon luminosity at the LHC and a logarithmic enhancement at high transverse momentum in the $`qg`$ and $`\overline{q}g`$ real emissions subprocesses are primarily responsible . In the channels which exhibit radiation zero behaviour (i.e. $`W\gamma `$ and $`WZ`$ ), the Born contribution is suppressed and NLO corrections are even larger . Since the $`O(\alpha _s)`$ subprocesses responsible for the enhancement at large transverse momentum do not involve TGCs, the overall effect of NLO corrections is a spoiling of sensitivity to anomalous TGCs.
##### Jet veto
Distributions obtained by vetoing hard jets in the central rapidity region for one possible choice of jet definition ($`p_T(\text{jet})>30`$ GeV, $`|\eta (\text{jet})|<3`$) are shown in Figure 34. The jet veto is effective in recovering the qualitative shape of the LO distributions including the approximate radiation zero in $`WZ`$ production (Figure 34, bottom left). The jet veto serves to recover anomalous TGC sensitivity which is otherwise lost when introducing NLO corrections. A 10-30% improvement in anomalous TGC coupling sensitivity limits in $`WZ`$ production can be achieved when a jet veto is applied as compared to the inclusive NLO case. These limits are often close to those obtained at LO. In general results derived at LO can be considered approximate zero jet results and their conclusions remain interesting. A jet veto also reduces the scale dependence of NLO results .
### 5.6 Determination of TGCs
At the LHC the measurement of TGCs will benefit from both the large statistics and the high centre-of-mass energy. The large available statistics will allow the use of multi-dimensional distributions to increase the sensitivity to the TGCs.
This section discusses the experimental observables sensitive to TGCs and describes the analysis methods employed to measure the TGCs.
#### 5.6.1 Experimental observables
The experimental sensitivity to the TGCs comes from the increase of the production cross section and the modification of differential distributions with non-standard TGCs. The sensitivity is enhanced at high centre-of-mass energies of the hard scattering process, more significantly for $`\lambda `$-type TGCs than for $`\kappa `$-type TGCs in the case of $`W\gamma `$ and $`WZ`$ production. As an example, the increase in the number of events with large di-boson invariant masses is a clear signature of non-standard TGCs as illustrated in Figure 35, where the invariant mass of the hard scattering is shown for $`W\gamma `$ events, simulated with a parametric description of the ATLAS detector, for the Standard Model and non-standard TGCs. A form factor of 10 TeV was used.
For the event generation employing non-standard values of the TGCs, leading order (LO) as well as next to leading order (NLO) calculations have been used (see Section 5.5). Limits on the TGCs can be obtained from event counting in the high invariant mass region. The disadvantage of such an approach alone is that the behaviour of the cross section as function of the TGCs makes it difficult to disentangle the contributions from different TGCs and even their sign (with respect to SM). It is therefore advantageous to combine it with information from angular distributions of the bosons and possibly their decay angles; this improves the sensitivity and improves the separation of contributions from different non-standard TGCs.
In general it is possible experimentally to reconstruct up to four (six) angular variables in the di-boson rest-frame describing an $`W\gamma `$ or $`Z\gamma `$ ($`WZ`$) event:
* Boson production angles, $`\mathrm{\Theta }`$ and $`\mathrm{\Phi }`$, of the di-boson system with respect to the beam-axis in the di-boson rest-frame.
* Decay angles of bosons, $`\theta _{1(2)}^{}`$ and $`\varphi _{1(2)}^{}`$, in the rest-frame of the decaying bosons.
The azimuthal boson production angle, $`\mathrm{\Phi }`$, has no sensitivity to the TGCs. In case of $`W\gamma /WZ`$, $`\mathrm{\Theta }`$ is the most sensitive kinematical variable. The enhanced sensitivity to the TGCs in $`WV`$ production is due to the vanishing of helicity amplitudes in the Standard Model prediction at $`\mathrm{cos}\mathrm{\Theta }1/3`$, affecting the small $`|\eta |`$ region . Non-standard TGCs may partially eliminate the radiation zero, although the zero radiation prediction is less significant when including NLO corrections . In $`Z\gamma `$ production, no radiation amplitude zero is present.
In contrast, the sensitivity to the TGCs from the decay angles is weak; the decay angles primarily serve as projectors of different helicity components, enhancing the sensitivity of other variables.
In the study presented here, several experimentally derived observables and combinations thereof have been studied to assess the possible sensitivity to the TGCs. For both ($`W\gamma `$, $`WZ`$) and ($`Z\gamma `$, $`ZZ`$) events the observables are very similar; for $`WZ`$, the $`Z`$ takes the role of the $`\gamma `$. The actual behaviour of the observables as function of the couplings and the energy is different between the processes, due to the different masses of the involved bosons.
One observable, the transverse momentum, $`p_T`$, of the $`\gamma `$ or $`Z`$ (depending on the di-boson process), which has traditionally been used at hadron colliders, has sensitivity from a combination of high mass event counting and the $`\mathrm{\Theta }`$ angular distribution. Figure 36 shows the enhancement of di-boson production cross section for large values of the photon transverse momentum in presence of non-standard couplings.
The distribution of $`p_T^{\gamma ,Z}`$ assuming an integrated luminosity of 30 fb<sup>-1</sup> is shown in Figure 37 for $`W\gamma `$ and $`WZ`$ events, simulated with a parametric detector simulation program, for the Standard Model and non-standard TGCs. The enhancement for non-standard TGCs at high $`p_T^{\gamma ,Z}`$ is clearly visible and, furthermore, the qualitative behaviour is the same for different TGCs.
For the statistics expected at the LHC, even after 3 years running at low luminosity, one may enhance the experimental sensitivity further by separating the different types of information in multi-dimensional distributions. For $`W\gamma `$ and $`WZ`$ di-boson production, two sets of variables have been studied (and the equivalent set for $`WZ`$): $`(m_{W\gamma },|\eta _\gamma ^{}|)`$, and $`(p_T^\gamma ,\theta ^{})`$, where $`|\eta _\gamma ^{}|`$ is the rapidity of $`\gamma `$ with respect to the beam direction in the $`W\gamma `$ system (equivalent to $`\mathrm{\Theta }`$), and $`\theta ^{}`$ is the polar decay angle of the charged lepton in the $`W`$ rest-frame. Both sets consist of one variable sensitive to the energy behaviour and one sensitive to the angular information. For $`|\eta _\gamma ^{}|`$ and $`\theta ^{}`$, a complete reconstruction of the $`W`$ is necessary. The momentum of the $`W`$ can be reconstructed by using the $`W`$ mass as a constraint and assuming that the missing transverse energy is carried away by the neutrino. This leads to a two-fold ambiguity in the reconstruction. Alternatively, $`|\eta _\gamma ^{}|`$, may be approximated by the rapidity difference between the lepton from the $`W`$ and the $`\gamma `$. Distributions of $`|\eta _\gamma ^{}|`$ and $`\theta ^{}`$ are shown in Figure 38, for both the standard model expectation and different non-standard TGCs. The high sensitivity to the TGCs from $`|\eta _\gamma ^{}|`$ is due to the characteristic “zero radiation” gap. In contrast, the sensitivity to the TGCs from the decay polar angle, $`\theta ^{}`$, is weak.
#### 5.6.2 Analysis techniques for TGC determination
Depending on the available statistics and the dimensionality of the experimental distributions, different extraction techniques can be used in the determination of the TGCs.
One approach employed in this study determines the couplings by a binned maximum-likelihood fit to distributions of the observables, combined with the total cross section information. The likelihood function is constructed by comparing the fitted histogram with a reference histogram using Poisson probabilities. The reference distributions can be obtained for different values of the couplings by reweighting Monte Carlo events at generator level or equivalently using several Monte Carlo event samples generated for different values of the TGCs.
Although the expected number of events at the LHC will allow binning in two dimensions, a general multidimensional binned fit using all the TGC sensitive information will not be possible. In the latter case, an unbinned maximum likelihood fit to the observed information can be used, where the probability distribution functions can be constructed by Monte Carlo techniques. In the case of many dimensions, this approach can be time-consuming, but it may be advantageously combined with the reweighting technique. The information from the absolute prediction of the cross section can be included by the so-called “extended maximum likelihood” method .
### 5.7 Sensitivities at LHC
Sensitivity limits have been derived for the triple gauge-couplings $`WW\gamma `$ (ATLAS, CMS), $`WWZ`$ (ATLAS) and $`ZZ\gamma `$ (CMS). The analysis techniques used by ATLAS and CMS are described in Section 5.6. The ATLAS studies assume an integrated luminosity of $`𝑑t=30\text{ fb}\text{-1}`$, corresponding to three years of LHC low luminosity operation. CMS assumes $`100\text{ fb}\text{-1}`$, which is the expectation for one year of LHC high luminosity running.
CMS has performed its studies for a range of different form factor scales $`\mathrm{\Lambda }_{FF}`$, as motivated in Section 5.4. The plots in Figure 39 show the expected 95% CL limits on the anomalous $`WW\gamma `$ and $`ZZ\gamma `$ coupling parameters together with the corresponding unitarity limits. Only the displayed coupling is considered to deviate from the Standard Model. The points where the experimental curves turn asymptotic with respect to $`\mathrm{\Lambda }_{FF}`$ \- or are crossed by the unitarity limit - give an indication on the range of form factor scales accessible by the experiments. While the current Tevatron measurements probe the triple gauge-couplings up to form factors of $`\mathrm{\Lambda }_{FF}=0.75`$ TeV and around 2 TeV for $`ZZ\gamma `$ and ($`WW\gamma ,WWZ`$), respectively , the LHC experiments will be able to study far smaller structures with scales up to 10 TeV, assuming an integrated luminosity of $`100\text{ fb}\text{-1}`$.
Multi-dimensional fits where several couplings are allowed to vary have also been performed . Here, the sensitivity limits extracted from the log likelihood curves form an ellipse for a particular confidence level. Figure 40 shows the typical $`WW\gamma `$ sensitivity contours in the two-dimensional CP-conserving $`(\kappa \times \lambda )`$ coupling space for a form factor scale of $`10`$ TeV.
Table 22 summarises the sensitivity limits obtained by ATLAS and CMS as reported in . In addition, ATLAS has performed a fit using the complete generator level phase space information . The results for this ideal case show that, as the high energy tails of the $`p_\gamma ^T`$ distributions exhibit a very strong sensitivity to the $`\lambda `$-like anomalous couplings, the additional information does not improve the limits on this type of couplings considerably. However, the $`\kappa `$-type couplings may profit from a more sophisticated data analysis.
From the numbers in Table 22, we expect an improvement in sensitivity by up to two (four) orders of magnitude for anomalous $`WW\gamma /WWZ`$ ($`ZZ\gamma `$) couplings, with respect to the current Tevatron limits. The strong increase in sensitivity is due to the pronounced high $`\widehat{s}`$ enhancement at the LHC, most prominently for $`ZZ\gamma `$ (see Section 5.4.2). A smaller choice of the form factor scale would cut off this enhancement and diminish the sensitivity considerably, as shown in the lower plots in Figure 39.
### 5.8 Backgrounds to $`𝑾𝜸`$
The $`W\gamma `$ signal has a very small cross section, compared to $`W+`$jet production for example, and can contain a significant amount of background. The dominant background to the $`W\gamma `$ signal is from $`W`$+jet production where the jet is misidentified as a photon, resulting in a fake signal. Radiative $`W`$ decay also contributes when the electron from the $`W`$ decay radiates a photon, and both $`t\overline{t}\gamma `$ and $`b\overline{b}\gamma `$ quark-gluon fusion processes can also produce a fake signal contributing to the background. $`Z\gamma `$ production and $`W`$($`\tau \nu `$)$`\gamma `$ also make a small contribution to the backgrounds.
Previous studies have shown that the $`W\gamma `$ signal will be observable at the LHC provided that the backgrounds can be suppressed. All the backgrounds were generated with PYTHIA 5.7 in conjunction with the CMSJET fast detector simulation for the CMS experiment.
#### 5.8.1 $`W+`$ jet and $`Wl\nu \gamma `$ backgrounds
The dominant background to the process $`ppW(e\nu )\gamma `$ arises from $`W+`$jet events where the jet decays electromagnetically and is reconstructed in the calorimeter as a photon. The probability for the jet to fluctuate into an isolated electromagnetic shower is small, but the large number of jets above 10 GeV in the $`W`$ sample guarantees that some jets will look identical to photons. Even if the jet is not misidentified as a photon, it is possible for a radiative decay of the $`W`$ to produce the same signature as the signal. If the lepton from the $`W`$ decay radiates a photon, an event signature of $`\gamma ,l,\nu `$ may be observed. Cuts must therefore be applied to reduce this background.
##### $`𝑾\mathbf{+}`$jet
Figure 41 shows the $`p_T`$($`\gamma `$) spectrum for misidentified photon from the $`W+`$jet background and the real photon from the $`W\gamma `$ signal. A photon isolation cut has been applied to both data sets. A rejection power of nearly 7 can be obtained with an efficiency loss of less than 5$`\%`$, by using an isolation area of $`\mathrm{\Delta }R`$ = 0.25 and a $`p_T`$ threshold of 2 GeV . A greater rejection power with a much smaller efficiency loss is available at low luminosity. Therefore an event is selected if the photon meets the isolation criteria and if it is within $`\eta `$ = $`\pm `$2.5. The isolation cut clearly makes it possible to observe the signal, especially at high $`p_T`$, however a cut at $`p_T`$($`\gamma `$) = 100 GeV further reduces the background. This would not harm the sensitivity to anomalous couplings greatly as the anomalies only manifest themselves at high $`p_T`$.
##### Radiative $`𝑾`$
One method of reducing the background of radiative $`W`$ decays is to make a cut on the invariant mass of the $`\gamma l\nu `$ system. For the $`W\gamma `$ signal, $`M(\gamma l\nu )`$ is always larger than $`M_W`$ if finite $`W`$ width effects are ignored.
However, the $`M(\gamma l\nu )`$ cannot be determined unambiguously as the four-momentum of the neutrino is unknown: even if the transverse momentum is correctly determined from the missing momentum in the event, there is no measurement of the missing longitudinal momentum. Therefore the cluster transverse mass, or minimum invariant mass, may be used instead . The transverse mass is independent of the longitudinal momenta of the parent particle and its decay products.
For $`W\gamma l\nu `$ the cluster transverse mass sharply peaks at $`M_W`$ and drops rapidly above the $`W`$ mass. Thus $`\gamma l\nu `$ events originating from $`W\gamma `$ production and radiative $`W`$ decays can be distinguished if $`M_T(\gamma l\nu )`$ is cut slightly above $`M_W`$ . Hence a cut at $`M_T(\gamma l\nu )>`$ 90 GeV should take into account the finite width of the $`W`$ whilst not significantly affecting the signal.
The $`W\gamma `$ signal produces the lepton and photon almost back-to-back. Ensuring that they are well separated will further reduce the radiative $`W`$ background. This can be done using the quantity $`\mathrm{\Delta }R=\sqrt{(\mathrm{\Delta }\varphi ^2+\mathrm{\Delta }\eta ^2)}`$. Leading order analysis of the signal and radiative background enabled a study of the optimum value of $`\mathrm{\Delta }R`$ to use for separation. Typically a cut at $`\mathrm{\Delta }R>0.5`$ is used to ensure separation, but increasing the separation to $`\mathrm{\Delta }R>0.7`$ makes little difference to the signal whilst greatly reducing the background.
In order to suppress the radiative $`W`$ background events, cuts of $`\mathrm{\Delta }R(\gamma ,l)>0.7`$ and $`M_T(\gamma l\nu )>90`$ GeV are used.
#### 5.8.2 Quark-Gluon fusion background
Quark-gluon fusion is important at the LHC because the rate is extremely high. There are lots of available gluons in the proton at relatively high $`x`$, and because the $`WW\gamma `$ reaction is suppressed in some regions of phase space.
##### $`𝒃\overline{𝒃}𝜸`$
At the LHC 10<sup>12</sup> $`b\overline{b}`$ events are expected for a years running at high luminosity. Although the $`b\overline{b}\gamma `$ events are not kinematically similar to the signal, the expected number of events is so large that the background will be a problem unless it is reduced by cuts.
The $`b\overline{b}\gamma `$ background was generated using the processes: $`q\overline{q}g\gamma `$, and $`q\overline{q}Z\gamma `$. Events were generated from $`\widehat{p}_T`$ = 500 GeV with a cross section of 1.055 pb. This parton-level requirement was for computational efficiency as only the very highest $`p_T`$ events contribute to the background. A cut on missing $`p_T`$ can be made at 50 GeV in order to reduce the $`b\overline{b}\gamma `$ background.
##### $`𝒕\overline{𝒕}𝜸`$
Since the $`M_t>M_W+M_b`$, $`t\overline{t}`$ events represent an irreducible background to $`W\gamma `$ pair production. $`t\overline{t}\gamma `$ production is a copious source of high $`p_T`$ photons in association with hard leptons and without cuts has a cross section, $`\sigma `$ 300 pb, of at least 3 orders of magnitude more than the $`W\gamma `$ signal . The subsequent decay of top quarks into a $`W`$ boson and a $`b`$ quark and also the $`W`$ decay into a $`f\overline{f}`$ pair provide the same event signature as the $`W\gamma `$ signal. Therefore, due to the very large top quark production cross section at LHC energies, the process $`ppt\overline{t}\gamma W\gamma +X`$ represents a potentially significant background.
Events were generated by the process $`q\overline{q}g\gamma `$ and looking for $`t\overline{t}`$ production. This method is very inefficient, 4 million events were generated and 489 $`t\overline{t}\gamma `$ events were produced, with 10 events passing all of the cuts. The $`t\overline{t}\gamma `$ events were generated from $`\widehat{p}_T`$ = 500 GeV (for the same reasons as $`b\overline{b}`$), with a cross section of 1.049 pb. The large cross section means that although only a few events pass the cuts, this background is a potential problem.
Studies for the SSC showed that the background can be reduced to a manageable level by requiring the photon to be isolated from the hadrons in the event, and by imposing a jet veto (i.e. by considering the exclusive reaction $`ppW\gamma +0`$ jets).
Since the top quark decays predominantly into a $`Wb`$ final state, $`t\overline{t}\gamma `$ events are characterised by a large hadronic activity which frequently results in one or several high-$`p_T`$ jets. If the second $`W`$ boson decays hadronically, up to four jets are possible. This observation suggests that the $`t\overline{t}\gamma `$ background may be suppressed by vetoing high-$`p_T`$ jets. Such a “zero jets” requirement has been demonstrated to be very useful in reducing the size of the NLO QCD corrections in $`ppW\gamma +X`$ at SSC energies . If the second $`W`$ in the $`t\overline{t}\gamma `$ events decays hadronically, the number of jets in $`ppt\overline{t}\gamma W\gamma +X`$ is generally larger than for leptonic $`W`$ decays, and the jet veto is more efficient.
Unfortunately the jet veto also drastically reduces the number of signal events. Only 10% of the signal survives the jet veto cut alone and only 4% survive all the cuts and the jet veto. This suggests that an alternative method for reducing this background needs to be found for the LHC.
ATLAS studied the possibility of exploiting the number of jets in the $`t\overline{t}\gamma `$ events by imposing a cut on the second jet in the event. The $`W\gamma `$ signal will not have a 2nd jet, or if it does, it is a misidentified jet and will be of very low $`p_T`$. The $`t\overline{t}\gamma `$ events will have up to four high $`p_T`$ jets in each event. By cutting all events where the $`p_T`$ of the second jet is greater that 25 GeV, the majority of the $`t\overline{t}\gamma `$ events will be eliminated without greatly affecting the signal.
#### 5.8.3 $`Z\gamma `$ background
There is a small background to $`e\nu \gamma `$ that comes from $`Z(ee)\gamma `$ events in which one of the electrons gives rise to significant missing energy (generally by entering a gap in the detector). As CMS is hermetic and the crystals of the ECAL are off-pointing with respect to the interaction point, this background is very small. ATLAS calculate this background to be $``$ 25 times smaller than the signal before any cuts are imposed. Thus the $`Z\gamma `$ background is assumed to be negligible.
#### 5.8.4 $`W(\tau \nu )\gamma `$ background
The final background to $`ppW(e\nu ,\mu \nu )\gamma `$ is $`ppW(\tau \nu )\gamma `$ where the $`\tau `$ lepton decays into an electron or muon. The background is very small because the decay of the tau lepton results in electrons or muons with significantly reduced $`p_T`$ and the kinematical threshold for an electron is 25 GeV. Previous studies at Fermilab have shown this background to be negligible .
#### 5.8.5 Summary of backgrounds
Table 23 shows a list of all the cuts proposed to reduce the backgrounds to the $`W\gamma `$ signal. Having chosen each cut to reduce an individual background, it is important to understand how each cut effects both the signal and the other backgrounds.
Table 24 shows the efficiency of the individual cuts on the signal and the backgrounds. The $`W+`$jet and radiative $`W`$ backgrounds are treated together.
#### 5.8.6 Conclusion
The backgrounds to the $`W\gamma `$ signal have been studied and cuts have been made in order to reduce the backgrounds to at least an order of magnitude less than the signal for $`p_T(\gamma )>200`$ GeV. The $`W+`$jet and radiative $`W`$ backgrounds have been well studied and understood and the cuts made reduce these significantly. The quark-gluon fusion backgrounds are not so well understood in this work since a less than optimal generator for $`t\overline{t}\gamma `$ was used. However, the cuts studied for this channel work well for the low statistic samples presented here. Further study of this background would be interesting.
Backgrounds to $`WZ`$ production have been studied briefly and are similar, within statistical errors, to those in the $`W\gamma `$ channel presented here.
## 6 VECTOR-BOSON FUSION AND SCATTERING <sup>19</sup><sup>19</sup>19Section coordinators: Z. Kunszt, R. Mazini, D. Rainwater
### 6.1 Searching for $`𝑽𝑽\mathbf{}𝑯\mathbf{}𝝉𝝉`$
#### 6.1.1 Introduction
The search for the Higgs boson and, hence, for the origin of electroweak symmetry breaking and fermion mass generation, remains one of the premier tasks of present and future high energy physics experiments. Fits to precision electroweak (EW) data have for some time suggested a relatively small Higgs boson mass, of order 100 GeV , hence we have studied an intermediate-mass Higgs, with mass in the $`110150`$ GeV range, beyond the reach of LEP at CERN and perhaps of the Fermilab Tevatron. Observation of the $`H\tau \tau `$ decay channel in weak boson fusion events at the Large Hadron Collider (LHC) is quite promising, both in the Standard Model (SM) and Minimal Supersymmetric Standard Model (MSSM). This channel has lower QCD backgrounds compared to the dominant $`Hb\overline{b}`$ mode, thus offering the best prospects for a direct measurement of a $`Hf\overline{f}`$ coupling.
At the LHC, despite the fact that the cross section for Higgs production by weak-boson fusion is significantly lower than that from gluon fusion (by almost one order of magnitude), it has the advantage of additional information in the event other than the decay products’ transverse momentum and their invariant mass resonance: namely, the observable quark jets. Thus one can exploit techniques like forward jet tagging to reduce the backgrounds. Another advantage is the different colour structure of the signal vs the background. Additional soft jet activity (minijets) in the central region, which occurs much more frequently for the colour-exchange processes of the QCD backgrounds , are suppressed via a central jet veto.
We have performed first analyses of intermediate-mass SM $`H\tau \tau `$ and of the main physics and reducible backgrounds at the LHC, considering separately the decay modes $`\tau \tau h^\pm l^{}/p_T,e^\pm \mu ^{}/p_T`$. These modes demonstrate the feasibility of Higgs boson detection in this channel with modest luminosity . We demonstrated that forward jet tagging, $`\tau `$ identification and reconstruction criteria alone yield a signal-to-background ($`S/B`$) ratio of approximately 1/1 or better. Additional large background suppression factors can be obtained with the minijet veto, achieving final $`S/B`$ ratios as good as 6/1, depending on the Higgs mass.
In the MSSM, strategies to identify the structure of the Higgs sector are much less clear. For large $`\mathrm{tan}\beta `$, the light neutral Higgs bosons may couple much more strongly to the $`T_3=1/2`$ members of the weak isospin doublets than its SM analogue. As a result, the total width can increase significantly compared to a SM Higgs of the same mass. This comes at the expense of the branching ratio $`BR(h\gamma \gamma )`$, the cleanest Higgs discovery mode, possibly rendering it unobservable over much of MSSM parameter space and forcing consideration of other observational channels. Instead, since $`BR(h\tau \tau )`$ is enhanced slightly, we have examined the $`\tau `$ mode as an alternative .
#### 6.1.2 Simulations of signal and backgrounds
The analyses used full tree-level matrix elements for the weak boson fusion Higgs signal and the various backgrounds. Extra minijet activity was simulated by adding the emission of one extra parton to the basic signal and background processes, with the soft singularities regulated via a truncated shower approximation (TSA) .
We simulated $`pp`$ collisions at the LHC, $`\sqrt{s}=14`$ TeV. For all QCD effects, the running of the strong-coupling constant was evaluated at one-loop order, with $`\alpha _s(M_Z)=0.118`$. We employed CTEQ4L parton distribution functions throughout. The factorisation scale was chosen as $`\mu _f=`$ min($`p_T`$) of the defined jets, and the renormalisation scale $`\mu _r`$ was fixed by $`(\alpha _s)^n=_{i=1}^n\alpha _s(p_{T_i})`$. Detector effects were considered by including Gaussian smearing for partons and leptons according to ATLAS expectations .
At lowest order, the signal is described by two single-Feynman-diagram processes, $`qq`$ $`qq(WW,`$ $`ZZ)`$ $``$ $`qqH`$, i.e. $`WW`$ and $`ZZ`$ fusion where the weak bosons are emitted from the incoming quarks . From a previous study of $`H\gamma \gamma `$ decays in weak boson fusion , we know several features of the signal which we could exploit directly here: the centrally produced Higgs boson tends to yield central decay products (in this case $`\tau ^+\tau ^{}`$), and the two quarks enter the detector at large rapidity compared to the $`\tau `$’s and with transverse momenta in the 20-80 GeV range, thus leading to two observable forward tagging jets.
We considered separately the cases of one $`\tau `$ decaying leptonically ($`e`$,$`\mu `$) and the other decaying hadronically (with a combined branching fraction of $`45\%`$), and both decaying leptonically but with different flavour ($`e\mu `$ or $`\mu e`$, with a combined branching fraction of $`6.3\%`$). Our analyses critically employed transverse momentum cuts on the charged $`\tau `$-decay products and, hence, some care was taken to ensure realistic momentum distributions. Because of its small mass, we simulated $`\tau `$ decays in the collinear and narrow-width approximations and with decay distributions to $`\pi `$,$`\rho `$,$`a_1`$ , adding the various hadronic decay modes according to their branching ratios. We took into account the anti-correlation of the $`\tau ^\pm `$ polarisations in the decay of the Higgs.
##### Lepton-hadron mode
Positive identification of the hadronic $`\tau ^\pm h^\pm X`$ decay requires severe cuts on the charged hadron isolation. We based our simulations on the possible strategies analysed by Cavalli et al. . Considering hadronic jets of $`E_T>40`$ GeV in the ATLAS detector, they found non-tau rejection factors of 400 or more while true hadronic $`\tau `$ decays are retained with an identification efficiency of $`26\%`$.
Given the $`H`$ decay signature, the main physics background to the $`\tau ^+\tau ^{}jj`$ events of the signal arises from real emission QCD corrections to the Drell-Yan process $`q\overline{q}(Z,\gamma )\tau ^+\tau ^{}`$, dominated by $`t`$-channel gluon exchange. All interference effects between virtual photon and $`Z`$-exchange were included, as was the correlation of $`\tau ^\pm `$ polarisations. The $`Z`$ component dominates, so we call these processes collectively the “QCD $`Zjj`$” background.
An additional physics “EW $`Zjj`$” background arises from $`Z`$ and $`\gamma `$ bremsstrahlung in (anti)quark scattering via $`t`$-channel electroweak boson exchange, with subsequent decay $`Z,\gamma \tau ^+\tau ^{}`$. Naively, this EW background may be thought of as suppressed compared to the analogous QCD process. However, the EW background includes electroweak boson fusion, $`VV\tau ^+\tau ^{}`$, which has a momentum and colour structure identical to the signal and thus cannot easily be suppressed via cuts.
Finally, we considered reducible backgrounds, i.e. any event that can mimic the $`Hjj`$ signature of a hard, isolated lepton and missing $`p_T`$, a hard, narrow $`\tau `$-like jet, and two forward tagging jets. Thus we examined $`W+jets`$, where the $`W`$ decays leptonically ($`e`$,$`\mu `$) and one jet fakes a hadronic $`\tau `$, and $`b\overline{b}+jets`$, where one $`b`$ decays leptonically and either a light quark or $`b`$ jet fakes a hadronic $`\tau `$. We neglected other sources like $`t\overline{t}`$ events which had previously been shown to give substantially smaller backgrounds .
Fluctuations of a parton into a narrow $`\tau `$-like jet are considered with probability $`0.25\%`$ for gluons and light-quark jets and $`0.15\%`$ for $`b`$ jets (which may be considered an upper bound) .
In the case of $`b\overline{b}+jj`$, we simulated the semileptonic decay $`bl\nu c`$ by multiplying the $`b\overline{b}jj`$ cross section by a branching factor of 0.395 and implementing a three-body phase space distribution for the decay momenta to estimate the effects of lepton isolation cuts. We normalised our resulting cross section to reproduce the same factor 100 reduction found in .
##### Dual lepton mode
For the dilepton mode, we consider decay only to $`e,\mu `$ pairs to completely eliminate the backgrounds from real $`Z`$ production decaying directly to $`ee`$ or $`\mu \mu `$. Tau decays were performed in the same manner as in the lepton-hadron channel. We again considered QCD and EW $`Zjj;Z\tau \tau `$ production as the physics backgrounds.
We calculated the primary contributions from reducible backgrounds by considering all significant sources of two $`W`$’s, which decay leptonically to form the signature $`e,\mu `$, and two forward jets. This consists of $`t\overline{t}+jets`$, as well as both QCD and EW $`WWjj`$ production. As with the EW $`Zjj`$ case, EW $`WWjj`$ processes contain an electroweak boson fusion component kinematically similar to the signal, and so cannot be ignored.
We also considered $`b\overline{b}jj`$ production, with each $`b`$ decaying semileptonically simulated by implementing the $`VA`$ decay distributions of the $`b`$-quarks in the collinear limit, and multiplying the resultant cross section by a branching fraction 0.0218 (for the $`e,\mu `$ or $`\mu ,e`$ final states).
Finally, we considered the overlapping contribution from the signal itself in the decay mode $`HWWe\mu /p_T`$, which can be significant above $`M_H130`$ GeV.
#### 6.1.3 Standard Model analysis
The basic acceptance requirements must ensure that the two jets and two $`\tau `$’s are observed inside the detector (within the hadronic and electromagnetic calorimeters, respectively), and are well-separated from each other:
$`p_{T_j}20\mathrm{GeV},|\eta _j|5.0,\mathrm{\Delta }R_{jj}0.7,`$
$`|\eta _\tau |2.5,\mathrm{\Delta }R_{j\tau }0.7.`$ (65)
Tau-tau separation and tau decay product $`p_T`$ requirements are slightly different for the two signatures and are discussed separately below.
The $`Hjj`$ signal is characterised by two forward jets with large invariant mass, and central $`\tau `$ decay products. The QCD backgrounds have a large gluon-initiated component and thus prefer lower invariant tagging jet masses. Also, their $`\tau `$ and $`W`$ decay products tend to be less central. Thus, to reduce the backgrounds to the level of the signal, we required tagging jets with a combination of large invariant mass, far forward rapidity, and high $`p_T`$, as well as $`\tau `$ decay products central with respect to the tagging jets :
$`\eta _{j,min}+0.7<\eta _{\tau _{1,2}}<\eta _{j,max}0.7,\eta _{j_1}\eta _{j_2}<0,`$
$`\mathrm{\Delta }\eta _{tags}=|\eta _{j_1}\eta _{j_2}|4.4,m_{jj}>m_{jj_{min}},`$ (66)
where $`m_{jj_{min}}`$ is chosen slightly differently for the two scenarios, as discussed below.
##### Lepton-hadron mode
Here we required two additional cuts to form the tagging jet signature:
$$p_{T_j}>40,20\mathrm{GeV},\mathrm{\Delta }R_{\tau \tau }0.7.$$
(67)
That is, the $`p_T`$ requirement on the tagging jets is staggered, and as one tau decay is hadronic, it must have a large separation from the leptonic tau.
Triggering the event via the isolated $`\tau `$-decay lepton and identifying the hadronic $`\tau `$ decay as discussed in requires sizable transverse momenta for the observable $`\tau `$ decay products: $`p_{T_{\tau ,lep}}>20\mathrm{GeV}`$ and $`p_{T_{\tau ,had}}>40\mathrm{GeV}`$. It is possible to reconstruct the $`\tau `$-pair invariant mass from the observable $`\tau `$ decay products and the missing transverse momentum vector of the event . The $`\tau `$ mass was neglected and collinear decays assumed, a condition easily satisfied because of the high $`\tau `$ transverse momenta required. The $`\tau `$ momenta were reconstructed from the charged decay products’ $`p_T`$ and missing $`p_T`$ vectors. We imposed a cut on the angle between the $`\tau `$ decay products to satisfy the collinear decay assumption, $`\mathrm{cos}\theta _{lh}>0.9`$, and demanded a physicality condition for the reconstructed $`\tau `$ momenta (unphysical solutions arise from smearing effects); that is, the fractional momentum $`x_\tau `$ a charged decay observables takes from its parent $`\tau `$ cannot be negative. Additionally, the $`x_{\tau _l}`$ distribution of the leptonically decaying $`\tau `$-candidate is softer for real $`\tau `$’s than for the reducible backgrounds, because the charged lepton shares the parent $`\tau `$ energy with two neutrinos. Cuts $`x_{\tau _l}<0.75`$ and $`x_{\tau _h}<1`$ proved very effective in suppressing the reducible backgrounds.
Our Monte Carlo predicted a $`\tau `$-pair mass resolution of 10 GeV or better, so we chose $`\pm 10`$ GeV mass bins for analysing the cross sections. To further reduce the QCD backgrounds, which prefer low invariant masses for the tagging jets, we required $`m_{jj}>1`$ TeV. Additionally, the $`Wj+jj`$ background exhibits a Jacobian peak in its $`m_T`$ distribution ; hence a cut $`m_T(l,/p_T)<30`$ GeV largely eliminates this background.
Finally, to compensate for overall rate loss based on ATLAS and CMS expected detector ID efficiencies, we apply a factor 0.86 to the cross section for each tagging jet, and a factor 0.95 for the charged lepton.
Using all these cuts together, although not in a highly optimised combination, we expect already a signal to background ratio of 2/1 with a signal cross section of 0.4 fb for $`M_H=120`$ GeV.
A probability for vetoing additional central hadronic radiation was obtained by measuring the fraction of events that have additional radiation in the central region, between the tagging jets, with $`p_T`$ above 20 GeV, using the matrix elements for additional parton emission. This minijet veto reduces the signal by about $`15\%`$, but eliminates typically $`70\%`$ of the QCD backgrounds; the EW $`Zjj`$ background is reduced by about $`20\%`$, indicating the presence of both boson bremsstrahlung and weak boson fusion effects. Because the veto probability for QCD backgrounds is found to be process independent, we applied the same value to the $`bb+jj`$ background.
Table 25 summarises the signal and various background cross sections at progressive levels of the cuts, ID efficiencies and minijet veto as described above, for the case $`M_H=120`$ GeV. Table 26 gives the expected numbers of events for 60 fb<sup>-1</sup> integrated luminosity (low luminosity running) at the LHC.
It is possible to isolate a virtually background-free $`qqqqHjj\tau \tau `$ signal at the LHC, leading to a $`5\sigma `$ observation of a SM Higgs boson with a mere 60 fb<sup>-1</sup> of data. The expected purity of the signal is demonstrated in Figure 43 showing the reconstructed $`\tau \tau `$ invariant mass for a SM Higgs of 120 GeV after all cuts, particle ID efficiency factors and a minijet veto have been applied. While the reducible $`Wj+jj`$ and $`b\overline{b}+jj`$ backgrounds are the most complicated and do require further study, they appear to be easily manageable.
##### Dual lepton mode
For this signature, we simulated tau decays as before, but with both decaying to final-state leptons. As this would form a different final state in experiment, to form the basic tagging jet signature we require the cuts of Equations 6.1.3 and 6.1.3 as before, but additionally a minimum separation of the charged leptons somewhat less than for the lepton-hadron scenario, $`\mathrm{\Delta }R_{\tau \tau }0.4`$. To be able to trigger on the leptons, we require them to have minimum transverse momentum $`p_{T_l}>10`$ GeV. In the LHC experiments, this may be slightly higher for electrons and slightly lower for muons, but we do not make the distinction here.
Both the $`t\overline{t}+jets`$ and $`b\overline{b}jj`$ backgrounds are about three orders of magnitude larger than the signal, but the contribution from $`b\overline{b}jj`$ may be reduced by a cut on missing transverse energy, $`/p_T>30`$ GeV, and that from $`t\overline{t}+jets`$ may be severely restricted by vetoing additional jets in the central region between the tagging jets, which even before considering additional gluon radiation (minijets) may come from the decays of central final-state $`b`$-quarks. We veto all events with a central $`b`$ with $`p_T>20`$ GeV. This provides approximately a factor 17 in reduction of the top quark background, which may be substantially improved to even lower $`p_T`$ threshold via a $`b`$-tag, which we cannot simulate.
As the dual lepton final state has a lower overall branching ratio than the lepton-hadron case, we retained more overall rate by making a looser cut on the tagging jet invariant mass, $`m_{jj}>800`$ GeV. This cut was still necessary to reduce the QCD backgrounds.
Our Monte Carlo again predicted an excellent $`\tau `$-pair mass resolution, so we retain the mass binning of $`\pm 10`$ GeV. We also rejected non-tau’s as in the lepton-hadron case, although our exact cut was somewhat differently defined:
$`x_{\tau _1},x_{\tau _2}>0,x_{\tau _1}^2+x_{\tau _2}^2<1.`$
Finally, we found that a cut on the maximal separation of the two charged leptons is very useful in reducing the heavy quark backgrounds: $`\mathrm{\Delta }R_{e\mu }<2.6`$.
Efficiency factors for detection are the same as in the previous case, although with two final-state leptons an extra factor 0.95 was taken into account. A minijet veto was applied as before, although other analyses we have performed suggest the survival probabilities change slightly due to the lower hardness of the event, which is strongly correlated with $`m_{jj}`$ (see Table 27).
Table 27 outlines the cross sections of signal and background for progressive levels of cuts as described above, for the case $`M_H=120`$ GeV. Table 28 gives the expected numbers of events for 60 fb<sup>-1</sup> integrated luminosity (low luminosity running) at the LHC.
Although the dual lepton channel does not appear to be able to achieve quite as high an $`S/B`$ ratio as the lepton-hadron channel, it is still better than 1/1 over much of the mass range of interest, which is also clearly evident in the tau pair invariant mass plot of Figure 44. Furthermore, the independent statistical significance of this channel is as good as that found for the lepton-hadron case.
#### 6.1.4 MSSM analysis
The production of CP even Higgs bosons in weak boson fusion is governed by the $`hWW,HWW`$ couplings, which are suppressed by factors $`\mathrm{sin}(\beta \alpha ),\mathrm{cos}(\beta \alpha )`$, respectively , compared to the SM case. Their branching ratios are modified with slightly more complicated factors. One can simply multiply SM cross section results from our analysis by these factors to determine the observability of $`H\tau \tau `$ in MSSM parameter space. We used a renormalisation group improved next-to-leading order calculation, which allows a light Higgs mass up to $`125`$ GeV, and examined two trilinear term mixing cases, no mixing and maximal mixing .
Varying the pseudoscalar Higgs boson mass $`M_A`$, one finds that $`M_h`$, $`M_H`$ each approach a plateau for the case $`M_A\mathrm{},0`$, respectively. Below $`M_A120`$ GeV, the light Higgs mass will fall off linearly with $`M_A`$, while the heavy Higgs will approach $`M_H125`$ GeV, whereas above $`M_A120`$ GeV, the light Higgs will approach $`M_h125`$ GeV and the heavy Higgs mass will rise linearly with $`M_A`$. The transition region behaviour is very abrupt for large $`\mathrm{tan}\beta `$, such that the plateau state will go to $`125`$ GeV almost immediately, while for small $`\mathrm{tan}\beta `$ the transition is much softer and the plateau state reaches the limiting value via a more gradual asymptotic approach.
With reasonable integrated luminosity and combination of the lepton-hadron and dual-lepton channels, 40 fb<sup>-1</sup> in the worst case, it will be possible to observe at the $`5\sigma `$ level either $`h`$ or $`H`$ decays to $`\tau `$ pairs when they are in their respective plateau region, with the possibility of some overlap in a small region of $`M_A`$, as shown in Figure 45. Very low values of $`\mathrm{tan}\beta `$ would be unobservable, but already excluded by LEP2; there should be considerable overlap between this mode at the LHC and the LEP2 excluded region. Furthermore, a parton shower Monte Carlo with full detector simulation should be able to optimise the analysis so that much less data is required to observe or exclude the MSSM Higgs.
#### 6.1.5 Conclusions
The production of a neutral, CP even Higgs via weak boson fusion and decay $`H\tau \tau `$ at the LHC has been studied for the Standard Model and MSSM, utilising parton level Monte Carlo analyses. Each of the decay channels $`\tau \tau h^\pm l^{}/p_T,e^\pm \mu ^{}/p_T`$ independently allows a $`5\sigma `$ observation of a Standard Model Higgs with an integrated luminosity of about 60 fb<sup>-1</sup> or less, and provides a direct measurement of the $`H\tau \tau `$ coupling. For the MSSM case, a highly significant signal for at least one of the Higgs bosons with reasonable luminosity is possible over the entire physical parameter space which will be left unexplored by LEP2. Only 40 fb<sup>-1</sup> of data is required after combining the two channels. We conclude that this mode provides a no-lose strategy for seeing at least one of the CP even neutral MSSM Higgs bosons.
### 6.2 Searching for $`𝑽𝑽\mathbf{}𝑯\mathbf{}𝑾𝑾`$
In the previous section, vector-boson fusion forming a Higgs which then decays to two $`\tau `$’s was identified as a valuable process by which to find a Higgs boson in the mass range 110 to 150 GeV. Rainwater and Zeppenfeld have shown that a heavier Higgs in the range 130 to 200 GeV could be found by looking for the process $`VVHWWe^\pm \mu ^{}\overline{)}p_T`$ . As for the lighter Higgs, the forward jet tagging is a powerful tool for removing background ($`W`$ pairs, $`t\overline{t}`$ and $`Z\tau \tau `$ accompanied by jets). This approach appears more promising than the a search for an inclusive $`HWWe^\pm \mu ^{}\overline{)}p_T`$ signal, yielding a significant result with $`5`$ fb<sup>-1</sup>.
Work has started in the context of the Workshop to investigate this with fast detector simulation, but has not yet been completed.
### 6.3 The strongly interacting symmetry breaking sector
One possible scenario for the spontaneous breaking of the electroweak (EW) symmetry is a strongly interacting symmetry breaking sector (SBS), which generically is formed by new particles with strong interactions at the TeV scale. This sector should provide a global $`SU(2)_L\times SU(2)_R`$ spontaneous symmetry breaking down to the custodial $`SU(2)_{L+R}`$ subgroup, thus triggering the Standard Model spontaneous breaking from the $`SU(2)_L\times U(1)_Y`$ gauge-symmetry down to $`U(1)_{\mathrm{em}}`$. This is the minimal symmetry pattern ensuring that $`\rho 1+O(g^2)`$.
By assuming that the new states appear at the TeV scale, we are only left, at low energies, with the three massless Goldstone Bosons (GB) associated to the global symmetry breaking. We will refer to this scenario as the minimal strongly interacting symmetry breaking sector (MSISBS). In this case, the low-energy EW interactions can be well described with the Electroweak Chiral Lagrangian (EChL) , which is an $`SU(2)\times U(1)`$ gauge-invariant effective field theory that couples the GB to the gauge-bosons and fermions, without any further assumptions than those just described. The EChL, inspired in Chiral Perturbation Theory , is organised as a derivative (momentum) expansion, with a set of effective operators of increasing dimension. Although the lowest-order Lagrangian is common to all models satisfying the minimal assumptions, at higher orders each effective operator has a coefficient, whose different values will account for different underlying symmetry breaking mechanisms. Within this approach it is possible, not only to calculate at tree level, but to include loops whose divergences will be absorbed in the coefficients of operators of higher dimension, thus yielding finite results order by order in the calculations. The values of these renormalised parameters are expected in the $`10^3`$ to $`10^2`$ range.
As far as physics at the LHC is concerned, the most characteristic feature of a strong SBS is the enhanced production of longitudinal gauge-boson pairs. We will review the EChL amplitudes for these processes. However, the EChL perturbative predictions can only describe EW physics at low energies, well below the mass of the heavy states. Indeed, any amplitude calculated with the EChL is obtained as a truncated series in powers of the external momenta. Hence, it will always violate unitarity bounds at high enough energies. In addition, it cannot reproduce any pole associated to new resonant states. Consequently, in order to apply this formalism to study strong SBS phenomenology at the LHC, we have several ways to proceed:
1. Perform studies strictly within the EChL, but restricted to subprocess energies below 1.5 TeV and to very small chiral parameters.
2. Enlarge the EChL introducing explicitly the heavy resonances of each particular model, but this adds new unknown parameters, namely the mass and the width of each resonance.
3. Follow a more model-independent approach, by unitarising the EChL amplitudes and generating heavy resonances from the information contained in the chiral coefficients.
In the last approach, it is possible to describe the different resonant scenarios with just two chiral parameters. Finally we present a study of the LHC sensitivity reach within this parameter space, using the signal of the cleanest leptonic decays of $`ZZ`$ and $`WZ`$ pairs.
#### 6.3.1 Effective Chiral Lagrangian description of electroweak interactions
The EChL provides a phenomenological description of EW interactions when the SBS is strongly-interacting. The only degrees of freedom at low energies are the GBs associated to the $`SU(2)_L\times SU(2)_RSU(2)_{L+R}`$ global symmetry breaking, which are coupled to the EW gauge and fermion fields in an $`SU(2)_L\times U(1)_L`$ invariant way. Customarily, the GBs, $`\omega ^a`$ with $`a=1,2,3`$, are gathered in an $`SU(2)`$ matrix $`U=\mathrm{exp}\left(i\omega ^a\tau ^a/v\right)`$, where $`\tau ^a`$ are the Pauli matrices and $`v=246\text{GeV}`$. The C and P invariant effective bosonic operators up to dimension four are (see the appendix for other notations)
$`_{\mathrm{EChL}}`$ $`=`$ $`{\displaystyle \frac{v^2}{4}}\mathrm{Tr}(D_\mu U(D^\mu U)^{})+a_0{\displaystyle \frac{g^2v^2}{4}}[\mathrm{Tr}(TV_\nu )]^2+a_1{\displaystyle \frac{igg^{}}{2}}_{\mu \nu }\mathrm{Tr}(T𝒲^{\mu \nu })`$ (68)
$`+`$ $`a_2{\displaystyle \frac{ig^{}}{2}}_{\mu \nu }\mathrm{Tr}(T[V^\mu ,V^\nu ])+a_3g\mathrm{Tr}(𝒲_{\mu \nu }[V^\mu ,V^\nu ])+a_4[\mathrm{Tr}(V_\mu V_\nu )]^2`$
$`+`$ $`a_5[\mathrm{Tr}(V_\mu V^\mu )]^2+a_6\mathrm{Tr}(V_\mu V_\nu )\mathrm{Tr}(TV^\mu )\mathrm{Tr}(TV^\nu )+a_7\mathrm{Tr}(V_\mu V^\mu )[\mathrm{Tr}(TV^\nu )]^2`$
$`+`$ $`a_8{\displaystyle \frac{g^2}{4}}[\mathrm{Tr}(T𝒲_{\mu \nu })]^2+a_9{\displaystyle \frac{g}{2}}\mathrm{Tr}(T𝒲_{\mu \nu })\mathrm{Tr}(T[V^\mu ,V^\nu ])+a_{10}[\mathrm{Tr}(TV_\mu )\mathrm{Tr}(TV_\nu )]^2`$
$`+`$ $`\text{e.o.m. terms }+\text{standard YM terms}`$
where we have defined $`TU\tau ^3U^{}`$ and $`V_\mu (D_\mu U)U^{}`$, as well as
$`D_\mu U`$ $``$ $`_\mu Ug𝒲_\mu U+g^{}U_\mu ,𝒲_\mu {\displaystyle \frac{i}{2}}\stackrel{}{W}_\mu \stackrel{}{\tau },_\mu {\displaystyle \frac{i}{2}}B_\mu \tau ^3,`$
$`𝒲_{\mu \nu }`$ $``$ $`_\mu 𝒲_\nu _\nu 𝒲_\mu g[𝒲_\mu ,𝒲_\nu ],_{\mu \nu }_\mu _\nu _\nu _\mu .`$ (69)
The “e.o.m.” terms refer to operators that can be removed using the equations of motion and the “standard YM terms” are the usual Yang Mills Lagrangian together with the gauge-fixing and Faddeev-Popov terms.
The first operator in Equation 68, which provides the $`W`$ and $`Z`$ masses, has dimension two and has the form of a gauged non-linear sigma model (NL$`\sigma `$M). Note that it is universal, since it only depends on $`v`$ \- that is why its predictions for longitudinal gauge-boson scattering amplitudes are called “Low Energy Theorems”. In contrast, the $`a_i`$ couplings will have different values depending on the underlying theory.
The gauge-boson observables are obtained from $`_{\mathrm{EChL}}`$ as a double expansion in $`p^n/(4\pi v)^n`$, $`p`$ being an external momentum, and in the gauge-couplings $`g`$ and $`g^{}`$. The lowest-order predictions are given by the tree level NL$`\sigma `$M, whereas the next order corrections are obtained with a one-loop calculation using the NL$`\sigma `$M vertices plus the tree level contributions of the other operators. The $`a_i`$ coefficients not only provide a model independent parametrisation of the unknown dynamics, but also some of them are used to absorb all the one-loop NL$`\sigma `$M divergences. This procedure could be carried out to any desired order, adding higher dimensional operators, thus yielding finite results order by order in the expansion.
In principle, the $`a_i`$ values for a particular scenario can be obtained by integrating out the heavy degrees of freedom. In fact, they have been determined for the particular cases of the SM with a heavy Higgs and for technicolor theories in the large $`N_{TC}`$ limit . In both cases, these couplings lie in the range $`10^2`$ to $`10^3`$, with either sign. They all have a constant contribution, but those needed in the renormalisation also have a logarithmic term.
#### 6.3.2 Present bounds on the chiral parameters
Let us now look at the present experimental constraints on the EChL parameters $`a_i`$ from low energy EW data. The best constraints come from the oblique radiative corrections, giving bounds on the $`a_0`$, $`a_1`$ and $`a_8`$ parameters that contribute to the gauge-bosons two-point functions up to order $`q^2`$. The EChL calculation of the $`S`$, $`T`$ and $`U`$ self-energy combinations give
$`S=16\pi \left[a_1(\mu )+\text{EChL loops}(\mu )\right],T={\displaystyle \frac{8\pi }{c_W^2}}\left[a_0(\mu )+\text{EChL loops}(\mu )\right],`$
$`U=16\pi \left[a_8(\mu )+\text{EChL loops}(\mu )\right]`$
Note that the $`a_i`$ have been renormalised to absorb the one-loop divergences from the NL$`\sigma `$M chiral loops, so that $`S`$, $`T`$ and $`U`$ are scale independent. Using the $`a_i`$ values for a heavy Higgs boson , the deviations of EW observables from the SM predictions at a reference value of the Higgs mass $`M_H`$ are
$`\mathrm{\Delta }SSS_{\mathrm{SM}}(M_H)=16\pi \left[a_1(\mu )+{\displaystyle \frac{1}{12}}{\displaystyle \frac{5/6\mathrm{log}M_H^2/\mu ^2}{16\pi ^2}}\right],`$
$`\mathrm{\Delta }TTT_{\mathrm{SM}}(M_H)={\displaystyle \frac{8\pi }{c_W^2}}\left[a_0(\mu ){\displaystyle \frac{3}{8}}{\displaystyle \frac{5/6\mathrm{log}M_H^2/\mu ^2}{16\pi ^2}}\right],\mathrm{\Delta }UUU_{\mathrm{SM}}(M_H)=16\pi a_8.`$
A global fit with $`M_H=300`$ GeV and $`m_t=175`$ GeV to the low energy EW data gives
$`\mathrm{\Delta }S=0.26\pm 0.14,\mathrm{\Delta }T=0.11\pm 0.16,\mathrm{\Delta }U=\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}0.26}\pm 0.24`$
which imply the following bounds for the three chiral couplings
$`a_1(1\text{TeV})=(6.8\pm 2.8)\times 10^3,a_0(1\text{TeV})=(4.3\pm 4.9)\times 10^3,a_8(1\text{TeV})=(4.9\pm 4.7)\times 10^3.`$
Other studies agree with these values . These data already disfavour the SM with a heavy Higgs boson and set strong constraints in models with a dominance of vector resonances (like technicolor). With further assumptions on the underlying SBS dynamics, the latter give a negative contribution to $`a_1`$. However, the precision EW measurements leave room for an strong SBS .
Further constraints come from the three-point functions, whose anomalous electroweak effective couplings were traditionally parametrised in terms of $`g_1^\gamma ,g_1^Z,\kappa _\gamma ,\kappa _Z,\lambda _\gamma `$ and $`\lambda _Z`$. A one-loop EChL calculation of these vertices gives
$`g_1^\gamma 1`$ $`=`$ $`0+\text{EChL loops},g_1^Z1={\displaystyle \frac{g^2}{c_W^2}}a_3+\text{EChL loops}(\mu )`$
$`\kappa _\gamma 1`$ $`=`$ $`g^2(a_2a_3a_1+a_8a_9)+\text{EChL loops},\lambda _\gamma =0`$
$`\kappa _Z1`$ $`=`$ $`g^2(a_8a_3a_9)+g^2(a_1a_2)+\text{EChL loops}(\mu ),\lambda _Z=0`$
There are several analyses that constrain these chiral couplings from LEP and Tevatron data. Ignoring the loops from the NL$`\sigma `$M, we get the following values from present LEP data (the Tevatron precision is comparable) $`\lambda _\gamma =0.037{\scriptscriptstyle \genfrac{}{}{0pt}{}{+0.035}{0.036}}`$,
$`\kappa _\gamma 1`$ $`=`$ $`0.038\genfrac{}{}{0pt}{}{+0.079}{0.075},a_2a_3a_1+a_8a_9=0.088\genfrac{}{}{0pt}{}{+0.184}{0.174},`$
$`g_1^Z1`$ $`=`$ $`0.010\pm 0.033a_3=0.018\pm 0.059.`$
Finally, some indirect bounds on quartic couplings have also been found . These indirect estimates come from loops containing $`a_i`$ vertices, but do not include 2-loop diagrams from the NL$`\sigma `$M. They find bounds on $`a_i`$ for $`i=4,5,6,7,10`$ ranging from $`10^1`$ to $`10^2`$.
In summary, the present data on the oblique EW corrections already sets significant bounds on the $`a_0,a_1`$ and $`a_8`$ chiral parameters, but there is not much sensitivity yet to those chiral parameters that contribute to the three or four-point functions. We will see next how, at the LHC, the situation will improve significantly.
#### 6.3.3 The Effective Chiral description at the LHC
At the next generation of colliders, we will be probing the $`W`$ and $`Z`$ interactions at TeV energies. As long as we are only considering the GBs and no other fundamental fields up to the TeV scale, we expect the self-interactions of longitudinal gauge-bosons, $`V_L`$, to become strong at LHC energies. This can be easily understood since, intuitively, longitudinal gauge-bosons are nothing but the GBs, which interact strongly. This intuitive statement is rigorously given in terms of on-shell amplitudes and is known as the Equivalence Theorem (ET),
$$A(V_L^a,V_L^b,V_L^c\mathrm{}\text{Other fields})A(\omega ^a\omega ^b\omega ^c\mathrm{}\text{Other fields})+O\left(M_W^2/\sqrt{s}\right),$$
(70)
which holds for any spontaneously broken non-Abelian theory. Indeed, it was first derived for the SM . Its usefulness is twofold: it relates the pure SBS fields with the observables, but also the calculations can now be performed in terms of scalars instead of gauge-bosons, at least in the high energy limit $`s>>M_W^2`$. At first sight it may seem that the ET is incompatible with the use of the EChL, since an effective theory is a low energy limit. Nevertheless, the ET can still be applied with the EChL, only at leading order in $`g`$ and $`g^{}`$, if we only consider energies below 1.5 TeV and small chiral parameters .
Hence, in a first approximation, we will simplify the high energy description of the strong SBS by neglecting EW corrections. Thus, due to our assumption that $`SU(2)_{L+R}`$ is preserved in the SBS, only the operators that respect custodial symmetry once the gauge-symmetries are switched off will be relevant in this regime. These are the universal term and the operators with $`a_i`$ couplings for $`i=3,4,5`$.
At the LHC, the two most relevant processes of $`V_LV_L`$ production are the scattering of two longitudinal vector-bosons in fusion reactions and the $`V_L`$ pair production from $`q\overline{q}`$ annihilation. Through the ET, they are identified with GB elastic scattering and $`q\overline{q}\omega \omega `$, respectively. Customarily, GB elastic scattering is described in terms of partial wave amplitudes of definite angular momentum, $`J`$, and isospin, $`I`$, associated to the custodial $`SU(2)_{L+R}`$ group. With the EChL, these partial waves, $`t_{IJ}`$ are obtained as
$$t_{IJ}(s)=t_{IJ}^{(2)}(s)+t_{IJ}^{(4)}(s)+\mathrm{},$$
(71)
where the superscript refers to the corresponding power of momenta. They are given by
$`t_{00}^{(2)}`$ $`=`$ $`{\displaystyle \frac{s}{16\pi v^2}},t_{00}^{(4)}={\displaystyle \frac{s^2}{64\pi v^4}}\left[{\displaystyle \frac{16(11a_5+7a_4)}{3}}+{\displaystyle \frac{101/950\mathrm{log}(s/\mu ^2)/9+4i\pi }{16\pi ^2}}\right],`$
$`t_{11}^{(2)}`$ $`=`$ $`{\displaystyle \frac{s}{96\pi v^2}},t_{11}^{(4)}={\displaystyle \frac{s^2}{96\pi v^4}}\left[4(a_42a_5)+{\displaystyle \frac{1}{16\pi ^2}}\left({\displaystyle \frac{1}{9}}+{\displaystyle \frac{i\pi }{6}}\right)\right],`$
$`t_{20}^{(2)}`$ $`=`$ $`{\displaystyle \frac{s}{32\pi v^2}},t_{20}^{(4)}={\displaystyle \frac{s^2}{64\pi v^4}}\left[{\displaystyle \frac{32(a_5+2a_4)}{3}}+{\displaystyle \frac{273/5420\mathrm{log}(s/\mu ^2)/9+i\pi }{16\pi ^2}}\right].`$ (72)
Note that, within our approximations, the above amplitudes only depend on $`a_4`$ and $`a_5`$. The projection in angular momentum has been defined, from the definite $`I`$ amplitude $`T_I`$, as
$$t_{IJ}=\frac{1}{64\pi }_1^1d(\mathrm{cos}\theta )P_J(\mathrm{cos}\theta )T_I(s,t).$$
(73)
The $`V_LV_L`$ production from $`q\overline{q}`$ annihilation, is very important since vector resonances can also couple to this channel. By means of the ET, we are thus interested in $`q\overline{q}\omega \omega `$. As far as GBs couple to quarks proportionally to their mass, the only relevant contribution comes from the $`s`$-channel annihilation through a vector-boson. In practice, for the $`WZ`$ final state, the $`W\omega z`$ interaction is described as $`gF_V(s)`$, by means of a vector form factor, $`F_V(s)`$, which is obtained from the EChL as
$$F_V(s)=1+F_V^{(2)}(s)+\mathrm{}\text{with}F_V^{(2)}(s)=\frac{s}{(4\pi v)^2}\left[64\pi ^2a_3(\mu )\frac{1}{6}\mathrm{log}\frac{s}{\mu ^2}+\frac{4}{9}+i\frac{\pi }{6}\right]$$
(74)
Let us then review the studies of the LHC sensitivity to the chiral parameters via these two processes.
#### 6.3.4 Non-resonant studies for LHC
The EChL formalism has been applied to study the LHC sensitivity to different non-resonant SBS sectors in . We summarise in Table 29 the results from where the expected number of gold-plated $`ZZ`$ and $`WZ`$ from $`VV`$-fusion and $`q\overline{q}`$-annihilation was calculated for values of the custodial preserving $`a_3,a_4`$ and $`a_5`$ parameters in the 10<sup>-2</sup> to 10<sup>-3</sup> range. Since for values of $`a_4`$ or $`a_55\times 10^3`$ *unitarity violations cannot be ignored at energies beyond* 1.5 TeV, these studies only include events in the region of low invariant mass $`V_LV_L`$ pair, i.e. $`M_{VV}1.5`$ TeV. The rest of kinematical cuts are similar to those given in Equation 81. To illustrate the agreement between these kinds of studies, we give in Table 29 other estimates of the $`a_i`$ bounds attainable at the LHC.
It will be very difficult to detect these non-resonant signals over the continuum background, since they just give small enhancements in the high energy region of the $`M_{VV}`$ and $`p_T`$ distributions. There is a general agreement that, although the present bounds could be significantly improved, with these non-resonant studies, the LHC would be hardly sensitive to values of the chiral parameters down to the $`10^3`$ level. Like-sign $`W^\pm W^\pm `$ production may be better in these channels .
Obviously, these studies do not describe one of the most characteristic features of strong interactions: resonances. Moreover, they are limited to moderate energies due to the unitarity violations mentioned already. These caveats can be overcome by means of unitarisation procedures which we explain next.
#### 6.3.5 Unitarisation and resonances in the SBS
In terms of the partial waves defined in Equation 72, the elastic $`V_LV_L`$ scattering unitarity condition, (basically, the Optical Theorem) for physical values of $`s`$, is
$$\text{Im}t_{IJ}(s)=t_{IJ}(s)^2\text{Im}\frac{1}{t_{IJ}(s)}=1,t_{IJ}(s)=\frac{1}{\text{Re}t_{IJ}^1(s)i}.$$
(75)
Hence we only have to use the EChL to approximate
$$\text{Re}t_{IJ}^1=(t_{IJ}^{(2)})^1[1\text{Re}t_{IJ}^{(4)}/t_{IJ}^{(2)}+\mathrm{}].$$
(76)
But since the EChL amplitudes satisfy elastic unitarity perturbatively, i.e.
$$\text{Im}t_{IJ}^{(4)}(s)=t_{IJ}^{(2)}(s)^2\frac{\text{Im}t_{IJ}^{(4)}(s)}{t_{IJ}^{(2)}(s)^2}=1,$$
(77)
we find
$$t_{IJ}(s)=\frac{t_{IJ}^{(2)}}{1t_{IJ}^{(4)}/t_{IJ}^{(2)}}$$
(78)
This is the $`O(p^4)`$ Inverse Amplitude Method (IAM), which has given remarkable results describing meson interactions, which have a symmetry breaking pattern almost identical to our present case . Note that it respects strict elastic unitarity, while keeping the correct EChL low energy expansion. Furthermore, the extension of Equation 78 to the complex plane can be justified using dispersion theory . In particular, it has the proper analytical structure and, eventually, poles in the second Riemann sheet for certain $`a_4`$ and $`a_5`$ values, that can be interpreted as resonances. Thus, EChL+IAM formalism can describe resonances without increasing the number of parameters and respecting chiral symmetry and unitarity.
The EChL+IAM has already been applied to the SBS to study some specific choices of $`a_4`$ and $`a_5`$ that mimic models with vector or scalar resonances. The LHC sensitivity to resonances parametrised with $`a_4`$ and $`a_5`$ was first studied in and , and more recently in . A map of these resonances in the $`(a_4,a_5)`$ space was first obtained in . We show in Figure 46 the vector and scalar neutral resonances expected in the $`(a_4,a_5)`$ parameter space. As far as we expect $`a_4`$ and $`a_5`$ to lie between $`10^2`$ and $`10^3`$, we scan only that range. Furthermore, the poles of the IAM amplitudes will give us the positions and widths of the resonances. Note that, from Equation 72 within our approximations, the $`I=J=1`$ and $`I=J=0`$ channels only depend on the $`a_42a_5`$ and $`7a_4+11a_5`$ combinations, respectively. Thus the straight lines that keep these combinations constant have the same physics in the corresponding channel. We give several examples in the tables within the figure. The fact that each IAM amplitude depends only on one combination of $`a_i`$ implies that their mass and width are related by the KSFR relation . In addition, we locate five points that we will use later as illustrative examples. The white area means that no resonances or saturation of unitarity is reached below $`4\pi v3\text{TeV}`$, which we expect to be the region of applicability for our approach.
We do not give results for the $`I=2,J=0`$ channel since we do not expect any heavy resonance with our minimal assumptions. Intuitively this occurs because the $`I=2,J=0`$ channel is repulsive.
The general resonance spectrum of the MSISBS is gathered in the last plot of Figure 46 . Depending on $`a_4`$ and $`a_5`$, we find one scalar resonance ($`S`$), one vector resonance ($`V`$), two resonances ($`S,V`$), a resonance and a doubly charged wide saturation effect ($`W_2`$) or even no resonances below 3 TeV (white area). For illustration, we have included points for some simple and familiar scenarios: minimal technicolor models with 3 and 5 technicolors ($`TC3`$ and $`TC5`$), and the heavy Higgs SM case, with a tree level mass of 1000 and 1200 GeV ($`H1000`$ and $`H1200`$). The black region is excluded by the constraints on the $`I=2,J=0`$ wave . In the dark “Light Resonances” areas (lighter than 700 GeV), our results should be interpreted cautiously. Outside these areas, we estimate that the predictions of Figure 46 are reliable within $``$ 20% .
Once we have the general spectrum, our aim is to study to what extent the LHC is sensitive to different resonant scenarios via $`V_LV_L`$ production. For that purpose, we cannot forget the unitarisation of $`q\overline{q}V_LV_L`$, since we expect the final state to re-scatter strongly, in particular when there is a resonance in the $`I=J=1`$ elastic channel. This effect can be parametrised in terms of a vector form factor, $`F_V`$. Again, the $`F_V`$ obtained from the EChL does not satisfy exactly its unitarity condition
$$\text{Im}F_V(s)=F_V(s)t_{11}^{}(s),$$
(79)
which implies that the phases of $`F_V`$ and $`t_{11}`$ should be the same (Watson’s Final State Theorem). Moreover, the poles of $`F_V`$ should be those of $`t`$. Hence, we can relate the combination of $`a_i`$ that appears in the perturbative expansion of $`F_V`$ (Equation 74) with $`a_42a_5`$. All in all, it is possible to unitarise $`F_V`$ using only the $`t_{11}`$ EChL result, as follows :
$$F_V\frac{1}{1t_{11}^{(4)}/t_{11}^{(2)}}.$$
(80)
In summary, $`F_V`$ is determined just by $`a_42a_5`$, and we can still use the map of resonances in Figure 46.
#### 6.3.6 Study of the LHC sensitivity to the resonance spectrum of the strong SBS
We will restrict the study to $`ZZ`$ and $`WZ`$ production, assuming that their gold-plated decays, $`ZZ4l`$ and $`WZl\nu ll`$ (with $`l=e,\mu `$) can be identified and reconstructed with a 100% efficiency. We do not consider like-sign $`W^\pm W^\pm `$ production, since, as we have seen, we do not expect $`I=2`$ resonances.
To evaluate $`VV`$ fusion processes, we use the leading-order Effective-$`W`$ Approximation (EWA) . Non-fusion diagrams are not included since they are expected to be small in our kinematic region. We also use the CTEQ4 parton distribution functions at $`Q^2=M_W^2`$ for $`VV`$ fusion and at $`Q^2=s`$ for $`q\overline{q}`$ annihilation and $`gg`$ fusion, with $`\sqrt{s}`$ being the centre of mass energy of the parton pair. More detail can be found in .
Since we do not consider final $`W`$ and $`Z`$ decays, the cuts are set directly on the gauge-boson variables. A first criterion to enhance the strong $`V_LV_L`$ signal over the background is to require high invariant mass $`M_{VV}`$ and small rapidities. We have applied the following set of minimal cuts:
$$500\mathrm{GeV}M_{V_1V_2}10\mathrm{TeV},|y_{\mathrm{lab}}(V_1)|,|y_{\mathrm{lab}}(V_2)|2.5,p_T(V_1),p_T(V_2)200\mathrm{GeV},$$
(81)
which are also required by our approximations, mainly by the ET. An additional invariant mass cut around each resonance will be imposed later.
The $`ZZ`$ production signal occurs through the $`W_L^+W_L^{}Z_LZ_L`$ and $`Z_LZ_LZ_LZ_L`$ fusion processes. In addition, we have included the following backgrounds
$`q\overline{q}ZZ,(61\%),W^+W^{}ZZ,(18\%),ggZZ,(21\%)`$
where we also give their relative contribution to the total background with the minimal cuts. The continuum from $`q\overline{q}`$ annihilation has only tree level SM formulae, which is probably too optimistic since the NLO QCD corrections can enhance significantly the tree level cross sections. The second background is calculated in the SM at tree level, with at least one transverse weak boson. Finally, the one-loop $`ggZZ`$ amplitude has been taken from .
For $`W^\pm Z`$ final states, two processes contribute to the signal: $`W_L^\pm Z_LW_L^\pm Z_L`$ and $`q\overline{q}^{}W_L^\pm Z_L`$, whereas the backgrounds, calculated at tree level within the SM, are
$`W^\pm ZW^\pm Z,(18\%),\gamma ZW^\pm Z,(15\%),q\overline{q}^{}W^\pm Z,(67\%).`$
The $`W^\pm ZW^\pm Z`$ amplitudes have at least one transverse boson and exclude the Higgs contribution. In the $`q\overline{q}^{}W^\pm Z`$ background, we have excluded the amplitude with a $`V_LV_L`$ pair, which is part of the signal. The QCD corrections to $`q\overline{q}^{}`$ annihilation would give an enhancement in both the signal and the background, so we expect that they will not modify considerably our estimates of the statistical significance of vector resonance searches. We have not studied the $`t\overline{t}`$ background since it can be efficiently suppressed after imposing kinematic constraints and isolation cuts to high-$`p_T`$ leptons .
For illustrative purposes, let us first concentrate on the five representative points given in Figure 46. Points 1, 3 and 4 represent models containing a $`J=I=1`$ resonance with masses in the range 900-2000 GeV. Point 5 represents a model with a scalar resonance with mass 730 GeV and a width of 140 GeV. Finally, point 2 represents both a scalar and a vector resonance. The $`M_{VV}`$ distributions for these five models are shown in Figure 47, where we have plotted the signal on top of the background for gold-plated $`ZZ`$ and $`WZ`$ events, assuming an integrated luminosity of 100 fb<sup>-1</sup>. The vector resonances in points 1 to 4 can be seen as peaks in the distribution of final $`WZ`$ pairs. The scalar resonances in points 2 and 5 give small enhancements of $`ZZ`$ pairs. Note that as both $`a_4`$ and $`a_5`$ tend to 0, the resonances become heavier and broader, yielding a less significant signal. It seems evident that it will be much harder to detect scalar than vector resonances. The reasons are that scalars are wider, they are not produced with a significant rate from $`q\overline{q}`$ annihilation, and there is a smaller rate of $`ZZ`$ production from $`VV`$ fusion. Furthermore, the $`ZZ`$ branching ratio to leptons is smaller that that of $`WZ`$.
The contributions to signal and background for $`WZ`$ and $`ZZ`$ production at these representative points are given in Table 30. In order to enhance the signal to background ratio, we have optimised the $`M_{VV}`$ cut, keeping events within approximately one resonance width around the resonance mass (see the second column of these tables). From the $`WZ`$ results, it is clear that the LHC will have a very good sensitivity to light vector resonances, due to the $`q\overline{q}^{}`$-annihilation, which dominates by far the $`VV`$-fusion process. As the vector resonance mass increases, the $`q\overline{q}`$ contribution is damped faster than that of $`VV`$ fusion, and both signals become comparable for vector masses around 2 TeV. Let us remark that, in $`ZZ`$ production, there is only strong interaction signal in $`VV`$ fusion, and therefore to tag forward jets is always convenient in this final state in order to reject non-fusion processes. This is not the case, however, for vector resonance searches since it is mostly due to $`q\overline{q}`$ annihilation. In these tables, we have also estimated the statistical significance, $`\text{Signal}/\sqrt{\text{Bkgd}}`$, assuming integrated luminosities of 100 and 400 fb<sup>-1</sup>. In $`ZZ`$ final states, we also give the significance assuming perfect forward jet-tagging.
Finally, we also show in Figure 47 the regions of the $`(a_4,a_5)`$ space accessible at the LHC, giving 3 and 5$`\sigma `$ contours and assuming integrated luminosities of 100 and 400 fb<sup>-1</sup>. In terms of resonance mass reach limits, we find that with 100 fb<sup>-1</sup>, scalar resonances could be discovered (5$`\sigma `$) in gold-plated $`ZZ`$ events up to a mass of 800 GeV with forward jet-tagging. Vector resonances could be discovered using gold-plated $`WZ`$ events up to a mass of 1800 GeV. These numbers are in good agreement with more realistic studies of particular cases. We can also see that there is a central region in the $`(a_4,a_5)`$ space that does not give significant signals in gold-plated $`ZZ`$ and $`WZ`$ events. This region corresponds to models in which either the resonances are too heavy or there are no resonances in the SBS and the scattering amplitudes are unitarised smoothly. It is a key issue as to whether this type of non-resonant $`V_LV_L`$ signal could be probed at the LHC. It has been argued that doubly-charged $`WW`$ production could be relevant to test this non-resonant region. But non-resonant $`VV`$ distributions would only have slight enhancements at high energies, and a very accurate knowledge of the backgrounds and the detector performance would be necessary in order to establish their existence.
#### 6.3.7 Appendix
### 6.4 Vector-boson scattering
The search for a fundamental scalar particle which would be responsible for electroweak symmetry breaking has so far proven unsuccessful. While the existence of a light Standard Model (SM) Higgs alone would be consistent with all precision electroweak measurements, the well known hierarchy problems make the theory unsatisfactory. The model makes ad hoc assumptions about the shape of the potential, responsible for electroweak symmetry breaking, and provides no explanation for the values of the parameters. Although supersymmetry is an appealing alternative, no indication exists, yet, of its validity. Therefore, in the absence of a low mass Higgs particle, a strongly coupled theory must be considered. The study of electroweak symmetry breaking will require measurements of the production rate of pairs of longitudinal gauge-bosons, since they are the Goldstone bosons of the symmetry breaking process. It will also be essential to search for the presence of resonances which regularise the vector-boson scattering cross-section. Scalar resonances occur in models with a heavy SM Higgs boson, and vector resonances, in charged or neutral channels, are also predicted in dynamical theories, such as technicolor.
In this section, different channels for scattering of high energy gauge-bosons at the LHC are considered These include heavy Higgs production and resonant $`WZ`$ as well as non-resonant $`WZ`$ and $`W^+W^+`$ production in the Chiral Lagrangian model. High mass gauge-boson pair production in a multi-scale technicolor model is also examined. The possibility of making such measurements at the LHC is evaluated.
#### 6.4.1 Heavy Higgs signal
It is now generally believed that a SM Higgs should be light, its mass being bound by requirements of vacuum stability and by the validity of the SM to high scales in perturbative calculations . The parameters of the Higgs used in this study were calculated at tree level. One should note that in NNLO, the resonance saturates . Nevertheless, the search for such a resonance at the LHC can serve as a testing ground for the measurement of the production of high mass longitudinal gauge-boson pairs or for the search of a generic resonance. The $`HWWl\nu jj`$ channel is presented in this section as an example of a typical analysis of a heavy Higgs signal. In fact, $`V_LV_L`$ fusion is also detectable in the case of a heavy Higgs resonance, through the processes $`HZZ`$, up to $`M_H800`$ GeV. Simultaneous detection of a heavy Higgs in other signals would not only confirm the discovery but also provide additional information on the Higgs couplings, which are essential for determining the nature of the resonance.
##### $`𝑯\mathbf{}𝑾𝑾\mathbf{}𝒍𝝂𝒋𝒋`$
In the vector-boson fusion process of Higgs production, $`qqqqH`$, the rate for this channel is sufficient to be observed at low luminosity with a very distinctive signature :
* A high-$`p_T`$ central lepton ($`|\eta _l|<`$2).
* A large $`E_T^{miss}`$.
* Two high-$`p_T`$ jets from the $`Wjj`$ decay in the central region and close-by in space ($`\mathrm{\Delta }R0.4`$) arising from the large boost of the $`W`$ boson.
* Two tag jets in the forward regions ($`|\eta _j|>2`$).
* No extra jet in the central region (central jet veto).
The main backgrounds are:
* $`W`$+jet which gives the largest contribution but also suffers from significant theoretical uncertainties due to higher-order corrections .
* $`t\overline{t}l\nu bjj\overline{b}`$, with the presence of a real $`Wjj`$ decay, but also additional hadronic activity from the $`b`$-jets in the central region.
* $`WWl\nu jj`$ continuum production, which has a much lower rate but is irreducible in the central region.
In addition to central jet veto and forward tag jets cuts, other cuts (high-$`p_T`$ cuts) have been used to optimise the statistical significance of the signal. They are:
* Lepton cuts: $`p_T^l`$, $`E_T^{miss}>`$ 100 GeV, $`p_T^{Wl\nu }>`$ 350 GeV.
* Jet cuts: two jets reconstructed within $`\mathrm{\Delta }R=0.2`$ with $`p_T>`$ 50 GeV and $`p_T^{Wjj}>`$ 350 GeV.
* $`W`$ mass window: $`m_{jj}=m_W\pm 2\sigma `$, where $`\sigma `$ is the resolution on $`m_{jj}`$.
Table 32 shows the number of events resulting from this selection, for an integrated luminosity of 30 fb<sup>-1</sup>, for $`M_H=`$ 1 TeV and $`M_H=`$ 800 GeV as evaluated with the ATLAS fast simulation program (ATLFAST, ). A significant signal remains above background. Variation of the $`E_{tag}`$ cut provides the possibility to compare the shape and cross section of the resonance production to the expected parameters of the Higgs signal (see Figure 48).
The $`HZZll\nu \nu `$ and $`HZZlljj`$ channels in ATLAS have also been studied over most of the mass range from 300 Gev to 1 TeV. It has been shown that forward jet tagging ($`2<|\eta _j|<5`$), is a powerful method for rejecting background and selecting $`qqqqH`$ production, i.e. the vector-boson fusion process.
#### 6.4.2 Strong vector-boson scattering
##### Chiral Lagrangian model
In the Chiral Lagrangian model , the form of the Lagrangian is only constrained by symmetry considerations which are common to any strong electroweak symmetry breaking sector. Differences among underlying theories appear through the values of the parameters of the Chiral Lagrangian. Within the chiral approach, the low-energy Lagrangian is built as an expansion in derivatives of the Goldstone boson fields. There is only one possible term with two derivatives which respects $`SU(2)_{L+R}`$ symmetry:
$$^{(2)}=\frac{v^2}{4}\mathrm{Tr}(D_\mu UD^\mu U^{})$$
where $`D_\mu U`$ = $`_\mu U`$ \- $`W_\mu U+UB_\mu `$, $`W_\mu =ig\sigma ^aW_\mu ^a/2`$, $`B_\mu =ig\sigma ^3B_\mu /2`$.
The dependence on the different models appears at next order through two phenomenological parameters $`L_1`$ and $`L_2`$:
$$^{(4)}=L_1(\mathrm{Tr}(D_\mu UD^\mu U^{}))^2+L_2(\mathrm{Tr}(D_\mu UD^\nu U^{}))^2$$
The $`SU(2)_{L+R}`$ symmetry allows us to define a weak isospin $`I`$. The $`W_LW_L`$ scattering can then be written in terms of isospin amplitudes, exactly as in low energy hadron physics. We assign isospin indices as follows:
$$W_L^aW_L^bW_L^cW_L^d$$
where $`W_L`$ denotes either $`W_L^\pm `$ or $`Z_L`$, where $`W_L^\pm =(1/\sqrt{2})`$ $`(W_L^1iW_L^2)`$ and $`Z_L=W_L^3`$. The scattering amplitude is given by:
$$(W_L^aW_L^bW_L^cW_L^d)A(s,t,u)\delta ^{ab}\delta ^{cd}+A(t,s,u)\delta ^{ac}\delta ^{bd}+A(u,t,s)\delta ^{ad}\delta ^{bc}$$
where $`a,b,c,d`$ =1,2,3 and $`s,t,u`$ are the usual Mandelstam kinematical variables.
In this approach it is possible to compute the function $`A(s,t,u)`$ in $`𝒪(p^4)`$ :
$`A(s,t,u)`$ $`=`$ $`{\displaystyle \frac{s}{v^2}}+{\displaystyle \frac{1}{4\pi v^4}}(2L_1s^2+L_2(t^2+u^2))`$
$`+{\displaystyle \frac{1}{16\pi ^2v^4}}\left({\displaystyle \frac{t}{6}}(s+2t)\mathrm{log}({\displaystyle \frac{t}{\mu ^2}}){\displaystyle \frac{u}{6}}(s+2u)\mathrm{log}({\displaystyle \frac{u}{\mu ^2}}){\displaystyle \frac{s^2}{2}}\mathrm{log}({\displaystyle \frac{s}{\mu ^2}})\right)`$
The values of $`L_1`$ and $`L_2`$ depend on the model, but are expected to be in the range $`10^2`$ to $`10^3`$.
The usual Chiral Lagrangian approach does not respect unitarity at high energies. The Inverse Amplitude Method (IAM) , which is based on the assumption that the inverse of the amplitude has the same analytic properties as the amplitude itself, has been very successful at describing low energy hadron scattering. The most interesting feature of this approach is that it allows us to describe different reactions by using only the two parameters $`L_1`$ and $`L_2`$.
In analogy to $`\pi \pi `$ scattering, there are three possible isospin channels $`I`$ = 0,1,2. At low energies, the states of lowest momentum $`J`$ are the most important, and thus only the $`a_{00}`$, $`a_{11}`$ and $`a_{20}`$ partial waves are considered. It is possible to reproduce, with the IAM model, the broad Higgs-like resonance in ($`I,J`$) = (0,0) channel as well as resonant and non-resonant scattering in the channel (1,1) by selecting appropriate values for $`L_1`$ and $`L_2`$. It has been shown that in the ($`I=1,J=1`$) channel there may exist narrow resonances up to 2500 GeV and this scattering only depends on the combination of ($`L_22L_1`$).
##### Resonant $`𝑾_𝑳𝒁_𝑳\mathbf{}𝑾_𝑳𝒁_𝑳`$ channel
As a reference for the IAM model, the process $`W_LZ_LW_LZ_L`$, with $`Zll`$ ($`l=e,\mu `$) and $`Wjj`$ is used . A modified version of PYTHIA 5.7 was used to generate $`V_LV_L`$ scattering processes for each value of $`L_1`$ and $`L_2`$. The simulation was done for two values of ($`L_22L_1`$) = 0.006 and 0.01, which yield $`\sigma \times BR`$ of 1.5 fb and 2.8 fb, with mass peaks at 1.5 TeV and 1.2 TeV respectively.
Irreducible background arises from continuum $`WZ`$ production and the main QCD background is from $`Z`$+jets production with two final state jets faking the $`W`$ decay if their invariant mass is close to $`m_W`$. $`t\overline{t}`$ production is potentially dangerous but is efficiently suppressed by a cut on the invariant mass of leptons from the $`W`$ decay . The following cuts were used for background rejection:
* Two isolated leptons with the same flavour and opposite charges in the region $`|\eta |<2.5`$ and $`p_T>100`$ GeV. Their invariant mass was required to lie in the region $`|m_{ll}m_Z|<6`$ GeV.
* Jets were reconstructed in a cone of width $`\mathrm{\Delta }R=0.2`$. Only two jets with $`p_T>50`$ GeV were allowed in the central region ($`|\eta |<2`$) and $`|m_{jj}m_W|<15`$ GeV was required. Only $`W`$ and $`Z`$ with $`p_T>200`$ GeV were kept.
* In the forward region ($`2<|\eta |<5`$), jets were reconstructed in a cone of width $`\mathrm{\Delta }R=0.5`$ and events were accepted only if jets with $`p_T>30`$ GeV and $`E_{jet}>500`$ GeV were present in each hemisphere.
The expected number of signal and background events after all cuts and for $`=100`$ fb<sup>-1</sup> are presented in Table 33. The mass spectra obtained after all cuts (Figure 49) shows a clear peak with a width of 75 GeV (100 GeV) for the 1.2 TeV (1.5 TeV) resonance and 14 (8) signal events in the window $`|m_{WZ}m_V|<2\sigma `$. The contribution from irreducible backgrounds is negligible and is below 0.05 events inside the mass window. It is clear that such a narrow resonance could be detected easily after a few years of high luminosity.
##### Non-resonant channels
If nature does not provide resonances in $`V_LV_L`$ scattering, the measurement of cross sections at high mass for non-resonant channels becomes the only probe for the mechanism of regularisation of the cross section. It would then be essential to understand very well the magnitude and energy dependence of backgrounds. Those channels can be particularly important since it has been shown that a complementary relationship exits between resonant and non-resonant processes . Both $`W_LZ_L`$ and $`W_LW_L`$ scattering have been studied within the ATLAS framework.
###### $`𝑾_𝑳𝒁_𝑳\mathbf{}𝑾_𝑳𝒁_𝑳`$
The non-resonant $`W_LZ_LW_LZ_L`$ process, with $`Zll`$ and $`Wl\nu `$ ($`l=e,\mu `$), was incorporated in PYTHIA and used with two values of $`L_1`$: 0.003 and 0.01, leading to $`\sigma \times BR=`$ 0.19 fb and 0.11 fb respectively. The main features of the signal are:
* The presence of two high-$`p_T`$ leptons of same flavour and opposite charge in the barrel region, having an invariant mass consistent with the mass of the $`Z`$ boson.
* One additional high-$`p_T`$ lepton in the barrel region.
* Significant missing momentum in the event due to the presence of a neutrino.
* The presence of energetic jets in the forward region.
The main irreducible background, coming from continuum $`WZ`$ production, was generated by PYTHIA with $`\sigma \times BR=`$ 13.5 fb. The main reducible background is the QCD process $`Zt\overline{t}`$ where one of the $`W`$ bosons from a $`t`$-quark decays into a lepton and an anti-neutrino. The value of $`\sigma \times BR`$ of this process is 26.3 fb. A less important contribution comes from $`ZZ`$ production with $`\sigma \times BR=`$ 1.52 fb. These different backgrounds were rejected with a high efficiency by using the following cuts:
* Two isolated leptons of same flavour and opposite charge were required in the central region with $`p_T>30`$ GeV and invariant mass satisfying $`|m_{ll}m_Z|<6`$ GeV. One additional lepton was required.
* A missing momentum of at least 75 GeV.
* At least one jet with $`p_T>40`$ GeV and $`E_{jet}>500`$ GeV should be present in the forward region.
In order to analyse $`WZ`$ scattering in the high-mass region, the transverse mass $`M_T`$
$$M_T^2=\left[\sqrt{M^2(lll)+p_T^2(lll)}+\overline{)}p_T\right]^2\left[\stackrel{}{p}_T(lll)+\overline{)}\stackrel{}{p}_T\right]^2$$
was used. $`M(lll)`$ and $`p_T(lll)`$ are the invariant mass and transverse momentum of the three charged leptons and $`\overline{)}p_T`$ is the missing momentum in the event. The transverse mass $`M_T`$ distribution for the $`W_LZ_L`$ scattering and for $`Zt\overline{t}`$ background, after the application of cuts, is shown in Figure 50. The number of signal and background events with the invariant mass of $`WZ`$ system larger then 600 GeV for an integrated luminosity of $`=500`$ fb<sup>-1</sup> and applying different cuts, are shown in Table 34. The $`ZZ`$ background is not shown since it is effectively removed by the requirement of missing transverse momentum.
###### Like-sign $`𝑾`$ pair production
$`W_L^+W_L^+`$ production has been extensively studied . As possible scenarios for this process by $`W_L^+W_L^+`$ scattering, the following are considered:
* A $`t`$-channel exchange of a Higgs with $`M_H`$ = 1 TeV, ($`W_LW_L`$ only), simulated with PYTHIA with $`\sigma \times BR`$ = 1.33 fb (the same parameters of the resonance as in Section 6.4.1 were used).
* The K-matrix unitarised amplitude $`a_{IJ}^K=\frac{Re(a_{IJ})}{1iRe(a_{IJ})}`$, where $`a_{IJ}`$ is the low-energy theorem amplitude, proportional to $`s`$. This model is constructed to satisfy explicitly elastic unitarity and would yield the maximum expected signal. The $`\sigma \times BR`$ = 1.12 fb.
* A Chiral Lagrangian model, as in the $`WZ`$ resonant channel, with the same parameters: $`L_1`$ = 0, and $`L_2`$ = 0.006 or 0.01, leading to $`\sigma \times `$ BR = 0.484 and 0.379 fb, respectively.
Backgrounds from continuum $`WW`$ bremsstrahlung produce mostly transverse $`W`$’s. Other backgrounds include processes involving non-Higgs exchange, as well as QCD processes of order $`\alpha \alpha _s`$ in amplitude, with gluon exchange and $`W`$ bremsstrahlung from interacting quarks. The effects of $`Wt\overline{t}`$ and $`WZ`$ backgrounds are also considered. The signal was generated with PYTHIA 6.2 and backgrounds were incorporated into PYTHIA from a Monte Carlo generator based on Barger’s work , which takes into account all diagrams. The contribution from electroweak processes not involving the Higgs were estimated by assuming a low-mass Higgs ($`M_H=`$ 100 GeV).
An analysis was performed using the fast ATLAS detector simulation (ATLFAST), with parameters set for high luminosity. The following leptonic cuts were first applied:
1. Two positively charged isolated leptons in the central region ($`p_T>40`$ GeV and $`|\eta |<1.75`$) must be identified. They will satisfy the trigger requirement.
2. The opening angle between the two leptons, in the transverse plane, must satisfy: $`\mathrm{cos}\mathrm{\Delta }\varphi <0.5`$. This cut selects preferentially events with longitudinal $`W`$’s which have high $`p_T`$. The invariant mass of the two leptons was further required to satisfy $`m_{ll}>100`$ GeV. This latter cut eliminates few events in the low $`m_{ll\nu \nu }`$ region.
At the jet level, backgrounds can be reduced by requiring that:
1. No jet having $`p_T>50`$ GeV be present in the central region ($`|\eta |<2`$). This reduces significantly the background from the $`Wt\overline{t}`$ process.
2. Two jets must be present in the forward and backward regions: $`\eta >2`$ and $`\eta <2`$, with energies $`>`$ 300 GeV.
3. A lower $`p_T`$ was required for the forward jets: $`p_T<`$ 150 GeV for the first and $`p_T<`$ 90 GeV for the second.
Figure 51 shows expected mass distribution of the $`ll\nu \nu `$ system, for an integrated cross section of 300 fb<sup>-1</sup>, after all cuts were applied, accounting only for transverse momentum. No correction was made for pile-up effects in jet tagging or central jet veto. If one counts only events with $`m_{ll\nu \nu }>`$ 400 GeV, a significant signal to background ratio is obtained (see Table 35). As expected, the K-matrix scenario gives the highest signal \- this could be observable after a few years of high luminosity running. By contrast, it was shown in Section 6.4.2 that if the $`\rho `$ resonance is itself clearly observable in the resonant channel, then the signal will be very low. The major remaining background, especially at low values of $`m_{ll\nu \nu }`$, is from continuum transverse $`W`$ pairs. Note that only a $`W_L^+W_L^+`$ signal was searched for in this analysis. Combining the results with $`W_L^{}W_L^{}`$ would add approximately one-half to one-third of the signal and backgrounds. The Chiral Lagrangian model, with its parameters leading to a resonance in the $`WZ`$ system, would yield a very weak signal in the $`W^+W^+`$ channel, confirming the complementarity relationship between those two channels .
#### 6.4.3 Technicolor
Technicolor (TC) provides a framework for dynamical electroweak symmetry breaking . It assumes the existence of techni-fermions possessing a technicolor charge and interacting strongly at high scale. Chiral symmetry is broken by techni-quark condensates giving rise to Goldstone bosons, the techni-pions, which are the longitudinal degrees of freedom of the $`W`$ and $`Z`$ gauge-bosons. TC has been extended (extended TC, or ETC) to allow the generation of fermion masses . In order to account for the absence of FCNCs, the coupling constant is required to “walk”, rather than “run”. To achieve a walking $`\alpha _{TC}`$, multi-scale TC models contain several representations of the fundamental family, and lead to the existence of techni-hadron resonances accessible at LHC energies. Such models are constrained by precision electroweak data , but not necessarily excluded . However, the constraints from those data make it unnatural to have a large top quark mass. In top-colour-assisted TC (TC2) models , the top quark arises in large part from a new strong top-colour interaction, which is a separate broken gauge-sector.
The possible observation of TC resonances using the ATLAS detector is described in . In particular, the search for a ($`I`$=1, $`J`$=1) techni-rho resonance, a techni-pion and a techni-omega has been performed. Although certain models, with a given set of parameters, are used as reference, the signals studied can be considered generic in any model which predicts resonances. The model adopted here is that of multi-scale TC , with the TC group $`SU(N_{TC})`$ where $`N_{TC}`$ = 4 and two isotriplets of techni-pions. The longitudinal gauge-boson and the techni-pions mix
$$|\mathrm{\Pi }_T>=\mathrm{sin}\chi |W_L>+\mathrm{cos}\chi |\pi _T>$$
with a mixing angle which has a value $`\mathrm{sin}\chi =1/3`$. The decay constant of the mixed state is $`F_T=F_\pi \mathrm{sin}\chi =82`$ GeV and the charge of the up-type (down-type) techni-fermion is $`Q_U=1`$ ($`Q_D=0`$). This model is incorporated in PYTHIA 6.1. The decay channels of $`\rho _T`$ depend on the assumed masses of the techni-particles. Some mass scenarios have been considered to be representative of what one may expect to probe at the LHC and it is also assumed that the $`\pi _T`$ coupling to the top quark is very small, as may be expected in TC2 models. The following sections present an example showing a typical analysis for extracting TC signals. More channels and an extensive description can be found in .
##### $`𝝆_𝑻^\mathbf{\pm }\mathbf{}𝑾^\mathbf{\pm }𝒁\mathbf{}𝒍^\mathbf{\pm }𝝂𝒍^\mathbf{+}𝒍^{\mathbf{}}`$
This decay could be the cleanest channel for the techni-rho detection and complements the study shown in Section 6.4.2. The good efficiency of the ATLAS and CMS detectors for lepton detection and missing transverse energy measurement will provide good identification of the $`W`$ and $`Z`$ bosons. Table 36 shows the parameters for the various sets of events which were generated. For each set, $`10^4`$ events were generated and the signal was normalised to three years of low luminosity running at the LHC (30 fb<sup>-1</sup>). The branching ratios quoted include a preselection on the transverse mass ($`\widehat{m}>`$ 150, 300, 600 GeV for $`m_{\rho _T^\pm }`$= 220, 500 and 800 GeV respectively).
The only background which needs to be considered is the continuum production of $`WZ`$ gauge-bosons, with $`\sigma =21`$ pb. The cuts which were applied are:
* At least three charged leptons were required (with $`E_T>20`$ GeV for electrons and $`E_T>6`$ GeV for muons), two of which must have the same flavour and opposite charge.
* The invariant mass of the lepton pair with the same flavour and opposite sign should be close to that of the $`Z`$: $`|m_{l^+l^{}}m_Z|<5`$ GeV.
* The longitudinal momentum of the neutrino is calculated (with a 2-fold ambiguity) from the missing transverse energy and the momentum of the unpaired lepton assuming an invariant mass $`m_{l\nu }=m_W`$. Once the $`W`$ and $`Z`$ were reconstructed, their transverse momentum was required to be larger than 40 GeV.
* Only events for which the decay angle with respect to the direction of the $`WZ`$ system ($`\rho _T`$) in its rest frame was $`|\mathrm{cos}\widehat{\theta }|<0.8`$ were accepted.
The significance ($`S/\sqrt{B}`$) of the signal ($`S`$) above the background ($`B`$) is shown in Table 36. The number of signal and background events was counted in mass regions around the $`\rho _T`$ peak: 210 to 240, 460 to 560 and 740 to 870 for $`m_{\rho _T}=`$220, 500 and 800 GeV respectively. No evident signal can be observed for cases (b), (e) and (f) (see Figure 52), principally because the $`\rho _T`$ resonance is too wide.
The Authors would like to thank M. Chanowitz, K. Lane, M. Mangano, J.R. Peláez, S.R. Slabospitsky and P. Savard for their technical help with some Monte Carlo generators and for fruitful discussions.
### 6.5 The degenerate BESS Model at the LHC
It is well known that naïve Dynamical Symmetry Breaking (DSB) models like standard QCD-scaled technicolor generally tend to provide large corrections to electroweak precision observables. New physics effects are naturally small if decoupling holds. In fact in this case the corrections to electroweak observables are power suppressed in the limit in which the masses of the new particles are made large. It is thus a natural question as to whether examples of DSB models with decoupling do exist.
Here we will focus on a scheme of DSB, called degenerate BESS (D-BESS) in which decoupling is naturally satisfied in the low energy limit. The model predicts the existence of two triplets of new resonances corresponding to the gauge-bosons of an additional gauge-symmetry $`SU(2)_LSU(2)_R`$. The global symmetry group of the theory is $`(SU(2)_LSU(2)_R)^3`$ breaking down spontaneously to $`SU(2)_D(SU(2)_LSU(2)_R)`$ and giving rise to nine Goldstone bosons. Six of these give mass to the new gauge-bosons, which turn out to be degenerate. As soon as we perform the gauging of the subgroup $`SU(2)_LU(1)_Y`$, the three remaining Goldstone bosons disappear giving masses to the SM gauge-bosons.
What makes the model so attractive is the fact that, due to the degeneracy of the masses and couplings of the extra gauge-bosons $`(L^\pm ,L_3,R^\pm ,R_3)`$, it decouples, so all the deviations in the low-energy parameters from their SM values are strongly suppressed. Also, the degeneracy is protected by the additional “custodial” symmetry $`(SU(2)_LSU(2)_R)`$. The deviations from the SM predictions come from the mixing of $`(𝐋_\mu ,𝐑_\mu )`$ with the standard gauge-bosons. In order to compare with the experimental data, radiative corrections have to be taken into account. Since the model is an effective parametrisation of a strongly interacting symmetry breaking sector, one has to introduce a UV cut-off $`\mathrm{\Lambda }`$. We neglect the new physics loop corrections and assume for D-BESS the same radiative corrections as for the SM with $`M_H=\mathrm{\Lambda }=1`$ TeV . The 95% CL bounds on the parameter space of the model coming from the precision electroweak data can be expressed by the following approximated relation: $`M`$(TeV)$`2.4g/g^{\prime \prime }`$, where $`M`$ is the common mass of the new resonances, $`g`$ and $`g^{\prime \prime }`$ are the standard $`SU(2)_L`$ and the new strong gauge-couplings respectively. Therefore one has a large allowed region available for the model even for the choice $`M_H=\mathrm{\Lambda }=1`$ TeV - a value highly disfavoured by the fit within the SM . Also, the bounds on the D-BESS model from the direct search for new gauge bosons performed at Tevatron are very loose . This allows the existence of a strong electroweak sector at relatively low energies such that it may be accessible with accelerators designed for the near future. A peculiar feature of this strong electroweak symmetry breaking model is the absence of $`WW`$ enhancement due to the absence of direct couplings of the new resonances to the longitudinal weak gauge-bosons. For this reason, the gold plated channels to consider for discovering $`(𝐋_\mu ,𝐑_\mu )`$ are the fermionic ones.
Here we have considered the production of these new resonances at the LHC for the following configuration $`\sqrt{s}=14`$ TeV and $`=10^{34}`$ cm<sup>-2</sup>sec<sup>-1</sup> and for the electron channel decay (the muon channel was studied in ). The events were generated using PYTHIA Monte Carlo (version 6.136) . Only the Drell-Yan mechanism for production was considered since it turns out to be the dominant one. We have analysed the production of the charged resonances in $`ppL^\pm ,W^\pm e\nu _e`$ ($`R^\pm `$ are completely decoupled) and neutral ones in $`ppL_3,R_3,Z,\gamma e^+e^{}`$. The signal events were compared with the background from SM production. We have performed a rough simulation of the detector, in particular, assuming a $`2\%`$ smearing in the momenta of charged leptons and a resolution $`\mathrm{\Delta }E_T^{miss}=0.6\sqrt{E_T^{miss}}`$ in the missing transverse energy. In the neutral channel, we have assumed an error of $`2\%`$ in the reconstruction of the $`e^+e^{}`$ invariant mass, which includes bremsstrahlung effects . We have considered several choices of the model parameters, in the region allowed by the present bounds, and for each case we have selected cuts to maximise the statistical significance of the signal. In Figure 53 we show the transverse mass distributions for the signal and for the SM background for the case $`M=1`$ TeV (left) and $`M=2`$ TeV (right) and $`g/g^{\prime \prime }=0.1`$. The following cuts have been applied for $`M=1`$ TeV: $`|p_T^e|`$ and $`|p_T^{miss}|>0.3`$ TeV and $`M_T>0.8`$ TeV. The number of signal events per year is 3200, the corresponding background is of 1900 events. The corresponding statistical significance $`S/\sqrt{S+B}`$ for one year of running is 44. For $`M=2`$ TeV, the applied cuts are: $`|p_T^e|`$ and $`|p_T^{miss}|>0.7`$ TeV and $`M_T>1.8`$ TeV, resulting in $`S=108`$, $`B=46`$ and $`S/\sqrt{S+B}=8.7`$.
In Figure 54, we show the results of our simulation for the same choice of the parameters as in Figure 53 for the neutral channel. The following cuts have been applied for $`M=1`$ TeV: $`|p_T^{e^+}|`$ and $`|p_T^e^{}|>0.3`$ TeV and $`M_{e^+e^{}}>0.8`$ TeV. The number of signal events per year is 620, the background is of 1200 events with a corresponding statistical significance of 15. For $`M=2`$ TeV, the cuts are: $`|p_T^{e^+}|`$ and $`|p_T^e^{}|>0.7`$ TeV and $`M_{e^+e^{}}>1.8`$ TeV, resulting in $`S=24`$, $`B=30`$ and $`S/\sqrt{S+B}=3.3`$. It turns out that the cleanest signature is in the neutral channel, but the production rate is lower than for the charged one. Also we observe that, due to the fact that the D-BESS resonances are almost degenerate ($`\mathrm{\Delta }M/M(g/g^{\prime \prime })^2`$), it will be impossible to disentangle $`L_3`$ and $`R_3`$ which both contribute to the peak of the signal in Figure 54.
Our conclusion is that the LHC will be able to discover a strong electroweak resonant sector as described by the degenerate BESS model for masses up to 2 TeV - in some cases with very significant numbers of events. Furthermore, if no deviations from the SM predictions are seen within the statistical and systematic errors, the LHC with $`L=100`$ fb<sup>-1</sup> will put a 95% CL bound $`g/g^{\prime \prime }<0.040.06`$ for $`0.5<M(`$TeV$`)<2`$ .
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# 1 Introduction
## 1 Introduction
One of the exactly soluble model in quantum field theory is the quantum electrodynamics of the massless fermions in 1+1 dimensions, which is known as Schwinger model . The massive Schwinger model, describing the electromagnetic interaction of a massive Dirac field, is no longer exactly soluble, however all non-trivial features of the massless model continue to hold for small fermion mass limit . The Schwinger model may serve as a laboratory to study some important features of particle physics, present also in higher dimensional theories, such as screening and quark confinement which are some of the most important problems in particle physics. For example it has been proposed that the infrared behavior of QCD<sub>4</sub> may be responsible for the confinement of quarks and gluons. But the concept of infrared slavery, i.e. the increase of potential between colored objects with separation, could not be verified to be true using the perturbation methods (because of infrared singularities) and must be studied nonperturbatively. These kinds of calculations can be done in an equivalent two-dimensional model.
As it is well known, the gauge field of the massless Schwinger model can be made massive by the standard Higgs mechanism and the Coulomb force is replaced by a finite range force. Then by introducing two opposite static external charges $`q`$ and $`\overline{q}`$, one can see that the potential tends to some constant for large separation of $`q\overline{q}`$ pairs, reflecting the screening of these charges by the induced vacuum polarization. On the other hand in the massive Schwinger model, a semiclassical analysis reveals a linear $`q\overline{q}`$ potential. In this case, by computing the Wilson loop for widely separated charges, within the framework of Euclidean path integral and mass perturbation (for small masses), one can see that integer external probe charges are completely screened whereas a linearly potential is formed between widely separated non–integer charges .
A particular intriguing and interesting case occurs when the two dimensional surface, on which the model is defined, is a curved space-time. (Similar investigations for pure Yang-Mills theories on arbitrary two dimensional compact Riemann surfaces have been done in several papers, see for example .) These models are useful for better understanding the confinement and screening mechanisms in curved space-time and can be viewed as a first step to study these phenomena in the presence of quantum gravity. Moreover, they may have application in string theory and quantum gravity coupled to nonconformal matter (note that the kinetic term of the gauge field spoils the conformal invariance of the theory).
The Schwinger model has been studied on different non flat surfaces, for example on closed Riemann surfaces of genus $`g2`$ , on torus , and on sphere . Also the Green function of the gauge field of the Schwinger model has been calculated on the Poincare disk in . Moreover, in , the authors have considered a $`D`$–dimensional hyperboloid with negative curvature, embedded in $`(D+1)`$–dimensional Minkowski space, and by considering the behavior of the gauge and matter fields near the boundary, they have chosen the solutions with suitable behavior. In this way, they have used the negative curvature space–time as a regulator for interacting Euclidean quantum field theories. However, the confinement and screening properties of the Schwinger model have not yet been studied on curved space–time. In , where the bosonization procedure of the Schwinger model in curved space has been discussed, it has been mentioned that this model continues to exhibit screening or confinement of the charges associated to the electromagnetic field on conformally flat spaces. As we will show, it is not true at least for the Poincare half plane. In , the authors have argued that as the perimeter and the area in hyperboloid space–times are proportional for large loops, one can not simply distinguish between different phases by only considering the Wilson loop dependence on area or perimeter. As we will show, this is right, i.e. by explicit computation of effective static potential between a quark and antiquark, we show that despite the different behavior of the Wilson loop of the Schwinger model in $`e0`$ and $`e=0`$ (the first has perimeter behavior, while the second has area behavior), both have a common phase structure.
In this paper we want to study the confining behavior of the Schwinger model on the Poincare half plane. This is an interesting case because it can illustrate the effects of the boundary and the metric of the space–time on the confinement feature of the Schwinger model. Other property of the Poincare half plane is that its metric is independent of one of the coordinates, so one can obtain the static potential of the external charges in terms of the spatial geodesic distance.
The paper is organized as follows. In section 2, following the method used in , we obtain an expression for the potential between the external charges by integrating out the fermionic degrees of freedom. We discuss the confining and screening like behaviors of the system and point out the differences of these features with respect to the flat case. We justify our results by calculating the expectation value of the Wilson loop. We also derive the bosonization rules for the Schwinger model on the Poincare half plane. In section 3 we consider the massive Schwinger model. Using the bosonization method and by solving the equations of motion of the gauge and matter fields, we obtain a perturbative expression for the interaction energy of the probe charges.
Note that in this paper we do not consider the nontrivial topologically sectors of the gauge fields.
## 2 Massless Schwinger model on the Poincare half plane and its confining behavior
The Poincare half plane, $`H=\{(x,t),x>0\}`$, is a non–compact Riemann surface equipped with the metric $`ds^2=(dx^2+dt^2)r^2/x^2`$ and the symplectic area form $`\sqrt{g}d^2x=(dxdt)r^2/x^2`$. $`r`$ is a scale parameter of the Poincare plane and is related to the scalar curvature by $`R=2/r^2`$. This space is conformally related to the compact orientable Riemann surface $`\mathrm{\Sigma }_g`$ with genus $`g2`$, $`\mathrm{\Sigma }_g`$=H/G, where G is a discrete subgroup of PSL(2,R) (the isometry group of H). The geodesics of the Poincare half plane are semi–circles centered on the horizontal axis $`t`$ (which we take it as the time axis), and straight lines parallel to the vertical axis $`x`$. The geodesic distance between the points $`(x_1,t_1)`$ and $`(x_2,t_2)`$ on the semi circle is
$$L=r\mathrm{cosh}^1[1+\frac{(x_2x_1)^2+(t_2t_1)^2}{2x_1x_2}].$$
(1)
For the points $`(x_1,t)`$ and $`(x_2,t)`$, $`x_2>x_1`$, situated on the straight line the geodesic distance is given by
$$d=r\mathrm{ln}\frac{x_2}{x_1}.$$
(2)
The Schwinger model is defined by the action
$$S=\sqrt{g}d^2x[i\overline{\psi }\widehat{\gamma }^ae_a^\mu (_\mu ieA_\mu )\psi +\frac{1}{4}g^{\mu \rho }g^{\nu \lambda }F_{\mu \nu }F_{\rho \lambda }],$$
(3)
where $`e`$ is the charge of dynamical fermions, and $`\widehat{\gamma }^a`$ are anti–Hermitian matrices which in terms of Pauli matrices are $`\widehat{\gamma }^0=i\sigma _2`$ and $`\widehat{\gamma }^1=i\sigma _1`$. $`F_{\mu \nu }`$ is defined by $`F_{\mu \nu }=_\mu A_\nu _\nu A_\mu `$, $`g_{\mu \nu }=\eta _{\mu \nu }r^2/x^2`$, $`\eta _{\mu \nu }=\mathrm{diag}(1,1)`$, and $`\sqrt{g}=r^2/x^2`$. The zwei-beins fields ($`e_a^\mu ,e_\mu ^a`$) are defined through
$$g_{\mu \nu }=e_\mu ^ae_\nu ^b\eta _{ab},g^{\mu \nu }=e_a^\mu e_b^\nu \eta ^{ab}.$$
(4)
For the metric $`g_{\mu \nu }=\eta _{\mu \nu }r^2/x^2`$, we obtain
$$e_\mu ^a=\frac{r}{x}\delta _\mu ^a,e_a^\mu =\frac{x}{r}\delta _a^\mu ,$$
$$e^{\mu a}=\frac{x}{r}\eta ^{\mu a},e_{\mu a}=\frac{r}{x}\eta _{\mu a}.$$
(5)
The action (3) is invariant under change of coordinate system, frame rotation $`e_\mu ^a\mathrm{\Lambda }_b^ae_\mu ^b`$, $`\mathrm{\Lambda }`$ SO(2), and local gauge transformation, but it is not conformal invariant since the Maxwell field theory is conformal invariant only in four dimensions.
### 2.1 Bosonization
The classical equation of motion of the field $`A_\mu `$ is
$$\frac{1}{\sqrt{g}}_\nu \sqrt{g}F^{\nu \sigma }=J^\sigma =e\overline{\psi }\widehat{\gamma }^be_b^\sigma \psi ,$$
which yields
$$_\sigma \sqrt{g}\overline{\psi }\widehat{\gamma }^be_b^\sigma \psi =0.$$
Hence
$$\overline{\psi }\widehat{\gamma }^be_b^\sigma \psi =\alpha ϵ^{\sigma \nu }_\nu \mathrm{\Phi },$$
(6)
where $`\alpha `$ is a constant, $`ϵ^{\sigma \nu }=\widehat{ϵ}^{\sigma \nu }/\sqrt{g}`$ and $`\widehat{ϵ}^{01}=\widehat{ϵ}_{01}=1`$, $`\widehat{ϵ}^{10}=1`$. This relation is one of the bosonization rules for massless fermions in a two dimensional (conformally flat) space. In , it has been shown that by performing a fermionic change of variables, $`\psi =\chi /g^{1/8}`$ and $`\overline{\psi }=\overline{\chi }/g^{1/8}`$, the bosonization of the fermionic part of the action (3) is realized in a similar method as in the flat case. On the other hand the bosonization rules on the half plane, R$`{}_{}{}^{+}\times `$R, is the same as the complete plane . Therefore on the Poincare half plane the bosonization rules are
$$i\overline{\psi }\gamma ^\mu _\mu \psi =\frac{1}{2}g^{\mu \nu }_\mu \mathrm{\Phi }_\nu \mathrm{\Phi },$$
$$\overline{\psi }\gamma ^\mu \psi =\frac{i}{\sqrt{\pi }}ϵ^{\mu \nu }_\nu \mathrm{\Phi },$$
(7)
$$\overline{\psi }\psi =\frac{1}{g^{1/4}}\overline{\chi }\chi =\frac{1}{g^{1/4}}\mathrm{\Sigma }\mathrm{cos}(2\sqrt{\pi }\mathrm{\Phi }),$$
in which $`\gamma ^\mu =\widehat{\gamma }^ae_a^\mu `$, and $`\mathrm{\Sigma }`$ is a $`c`$–number which depends on the normal ordering of the composite operator $`\overline{\psi }\psi `$. To determine $`\mathrm{\Sigma }`$, we proceed as .
On one hand, the bosonization of the composite operator $`\overline{\chi }\chi `$ is the same as in the flat case, that is $`\overline{\chi }\chi =\mathrm{\Sigma }N_\mu \mathrm{cos}(2\sqrt{\pi }\mathrm{\Phi })`$, where $`N_\mu `$ is the normal ordering with respect to the mass $`\mu =e/\sqrt{\pi }`$. Hence
$$<\overline{\chi }\chi (\xi _1)\overline{\chi }\chi (\xi _2)>=\mathrm{\Sigma }^2<N_\mu \mathrm{cos}[2\sqrt{\pi }\mathrm{\Phi }(\xi _1)]N_\mu \mathrm{cos}[2\sqrt{\pi }\mathrm{\Phi }(\xi _2)]>$$
$$=\frac{\mathrm{\Sigma }^2r^2}{x_1x_2}\mathrm{cosh}[4\pi D(\xi _1,\xi _2)],$$
(8)
where $`\xi _1=(x_1,t_1)`$, $`\xi _2=(x_2,t_2)`$, are two points on the upper half plane, and $`D(\xi _1,\xi _2)`$ is the bosonic propagator
$$D(\xi _1,\xi _2)=\frac{1}{2\pi }Q_l(1+\frac{2|\xi _1\xi _2|^2}{4x_1x_2}),$$
(9)
computed from the Lagrangian
$$L=\frac{1}{2}g^{\mu \nu }\sqrt{g}_\mu \mathrm{\Phi }_\nu \mathrm{\Phi }+\frac{1}{2}\mu ^2\sqrt{g}\mathrm{\Phi }^2.$$
(10)
$`Q_l`$ is the Legendre function of the second kind and $`l=(1+\sqrt{1+4\mu ^2r^2})/2`$. The appearance of the metric dependent term $`r^2/(x_1x_2)`$ in eq.(8) is related to the renormalization of vertex operators on the curved space–time . In the limit $`\xi _1\xi _2`$, we use the asymptotic behavior of $`Q_l`$ and obtain
$$<\overline{\chi }\chi (\xi _1)\overline{\chi }\chi (\xi _2)>=\frac{\mathrm{\Sigma }^2r^2}{2x_1x_2}\mathrm{exp}[\mathrm{ln}\frac{|\xi _1\xi _2|^2}{4x_1x_2}2\gamma 2\mathrm{\Psi }(l+1)],$$
(11)
where $`\gamma `$ is the Euler constant, and $`\mathrm{\Psi }`$ is the digamma function.
On the other hand, in the limit $`\xi _1\xi _2`$, we have
$$<\overline{\chi }\chi (\xi _1)\overline{\chi }\chi (\xi _2)>=\frac{1}{2\pi ^2|\xi _1\xi _2|^2}.$$
(12)
Note that this relation is the same as one in flat space–time. The reason of this equality lies in the fact that in the limit $`\xi _1\xi _2`$, all the $`A_\mu `$–dependent terms in evaluating $`<\overline{\chi }\chi (\xi _1)\overline{\chi }\chi (\xi _2)>`$ are canceled out , and this calculation reduces to one in the free fermion model, i.e. without gauge field, on a flat Euclidean space–time, described by the action
$$S_{\mathrm{free}}=d^2x(i\overline{\chi }\widehat{\gamma }^a_a\chi ).$$
(13)
Comparing (11) and (12) we obtain
$$\mathrm{\Sigma }=\frac{1}{2\pi r}\mathrm{exp}[\gamma +\mathrm{\Psi }(l+1)],$$
(14)
which differs from the result obtained for the complete flat plane: $`(e/2\pi ^{3/2})\mathrm{exp}(\gamma )`$ . This difference is due to the presence of the curvature which modify the Green function of the gauge fields appeared in the fermionic two–point functions . In the limit $`R0`$ ($`r\mathrm{}`$), using $`lim_x\mathrm{}\mathrm{\Psi }(x)=\mathrm{ln}(\mathrm{x})`$, we obtain the same $`\mathrm{\Sigma }`$ as the flat case.
### 2.2 Confinement: the effective action approach
In order to investigate the confining behavior of the action (3), we will obtain the equation of motion of the gauge field derived from the corresponding effective action. Using (7), the bosonic version of (3) is
$$S=\sqrt{g}d^2x(\frac{1}{2}g^{\mu \nu }_\mu \mathrm{\Phi }_\nu \mathrm{\Phi }\frac{ie}{\sqrt{\pi }}ϵ^{\mu \nu }A_\mu _\nu \mathrm{\Phi }+\frac{1}{4}F^{\mu \nu }F_{\mu \nu }).$$
(15)
Integrating over the bosonic degrees of freedom we arrive at
$$S_{\mathrm{eff}}=(\frac{e^2}{2\pi }\sqrt{g}\frac{F}{\sqrt{g}}\frac{1}{\mathrm{\Delta }}\frac{F}{\sqrt{g}}+\frac{1}{2\sqrt{g}}F^2)d^2x,$$
(16)
in which $`\mathrm{\Delta }=(1/\sqrt{g})_\mu g^{\mu \nu }\sqrt{g}_\nu `$, and $`F=\widehat{ϵ}^{\mu \nu }_\mu A_\nu `$.
As an alternative method, this effective action can be also obtained by integrating out the fermionic degrees of freedom of the action (3). To do this, we should compute the determinant of the Dirac operator $`i\gamma ^\mu D_\mu =i\gamma ^\mu (_\mu ieA_\mu )`$,
$$D:=\mathrm{ln}\frac{D\overline{\psi }D\psi \mathrm{exp}(\sqrt{g}d^2x\overline{\psi }i\gamma ^\mu (_\mu ieA_\mu )\psi )}{D\overline{\psi }D\psi \mathrm{exp}(\sqrt{g}d^2x\overline{\psi }i\gamma ^\mu _\mu \psi )}=\mathrm{ln}\frac{\mathrm{det}i\gamma ^\mu D_\mu }{\mathrm{det}i\gamma ^\mu _\mu }.$$
(17)
Using the one–loop radiative correction of the two–point function of the gauge field and also by considering the requirement of the invariance of the theory under PSL(2,R), it can be shown that
$$D=\frac{e^2}{2\pi }\sqrt{g}d^2x\frac{F}{\sqrt{g}}\frac{1}{\mathrm{\Delta }}\frac{F}{\sqrt{g}}.$$
(18)
Adding the kinetic term of the gauge field, we arrive at (16).
In the gauge $`A_1=0`$ and in the static case $`dA_0/dt=0`$, the effective Lagrangian (16) becomes
$$_{\mathrm{eff}}=\frac{e^2}{2\pi }A_{0}^{}{}_{}{}^{2}+\frac{1}{2}\frac{x^2}{r^2}(\frac{dA_0}{dx})^2.$$
(19)
The above effective Lagrangian density shows that the photon gains a mass equal to $`e/\sqrt{\pi }`$, which can be interpreted as a peculiar two–dimensional version of the Higgs phenomenon.
Now following , if we introduce a static external charge distribution composed of a quark and an anti–quark with charges $`e_1=e^{}`$ and $`e_2=e^{}`$ at points $`\xi _1=(a,t)`$ and $`\xi _2=(b,t)`$, respectively, this Lagrangian becomes
$$=\frac{e^2}{2\pi }A_0^2+\frac{1}{2}\frac{x^2}{r^2}(\frac{dA_0}{dx})^2+J^0A_0\frac{r^2}{x^2},$$
(20)
where
$$J^0=\frac{i}{\sqrt{g}}\underset{n=1}{\overset{2}{}}e_n\delta ^2(\xi \xi _n)𝑑t_n$$
$$=i\frac{x^2}{r^2}e^{}[\delta (xb)\delta (xa)].$$
(21)
In this case, the equation of motion of the field $`A_0`$ is
$$\frac{d}{dx}\frac{x^2}{r^2}\frac{dA_0}{dx}\frac{e^2}{\pi }A_0=ie^{}[\delta (xb)\delta (xa)].$$
(22)
To find $`A_0(x)`$, we note that the Green function of the self adjoint operator $`P=\frac{d}{dx}\frac{x^2}{r^2}\frac{d}{dx}\frac{e^2}{\pi }`$ is
$$G_P(x,x^{})=\frac{r^2}{2l+1}\frac{x_<^l}{x_>^{l+1}},$$
(23)
where $`x_<`$ ($`x_>`$) is the smaller (larger) value of $`x`$ and $`x^{}`$. This Green function is the same as the Green function of the radial part of Poisson operator in spherical coordinates and satisfies the Dirichlet boundary condition at $`x=0`$ (the Poincare half plane has no boundary, but by boundary condition we mean the behavior of the fields near the horizontal axis). In the flat case limit, $`r\mathrm{}`$, $`l`$ leads to $`\mu r`$ and therefore the eq.(23) reduces to
$$\frac{(g(x)g(x^{}))^{1/8}}{2\mu }e^{\mu d(xx^{})},$$
(24)
where $`d`$ is the geodesic distance (2). By setting $`g_{\mu \nu }=\eta _{\mu \nu }`$, eq.(24) leads to the Green function of the flat case, i.e. $`(1/2\mu )e^{\mu d}`$. Using (23), we obtain
$$A_0(x)=ie^{}[G_P(xb)G_P(xa)]=\{\begin{array}{cc}\frac{ie^{}r^2}{2l+1}(\frac{b^l}{x^{l+1}}\frac{a^l}{x^{l+1}}),\hfill & b<x\hfill \\ \frac{ie^{}r^2}{2l+1}(\frac{x^l}{b^{l+1}}\frac{a^l}{x^{l+1}}),\hfill & a<x<b\hfill \\ \frac{ie^{}r^2}{2l+1}(\frac{x^l}{b^{l+1}}\frac{x^l}{a^{l+1}}),\hfill & x<a.\hfill \end{array}$$
(25)
To calculate the quark–antiquark energy, we must note that the Schwinger model on the Poincare half plane can be considered as the analytical continuation of the corresponding model on a Minkowskian space–time described by the metric $`ds^2=(r^2/x^2)(dt^2dx^2)`$. By ignoring the $`i`$ factors in eqs.(21) and (25) and substituting them back into $`^{\mathrm{Min}.}`$, which has the same form as (20), one can obtain the static external charges energy as $`U=^{\mathrm{Mink}.}𝑑x=^{\mathrm{Eucl}.}𝑑x`$ . In this way we find the interaction energy of the external charges as
$$U=\frac{1}{2}J^0A_0\frac{r^2}{x^2}𝑑x=\frac{e^2}{2l+1}\frac{r^2}{2a}(2e^{\frac{d}{r}(l+1)}+e^{\frac{d}{r}}+1).$$
(26)
For a detailed discussion on the relation of the Euclidean action with the Minkowskian static energy, see .
In the flat case, the eq.(22) is replaced by
$$\frac{d^2}{dx^2}A_0\frac{e^2}{\pi }A_0=ie^{}[\delta (xb)\delta (xa)],$$
(27)
which is invariant under translation ($`xx+c,c`$ R), hence the potential is only a function of charge separation, which is a translational invariant quantity. But in our case, (22) is not invariant under scale transformation (dilatation $`x\lambda x`$; $`\lambda `$ R), which leaves the distance $`d`$ invariant, hence the potential depends on both the distance $`d=r\mathrm{ln}(b/a)`$ and the position of the external charges.
For large separation, $`b>>a`$, the potential tends to
$$\underset{d\mathrm{}}{lim}U=\frac{r^2}{2a}\frac{e^2}{2l+1},$$
(28)
which indicates the screening like phenomenon: By fixing the position of one of the charges at an arbitrary point $`x=a`$, and moving the other charge, the potential increases linearly for small separation and tends to a finite value for large $`d`$. But the crucial point is that the geometry of the Poincare half plane is non–trivial, and a model defined in this space–time, may have different behaviors in different regions. For example, while the confinig phase is dominant in a region, the system may be in screening phase in another region. To see this, one must study the behavior of some external charge in this space, as a probe. Now as it is clear from (26), the system is in confining phase near the boundary $`x=0`$: for $`a0`$, we must have $`b=a+O(a^2)`$ in order to have a finite energy for the system, otherwise $`U\mathrm{}`$. This means that in the massless Schwinger model on the Poincare half plane, the confining phase, is dominant near the horizontal axis. This is related to the singularity of the metric at $`x=0`$.
On the flat plane, the Schwinger model is confining in the absence of dynamical fermions: The screening potential
$$U_{\mathrm{flat}}=\frac{e^2}{2\mu }(1e^{\mu |ba|}),$$
(29)
in the limit $`\mu 0`$, becomes $`(e^2/2)|ba|`$ which increases with the relative distance of the charges. In this limit, the effects of the fermionic vacuum polarization is switched off: the screening is replaced by the confining behavior of the system.
But on the Poincare half plane at $`\mu =0`$, that is when dynamical massless fermions are absent, the potential becomes
$$U=\frac{e^2r^2}{2a}(1e^{\frac{d}{r}}),$$
(30)
which has the same confining or screening nature as (26). The dynamical fermions can only decrease the amount of saturated energy. Hence the screening like (or confining) behavior of the Schwinger model depends on the vacuum polarization and the curvature of the space–time.
### 2.3 Confinement: the Wilson loop approach
Now it is interesting to obtain and interpret these results by computing the Wilson loop expectation value. The interaction of an external current density $`j^\mu `$ and the gauge field $`A_\mu `$ is described by the action $`S_{\mathrm{int}.}=\sqrt{g}j^\mu A_\mu d^2x`$. We assume that $`j^\mu `$ is produced by two external charges moving on a loop which is obtained as follows. Two charges $`e^{}`$ and $`e^{}`$ are created at the point $`(x,t)`$ and move apart in (Euclidean) time $`\tau `$ to points $`(a,t+\tau )`$ and $`(b,t+\tau )`$. Then they stay static at their positions for a period of time $`T`$, and after that come together to annihilate. In the limit $`T>>\tau `$, in which we are interested, this Wilson loop becomes a rectangle $`c`$ characterized by $`a,b`$, and $`T`$, on the Poincare half plane. The reason for choosing this kind of Wilson loop is that in the large $`T`$–limit, the expectation value of this Wilson loop becomes proportional to exp\[$`U(d)T`$\] ($`U(d)`$ is the static external charge potential) for time–independent metrics . (See also for the same calculations on a curved space–time.) The interaction term of this process is $`S_{\mathrm{int}}=ie^{}_c𝑑x^\mu A_\mu `$, and the expectation value of the corresponding Wilson loop is
$`<W_c[A]>=`$
$$\frac{DA_\alpha D\mathrm{\Phi }\delta (H[A_\alpha ])\mathrm{det}[\frac{\delta H[A^\lambda ]}{\delta \lambda }]\mathrm{exp}(ie^{}_cA_\mu 𝑑x^\mu )\mathrm{exp}[(\frac{1}{2}(_\mu \mathrm{\Phi })^2\frac{ie}{\sqrt{\pi }}F\mathrm{\Phi }\frac{1}{2}\frac{x^2}{r^2}F^2)d^2x]}{DA_\alpha D\mathrm{\Phi }\delta (H[A_\alpha ])\mathrm{det}[\frac{\delta H[A^\lambda ]}{\delta \lambda }]\mathrm{exp}[(\frac{1}{2}(_\mu \mathrm{\Phi })^2\frac{ie}{\sqrt{\pi }}F\mathrm{\Phi }\frac{1}{2}\frac{x^2}{r^2}F^2)d^2x]}.$$
(31)
$`H[A_\alpha ]=0`$ is the gauge–fixing condition and $`\lambda `$ parameterizes the gauge transformation $`A_\alpha ^\lambda =A_\alpha +_\alpha \lambda `$. One can show that by using the change of variables $`A(F,\eta )`$, $`\eta :=H[A]`$, the Jacobian of this transformation:
$$DA_\alpha =\mathrm{det}^1[\frac{\delta H[A^\lambda ]}{\delta \lambda }]D\eta DF,$$
(32)
cancels precisely against the ghost determinant . Thus
$$<W_c[A]>=\frac{DFD\mathrm{\Phi }\mathrm{exp}[(\frac{1}{2}(_\mu \mathrm{\Phi })^2\frac{1}{2}\frac{x^2}{r^2}F^2\frac{ie}{\sqrt{\pi }}F\mathrm{\Phi })d^2x]\mathrm{exp}(ie^{}_cA_\mu 𝑑x^\mu )}{DFD\mathrm{\Phi }\mathrm{exp}[(\frac{1}{2}(_\mu \mathrm{\Phi })^2\frac{1}{2}\frac{x^2}{r^2}F^2\frac{ie}{\sqrt{\pi }}F\mathrm{\Phi })d^2x]}.$$
(33)
Using the Stokes theorem
$$_cA_\mu 𝑑x^\mu =_D\eta (\xi )F(\xi )d^2x,$$
(34)
where $`c=D`$, and $`\eta (\xi )=\{\begin{array}{cc}1,\hfill & \xi D\hfill \\ 0,\hfill & \xi D.\hfill \end{array}`$, we arrive at
$$<W_c[A]>=\mathrm{exp}[\frac{e^2}{2}\frac{r^2}{x^2}\eta ^2(\xi )d^2\xi \frac{e^2e^2}{2\pi }\eta (\xi )G_W(\xi ,\xi ^{})\eta (\xi ^{})d^2\xi d^2\xi ^{}],$$
(35)
in which the Green function $`G_W(\xi ,\xi ^{})`$ satisfies
$$[\frac{x^2}{r^2}(\frac{d^2}{dx^2}+\frac{d^2}{dt^2})\frac{x^2}{r^2}\mu ^2\frac{x^2}{r^2}]G_W(\xi ,\xi ^{})=\delta ^2(\xi \xi ^{}).$$
(36)
If we insert the Fourier expansion
$$G_W(\xi ,\xi ^{})=\frac{1}{2\pi }_{\mathrm{}}^{\mathrm{}}f_k(x,x^{})e^{ik(tt^{})}𝑑k,$$
(37)
in eq.(36), the coefficients are found as following
$$f_k(x,x^{})=r^4(xx^{})^{3/2}I_{l+\frac{1}{2}}(kx_<)K_{l+\frac{1}{2}}(kx_>).$$
(38)
$`I_{l+\frac{1}{2}}`$ and $`K_{l+\frac{1}{2}}`$ are modified Bessel functions of the first and second kind, respectively. Performing the integration over $`k`$ we obtain
$$G_W(\xi ,\xi ^{})=\frac{r^4}{2\pi (xx^{})^2}Q_l[\frac{x^2+x^2+(tt^{})^2}{2xx^{}}],$$
(39)
Now as the functional integral in the massless Schwinger model is Gaussian, and the higher order correlators factorize into product of pair correlators $`<F(\xi )F(\xi ^{})>`$, we have
$$<\mathrm{exp}(ie^{}A_\mu 𝑑x^\mu )>=\mathrm{exp}[\frac{e^2}{2}_Dd^2\xi d^2\xi ^{}<F(\xi )F(\xi ^{})>].$$
(40)
Comparing (40) with (35) and (39), results
$$<F(\xi )F(\xi ^{})>=\delta ^2(\xi \xi ^{})\frac{r^2}{x^2}\frac{\mu ^2r^4}{2\pi (xx^{})^2}Q_l(\mathrm{cosh}\frac{L}{r}),$$
(41)
where $`L`$ is the geodesic distance between $`\xi `$ and $`\xi ^{}`$ (see eq.(1)). By considering the behavior of $`I_{l+\frac{1}{2}}(x)`$ at $`x0`$ ($`Q_l(\mathrm{cosh}\frac{L}{r})`$ at $`L=\mathrm{}`$), one can easily check that the behavior of $`f_k(x,x^{})`$ ($`<F(\xi )F(\xi ^{})>`$) is consistent with the Dirichlet boundary condition imposed on eq.(23).
In the flat case limit, one can show that $`lim_r\mathrm{}Q_l[\mathrm{cosh}(L/r)]=K_0(\mu L)`$, and by setting $`g_{\mu \nu }\eta _{\mu \nu }`$, the eq.(41) becomes
$$<F(\xi )F(\xi ^{})>_{\mathrm{flat}}=\delta ^2(\xi \xi ^{})\frac{\mu ^2}{2\pi }K_0(\mu |\xi \xi ^{}|),$$
(42)
which is the strength fields correlator on the flat space–time . In the absence of dynamical fermions ($`\mu =0`$), eq.(42) reduces to $`\delta ^2(\xi \xi ^{})`$ and one obtains the area law for the Wilson loop, which is a characteristic of a confining potential. In the presence of dynamical fermions, since $`K_0(\mu |\xi \xi ^{}|)`$ decays exponentially as $`e^{\mu |\xi \xi ^{}|}`$, the correlator exhibits the finite correlation length (related physically to the screening effect), and the perimeter law is arisen for large contour . On the Poincare half plane, $`K_0`$ is replaced by $`Q_l(\mathrm{cosh}\frac{L}{r})`$ which decays as $`[\mathrm{cosh}(L/r)]^{l1}`$ for large $`L/r`$. Hence in this case we have also a finite correlation length for $`<F(\xi )F(\xi ^{})>`$ and, as we will show, the Wilson loop is perimeter dependent. In fact the area term arising from the delta function is canceled out by the corresponding term in the integration of $`Q_l`$.
Using
$$\underset{T\mathrm{}}{lim}\frac{1}{2\pi T}_0^Te^{ikt}𝑑t_0^Te^{ikt^{}dt^{}}=\delta (k),$$
(43)
and
$$f_0(x,x^{})=\frac{r^4}{2l+1}\frac{x_<^{l1}}{x_>^{l+2}},$$
(44)
one can obtain the following expression
$$U=\underset{T\mathrm{}}{lim}(\frac{1}{T}\mathrm{ln}W)=\frac{r^2}{2l+1}\frac{e^2}{2}[\frac{1}{a}+\frac{1}{b}2(\frac{a^l}{b^{l+1}})],$$
(45)
for the static potential between external charges which is equal to one obtained in eq.(26). After some calculations, one can show that the Wilson loop for $`T>>b>>a`$ is
$$<W_c[A]>=\mathrm{exp}[\frac{r^2}{2l+1}\frac{e^2}{2}T(\frac{1}{a}+\frac{1}{b})+O(\frac{1}{T^\lambda })],$$
(46)
where using the hypergeometric representation of the Legendre function, $`\lambda `$ is found to be $`\lambda =2l+1`$.
But the perimeter of the large Wilson loop $`T>>b,a`$ is
$$_c\frac{r\sqrt{dx^2+dy^2}}{x}=T(\frac{r}{a}+\frac{r}{b}),$$
(47)
therefore eq.(46) shows that for a large contour, the perimeter law is satisfied, which is the same behavior as the flat space–time case.
In the quenched Schwinger model ($`\mu =0`$), the Wilson loop expectation value is
$$<W_c[A]>=\mathrm{exp}[\frac{e^2}{2}\eta ^2\frac{r^2}{x^2}𝑑x𝑑t]=\mathrm{exp}[\frac{e^2}{2}Tr^2(\frac{1}{a}\frac{1}{b})].$$
(48)
But we note that $`Tr^2(1/a1/b)`$ is nothing but the area of the rectangle bounded by the Wilson loop. Therefore on the Poincare half plane, the Wilson loop of the Schwinger model in the absence of dynamical fermions (that is the pure QED<sub>2</sub>) is equal to the exponential of the area, like the flat case. But in contrast to the flat case, the area is not proportional to the geodesic distance of the charges, and we can not conclude that the system is in confining phase. This result is consistent with our discussion after eq.(30).
## 3 Confining aspect of massive Schwinger model on the Poincare half plane
The massive Schwinger model, i.e. U(1) gauge theory with massive dynamical fermions of charge $`e`$ and mass $`m`$, is defined by the action
$$S=\sqrt{g}d^2x[i\overline{\psi }\widehat{\gamma }^ae_a^\mu (_\mu ieA_\mu )\psi +m\overline{\psi }\psi +\frac{1}{4}g^{\mu \nu }g^{\rho \sigma }F_{\mu \rho }F_{\nu \sigma }].$$
(49)
This model is not soluble even in the flat case, but in the limit $`m<<e`$, the physical quantities may be evaluated using perturbative expansion in fermions mass.
In conformally flat curved space–time, the bosonic form of the action (49) is
$$S=\sqrt{g}d^2x[\frac{1}{2}g^{\mu \nu }_\mu \mathrm{\Phi }_\nu \mathrm{\Phi }\frac{ie}{\sqrt{\pi }}ϵ^{\mu \nu }A_\mu _\nu \mathrm{\Phi }\frac{m\mathrm{\Sigma }}{g^{\frac{1}{4}}}\mathrm{cos}(2\sqrt{\pi }\mathrm{\Phi })+\frac{1}{4}g^{\mu \nu }g^{\rho \sigma }F_{\mu \rho }F_{\nu \sigma }],$$
(50)
where the constant $`\mathrm{\Sigma }`$ is given by (14). In fact this model is a Sine-Gordon model whose interaction is position dependent . The confining behavior of this system can be analyzed perturbatively by expanding the mass term in a power of $`\mathrm{\Phi }`$ . The equations of motion followed from the bosonized action (50), in the presence of external charges (21), are
$$\frac{d}{dx}\frac{x^2}{r^2}\frac{d}{dx}A_0+\frac{ie}{\sqrt{\pi }}\frac{d\mathrm{\Phi }}{dx}=ie^{}[\delta (xb)\delta (xa)],$$
$$\frac{d^2\mathrm{\Phi }}{dx^2}+\frac{ie}{\sqrt{\pi }}\frac{dA_0}{dx}+\frac{2\sqrt{\pi }m\mathrm{\Sigma }r}{x}\mathrm{sin}(2\sqrt{\pi }\mathrm{\Phi })=0,$$
(51)
where we have used, as before, the Coulomb gauge $`A_1=0`$ and $`dA_0/dt=0`$. Using the approximation $`\mathrm{sin}(2\sqrt{\pi }\mathrm{\Phi })2\sqrt{\pi }\mathrm{\Phi }`$, and assuming that the field $`\mathrm{\Phi }`$ is a slowly varying field , we arrive at
$$\frac{d}{dx}(\frac{x^2}{r^2}+\frac{e^2}{4\pi ^2m\mathrm{\Sigma }r}x)\frac{dA_0}{dx}=ie^{}[\delta (xb)\delta (xa)],$$
(52)
$$\frac{4\pi m\mathrm{\Sigma }r}{x}\mathrm{\Phi }+\frac{ie}{\sqrt{\pi }}\frac{dA_0}{dx}=0,$$
with solutions
$$A_0(x)=\{\begin{array}{cc}0,\hfill & x>b\hfill \\ i\frac{4\pi ^2m\mathrm{\Sigma }e^{}r}{e^2}[\mathrm{ln}(1+\frac{e^2}{4\pi ^2m\mathrm{\Sigma }}\frac{r}{b})\mathrm{ln}(1+\frac{e^2}{4\pi ^2m\mathrm{\Sigma }}\frac{r}{x})],\hfill & a<x<b\hfill \\ i\frac{4\pi ^2m\mathrm{\Sigma }e^{}r}{e^2}[\mathrm{ln}(1+\frac{e^2}{4\pi ^2m\mathrm{\Sigma }}\frac{r}{b})\mathrm{ln}(1+\frac{e^2}{4\pi ^2m\mathrm{\Sigma }}\frac{r}{a})],\hfill & x<a\text{,}\hfill \end{array}$$
(53)
and
$$\mathrm{\Phi }(x)=\{\begin{array}{cc}0,\hfill & x>b\hfill \\ (\frac{e^{}}{e})\frac{e^2}{4\pi ^{\frac{3}{2}}m\mathrm{\Sigma }}(\frac{1}{\frac{x}{r}+\frac{e^2}{4\pi ^2m\mathrm{\Sigma }}}),\hfill & a<x<b\hfill \\ 0,\hfill & x<a\text{.}\hfill \end{array}$$
(54)
Therefore we find the potential $`U=\frac{1}{2}\rho A_0𝑑x`$ as
$$U=2\pi ^2m\mathrm{\Sigma }(\frac{e^{}}{e})^2r[\mathrm{ln}(1+\frac{e^2}{4\pi ^2m\mathrm{\Sigma }}\frac{r}{a})\mathrm{ln}(1+\frac{e^2}{4\pi ^2m\mathrm{\Sigma }}\frac{r}{b})].$$
(55)
By fixing $`a`$ and increasing the separation of the charges, $`U`$ increases and finally achieves the limiting value $`2\pi ^2m\mathrm{\Sigma }(e^{}/e)^2r\mathrm{ln}(1+\frac{e^2r}{4\pi ^2m\mathrm{\Sigma }a})`$, which shows that the system is in the screening phase. When one of the charges is located near the horizontal axis, $`a0`$, the eq.(55) goes to infinity, unless $`ba`$. So the system is in confining phase at $`x0`$, as we expect. On the other hand, in order to satisfy the conditions that we have assumed for the field $`\mathrm{\Phi }`$ (to be small and small varying), we must take
$$(\frac{x}{r}+\frac{e^2}{4\pi ^2m\mathrm{\Sigma }})>>(\frac{e^{}}{e})\frac{e^2}{4\pi ^2m\mathrm{\Sigma }}.$$
(56)
By expanding $`U`$ in terms of $`r/a`$ and $`r/b`$, we obtain
$$U=\frac{e^2}{2}r^2(\frac{1}{a}\frac{1}{b})+O(\frac{r^2}{a^2},\frac{r^2}{b^2}).$$
(57)
For large $`a/r`$ and $`b/r`$, $`U`$ is proportional to the area of the Wilson loop characterized by $`a,b,`$ and $`T`$. In the flat case, this behavior is interpreted as a sign of confinement but, as we have discussed earlier, this is not true for the Poincare half plane.
Finally if we consider the small fermion mass limit $`m<<e`$, and also $`e^{}<<e`$, the eq.(55) reduces to
$$U=2\pi ^2m\mathrm{\Sigma }(\frac{e^{}}{e})^2r\mathrm{ln}(\frac{b}{a}),$$
(58)
which is comparable with the corresponding result in the flat case, after substituting $`r`$ln$`(b/a)(ba)`$ and $`\mathrm{\Sigma }`$(Poincare) $`\mathrm{\Sigma }`$ (flat) . Note that the potential (58) is proportional to the geodesic distance $`d=r`$ln $`(b/a)`$, but this is not a sign of confinement as in the flat case. To see this, note that if one fixes the position of the first charge at $`x=a`$ and moves the other charge to a large distance ($`b/r\mathrm{}`$), the eq.(55) reduces to (for $`m<<e`$, $`e^{}<<e`$)
$$U(\frac{b}{r}\mathrm{})=e^2r^2\frac{1}{2a},$$
(59)
which has a finite value, and the system is again in the screening (and not confining) phase.
## 4 Conclusion
Let us summarize the main results of the paper:
1- In $`m=0`$, the Schwinger model on flat space–time is in screening phase, but on the Poincare half plane, the system is in confining phase in $`x0`$ region, and in screening phase in regions far enough from the horizontal axis ( after eq.(28)),
2- In $`m=0`$ and $`e=0`$, the model is confining in flat case (after eq.(29)), but on the Poincare half plane, the phase depends on the region under study (after eq.(30)).
3- In $`m=0`$, the Wilson loop obeys the perimeter (area)–law for $`e0`$ ($`=0`$) in both the flat (after eq.(42)) and the Poincare (eqs.(47) and (48)) cases. But on the Poincare half plane, the area dependence does not indicate the confining phase (in contrast to the flat case) (after eq.(48)).
4- In $`m0`$, the Schwinger model on the flat space–time is in confining phase but in the Poincare case, the model is in screening phase (after eqs.(55) and (59)).
Acknowledgement We would like to thank A. Arfaei for useful discussion and the research council of the University of Tehran for partial financial support.
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# 1 Introduction
## 1 Introduction
Recent years have seem dramatic progress in the understanding of non-perturbative aspects of string theory . With these studies has come the realization that solitonic extended objects, other than just strings, play an essential role. An important object in these investigations has been Dirichlet branes . D-branes are non-perturbative states on which open string can live, and to which various closed strings including Ramond-Ramond states can couple.
Type II string theories have two kind of D-branes, BPS D$`p`$-branes for $`p`$ even(odd) and non-BPS D$`p`$-branes for $`p`$ odd(even) in IIA(IIB) theory. The BPS branes are stable solitons which break half of the space-time supersymmetries and their dynamics are properly described in field theory by Dirac-Born-Infeld(DBI) action (see also ) and Chern-Simons action . The non-BPS D$`p`$-branes on the other hand suffer from open string tachyonic mode whose mass causes the brane in a flat background to be unstable. However, there are other terms in the tachyon potential which makes it bounded from below. Consequently, the non-BPS branes decay to minimum of the tachyon potential. It has been conjectured that at the stationary point of the tachyon potential, the negative minimum energy of the tachyon potential plus the positive energy of the brane tension is exactly zero . Hence, the unstable non-BPS branes in flat space-time vacuum should decay to the true vacuum of the theory in which there is no branes. This conjecture was studied in by explicit calculation of the tachyon potential using the string field theory framework.
The dynamics of massless bosonic excitations of non-BPS D$`p`$-branes are suitablely described by the DBI action in field theory. This action has been generalized to the supersymmetric form to include the dynamics of massless fermionic modes of the branes as well . The RR fields of the type II theory have no coupling to the non-BPS D-branes through the usual Chern-Simons action. However, there is a non-vanishing coupling between the RR field and tachyon on the world-volume of the branes . The Chern-Simons action hence was modified in to incorporate this coupling. In the present paper, on the other hand, we are interested in generalizing the DBI action to take into account the dynamics of the tachyon field. We study this by explicit evaluation of some nontrivial disk S-matrix elements in the first quantized string theory. From these matrix elements we conjecture an extension for the DBI action which includes the tachyon field as well.
An outline of the paper is as follows. We begin in section 2 by expressing our conjecture for extension of the DBI action which includes dynamics of the tachyon field. Then we expand this action around a background B-flux to produce various couplings involving two tachyons and one gauge field, graviton, dilaton or Kalb-Ramond antisymmetric tensor. In Section 2.1 we transform the above couplings between commutative fields to their non-commutative counterparts. We do this because the disk S-matrix elements in the presence of the background B-flux with which we are going to compare our conjectured action in the subsequent section are corresponding to non-commutative open string fields . In Section 3, we evaluate the S-matrix elements and check their consistency with the proposed field theory couplings. In Section 4 we extend our proposed action for describing dynamics of a single non-BPS D$`p`$-brane to the case of the non-abelian theory of $`N`$ coincident branes using the general grounds of the symmetrized trace, and non-abelian gauge and invariance under T-duality transformations. Appendix contains our conventions and some useful comments on conformal field theory propagators and vertex operators used in our calculations.
## 2 Abelian action
The world-volume theory of a single non-BPS D-brane in type II theory includes a massless U(1) vector $`A_a`$, a set of massless scalars $`X^i`$, describing the transverse oscillations of the brane, a tachyonic state $`T`$ and their fermionic partners (see, e.g., ). The leading order low-energy action for the massless fields corresponds to a dimensional reduction of a ten dimensional U(1) Yang Mills theory. As usual in string theory, there are higher order $`\alpha ^{}=\mathrm{}_s^2`$ corrections, where $`\mathrm{}_s`$ is the string length scale. As long as derivatives of the field strengths (and second derivatives of the scalars) are small compared to $`\mathrm{}_s`$, then the action takes a Dirac-Born-Infeld form . To take into account the couplings of the massless open string states with closed strings, the DBI action may be extended naturally to include massless Neveu-Schwarz closed string fields,i.e., the metric, dilaton and Kalb-Ramond filed. In this case one arrives at the following world-volume action:
$`S`$ $`=`$ $`T_p{\displaystyle d^{p+1}\sigma e^\mathrm{\Phi }\sqrt{\mathrm{det}(P[G_{ab}+B_{ab}]+2\pi \alpha ^{}F_{ab}}}.`$
Here, $`F_{ab}`$ is the abelian field strength of the world-volume ordinary gauge field, while the metric and antisymmetric tensors are the pull-backs of the bulk tensors to the D-brane world-volume, e.g.,
$`P[G_{ab}]`$ $`=`$ $`G_{\mu \nu }{\displaystyle \frac{X^\mu }{\sigma ^a}}{\displaystyle \frac{X^\nu }{\sigma ^b}}`$ (1)
$`=`$ $`G_{ab}+2G_{i(a}_{b)}X^i+G_{ij}_aX^i_bX^j`$
where in the second line above we have used that fact that we are employing static gauge throughout the paper, i.e., $`\sigma ^a=X^a`$ for world-volume and $`X^i(\sigma ^a)`$ for transverse coordinates.
In order to extend this action to incorporate dynamics of the tachyonic mode as well, we shall evaluate some disk S-matrix elements in string theory and read from them various couplings involving the tachyon field. Our results indicate that the tachyon should appear in the following extension of DBI action:
$$S=T_pd^{p+1}\sigma e^\mathrm{\Phi }V(T)\sqrt{\mathrm{det}(P[G_{ab}+B_{ab}]+2\pi \alpha ^{}F_{ab}+2\pi \alpha ^{}_aT_bT)}$$
(2)
where the tachyon potential is $`V(T)=1+2\pi \alpha ^{}m^2T^2/2+O(T^4)`$ and the tachyon mass is $`m^2=1/2\alpha ^{}`$. In our conventions the tachyon field is dimensionless. The conjecture in is that the tachyon potential is zero at the minimum of the potential, i.e., $`V(T_0)=0`$. Hence, upon tachyon condensation at this point the abelian action (2) of the non-BPS brane becomes zero.
Note that appearance of the tachyon kinetic term and potential in (2) is similar in form to the kinetic term and potential of the transverse scalar fields in the non-abelian DBI action of $`N`$ coincident BPS D-branes . In this case though the kinetic term appears in the pull-back of the metric under the square root and the potential for scalar fields multiplies the square root in the DBI action(see eq. (4) for $`T=0`$).
We now continue backward, assuming the above action (2) and check its consistency with some S-matrix elements. To have nontrivial check, we shall evaluate disk amplitudes describing decay of two tachyons to dilaton, graviton or Kalb-Ramond antisymmetric tensor on the world-volume of a single non-BPS D$`p`$-brane with background B-flux. The amplitude describing the world-volume coupling of two tachyons to gauge field will be evaluated as well. Therefore, we begin by expanding (2) for fluctuations around the background $`G_{\mu \nu }=\eta _{\mu \nu }`$, $`B_{\mu \nu }=^{ab}\eta _{a\mu }\eta _{b\nu }`$, $`\mathrm{\Phi }=0`$ and $`T=0`$ to extract interactions expected from the proposed action (2). The fluctuations should be normalized as the conventional field theory modes which appear in the string vertex operators. As a first step, we recall that the graviton vertex operator corresponds to string frame metric. Hence, one should transform the Einstein frame metric $`G_{\mu \nu }`$ to the string frame metric $`g_{\mu \nu }`$ via $`G_{\mu \nu }=e^{\mathrm{\Phi }/2}g_{\mu \nu }`$. Now with our conventions for string vertex operators (see Appendix), the string mode fluctuations take the form
$`g_{\mu \nu }`$ $`=`$ $`\eta _{\mu \nu }+2\kappa h_{\mu \nu }`$
$`\mathrm{\Phi }`$ $`=`$ $`\sqrt{2}\kappa \varphi `$
$`B_{\mu \nu }`$ $`=`$ $`^{ab}\eta _{a\mu }\eta _{b\nu }2\kappa b_{\mu \nu }`$
$`T`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2\pi \alpha ^{}T_p}}}\tau `$
$`A_a`$ $`=`$ $`{\displaystyle \frac{1}{2\pi \alpha ^{}\sqrt{T_p}}}a_a`$
$`X^i`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{T_p}}}\lambda ^i.`$
With these normalizations, the pull back of the Einstein frame metric becomes:
$`P[G_{ab}]`$ $`=`$ $`\eta _{ab}(1+{\displaystyle \frac{\kappa }{\sqrt{2}}}\varphi )+2\kappa P[h_{ab}]+{\displaystyle \frac{1}{T_p}}(1+{\displaystyle \frac{\kappa }{\sqrt{2}}}\varphi )_a\lambda ^i_b\lambda _i+\mathrm{}`$
where the dots represents terms with two and more closed string fields.
Now it is straightforward, to expand eq. (2) using
$`\sqrt{\mathrm{det}(M_0+M)}`$ $`=`$ $`\sqrt{\mathrm{det}(M_0)}(1+{\displaystyle \frac{1}{2}}\mathrm{Tr}(M_0^1M)`$
$`{\displaystyle \frac{1}{4}}\mathrm{Tr}(M_0^1MM_0^1M)+{\displaystyle \frac{1}{8}}(\mathrm{Tr}(M_0^1M))^2+\mathrm{})`$
to produce a vast array of interactions. We are interested in the interactions quadratic in tachyon, and linear in massless open and closed string fluctuations. The appropriate Lagrangian are:
$`_{2,0}`$ $`=`$ $`c\left({\displaystyle \frac{1}{2}}m^2\tau ^2+{\displaystyle \frac{1}{2}}(V_S)^{ab}_a\tau _b\tau \right)`$ (3)
$`_{3,0}`$ $`=`$ $`{\displaystyle \frac{c}{2\sqrt{T_p}}}\left({\displaystyle \frac{1}{2}}m^2(V_A)^{ab}f_{ba}\tau ^2+{\displaystyle \frac{1}{2}}V^{ab}f_{ba}V^{cd}_c\tau _d\tau V^{ab}f_{bc}V^{cd}_d\tau _a\tau \right)=\mathrm{\hspace{0.17em}\hspace{0.17em}0}`$
$`_{2,1}`$ $`=`$ $`\kappa c((V^{ab}(h_{ba}b_{ba})+{\displaystyle \frac{1}{2\sqrt{2}}}(\mathrm{Tr}(V)4)\varphi )({\displaystyle \frac{1}{2}}(V_S)^{ab}_a\tau _b\tau +{\displaystyle \frac{1}{2}}m^2\tau ^2)`$
$`V^{ab}(h_{bc}b_{bc}+{\displaystyle \frac{1}{2\sqrt{2}}}\varphi \eta _{bc})V^{cd}_d\tau _a\tau ).`$
where we have dropped some total derivative terms in $`_{3,0}`$. In the above Lagrangian, $`f_{ab}=_aa_b_ba_a`$ and
$`c\sqrt{\mathrm{det}(\eta _{ab}+_{ab})}`$ $`,`$ $`V^{ab}\left((\eta +)^1\right)^{ab},`$ (4)
and $`V_S(V_A)`$ is symmetric(antisymmetric) part of the V matrix above. It is important to note that the antisymmetric matrix $`V_A`$ appears in total derivative terms in the Lagrangian $`_{3,0}`$ which involves only open string fields.
In the case that the background B-flux is zero, the coupling of Kalb-Ramond field to tachyon in the second line of $`_{2,1}`$ vanishes, and the graviton and dilaton couplings reduce to the natural coupling of these fields to the kinetic term of the tachyon field, i.e., $`e^\mathrm{\Phi }G^{ab}_aT_bT`$. In that way, there is no nontrivial coupling that confirms the conjectured action (2) is valid or not. So we continue our discussion for non-vanishing background B-flux.
The open string fields appearing in the DBI action (2) or (3) are ordinary commutative fields. Whereas, open string vertex operators in string theory with background B-flux correspond to non-commutative fields . Hence, in order to compare the couplings in (3) with corresponding S-matrix elements, one has to transform the commutative fields in (3) to their non-commutative variables.
### 2.1 Change of variables
In differential equation for transforming commutative gauge field to its non-commutative counterpart was found to be
$`\delta \widehat{A}_a(\theta )`$ $`=`$ $`{\displaystyle \frac{1}{4}}\delta \theta ^{cd}\left(\widehat{A}_c\widehat{F}_{ad}+\widehat{F}_{ad}A_c\widehat{A}_c_d\widehat{A}_a_d\widehat{A}_aA_c\right)`$ (5)
$`\delta \widehat{F}_{ab}(\theta )`$ $`=`$ $`{\displaystyle \frac{1}{4}}\delta \theta ^{cd}\left(2\widehat{F}_{ac}\widehat{F}_{bd}+2\widehat{F}_{bd}\widehat{F}_{ac}\widehat{A}_c(\widehat{D}_d\widehat{F}_{ab}+_d\widehat{F}_{ab})(\widehat{D}_d\widehat{F}_{ab}+_d\widehat{F}_{ab})\widehat{A}_c\right)`$
where the gauge field strength and $``$ product were defined to be
$`\widehat{F}_{ab}`$ $`=`$ $`_a\widehat{A}_b_b\widehat{A}_ai\widehat{A}_a\widehat{A}_b+i\widehat{A}_b\widehat{A}_a`$
$`=`$ $`_a\widehat{A}_b_b\widehat{A}_ai[\widehat{A}_a,\widehat{A}_b]_M`$
$`\widehat{f}(x)\widehat{g}(x)`$ $`=`$ $`e^{\frac{i}{2}\theta _{ab}_y^a_z^b}\widehat{f}(y)\widehat{g}(z)|_{y=z=x}.`$
These differential equations can be integrated perturbatively to find a relation between ordinary fields appearing in (3) and non-commutative fields corresponding to open string vertex operators. The result for abelian case that we are interested in is :
$`A_a`$ $`=`$ $`\widehat{A}_a{\displaystyle \frac{1}{2}}\theta ^{cd}\left(\widehat{A}_c^{}\widehat{F}_{ad}\widehat{A}_c^{}_d\widehat{A}_a\right)+O(\widehat{A}^3)`$
$`F_{ab}`$ $`=`$ $`\widehat{F}_{ab}\theta ^{cd}\left(\widehat{F}_{ac}^{}\widehat{F}_{bd}\widehat{A}_c^{}_d\widehat{F}_{ab}\right)+O(\widehat{A}^3)`$ (6)
where the commutative multiplication $`^{}`$ operates as
$`\widehat{f}(x)^{}\widehat{g}(x)`$ $`=`$ $`{\displaystyle \frac{\mathrm{sin}(\frac{1}{2}\theta _{ab}_y^a_z^b)}{\frac{1}{2}\theta _{ab}_y^a_z^b}}\widehat{f}(y)\widehat{g}(z)|_{y=z=x}.`$
In we verified by explicit calculation of S-matrix elements of one massless closed and two open string states that the transformation (6) is exactly reproduced by perturbative string theory. Appropriate transformations for scalar fields such as the tachyon can be read from (6), i.e.,
$`T`$ $`=`$ $`\widehat{T}+\theta ^{cd}\widehat{A}_c^{}_d\widehat{T}+\mathrm{}`$ (7)
$`_aT`$ $`=`$ $`_a\widehat{T}i[\widehat{A}_a,\widehat{T}]_M+\theta ^{cd}\left(\widehat{F}_{ca}^{}_b\widehat{T}+\widehat{A}_c^{}_a_d\widehat{T}\right)+\mathrm{}`$
where dots represent terms which involve more than two open string fields. They produce couplings between more than three fields upon replacing them into (3) in which we are not interested.
The differential equation (5) expresses infinitesimal variation of linear field , e.g., $`\delta \widehat{A}_a`$, in terms of infinitesimal variation of the non-commutative parameter, i.e., $`\delta \theta ^{cd}`$. Upon integration (6), this transforms the linear commutative gauge field in terms of nonlinear combination of non-commutative fields. However, we are interested in transforming quadratic combinations of commutative fields appearing in (3) in terms of non-commutative fields. Such a transformation, in principle, might be found from a differential equation alike (5) that expresses infinitesimal variation of the quadratic fields in terms of infinitesimal variation of non-commutative parameter. Upon integration, that would produced the desired transformation. In that way, one would find that not only the fields transform as in (7) but the multiplication rule between fields also undergo appropriate transformation. We are not going to find such a differential equations here. Instead, we simply note that the transformation for multiplication rule between two open string fields can be conjectured from the right hand side of eq. (6) to be
$`fg|_{\theta =0}`$ $``$ $`f^{}g|_{\theta 0}`$ (8)
for $`f`$ and $`g`$ being any arbitrary open string fields. This transformation rule was confirmed in by explicit evaluation of S-matrix elements of one massless closed and two open string states.
Now with the help of equation (7) and (8), one can transform the commutative Lagrangian (3) to non-commutative counterparts. In doing so, one should first using (8) replace ordinary multiplication of two tachyons by the $`^{}`$ multiplication. Then, using (7) the ordinary tachyon fields should be transformed to their non-commutative counterparts. The results are
$`\widehat{}_{2,0}`$ $`=`$ $`c\left({\displaystyle \frac{1}{2}}m^2\widehat{\tau }\widehat{\tau }+{\displaystyle \frac{1}{2}}(V_S)^{ab}_a\widehat{\tau }_b\widehat{\tau }\right)+{\displaystyle \frac{ic}{4\pi \sqrt{T_p}}}(V_S)^{ab}_a\widehat{\tau }[\widehat{a}_b,\widehat{\tau }]_M+\mathrm{}`$ (9)
$`\widehat{}_{2,1}`$ $`=`$ $`\kappa c((V^{ab}(h_{ba}b_{ba})+{\displaystyle \frac{1}{2\sqrt{2}}}(\mathrm{Tr}(V)4)\varphi )({\displaystyle \frac{1}{2}}(V_S)^{ab}_a\widehat{\tau }^{}_b\widehat{\tau }+{\displaystyle \frac{1}{2}}m^2\widehat{\tau }^{}\widehat{\tau })`$ (10)
$`V^{ab}(h_{bc}b_{bc}+{\displaystyle \frac{\varphi \eta _{bc}}{2\sqrt{2}}})V^{cd}_d\widehat{\tau }^{}_a\widehat{\tau })+\mathrm{}`$
where ellipsis represent terms which have more than three fields. Here we have dropped some total derivative terms which appeared in $`\widehat{}_{2,0}`$ and also replaced $`^{}`$ between two non-commutative tachyons in (9) with ordinary multiplication rule because the deference is some total derivative terms. In the eq. (10) on the other hand, the difference between $`^{}`$ and ordinary multiplication rules is not just a total derivative terms. Note that upon inserting the transformation (7) into (3), the antisymmetric matrix $`(V_A)^{ab}`$ appears in eq. (9) only in the $``$ product terms<sup>1</sup><sup>1</sup>1Note that our conventions set $`\theta ^{ab}=4\pi (V_A)^{ab}`$..
It is interesting to note that the symmetric part of the $`(\eta +)^1`$ matrix, i.e., $`V_S`$, appears as the metric in (9) and the antisymmetric part appears in the definition of $``$ product in the Moyal bracket. This is consistent with the conclusion reached in . In our discussion, however, the symmetric part $`V_S`$ appears naturally as a result of expanding the ordinary DBI action around the background B-flux, and the antisymmetric part $`V_A`$ appears as a result of transforming commutative fields to non-commutative variables. In the Lagrangian (10), on the other hand, which involves open and closed string fields, both symmetric and antisymmetric matrices appear in its different coupling terms.
## 3 Scattering Calculations
The couplings in (9) and (10) should be reproduced by disk S-matrix elements of string theory if the proposed action (2) is going to be valid. In this section, we work at the string theory side and evaluate these couplings using the conformal field theory technique. We begin with the evaluation of the coupling of two tachyons to gauge or scalar fields.
### 3.1 Open-Open-Open couplings
In the world-sheet conformal field theory framework, the coupling of two tachyons to a gauge or scalar field is described properly by the correlation of their corresponding vertex operators inserted at the boundary of the disk world-sheet, that is
$`A`$ $``$ $`(\zeta _3𝒢)_\mu {\displaystyle }dx_1dx_2dx_3:V_1(2k_1V^T,x_1)::V_1(2k_2V^T,x_2)::V_0^\mu (2k_3V^T,x_3):`$
where the vertex operators are given in the Appendix. Using the world-sheet conformal field theory technique, it is not difficult to perform the correlators above and show that the integrand is invariant under $`SL(2,R)`$. Gauging this symmetry by fixing the positions of the vertices at arbitrary points, one finds $`A(\tau _1,\tau _2,\lambda _3)=0`$ and
$`A(\tau _1,\tau _2,a_3)`$ $`=`$ $`{\displaystyle \frac{c\mathrm{sin}(\pi l)}{\pi \sqrt{T_p}}}(k_1V_S\zeta _3)`$ (11)
where we have defined $`l2k_1V^TVk_2=2k_1V_Ak_2`$. We have also normalized the amplitude by the appropriate coupling factor $`c/2\pi \sqrt{T_p}`$. The $`\mathrm{sin}(\pi l)`$ factor above arises basically from two different phase factors corresponding to two distinct cyclic orderings of the vertex operators. Each phase factor stems from the second term of the world-sheet propagator (29). Using the fact that our conventions set $`\theta ^{ab}=4\pi V_A^{ab}`$, it is not difficult to verify that the S-matrix element (11) is exactly reproduced by the second term in (9). At the same time, vanishing of $`A(\tau _1,\tau _2,\lambda _3)`$ is consistent with (9).
### 3.2 Closed-Open-Open amplitudes
The amplitudes describing decay of two open string tachyons to one massless closed string NSNS mode is given by the following correlation:
$$A(\epsilon _3D)_{\mu \nu }dx_1dx_2d^2z:V_0(2k_1V^T,x_1)::V_0(2k_2V^T,x_2)::V_1^\mu (p_3,z_3)::V_1^\nu (p_3D,\overline{z}_3):$$
where the closed string vertex operator inserted at the middle and open string vertex operators at the boundary of the disk world-sheet. Explicit form of the vertex operators in terms of world-sheet fields are given in the Appendix. Here again using appropriate world-sheet propagators, one can evaluate the correlations above and show that the integrand is $`SL(2,R)`$ invariant. We refer the reader to Refs. for the details of the calculations. Gauging the $`SL(2,R)`$ symmetry by fixing $`z_3=i`$ and $`x_1=\mathrm{}`$, one arrives at
$`A`$ $``$ $`2^{2s2}{\displaystyle 𝑑x\left((2s+1)\mathrm{Tr}(\epsilon _3D)\frac{8ik_2V^T\epsilon _3DVk_1}{xi}+\frac{8ik_1V^T\epsilon _3DVk_2}{x+i}\right)}`$
$`\times (xi)^{sl}(x+i)^{s+l}`$
where the integral is taken from $`\mathrm{}`$ to $`+\mathrm{}`$, and $`s=p_3V_Sp_3=1/22k_1V_Sk_2`$. This integral is doable and the result is
$`A`$ $`=`$ $`{\displaystyle \frac{i\kappa c}{2}}\left(a_1(s+l)a_2(sl)\right){\displaystyle \frac{\mathrm{\Gamma }(2s)}{\mathrm{\Gamma }(1sl)\mathrm{\Gamma }(1s+l)}}`$ (12)
where $`a_1`$ and $`a_2`$ are two kinematic factors depending only on the space time momenta and the closed string polarization tensor
$`a_1`$ $`=`$ $`4k_2V^T\epsilon _3DVk_1`$
$`a_2`$ $`=`$ $`(s+l)\mathrm{Tr}(\epsilon _3D)+4k_1V^T\epsilon _3DVk_2.`$
We have also normalized the amplitude (12) at this point by the coupling factor $`i\kappa c/2\pi `$. As a check of our calculations, we have inserted the dilaton polarization (31) into (12) and found that the result is independent of the auxiliary vector $`\mathrm{}^\mu `$. The amplitude (12) has the pole structure at $`m_{open}^2=n/\alpha ^{}`$<sup>2</sup><sup>2</sup>2We explicitly restore $`\alpha ^{}`$ here. Otherwise our conventions set $`\alpha ^{}=2`$. This does not have tachyon pole which is consistent with the fact that coupling of three tachyons is zero. In fact due to the world-sheet fermions in the tachyon vertex operator, coupling of any odd number of tachyons is zero. Hence, the world-volume of the non-BPS D$`p`$-branes has a $`Z_2`$ symmetry under which the tachyon changes sign.
#### 3.2.1 Massless poles
Given the general form of the string amplitude in eq. (12), one can expand this amplitude as an infinite sum of terms reflecting the infinite tower of open string states that propagate on the world-Volume of D-brane. In the domain where $`s0`$, the first term of the expansion representing the exchange of massless string states dominate. In this case the scattering amplitude (12) reduces to
$`A`$ $`=`$ $`{\displaystyle \frac{i\kappa c\mathrm{sin}(\pi l)}{4\pi s}}(a_1+a_2)+\mathrm{}`$ (13)
where dots represent contact terms and the infinite series of massive poles. Making the appropriate explicit choices of polarizations, we find
$`A_s(\tau _1,\tau _2,\varphi _3)`$ $`=`$ $`{\displaystyle \frac{i\kappa c\mathrm{sin}(\pi l)}{8\pi \sqrt{2}s}}({\displaystyle \frac{l}{2}}(\mathrm{Tr}(D)+2)4k_1VVk_2)+12`$ (14)
$`A_s(\tau _1,\tau _2,h_3)`$ $`=`$ $`{\displaystyle \frac{i\kappa c\mathrm{sin}(\pi l)}{4\pi s}}({\displaystyle \frac{l}{2}}\mathrm{Tr}(\epsilon _3D)4k_1V\epsilon _3^TVk_2)+12.`$
Here $`h_3`$ stands for both graviton and Kalb-Ramond antisymmetric tensors. In writing explicitly the above massless poles, one finds some terms which are proportional to $`s`$ as well. We will add these terms which have no contribution to the massless poles of field theory to the contact terms in (16). The amplitudes (14) should be reproduced in $`s`$-channel of field theory. They can be evaluated in field theory as
$`A_s^{}(\tau _1,\tau _2,\varphi _3)`$ $`=`$ $`(\stackrel{~}{V}_{\varphi _3a})^a(\stackrel{~}{G}_a)_{ab}(\stackrel{~}{V}_{a\tau _1\tau _2})^b`$
$`A_s^{}(\tau _1,\tau _2,h_3)`$ $`=`$ $`(\stackrel{~}{V}_{h_3a})^a(\stackrel{~}{G}_a)_{ab}(\stackrel{~}{V}_{a\tau _1\tau _2})^b`$ (15)
where the propagator and the vertices can be read from expansion of (2) in terms of non-commutative fields. They are
$`(\stackrel{~}{G}_a)^{ab}`$ $`=`$ $`{\displaystyle \frac{i}{c}}{\displaystyle \frac{(V_S^1)^{ab}}{s}}`$
$`(\stackrel{~}{V}_{\varphi _3a})^a`$ $`=`$ $`{\displaystyle \frac{\sqrt{T_p}\kappa c}{2\sqrt{2}}}\left({\displaystyle \frac{1}{2}}(\mathrm{Tr}(D)+2)p_3V_A^ap_3VV^a+V^aVp_3\right)`$
$`(\stackrel{~}{V}_{h_3a})^a`$ $`=`$ $`\sqrt{T_p}\kappa c\left({\displaystyle \frac{1}{2}}\mathrm{Tr}(\epsilon _3D)p_3V_A^ap_3V\epsilon _3^TV^a+V^a\epsilon _3^TVp_3\right)`$
$`(\stackrel{~}{V}_{a\tau _1\tau _2})^a`$ $`=`$ $`{\displaystyle \frac{c\mathrm{sin}(\pi l)}{2\pi \sqrt{T_p}}}k_1V_S^a+12.`$
In writing the above propagator, we have used the covariant gauge $`V_S^{ab}_a\widehat{A}_b=0`$. Replacing above propagator and vertices into (15), one finds exactly the string massless poles (14).
#### 3.2.2 Contact terms
Having examined in detail the massless poles of string amplitude, we now extract the low energy contact terms of the string amplitude (12). Expanding the gamma function appearing in this amplitude for $`s0`$, one will find
$`A`$ $`=`$ $`{\displaystyle \frac{i\kappa c}{2}}(({\displaystyle \frac{a_1+a_2}{2\pi }}){\displaystyle \frac{\mathrm{sin}(\pi l)}{s}}+({\displaystyle \frac{a_1a_2}{2}}){\displaystyle \frac{\mathrm{sin}(\pi l)}{\pi l}}`$
$`+(a_1+a_2){\displaystyle \frac{\mathrm{sin}(\pi l)}{\pi l}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\zeta (2n+1)l^{(2n+1)}+k^2O(s,l)).`$
The factor $`\mathrm{sin}(\pi l)/(\pi l)`$ appears for all the contact terms which is consistent with the transformation of multiplication rule in (8). The second term of the first line above is the contact term with minimum number of momentum in which we are interested, that is,
$`A_c`$ $``$ $`{\displaystyle \frac{i\kappa c\mathrm{sin}(\pi l)}{4\pi l}}(a_1a_2)`$ (16)
Inserting appropriate polarization (see Appendix) and adding the residue contact terms of the massless poles (13), one finds
$`A_c(\tau _1,\tau _2,\varphi _3)`$ $`=`$ $`{\displaystyle \frac{i\kappa c\mathrm{sin}(\pi l)}{8\pi \sqrt{2}l}}\left({\displaystyle \frac{s}{2}}(\mathrm{Tr}(D)+2)4k_1VVk_2\right)+12`$ (17)
$`A_c(\tau _1,\tau _2,h_3)`$ $`=`$ $`{\displaystyle \frac{i\kappa c\mathrm{sin}(\pi l)}{4\pi l}}\left({\displaystyle \frac{s}{2}}\mathrm{Tr}(\epsilon _3D)4k_1V\epsilon _3^TVk_2\right)+12.`$
where again $`h_3`$ stands for both graviton and Kalb-Ramond antisymmetric tensors. These contact terms are reproduced exactly by the Lagrangian in (10). The first terms in (17) by the terms in the first line of (10) and the second terms in (17) by the terms in the second line of (10). Note that, while the first terms in (17) can be reproduced in field theory by an action in which the tachyon kinetic term appears linearly like the one proposed in , the second terms in (17) can be reproduced only if the tachyon kinetic term appears non-linearly in the determinant under the square root in the BDI action, i.e., eq. (2). This ends our illustration of consistency between disk S-matrix elements and the proposed action (2).
## 4 Non-abelian action
In this section we extend the proposed action (2) for a non-BPS D$`p`$-brane to the case of N coincident non-BPS D$`p`$-brane where the world-volume theory involves a U(N) gauge theory. Our guiding principle in constructing such a non-abelian action is that the action should be consistent with the familiar rules of T-duality. This guideline has been recently employed by Myers to construct non-abelian DBI action. In this way, one should start with non-abelian action for D9-branes and then use some sort of T-duality transformations to convert the D9-brane action to non-abelian action for D$`p`$-branes. Therefore, we begin by extending the abelian action (2) to non-abelian action for non-BPS D9-branes. In this case there is no scalar field corresponding to the transverse direction of D9-branes. Hence, the non-abelian action may be constructed from the corresponding abelian case by simply extending the derivative of open string fields to its non-abelian covariant derivative , and a trace over the U(N) representations (see also ), that is,
$$S=T_9d^{10}\sigma \mathrm{Tr}\left(e^\mathrm{\Phi }V(T)\sqrt{\mathrm{det}(G_{\mu \nu }+B_{\mu \nu }+2\pi \alpha ^{}F_{\mu \nu }+2\pi \alpha ^{}D_\mu TD_\nu T)}\right)$$
(18)
where $`G_{\mu \nu }`$, $`B_{\mu \nu }`$ and $`F_{\mu \nu }`$ are the metric, antisymmetric tensor and the non-abelian gauge field strength, respectively, and
$`D_\mu T`$ $`=`$ $`{\displaystyle \frac{T}{\sigma ^\mu }}i[A_\mu ,T].`$
This action is still incomplete without a precise prescription for how the gauge trace should be implemented. We expect that a prescription similar to that given for Born-Infeld action should also be given here. That is, the gauge trace should be completely symmetric between all non-abelian expression of the form $`F_{\mu \nu }`$, $`D_\mu T`$ and individual $`T`$ appearing in the tachyon potential $`V(T)`$.
Now we generalize (18) to the action appropriate for non-BPS D$`p`$-branes for any $`p`$. To this end, we apply familiar T-duality transformations rules to the non-abelian D9-brane action (18). T-duality transformations in $`i=p+1,\mathrm{},9`$ directions of the D9-brane world-volume converts the D9-brane to D$`p`$-brane, the gauge fields in those direction to $`\stackrel{~}{A}_i=X^i/2\pi \alpha ^{}`$ and leaves unchanged the tachyon, i.e., $`\stackrel{~}{T}=T`$. The new scalar fields $`X^i`$ are now transverse coordinates of the new D$`p`$-brane. Under this transformation, the covariant derivative of tachyon becomes
$`\stackrel{~}{D}_i\stackrel{~}{T}`$ $`=`$ $`{\displaystyle \frac{T}{\sigma ^i}}{\displaystyle \frac{i}{2\pi \alpha ^{}}}[X^i,T].`$
Using the fact that we are always working in the static gauge, the first term on the right hand side becomes zero because of the assumption implicit in the T-duality transformations that all fields must be independent of the coordinates $`\sigma ^i`$. Now adding this transformation to the T-duality transformation of massless fields (see e.g., ), one has complete list of T-duality transformations for the fields appearing in (18);
$`\stackrel{~}{E}_{ab}=E_{ab}E_{ai}E^{ij}E_{jb}`$ , $`\stackrel{~}{E}_{ai}=E_{aj}E^{ji}`$
$`\stackrel{~}{E}_{ij}=E^{ij}`$ , $`E_{ia}=E^{ij}E_{ja}`$
$`e^{2\stackrel{~}{\varphi }}=e^{2\varphi }\mathrm{det}(E^{ij})`$ , $`\stackrel{~}{D}_i\stackrel{~}{T}={\displaystyle \frac{i}{2\pi \alpha ^{}}}[X^i,T]`$ (19)
$`\stackrel{~}{F}_{ab}=F_{ab}`$ , $`\stackrel{~}{F}_{ai}={\displaystyle \frac{1}{2\pi \alpha ^{}}}D_aX^i`$
$`\stackrel{~}{F}_{ij}={\displaystyle \frac{i}{(2\pi \alpha ^{})^2}}[X^i,X^j]`$ , $`\stackrel{~}{F}_{ia}={\displaystyle \frac{1}{2\pi \alpha ^{}}}D_aX^i`$
where we have defined $`E_{\mu \nu }G_{\mu \nu }+B_{\mu \nu }`$. Here $`E^{ij}`$ denotes the inverse of $`E_{ij}`$, i.e., $`E^{ik}E_{kj}=\delta ^i_j`$. Under above T-duality transformations the determinant in (18) becomes
$`\stackrel{~}{D}=\mathrm{det}\left(\begin{array}{ccc}& \stackrel{~}{E}_{ab}+2\pi \alpha ^{}F_{ab}+2\pi \alpha ^{}D_aTD_bT& \stackrel{~}{E}_{aj}+D_aX^jiD_aT[X^j,T]\\ \\ & \stackrel{~}{E}_{ib}D_bX^ii[X^i,T]D_bT& \stackrel{~}{E}_{ij}\frac{i}{2\pi \alpha ^{}}[X^i,X^j]\frac{1}{2\pi \alpha ^{}}[X^i,T][X^j,T]\end{array}\right)`$
Manipulating the matrix inside the determinant, one finds
$`\stackrel{~}{D}`$ $`=`$ $`\mathrm{det}(P[E_{ab}+E_{ai}(Q^1\delta )^{ij}E_{jb}]+2\pi \alpha ^{}F_{ab}+T_{ab})\mathrm{det}(Q^i{}_{j}{}^{})\mathrm{det}(E^{ij})`$ (20)
where now the definition of the pull-back above is the extension of (1) in which ordinary derivative is replaced by its non-abelian covariant derivative. Here the matrices $`Q^i_j`$ and $`T_{ab}`$ are defined to be
$`Q^i_j`$ $`=`$ $`\delta ^i{}_{j}{}^{}{\displaystyle \frac{i}{2\pi \alpha ^{}}}[X^i,X^k]E_{kj}{\displaystyle \frac{1}{2\pi \alpha ^{}}}[X^i,T][X^k,T]E_{kj}`$
$`T_{ab}`$ $`=`$ $`2\pi \alpha ^{}D_aTD_bTD_aT[X^i,T](Q^1)_{ij}[X^j,T]D_bT`$
$`iE_{ai}(Q^1)^i{}_{j}{}^{}[X^j,T]D_bTiD_aT[X^i,T](Q^1)_i{}_{}{}^{j}E_{jb}^{}`$
$`iD_aX^i(Q^1)_{ij}[X^j,T]D_bTiD_aT[X^i,T](Q^1)_{ij}D_bX^j`$
In equations (20) and (4), indices are raised and lowered by $`E^{ij}`$ and $`E_{ij}`$, respectively. Now replacing (20) into T-dual of (18) and using the transformation for dilaton field (19), one finds the final T-dual action
$`\stackrel{~}{S}`$ $`=`$ $`T_p{\displaystyle d^{p+1}\sigma }`$
$`\times \mathrm{Tr}\left(e^\mathrm{\Phi }V(T)V^{}(T,X^i)\sqrt{\mathrm{det}(P[E_{ab}+E_{ai}(Q^1\delta )^{ij}E_{jb}]+2\pi \alpha ^{}F_{ab}+T_{ab})}\right)`$
where we have defined $`V^{}(T,X^i)=\sqrt{\mathrm{det}(Q^i{}_{j}{}^{})}`$. This potential term is one for abelian case. If the tachyon field is set to zero, this action would get to the result of non-abelian action for N coincident BPS D$`p`$-branes . In this case, the prescription for the gauge trace is studied in . The trace is completely symmetric between $`F_{ab}`$, $`D_aX^i`$, $`i[X^i,X^j]`$ and individual $`X^i`$. The latter non-abelian field stems from non-abelian Taylor expansion of the closed string fields that appear in the DBI action . Natural extension of this prescription for the trace in the action (4) is that the trace should be completely symmetric between all non-abelian expressions of the form $`F_{ab}`$, $`D_aX^i`$, $`i[X^i,X^j]`$, $`X^i`$, $`D_aT`$, $`i[X^i,T]`$ and individual $`T`$ of the tachyon potential.
Acknowledgments
I would like to acknowledge useful conversation with R.C. Myers. This work was supported by University of Birjand and IPM.
## Appendix A Perturbative string theory with background field
In perturbative superstring theories, to study scattering amplitude of some external string states in conformal field theory frame, one usually evaluate correlation function of their corresponding vertex operators with use of some standard conformal field theory propagators . In trivial flat background one uses an appropriate linear $`\sigma `$-model to derive the propagators and define the vertex operators. In nontrivial D-brane background the vertex operator remain unchanged while the standard propagators need some modification. Alternatively, one may use a doubling trick to convert the propagators to standard form and give the modification to the vertex operators . In this appendix we would like to consider a D-brane with constant gauge field strength / or antisymmetric Kalb-Ramond field in all directions of the D-brane. The modifications arising from the appropriate linear $`\sigma `$-model appear in the following boundary conditions <sup>3</sup><sup>3</sup>3 Our notation and conventions follow those established in . So we are working on the upper-half plane with boundary at $`y=0`$ which means $`_y`$ is normal derivative and $`_x`$ is tangent derivative. And our index conventions are that lowercase Greek indices take values in the entire ten-dimensional space-time, e.g., $`\mu ,\nu =0,1,\mathrm{},9`$; early Latin indices take values in the world-volume, e.g., $`a,b,c=0,1,\mathrm{},p`$; and middle Latin indices take values in the transverse space, e.g., $`i,j=p+1,\mathrm{},8,9`$. Finally, our conventions set $`\mathrm{}_s^2=\alpha ^{}=2`$.:
$`_yX^ai^a{}_{b}{}^{}_{x}^{}X^b=\mathrm{\hspace{0.17em}\hspace{0.17em}0}`$ $`\mathrm{for}`$ $`a,b=0,1,\mathrm{}p`$
$`X^i=\mathrm{\hspace{0.17em}\hspace{0.17em}0}`$ $`\mathrm{for}`$ $`i=p+1,\mathrm{}9`$ (23)
where $`_{ab}`$ are the constant background fields, and these equations are imposed at $`y=0`$. The world-volume (orthogonal subspace) indices are raised and lowered by $`\eta ^{ab}(N^{ij})`$ and $`\eta _{ab}(N_{ij})`$, respectively. Now we have to understand the modification of the conformal field theory propagators arising from these mixed boundary conditions. To this end consider the following general expression for propagator of $`X^\mu (z,\overline{z})`$ fields:
$`<X^\mu (z,\overline{z})X^\nu (w,\overline{w})>`$ $`=`$ $`\eta ^{\mu \nu }\mathrm{log}(zw)\eta ^{\mu \nu }\mathrm{log}(\overline{z}\overline{w})`$ (24)
$`D^{\mu \nu }\mathrm{log}(z\overline{w})D^{\nu \mu }(\overline{z}w)`$
where $`D^{\mu \nu }`$ is a constant matrix. To find this matrix, we impose the boundary condition (23) on the propagator (24), which yields
$`\eta ^{ab}D^{ba}^{ab}^a{}_{c}{}^{}D_{}^{bc}`$ $`=`$ $`0`$ (25)
for the world-volume directions, $`D^{ij}=N^{ij}`$ for the orthogonal directions, and $`D^{ia}=0`$ otherwise. Now equation (25) can be solved for $`D^{ab}`$, that is
$`D_{ab}`$ $`=`$ $`2(\eta )_{ab}^{(1)}\eta _{ab}`$ (26)
$`=`$ $`2V_{ba}\eta _{ab}`$ (27)
where matrix $`V`$ is the dual metric that appears in the expansion of DBI action (4). Note that the $`D^{\mu \nu }`$ is orthogonal matrix, i.e., $`D^\mu {}_{\alpha }{}^{}D_{}^{\nu \alpha }=\eta ^{\mu \nu }`$.
Using two dimensional equation of motion, one can write the world-sheet fields in terms of right- and left-moving components. In terms of these chiral fields, closed NSNS and open NS vertex operators are
$`V^{\mathrm{NSNS}}`$ $`=`$ $`:V_n^{\mathrm{NS}}(X(z),\psi (z),\varphi (z),p)::V_m^{\mathrm{NS}}(\stackrel{~}{X}(\overline{z}),\stackrel{~}{\psi }(\overline{z}),\stackrel{~}{\varphi }(\overline{z}),p):`$
$`V^{\mathrm{NS}}`$ $`=`$ $`:V_n^{\mathrm{NS}}(X(x)+\stackrel{~}{X}(x),\psi (x)+\stackrel{~}{\psi }(x),\varphi (x)+\stackrel{~}{\varphi }(x),k):`$
where $`\psi ^\mu `$ is super partner of world-sheet field $`X^\mu `$ and $`\varphi `$ is world-sheet superghost field. The indices $`n,m`$ refer to the superghost charge of vertex operators, and $`p`$ and $`k`$ are closed and open string momentum, respectively. In order to work with only right-moving fields, we use the following doubling trick:
$`\stackrel{~}{X}^\mu (\overline{z})D^\mu {}_{\nu }{}^{}X_{}^{\nu }(\overline{z})`$ $`\stackrel{~}{\psi }^\mu (\overline{z})D^\mu {}_{\nu }{}^{}\psi _{}^{\nu }(\overline{z})`$ $`\stackrel{~}{\varphi }(\overline{z})\varphi (\overline{z}).`$ (28)
These replacements in effect extend the right-moving fields to the entire complex plane and shift modification arising from mixed boundary condition from propagators to vertex operators. Under these replacement, world-sheet propagator between all right-moving fields take the standard form except the following boundary propagator:
$`<X^\mu (x_1)X^\nu (x_2)>`$ $`=`$ $`\eta ^{\mu \nu }\mathrm{log}(x_1x_2)+{\displaystyle \frac{i\pi }{2}}^{\mu \nu }\mathrm{\Theta }(x_1x_2)`$ (29)
where $`\mathrm{\Theta }(x_1x_2)=1(1)`$ if $`x_1>x_2(x_1<x_2)`$. Note that the orthogonal property of the $`D`$ matrix is an important ingredient for writing the propagators in the standard form. The vertex operators under transformation (28) becomes
$`V^{\mathrm{NSNS}}`$ $`=`$ $`:V_n^{\mathrm{NS}}(X(z),\psi (z),\varphi (z),p)::V_m^{\mathrm{NS}}(DX(\overline{z}),D\psi (\overline{z}),\varphi (\overline{z}),p):`$
$`V^{\mathrm{NS}}`$ $`=`$ $`:V_n^{\mathrm{NS}}(X(x)+DX(x),\psi (x)+D\psi (x),2\varphi (x),k):.`$
The vertex operator for massless NSNS and NS states and open string tachyon are
$`V^{\mathrm{NSNS}}`$ $`=`$ $`(\epsilon D)_{\mu \nu }:V_n^\mu (p,z)::V_m^\nu (pD,\overline{z}):`$
$`V^{\mathrm{NS}}`$ $`=`$ $`(\zeta 𝒢)_\mu :V_n^\mu (2kV^T,x):`$
$`V^\tau `$ $`=`$ $`:V_n(2kV^T,x):`$ (30)
where $`𝒢^{ab}=(\eta ^{ab}+D^{ab})/2=V^{ba}`$ for gauge field, $`𝒢^{ij}=(\eta ^{ij}D^{ij})/2=N^{ij}`$ for scalar field and $`𝒢^{ai}=0`$ otherwise. The open string vertex operators in $`(0)`$ and $`(1)`$ pictures are
$`V_0^\mu (k,x)`$ $`=`$ $`\left(X^\mu (x)+ik\psi (x)\psi ^\mu (x)\right)e^{ikX(x)}`$
$`V_1^\mu (k,x)`$ $`=`$ $`e^{\varphi (x)}\psi ^\mu (x)e^{ikX(x)}`$
$`V_0(k,x)`$ $`=`$ $`ik\psi (x)e^{ikX(x)}`$
$`V_1(k,x)`$ $`=`$ $`e^{\varphi (x)}e^{ikX(x)}.`$
The physical conditions for the massless open string and tachyon are
$`\mathrm{massless}:`$ $`kV_Sk=0,kV_S\zeta =0`$
$`\mathrm{tachyon}:`$ $`kV_Sk={\displaystyle \frac{1}{4}}`$
and for massless closed string are $`p^2=0`$ and $`p_\mu \epsilon ^{\mu \nu }=0`$ where $`\epsilon `$ is the closed string polarization which is traceless and symmetric(antisymmetric) for graviton(Kalb-Ramond) and
$$\epsilon ^{\mu \nu }=\frac{1}{\sqrt{8}}(\eta ^{\mu \nu }\mathrm{}^\mu p^\nu \mathrm{}^\nu p^\mu ),\mathrm{}p=1$$
(31)
for the dilaton. Using the fact that $`D^{\mu \nu }`$ is orthogonal matrix, one finds the following identities:
$$𝒢𝒢^T=𝒢^S,(D𝒢^T)^{ab}=𝒢^{ab},(D𝒢^T)^{ij}=N^{ij}$$
where the $`𝒢^S`$ is symmetric part of the $`𝒢`$ matrix.
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# Absorbing-state phase transitions in fixed-energy sandpiles
## I INTRODUCTION
Sandpile models are one of the simplest examples of avalanche dynamics, a phenomenon of growing experimental and theoretical interest. In these models, grains of “energy” (sand) are injected into the system, while open boundaries allow the system to reach a stationary state, in which energy inflow (a kind of external drive) and outflow (dissipation) balance. In the limit of infinitely small external driving, the system displays a highly fluctuating, scale-invariant avalanche-like response: the hallmark of criticality.
Ten years after the introduction of the first sandpile automaton by Bak, Tang and Wiesenfeld (BTW), our understanding of its critical behavior remains frustratingly limited, although several variants of the original model have been studied intensively . Despite some remarkable exact results , and various renormalization group analyses , the tempting possibility of assigning these models their proper universality classes remains unfulfilled. Theoretical and numerical difficulties have likewise hampered the precise estimation of critical exponents. Only recently was the upper critical dimension $`d_c=4`$ established under some assumptions for the avalanche structure .
Originally, sandpile models were proposed as the paradigm of self-organized criticality (SOC), i.e., evolution to a critical state without tuning of parameters. For this reason, sandpile models were considered for a long time to inhabit a different world than that of standard critical phenomena. Later, several authors pointed out that, in fact, the SOC state can be ascribed to the presence of two infinitely separated time scales. The two time scales correspond to the external energy input or driving, and the microscopic evolution (“avalanches”). This time-scale separation (also called slow driving), effectively tunes the system to its critical point. What is the relation between critical states due to infinite time-scale separation and regular critical points? This question stimulated many theoretical studies aimed at elucidating the links among sandpile automata and models exhibiting nonequilibrium phase transitions, such as systems with absorbing states, interfaces in disordered media, the voter model, and branching processes.
In order to make the connections with other nonequilibrium phenomena more firm, and to establish universality classes, precise critical exponent values are needed. Unfortunately, critical exponents governing the deviation from criticality cannot be measured in slowly driven sandpiles, which are posed by definition at their critical point. Thus correspondences between sandpiles and other nonequilibrium phase transitions can be only partial and inconclusive. In order to overcome this conceptual difficulty, a different approach to sandpiles has been recently pursued. It consists in analyzing sandpiles with fixed energy, that is, in considering the same microscopic rules that define sandpile dynamics, but without driving and boundary dissipation. In this way the system is closed and thus the total energy is a conserved quantity, fixed by the initial condition, and can be identified as a (temperaturelike) control parameter. The system turns out to be critical only for a particular value of the energy density (equal to that of the stationary, slowly driven sandpile) and it is thus possible to study deviations from criticality. This approach to sandpiles suggests further analogies with systems with absorbing states and interfaces in disordered media .
The stationary state of standard sandpile models is reached through the balance between the input and loss processes, identified by the energy addition and dissipation rates $`h`$ and $`ϵ`$, respectively. Critical behavior is observed in the slow driving regime, in which the parameters $`h`$ and $`ϵ`$ are tuned to their critical values ($`h0`$ and $`ϵ0`$, with $`h/ϵ0`$). In this regime, the system jumps among absorbing configurations (in which activity is null) via avalanche-like rearrangements. Evidently, in absence of external driving, any sandpile model can fall into an absorbing configuration. The connection to absorbing state phase transitions is made more clear by defining closed, fixed-energy sandpiles in which $`h0`$ and $`ϵ0`$, and periodic boundary conditions are imposed. Since the dynamics admits neither input nor loss, the total energy $`E`$ is conserved, and the energy density $`\zeta =E/L^d`$ is a tuning parameter. In this case, if the energy density $`\zeta `$ is large enough, the system reaches a stationary state with sustained activity, i.e., it is in the active phase. On the contrary, for small energy values, the system relaxes with probability one into a frozen configuration, i.e., it is in the absorbing phase. Separating these two regimes is a critical point ($`\zeta =\zeta _c`$) with marginal propagation of activity.
Once it is appreciated that fixed-energy sandpiles exhibit a continuous transition to an absorbing state, the existence of a critical stationary state in the corresponding driven dissipative sandpile is easily understood. That is because energy is added only in the absence of activity ($`\zeta <\zeta _c`$) while dissipation occurs only in the presence of activity ($`\zeta >\zeta _c`$). Thus $`d\zeta /dt`$ is positive for $`\zeta <\zeta _c`$, and vice-versa, leaving $`\zeta _c`$ as the only possible stationary value for the energy density. (The condition that dissipation and hence activity be absent in the subcritical phase makes the absorbing nature of this phase an essential ingredient of SOC.) Since SOC means tuning a system to its critical point by means of an infinitely slow drive, it is natural to try to understand the critical behavior first in the simpler context of a fixed-energy model. But while many examples of absorbing-state phase transitions have been studied in detail in recent years, we will see that characterizing sandpile criticality, even in the fixed-energy formulation, is a nontrivial project.
In this paper we define and study fixed-energy sandpiles (FES) with various microscopic dynamics. In particular we analyze the BTW sandpile , the stochastic Manna model, and a model with random mixing of a (real-valued) energy: the shuffling model (full definitions are given in the following section). We show that all of these models exhibit an absorbing-state phase transition at a critical value $`\zeta _c`$ of the energy density. What distinguishes the sandpile from other models with absorbing states is that the control parameter $`\zeta `$ represents the global value of a conserved field. This phase transition is the basis of the critical behavior of driven self-organized sandpiles. The transition is also studied using mean-field approximations, which yield good qualitative predictions for the order parameter and transition points.
Using the insights provided by the connection with absorbing states, we discuss in detail the attempt to construct a field theory for sandpiles. The latter is a generalization of Reggeon field theory (RFT), the minimal continuum theory describing absorbing-state phase transitions. We also discuss an alternative approach that considers sandpiles from the perspective of linear interface models (LIM) in disordered media. Since continuum descriptions have proved to be of fundamental importance in understanding universality and critical behavior, we analyze in detail open questions and possible improvements of these theoretical approaches.
For all the models mentioned, we report results of simulations close to the critical point, and discuss them in terms of universality classes. Numerical results indicate three distinct critical behaviors, depending upon the microscopic dynamics of models. In particular, the BTW model defines a critical behavior per se, related to the deterministic nature of the dynamics. We find striking evidence of non-ergodicity in the BTW FES: an anomalous transient to the stationary state, and lack of self-averaging. Stochastic automata, such as the Manna model, have a critical behavior that is rather close to the one of linear interface depinning models. Finally, the shuffling model shows a critical behavior that could be compatible with the RFT universality class. However, the nonlocal dynamics of this model merits a detailed examination. It is also important to note that all models show a violation of certain scaling relations usually associated with absorbing-state phase transitions. This seems to point out the particular role of the conserved field in these systems. Finally, we discuss the numerical results in the perspective of the theoretical frameworks mentioned above.
The outline of this paper is as follows: after defining the models in Sec. II, we discuss the generalized RFT theory (Sec. III) and LIM approach (Sec. IV) to FES models. We analyze from a critical perspective the approximations and hypotheses involved in these approaches. In particular, we discuss the nature of the different noise terms; this turns out to be essential to the identification of universality classes. In Sec. V we present the results of extensive simulations in two dimensions, and analyze them in the perspective of absorbing-state transitions , and the LIM mapping, which focuses on the roughness of a suitably defined interface . We find differences between BTW, Manna and fully stochastic FES exponents that persist upon enlarging the system size. Sec. VI is concerned with the origins of these differences and possible improvements in the theoretical descriptions to capture the true critical behavior of FES models. A brief summary is provided in Sec. VII. Mean-field theory approaches at the one- and two-site levels are described in the Appendix.
## II Fixed-energy sandpiles
In this paper we consider three different sandpile models. All are defined on a $`d`$-dimensional hypercubic lattice ($`d=2`$ in this study); the configuration is specified by giving the energy, $`z_i`$, at each site. The energy may take integer or real values, depending on the model, but is nonnegative in all cases. The specific models are defined as follows.
BTW model : Each active site, i.e., with (integer) energy greater than or equal to the activity threshold $`z_{th}`$ ($`z_iz_{th}=2d`$), topples at unit rate, i.e., $`z_iz_iz_{th}`$, and $`z_jz_j+1`$ at each of the $`2d`$ nearest neighbors of $`i`$. The toppling rate is introduced in order to define a Markov process with finite transition rates between configurations that differ at a small number of sites. The next site to topple is selected at random from the set of active sites; this is the only stochastic element in the dynamics. (The initial configuration is, in general, random as well.) The BTW dynamics with parallel updating (all active sites topple at each update) is completely deterministic, and it has been possible to obtain many exact results for the driven sandpile in this case, due to the Abelian property. This property implies that the order in which active sites are updated is irrelevant in the generation of the final (inactive) configuration. Accordingly, it is reasonable to expect that sequential or parallel updating does not affect the qualitative behavior. The BTW model is the prototypical sandpile model, and has been the subject of extensive numerical studies . Despite the huge numerical effort devoted to the analysis of its critical behavior, the model presents scaling anomalies which have precluded a definitive characterization. The scattered numerical values of the avalanche critical exponents were recently interpreted in terms of multiscaling properties .
Manna sandpile : In this case $`z_{th}=2`$ regardless of the number of dimensions; the energy is again integer-valued. The two particles liberated when the site $`i`$ topples move independently to randomly chosen nearest neighbors $`j`$ and $`j^{}`$ (That is, $`j=j^{}`$ with probability $`1/2d`$) . This model has a stochastic dynamics, which still enjoys a “stochastic” Abelian property, as shown recently by Dhar . The Manna model has also been the subject of many numerical studies. Together with the BTW model, it has been at the center of the long debate over universality classes for (driven) sandpiles , that we will discuss in later sections. The Manna model, fortunately, has a regular scaling behavior. The most recent analyses provide a coherent picture of its critical properties and exponent values.
Shuffling model : This model has nonnegative real-valued energies. When a site $`i`$ topples, the energy $`Z=z_i+_{jNNi}z_j`$ at that site and its nearest neighbors is redistributed randomly amongst these five sites. That is, we generate random numbers $`\eta _1,\mathrm{}\eta _5`$, uniform on , and let $`z_jz_j^{}=\eta _jZ/(\eta _1+\mathrm{}+\eta _5)`$ ($`j=1,\mathrm{},5`$). Sites with energy $`z_j^{}z_{th}=2`$ topple with probability one. In addition, the nearest neighbors of the toppling site that have energy $`z_j^{}<z_{th}`$ also become active with probability $`z_j^{}/z_{th}`$. This model contains stochasticity in each ingredient of the dynamics, and for this reason can be considered a fully stochastic model. It is clearly non-Abelian: the final configuration depends dramatically upon the order in which sites are updated. The parallel-updating version studied in this work exhibits an interesting nonlocal dynamical effect. At each update, the energy around a site is shuffled among nearest-neighbor sites. If a nearest-neighbor (or next-nearest neighbor) pair of sites are both active, the energy at a certain site or sites will be shuffled twice within a single time step. For larger aggregates of active sites, the reshuffling may involve the same site several times. In particular, energy can be transported over large distances by consecutive shuffling events along the front of active sites. This non-locality will create a mixing effect in the energy transport that one expects to influence the critical behavior.
In the present paper, we study the Manna and shuffling models with the parallel updating customarily used in sandpile automata. The BTW model is implemented using random sequential dynamics, with each active site having a toppling rate of unity. The next site to topple is chosen at random from a list of active sites, which must naturally be updated following each toppling event. The time increment associated with each such event is $`\mathrm{\Delta }t=1/N_A`$, where $`N_A`$ is the number of active sites. This is the mean waiting-time to the next event, if we were to choose sites blindly, instead of using a list. (In this way, $`N_A`$ sites topple per unit time, just as in a simultaneously updated version of the model.) Since the BTW model is Abelian, the choice of updating (parallel versus sequential) should be irrelevant to the asymptotic critical properties. This has been tested in independent simulations using parallel dynamics.
In a FES, the energy density $`\zeta `$ is fixed in the initial condition. The latter is generated by distributing $`\zeta L^d`$ particles randomly among the $`L^d`$ sites, yielding an initial (product) distribution that is spatially homogeneous and uncorrelated. Once the particles have been placed the dynamics begins. The condition to have at least one active site in the initial configuration is trivially satisfied on large lattices, for the $`\zeta `$ values of interest, i.e., close to the critical value. (For large $`L`$, the initial height at a given site is essentially a Poisson random variable, and the probability of having no active sites is exponentially decreasing with the lattice size). It is worth remarking that while the initial conditions are statistically homogeneous, the energy density is not perfectly smooth. For $`1lL`$, the energy density on a set of $`l^d`$ sites is essentially a Gaussian random variable with mean $`\zeta `$ and variance $`l^d`$. The initial value of the critical-site density $`\rho _c`$ (sites that become active upon receiving energy), moreover, is generally far from its stationary value, complicating relaxation to the steady state.
If after some time the system falls into a configuration with no active sites, the dynamics is permanently frozen, i.e., the system has reached an absorbing configuration. We shall see that as we vary $`\zeta `$, fixed-energy sandpiles show a phase transition separating an absorbing phase (in which the system always encounters an absorbing configuration), from an active phase possessing sustained activity . This is a continuous phase transition, at which the system shows critical behavior. The order parameter is the stationary average density of active sites $`\rho _a`$, which equals zero for $`\zeta <\zeta _c`$, and follows a power law, $`\rho _a(\zeta \zeta _c)^\beta `$, for $`\zeta >\zeta _c`$. The correlation length $`\xi `$ and relaxation time $`\tau `$ both diverge as $`\zeta \zeta _c`$; their critical behavior is characterized by the exponents $`\nu _{}`$ and $`\nu _{}`$, defined via $`\xi |\zeta \zeta _c|^\nu _{}`$ and $`\tau |\zeta \zeta _c|^\nu _{}`$, respectively. The dynamical critical exponent is defined via $`\tau \xi ^z`$, which implies $`z=\nu _{}/\nu _{}`$. The exponents $`\beta `$, $`\nu _{}`$ and $`\nu _{}`$ define the stationary critical behavior at the absorbing-state phase transition . In the vicinity of the critical point, where $`\xi `$ is very large, the actual characteristic length of the system is the lattice size $`L`$. We shall see that the application of finite-size scaling allows us to locate the critical point as well as estimate critical exponents.
## III Sandpiles as systems with absorbing states
In this section we discuss a recently proposed phenomenological field theory of sandpiles . Our main goal is to clarify the connection between fixed-energy sandpiles and Reggeon field theory (RFT), which is the minimal field theory describing absorbing-state phase transitions (whose prototypical examples are directed percolation (DP) and contact processes (CP)).
In Ref. we proposed a Langevin description for sandpiles by considering the mean-field description of sandpiles reported in Ref., and introducing spatial dependence and fluctuations. This allows a derivation that is based on the microscopic dynamics of sandpile automata, but involves several approximations.
Here we show how to write down a general Langevin description of sandpiles by using very general symmetry considerations. This results in a complete description, but one that is not easy to deal with, unless the proper approximations are introduced. After the introduction of some specific assumptions regarding noise terms, we recover the results of Ref. . On the other hand, the present more general treatment indicates possible modifications that may be needed for a complete characterization of sandpile models.
In sandpiles, the order parameter is $`\rho _a`$, the density of active sites (i.e., whose height $`zz_c`$) ; if at a given time $`\rho _a(𝐱)=0`$ for all x, the system has reached an absorbing configuration. The only dynamics in the model is due the field $`\rho _a(𝐱)`$, which is coupled to the local energy density, $`\zeta (𝐱,t)`$, which enhances or depresses the generation of new active sites. We therefore consider the dynamics of the local order-parameter field $`\rho _a(𝐱,t)`$ in a coarse-grained description, bearing in mind that the energy density $`\zeta (𝐱,t)`$ is a conserved field. Note that both $`\rho _a(𝐱,t)`$ and $`\zeta (𝐱,t)`$ are nonnegative. The most general dynamical equation that imposes local conservation of energy is
$$\frac{\zeta (𝐱,t)}{t}=^2(f_\zeta [\{\rho _a\},\{\zeta \}])+[g_\zeta (\{\rho _a\},\{\zeta \})\stackrel{}{\eta }(𝐱,t)],$$
(1)
where $`f_\zeta `$ and $`g_\zeta `$ are functionals of $`\rho _a`$ and $`\zeta `$. Conservation is enforced by the $`^2`$ term and the standard form of conserving noise, as for example in Cahn-Hilliard-type equations ($`\stackrel{}{\eta }`$ is a $`d`$-component vectorial noise). The dynamical equation for the density of active sites can be written analogously as
$$\frac{\rho _a(𝐱,t)}{t}=f_a(\{\rho _a\},\{\zeta \})+g_a(\{\rho _a\},\{\zeta \})\eta (𝐱,t),$$
(2)
where $`f_a`$ and $`g_a`$ are functionals of $`\rho _a`$ and $`\zeta `$ and $`\eta (𝐱,t)`$ is an uncorrelated Gaussian noise. We note that $`\eta `$ is a nonconserved noise: the active-site density is not a conserved quantity. The functionals $`f_a`$ and $`f_\zeta `$, and variances $`g_a^2`$ and $`g_\zeta ^2`$ appearing on the r.h.s. of Eqs. (1) and (2) are analytic functions (polynomials) of the local densities and (in principle) their spatial derivatives.
The right-hand-sides of Eq.s (1) and (2) must vanish when $`\rho _a=0`$ (if they did not, the state $`\rho _a=0`$ would not be absorbing!). This implies that none of the functionals $`f_a`$, $`g_a^2`$, $`f_\zeta `$, and $`g_\zeta ^2`$ contain terms independent of $`\rho _a`$; they are functions of $`\rho _a(𝐱,t)`$ and the product $`\zeta (𝐱,t)\rho _a(𝐱,t)`$. In this way activity is sustained only if $`\rho _a(𝐱,t)>0`$. It is convenient at this point to introduce a reference value $`\zeta _0`$ of $`\zeta `$ (for instance the global average energy), and expand the term $`\zeta \rho _a`$ about $`\zeta _0`$. Introducing $`\mathrm{\Delta }\zeta (𝐱,t)\zeta (𝐱,t)\zeta _0`$ we can express all the functionals as functions of $`\mathrm{\Delta }\zeta (𝐱,t)\rho _a(𝐱,t)`$, where all terms of the form $`\zeta _0[\rho _a(𝐱,t)]^n`$ are absorbed into the coefficient of $`[\rho _a(𝐱,t)]^n`$, $`\zeta _0`$ being constant.
In order to write the various functionals more explicitly, we have to consider the symmetry of the lattice in question. For isotropic models the system is inversion-symmetric under $`𝐱𝐱`$, so that odd powers of gradients, such as $`\rho _a`$, are forbidden. This leaves us with functionals such as
$$f_a(\{\rho _a\},\{\zeta \})=D_a^2\rho _a(𝐱,t)r\rho _a(𝐱,t)+\mu \rho _a(𝐱,t)\mathrm{\Delta }\zeta (𝐱,t)b\rho _a^2(𝐱,t)+....$$
(3)
where $`D_a,r,\mu `$ and $`b`$ are constants whose connection with the microscopic dynamics will be clarified below. The functionals $`f_\zeta `$, $`g_a`$ and $`g_\zeta `$ have similar forms. If we do not want to deal with an infinite set of power and derivative terms in $`\rho _a(𝐱,t)`$ and $`\mathrm{\Delta }\zeta (𝐱,t)`$, we have to identify the relevant terms from the renormalization group point of view. This can be done via power counting analysis at the upper critical dimension. This implies the knowledge of the noise term, i.e., we have to decide the terms to retain in $`g_a`$ and $`g_\zeta `$. The most relevant term is the linear one, corresponding to $`g_ag_\zeta \rho _a^{1/2}(𝐱,t)`$ . In RFT, the rationale for the noise variance being proportional to the local order parameter is that the numbers of elementary (birth and death) events in a given space-time cell are Poissonian random variables, so the variance is equal to the expected value. That the noise term for sandpile models has the same form as in RFT is by no means guaranteed. For instance, the BTW model is fully deterministic, and the nontrivial assumption that at the coarse-grained level it is described by a time-dependent noise should be tested. Further, the fact that the field $`\zeta (𝐱,t)`$ is conserved could affect the noise form. In fact, it is well known that additional symmetries on the fields can change the noise form. In the absence of an exact derivation of the noise terms, we proceed by showing the Langevin description resulting from the choice of a RFT-like noise.
Assuming RFT-like noise terms, the activity equation takes the form
$`{\displaystyle \frac{\rho _a(𝐱,t)}{t}}`$ $`=`$ $`D_a^2\rho _a(𝐱,t)r\rho _a(𝐱,t)b\rho _a^2(𝐱,t)`$ (4)
$`+`$ $`\mu \rho _a(𝐱,t)\mathrm{\Delta }\zeta (𝐱,t)+\eta _a(𝐱,t),`$ (5)
where $`\eta _a=\rho _a^{1/2}\eta `$. Here we have retained only relevant terms with respect to the noise considered. In mean-field theory the critical point corresponds to $`r=r_c=0`$; we expect fluctuations to renormalize $`r_c`$ to a nonzero value. In any case, the value of $`r`$ depends on $`\zeta _0`$, i.e., the energy density $`\zeta _0`$ plays the role of a (temperaturelike) control parameter.
The evolution of $`\mathrm{\Delta }\zeta (𝐱,t)`$ is governed only by the most relevant term in the functional $`f_\zeta `$, that is, the one linear in $`\rho _a`$. The equation may be integrated formally to yield
$$\mathrm{\Delta }\zeta (𝐱,t)=\mathrm{\Delta }\zeta (𝐱,0)+_0^t𝑑t^{}\left[D_\zeta ^2\rho _a(𝐱,t^{})+\left(\sqrt{\rho _a(𝐱,t^{})}\stackrel{}{\eta }\right)\right].$$
(6)
Substituting this into Eq. (5) and disregarding irrelevant higher order terms, the proposed Langevin equation for fixed-energy sandpiles becomes :
$`{\displaystyle \frac{\rho _a(𝐱,t)}{t}}`$ $`=`$ $`D_a^2\rho _a(𝐱,t)r(𝐱)\rho _a(𝐱,t)b\rho _a^2(𝐱,t)`$ (7)
$`+`$ $`w\rho _a(𝐱,t){\displaystyle _0^t}𝑑t^{}^2\rho _a(𝐱,t^{})+\sqrt{\rho _a}\eta (𝐱,t).`$ (8)
$`\eta `$ is a Gaussian white noise whose only non-vanishing cumulants are $`\eta (𝐱,t)\eta (𝐱^{},t^{})=D\delta (𝐱𝐱^{})\delta (tt^{})`$, $`c,b`$ and $`w`$ are fixed parameters, and the coefficient of the linear term,
$$r(𝐱)=r\mu \mathrm{\Delta }\zeta (𝐱,0),$$
(9)
inherits its spatial dependence from the initial energy distribution $`\mathrm{\Delta }\zeta (𝐱,0)`$. Observe that $`b`$ has to be positive to ensure stability; $`w>0`$ follows from the diffusion coefficient $`D_\zeta >0`$. This equation recovers the result obtained in Ref. ; we refer the reader interested in a more phenomenological approach to that paper.
We find, by standard power-counting analysis, that the upper critical dimension of this theory is $`d_c=4`$ . Above $`d_c`$, a qualitatively correct mean-field description is obtained by dropping the noise and gradient terms and replacing $`\zeta (𝐱,0)`$ by the spatially uniform $`\zeta =\zeta _0`$, yielding:
$$_t\rho _a(t)=\overline{r}\rho _a(t)\overline{b}\rho _a^2(t).$$
(10)
The critical point, $`\zeta =\zeta _c`$, corresponds to $`\overline{r}=0`$. Above $`\zeta _c`$, we have an active stationary state with $`\rho _a(\zeta \zeta _c)^\beta `$ with $`\beta =1`$; for $`\zeta <\zeta _c`$, the system falls into an absorbing configuration in which $`\rho _a=0`$. Other MF critical exponents can be calculated as well.
The present Langevin equation resembles RFT, except for the spatial dependence of $`r`$ and the non-Markovian term. Both stem from the interaction between activity with the energy background. Let us present here some comments on these two extra terms.
Quenched disorder: in the absence of the memory term, and for generic initial conditions, $`\mathrm{\Delta }\zeta (𝐱,0)const.`$, Eq. (8) is the field theory of directed percolation with quenched disorder. Disorder is known to be a relevant perturbation in DP below $`d_c=4`$ . On the other hand, the memory and spatially-dependent linear terms together represent coupling to the energy density, which is not quenched-in, but relaxes via the diffusion of activity \[see Eq. (9)\]. Thus the effect of a spatially-dependent $`r`$, in the present context, is not that of quenched disorder. In fact, we expect the physical effects of quenched disorder, and the present coupling to a conserved energy density (frozen only in the absence of activity), to be quite different. A handwaving argument to justify this assertion is the following: In the active stationary state, close to the critical point, activity will tend to be localized at any given moment, and a given point x will experience bursts of activity interspersed amongst dormant intervals. As activity alternately enters and vanishes from the neighborhood of x, the positive and negative contributions to the Laplacian memory term in Eq. (8) will largely cancel, and so this term will be dominated by the most recent changes in the state of the region. Thus the initial spatial variation in $`r(𝐱,0)`$ will effectively be forgotten in the stationary state. Another way to see this is to note that the effects of quenched disorder are found in a non-Markovian version of the contact process (using the so-called “run-time statistics”) in which the creation rate at site $`i`$ is $`\lambda _i(t)=(c_i+a)/(n_i+a+1)`$, where $`a`$ is a parameter, and $`c_i`$ represents the number of creation events out of $`n_i`$ total events at site $`i`$, up to time $`t`$. Evidently, sites which by chance have enjoyed a larger fraction of creation events in the past are likely to continue to do so, mimicking a quenched random creation rate. In the present case, the effective creation rate ($`r(x)`$) is $`\lambda (x,t)=\lambda (x,0)+w_0^t𝑑t^{}^2\rho _a(x,t^{})`$. Now, regions with $`\rho _a`$ larger than $`<\rho _a>`$ tend to have $`^2\rho _a<0`$. Thus the non-Markovian term provides a stabilizing, negative feedback on the creation rate. Regions currently experiencing above-average activity will be harder to excite in the near future. (Note however, that $`r(x,t)𝑑x`$ is time independent, since $`^2\rho _adx=0`$.) While the non-Markovian term effectively erases the initial distribution $`r(x,0)`$, we do expect the spatial dependence of $`r`$ to play an important role when we consider avalanches, i.e., the spread of activity from a localized seed, in a nonuniform energy density.
Non-Markovian term: As we have just discussed, this term enables the theory to forget the quenched, stochastic reproduction rate $`r(𝐱,0)`$. Naively, its associated coefficient, $`w`$, has the same dimensionality as $`b`$ and $`D`$, which are the two marginal parameters of RFT at its upper critical dimension, $`d_c=4`$. Below $`d_c`$ we expect the critical fixed point to be renormalized to $`r=r^{}`$, defining a renormalized $`\zeta _c`$ and nontrivial critical exponents. If the non-Markovian term is irrelevant, the field theory would be governed at criticality by the RFT fixed point. In $`d=2`$ the RFT critical behavior is characterized by $`\beta 0.58`$, $`\nu _{}0.73`$ and $`z1.77`$. We shall see in the following sections that numerical results are not compatible with this picture in the BTW and Manna case. This calls for a full RG analysis of Eq. (8). Unfortunately, this is a very dificult task because of primitive divergencies appearing in the perturbative approaches. A discussion of the RG treatment of the present field theory will be reported elsewhere.
Possible modifications and generalizations of Eq. (8), and their implications for critical behavior, will be discussed in later sections. Finally, a microscopic derivation of the field theory would ensure that the conservation symmetry has been properly taken into account in the present phenomenological approach.
## IV Sandpiles as interfaces in random media
A connection between sandpiles and interfaces moving in disordered media can be obtained by defining a variable $`H(i,t)`$ that counts the number of topplings (instances of activity) at site $`i`$ up to time $`t`$. This variable defines a growing surface in a $`d+1`$ dimensional space. The interface is said to be in the pinned phase if its disorder-average velocity $`<_tH(i,t)>`$ is null; a finite velocity marks the moving phase. It is then easy to recognize that the pinned phase in interface models is completely analogous to an absorbing state, while the moving phase corresponds to an active state . To make this correspondence more precise let us note that a nonzero interface velocity is only possible if active sites are present in the system; equivalently we can notice that $`_tH(i,t)=\rho _a(i,t)`$, so in either representation the dynamically active phase is restricted to the regime with nonvanishing $`\rho _a(x,t)`$. In this way it is evident that pinned (unpinned) and absorbing (active) states are just two ways of looking at the same physical situation. The connection between driven sandpiles and interfaces was first proposed by Narayan and Middleton and by Paczuski and Boettcher and recently generalized by Lauritsen and Alava who provided a direct mapping between the BTW model and a linear interface with quenched disorder. In the following we adapt their approach to fixed-energy sandpiles.
Let $`H_i(t)`$ be the number of topplings at site $`i`$ up to time $`t`$, and $`z_i(t)`$ the energy at $`i`$ at time $`t`$. The latter is evidently the difference between the inflow and the outflow of energy at site $`i`$ in the past. The outflow is given by $`2dH_i(t)`$, since in each toppling $`2d`$ particles are expelled from the site. There are two contributions to the inflow, the first being the energy $`z_i(0)`$ present at time $`t=0`$. The second comes from topplings of the nearest-neighbor sites, and can be expressed as $`_{NN}H_j(t)`$. Summing the above contributions we obtain:
$`z_i(t)`$ $`=`$ $`z_i(0)+{\displaystyle \underset{jNNi}{}}H_j(t)2dH_i(t)`$ (11)
$`=`$ $`z_i(0)+_D^2H_i(t),`$ (12)
where $`_D^2`$ stands for the discretized Laplacian.
Since sites with $`z_i(t)>z_c=2d1`$ topple at unit rate, the dynamics of the height follows
$$\frac{dH_i(t)}{dt}=\mathrm{\Theta }[z_i(0)+_D^2H_i(t)z_c],$$
(13)
where $`dH_i(t)/dt`$ is a shorthand notation for the rate at which the integer-valued variable $`H_i(t)`$ jumps to $`H_i(t)+1`$, and $`\mathrm{\Theta }(x)=1`$ for $`x>0`$ and is zero otherwise. Since $`z_i(t)`$ takes integer values, the smallest argument of the $`\mathrm{\Theta }`$-function yielding a nonzero toppling rate is unity. If we replace $`\mathrm{\Theta }(x)`$ by $`x`$, and assume this change to be irrelevant for critical properties , then the BTW FES is mapped onto a discretized Edward Wilkinson (EW) equation with quenched disorder, represented by the fluctuations in the $`z_i(0)`$ term. A noise term of this kind, which varies from site to site, but is time-independent, is referred to as columnar noise in the field of interface dynamics .
To understand the phenomenology of Eq. (13), let us define the average initial energy as $`f=z_i(0)`$. There are three different possibilities.
* If $`f`$ is small then with probability one the system is eventually pinned by disorder.
* If $`f`$ is large enough, the system has a finite velocity and keeps moving indefinitely.
* Separating these two regimes is a critical point marking the depinning transition.
Thus the phase transition in the BTW FES is analogous to a depinning transition. If the caveat noted above regarding the replacement $`\mathrm{\Theta }(x)x`$ turns out to be unimportant, then the transition should show the same scaling properties as depinning in the Edward-Wilkinson equation with columnar noise.
How are these results changed for the Manna model? For the outflow at site $`i`$ we now have $`2H_i(t)`$, since only two particles are transferred in each toppling event. The total input is the sum of the initial energy, $`z_i(0)`$, and a stochastic contribution $`I_i(t)`$ associated with topplings at the nearest neighbors of $`i`$:
$$I_i(t)=\underset{jNNi}{}\underset{\tau =1}{\overset{H_j(t)}{}}\eta _{i,j}(\tau ),$$
(14)
where the $`\eta _{i,j}(\tau )`$ are a set of independent (for $`i`$ fixed!), identically distributed random variables that specify the number of particles (0, 1, or 2) received by site $`i`$ at the $`\tau `$-th toppling of site $`j`$. Thus
$$\eta _{i,j}(\tau )=\{\begin{array}{cc}0\hfill & \text{with probability }(11/2d)^2\hfill \\ 1\hfill & \text{with probability }(11/2d)/d\hfill \\ 2\hfill & \text{with probability }(1/2d)^2\hfill \end{array}$$
(15)
Of course, the variables associated with different acceptor sites $`i`$ are highly correlated, since $`_i\eta _{i,j}(\tau )=2`$. $`\eta _{i,j}(\tau )`$ has mean $`1/d`$ and variance $`(11/2d)/d`$. It is convenient to introduce $`\xi _{i,j}(\tau )\eta _{i,j}(\tau )1/d`$, which has zero mean, the same variance as $`\eta _{i,j}(\tau )`$, and obeys $`_i\xi _{i,j}(\tau )=0`$. We may now write the analog of Eq. (12) for the Manna model:
$$z_i(t)=z_i(0)+\frac{1}{d}_D^2H_i(t)+\underset{jNNi}{}\underset{\tau =1}{\overset{H_j(t)}{}}\xi _{i,j}(\tau ).$$
(16)
To obtain a simple EW-like equation for the height in the Manna model, we must (1) ignore the correlations between noise terms associated with different sites, and (2) imagine that the noise is updated when site $`i`$ itself, rather than one of its neighbors, topples; we will denote the noise term as $`\xi _i(H)`$. Under these assumptions we may write
$$\frac{dH_i(t)}{dt}=\{\begin{array}{cc}1,\hfill & \text{if}z_i(0)+\frac{1}{d}_D^2H_i(t)+\xi _i(H)2\hfill \\ 0,\hfill & \text{ otherwise}.\hfill \end{array}$$
(17)
We have obtained an EW-like equation with quenched as well as columnar disorder, the so-called linear interface model. This last equation has been studied extensively both theoretically and numerically. If the previously discussed approximations are irrelevant, the Manna model should belong to the LIM universality class . The fact that the correlations between the noise terms are short range argues in favor of this conclusion .
The shuffling model deserves a particular note. In fact, it is not obvious that we can write an exact relation between the (continuous-valued) energy $`z_i(t)`$ and the height (number of topplings) $`H_i(t)`$. It is possible to write an interface equation for the shuffling model if we introduce some phenomenological constants and approximations beyond those used in the Manna case. We do not report the full derivation, which finally leads to an equation of the form of Eq. (17).
We have seen that two issues remain unresolved:
i) Whether the approximations involved in the Manna and shuffling cases change the critical behavior from the LIM universality class.
ii) Whether the various models are in the same universality class, since even if the approximations in i) are irrelevant, the Manna equation involves quenched as well as columnar noise, while only the latter appears in the BTW equation.
In order to answer the above questions analytically, a more rigorous study of the noise terms appearing in the interface equations is needed. This is analogous to the Langevin description of the previous section. We caution however that this analogy does not imply that it is easy, or even possible, to translate equations or results from one language to the other. For example, to the best of our knowledge, no one has succeeded in writing down an interface-like equation equivalent to RFT.
From a numerical point of view it is possible to measure various exponents characterizing the behavior of moving interfaces. Many of these exponents can be related to those measured in the context of absorbing-state phase transitions. It appears clear from the previous discussion that the driving force in the interface picture is equivalent to the energy density $`\zeta `$. This is the control parameter, and the exponents $`z`$ and $`\nu _{}`$ are the same in both pictures. Moreover, the order parameter exponent $`\beta `$ is equivalent to the interface velocity exponent usually measured in interface depinning models. More interestingly, associated with the interface picture are new exponents, related to the interface roughness, defined as :
$$W^2(L,t)=\frac{1}{L^d}<\underset{i}{}(H_i(t)\overline{H(t)})^2>$$
(18)
where $`\overline{H(t)}=L^d_iH_i(t)`$ and the $`<>`$ brackets represent an average over different realizations. In general one expects $`W^2`$ to exhibit an $`L`$-independent, power-law growth regime prior to saturating, that is
$$W^2(t,L)\{\begin{array}{cc}t^{2\beta _W},\hfill & tt_\times \hfill \\ L^{2\alpha },\hfill & tt_\times ,\hfill \end{array}$$
(19)
where the crossover time $`t_\times L^z`$. The limiting behaviors described above follow from the dynamic scaling property,
$$W^2(t,L)=L^{2\alpha }𝒲(t/L^z),$$
(20)
where the scaling function $`𝒲(x)x^{2\beta _W}`$ for small $`x`$, and attains a constant value for $`x\mathrm{}`$. The dynamic exponent thus satisfies the scaling relation $`z=\alpha /\beta _W`$ (first proposed by Family and Viseck ). We expect a data collapse for different system sizes in a plot of $`L^{2\alpha }W^2(t,L)`$ versus $`t/L^z`$. The roughness exponents are related via scaling relations to the other critical exponents. One may show, for example, that $`\beta _W=1\theta `$. To see this, note that in the power-law growth regime, for which the correlation length $`\xi (t)L`$, growth events in different regions are uncorrelated. Given the scaling property of the single-site height probability, $`P[H_i(t)]=f[H_i(t)/\overline{H(t)}]`$, we have $`W^2(t)=\mathrm{var}[H_i(t)][\overline{H(t)}]^2`$. Since $`\overline{H(t)}`$ is simply the integrated activity, $`\overline{H(t)}=_0^t𝑑t^{}\rho _a(t^{})t^{1\theta }`$, yielding $`\beta _W=1\theta `$.
At this point it is well to raise a caution regarding the naive application of scaling laws such as those mentioned in the preceding paragraph. Recent numerical studies have revealed that many growth models may exhibit anomalous roughening, i.e., the local width (calculated on ‘windows’ of size $`l<<L`$) scales with an exponent, $`\alpha _{loc}`$, other than $`\alpha `$. In these cases, simple scaling a la Family-Viscek does not hold. Technically this corresponds to the following situation: $`W(l,t)t^{\beta _W}_A(l/t^{1/z})`$, with an anomalous scaling function given by:
$$_A(u)\{\begin{array}{ccc}u^{\alpha _{loc}}\hfill & \mathrm{if}& u1\hfill \\ \mathrm{const}.\hfill & \mathrm{if}& u1\hfill \end{array},$$
(21)
it is only for $`\alpha _{loc}=\alpha `$ that usual self-affine scaling is recovered. This phenomenon has recently been elucidated by López (see and references therein). In general it originates from an additional correlation length, shorter than the system size, that enters as a relevant parameter in scaling equations, destroying self-affinity. In practical terms, it is important to observe that in the presence of anomalous roughening, if due attention is not paid (i.e., if scaling relations are naively assumed to hold), one can measure different correlation-time exponents depending on the type of experiment one performs. Let us finally point out that the linear interface model, at least in $`d=1`$, exhibits anomalous roughening , and therefore some of the scaling anomalies we observe could be ascribed to effects of this nature. This is an issue that certainly deserves further study.
## V Simulation results
In this section we present numerical simulations of FES models. All three FES models studied here exhibit a critical point; for large enough values of $`\zeta `$ the active site density (in the infinite-size limit) has a nonzero stationary value. In order to study the critical point and the scaling behavior of the active state in simulations of finite systems, we must study the quasistationary state that describes the statistical properties of surviving trials. The finite system size $`L`$, in fact, introduces a correlation length so that even above the critical point some initial configurations lead to an absorbing state. In practice, we compute average properties over a set of $`N_{samp}`$ independent trials, each using a different initial configuration ($`N_{samp}`$ ranges from $`10^3`$ to $`10^5`$ depending on the lattice size). Quasistationary properties are calculated from averages restricted to surviving trials. The active-site density exhibits the usual finite-size rounding in the neighborhood of the transition point; only in the limit $`L\mathrm{}`$ does the transition become sharp. For this reason, finite-size scaling is a fundamental tool in the location of the critical point as well as the calculation of critical exponents .
### A The Manna FES model
We performed simulations of the Manna fixed-energy sandpile in the version in which the two particles liberated when a site topples move independently to randomly chosen nearest neighbors. We studied lattices ranging from $`L=32`$ to $`1024`$ sites on a side, using homogeneous, random initial configurations as described in Sec. II.
After a transient whose duration depends on the system size $`L`$ and on $`\mathrm{\Delta }\zeta \zeta _c`$, the surviving sample averages reach a steady value. In Fig. (1) we show how the density of active sites approaches a mean stationary value $`\overline{\rho _a}(\mathrm{\Delta },L)`$. At a continuous transition to an absorbing state, the order parameter ($`\rho _a`$ in this instance) is expected to follow the finite-size scaling form
$$\overline{\rho _a}(\mathrm{\Delta },L)=L^{\beta /\nu _{}}(L^{1/\nu _{}}\mathrm{\Delta }),$$
(22)
where $``$ is a scaling function with $`(x)x^\beta `$ for large $`x`$, since for large enough $`L>>\xi \mathrm{\Delta }^\nu _{}`$ we must have $`\overline{\rho _a}\mathrm{\Delta }^\beta `$. To locate $`\zeta _c`$ we study the stationary active-site density as a function of system size. When $`\mathrm{\Delta }=0`$ we have that $`\overline{\rho _a}(0,L)L^{\beta /\nu _{}}`$; for $`\mathrm{\Delta }>0`$, by contrast, $`\overline{\rho _a}`$ approaches a stationary value, while for $`\mathrm{\Delta }<0`$ it falls off as $`L^d`$. Only at the critical point do we obtain a nontrivial power law, which allows us to locate the critical value $`\zeta _c`$. In Fig. 2 we observe power-law scaling for $`\zeta =0.71695`$, but clearly not for 0.7170 or 0.7169, allowing us to conclude that $`\zeta _c=0.71695(5)`$. (Figures in parenthesis denote statistical uncertainties.) The associated exponent ratio is $`\beta /\nu _{}=0.78(2)`$.
Next we consider the scaling behavior of the active-site density away from the critical point. The finite-size scaling form of Eq. (22) implies that a plot of $`\rho L^{\beta /\nu _{}}\overline{\rho _a}`$ versus $`xL^{1/\nu _{}}\mathrm{\Delta }`$ will show a data collapse for systems of different sizes. In practice, we determine the horizontal and vertical shifts (i.e., in a log-log plot of $`\rho _a`$ versus $`\mathrm{\Delta }`$) required for a data collapse. In Fig. 3, the best data collapse for $`L48`$ is obtained with $`\beta /\nu _{}=0.78(2)`$ and $`1/\nu _{}=1.22(2)`$. These values correspond to an exponent $`\beta =0.64(2)`$. This is recovered also by a direct fitting of the scaling function $`(x)`$ for large $`x`$ (see Fig. 3). A good estimate of $`\beta `$ can be also obtained by looking at the scaling of the stationary density with respect to $`\mathrm{\Delta }`$ for the largest possible sizes $`L`$. In this case if $`\mathrm{\Delta }>0`$ and $`L>>\xi `$ we have the scaling behavior $`\overline{\rho _a}\mathrm{\Delta }^\beta `$. In Fig. 4, we show the active site density as a function of $`\mathrm{\Delta }`$ for $`L=1024`$. The resulting power-law behavior yields $`\beta =0.64(1)`$, where the error is dominated by the uncertainty in the critical point $`\zeta _c`$.
To determine the dynamical exponent $`z=\nu _{||}/\nu _{}`$ we study the probability $`P(t)`$ that a trial has survived up to time $`t`$. The latter appears to decay, for long times, as $`P(t)\mathrm{exp}(t/\tau _P)`$. At the critical point, the characteristic decay time $`\tau _P`$ is a power-law function of the only characteristic length in the system, the system size $`L`$. Thus, we have $`\tau _P(L)L^z`$ for $`\mathrm{\Delta }=0`$. An estimate of $`\tau _P(L)`$ can be obtained by direct fitting of the exponential tail of $`P(t)`$, or by the time required for the survival probability to decay to one half. In Fig. 5 we report the behavior of $`\tau (L)`$ close to the critical point. Power-law behavior is recovered at the critical point, yielding $`z=1.57(4)`$. (The error bar is again dominated by the uncertainty in the critical value $`\zeta _c`$.) As a further consistency check we considered the density $`\rho _{a,all}(t,L)`$, that is, the active-site density averaged over all trials, including those that have reached the absorbing state $`\rho _a=0`$. Assuming that the time dependence involves a single characteristic time that scales as $`L^z`$, we write at the critical point $`\mathrm{\Delta }=0`$
$$\rho _{a,all}(t,L)=t^\theta g(tL^z)$$
(23)
where $`g(x)`$ is a constant for $`x1`$ and decays faster than any power law for $`x1`$. A data collapse can be obtained by plotting $`\rho _{all}=\rho _{a,all}(t,L)t^\theta `$ versus $`x=tL^z`$. The best data collapse is obtained with $`\theta =0.42(1)`$ and $`z=1.56(3)`$; it is shown in Fig. 6. This result confirms that the dynamical exponent is in the range $`z1.551.6`$. An exponent $`\theta =0.42(1)`$ is found also in the decay of the active-site density $`\rho _a(t)`$ averaged only over the surviving trials (see Fig. 1). In simple absorbing-state transitions, the latter exponent is consistent with the usual scaling relation $`\theta =\beta /\nu _{||}`$, obtained by assuming, for $`\mathrm{\Delta }=0`$, the simple scaling behavior $`\rho _a(t)=L^{\beta /\nu _{}}y(tL^z)`$, with $`y(x)=const`$ for $`x\mathrm{}`$. In the Manna FES model, this simple scaling behavior is not observed, and the relaxation of the order parameter shows qualitatively different scaling regimes. In particular, $`\rho _a(t)`$ exhibits a sharp drop (which seems to grow steeper with increasing $`L`$) just before entering the final approach to $`\overline{\rho _a}`$ (see Fig. 1). Accordingly, the exponent $`\theta `$ violates the usual scaling relation, and it is impossible to obtain a good data collapse with simple scaling forms. This is probably due to the introduction of an additional characteristic length that defines the relaxation to the quasistationary state (we are presently studying the possible relation between this effect and anomalous roughening). Moreover, it is not clear if the choice of initial conditions plays a role in this peculiar behavior. A more detailed study of the relaxation to the stationary state is required in order to understand the origin of these scaling anomalies, which appear in all the sandpile models analyzed in this paper, as well as in the one-dimensional Manna FES .
The interface mapping described in Sec. IV prompted us to study the dynamics of the mean width $`W(t,L)`$ \[see Eq.(18)\]. We studied the evolution of the width at $`\zeta _c`$, in systems of size $`L=128`$ to 800. Unfortunately, we were not able to reach the complete saturation regime of the roughness, which would afford an independent estimate of the exponent $`\alpha `$. This is due to the exponential decay of the survival probability at very large times. As shown in Fig. 7, we obtain a good collapse using the values $`\alpha =0.80(3)`$ and $`z=1.57(2)`$. Following Eq. (19), the short-time behavior of $`W(t,L)`$ gives an exponent $`\beta _W=0.51(1)`$. This exponent, however, shows a systematic increase with the system size $`L`$. In particular, for large sizes ($`L512`$) it seems that a simple power-law regime is not adequate to represent the temporal behavior of the interface width. Note also that the scaling relation $`\theta +\beta _W=1`$, satisfied to within uncertainty for the other models considered, is violated in the Manna case: $`\theta +\beta _W=0.93(2)`$. It appears that some of the anomalies affecting the temporal scaling of surviving trials could be influencing the estimates of the roughness exponents. Also in this case, further studies, for example of the local roughness, are needed for a direct comparison with other interface growth models.
In summary, numerical results show clear evidence of the critical behavior usually observed in absorbing phase transitions. Critical exponents and a discussion about universality classes will be provided in the next section. Finally, we note that the Manna sandpile does not exhibit the strong nonergodic effects reported below for the BTW model.
### B The BTW FES model
In Ref. preliminary results on the two-dimensional BTW model were reported; here we present a more detailed study, including considerably larger lattices. To study stationary properties, we performed, for each system size $`L`$ = 20, 40,…1280, and energy density $`\zeta `$, $`N_{samp}`$ independent trials (ranging from $`5\times 10^4`$ for $`L=20`$ to 1600 for $`L=1280`$), each extending up to a maximum time $`t_{max}`$. The latter, which ranged from 800 for $`L=20`$ to $`3\times 10^5`$ for $`L=1280`$, was sufficient to probe the stationary state. An overall idea of the dependence of the active-site density on $`\zeta `$ can be gotten from Fig. 8, which compares simulation results with the pair approximation derived in the Appendix.
The simulations reported in Ref. , which extended to systems of linear dimension $`L=160`$, permitted us to conclude that $`\zeta _c=2.1250(5)`$ . We first discuss the results of simulations performed at $`\zeta _c`$. Figure 9 shows the relaxation of the active- and critical-site densities at $`\zeta _c`$; note the non-monotonic approach to the limiting values. The inset shows that there is a deterministic, linear relation between the two densities during the relaxation process: for $`\zeta =\zeta _c`$, a least-squares fit yields $`\rho _c=\rho _{c,cr}C\rho _a`$, where $`C=1.368`$ and $`\rho _{c,cr}=0.4459`$ is the critical site density at $`\zeta _c`$ in the limit $`L\mathrm{}`$ (for which $`\rho _a`$ naturally falls to zero). We note that this relation is independent of system size $`L`$ and of sample-to-sample variations (for the same $`L`$); all that changes is the portion of the line filled in by the data. For off-critical values of the energy density, the active- and critical-site densities follow a different linear trend .
In Fig. 10 we plot $`\overline{\rho _a}(\zeta _c,L)`$ and the excess critical-site density $`|\overline{\rho _c}(\zeta _c,L)\zeta _{c,cr}|`$ (overbars denote mean stationary values), versus $`L`$ on log scales, anticipating that these decay $`L^{\beta /\nu _{}}`$. The apparent power-law behavior for small $`L`$ is followed, for larger $`L`$, by an approach to a larger exponent. For $`L320`$ we obtain the estimates of $`\beta /\nu _{}=0.78(3)`$ and 0.77(2) from the active-site and critical-site density, respectively, but it is clear that the slope of this plot has not stabilized even for $`L=1280`$.
Next we consider the relaxation time at $`\zeta _c`$. There are two independent quantities whose relaxation is readily monitored: the survival probability $`P(t)`$ and the active-site density $`\rho _a(t)`$. (Given the strict linear relationship between $`\rho _a`$ and $`\rho _c`$, we cannot treat the latter as an independent dynamical variable; not surprisingly, the two yield essentially the same relaxation times.) We studied four different relaxation times; the first two are associated with the survival probability $`P(t)`$. This quantity decays slowly at first, then enters a regime of roughly exponential decay, after which it attains a nearly constant value $`P_P`$. (While $`P(t)`$ appears to decay very slowly after attaining $`P_P`$, the relaxation times we study here are for the approach to $`P_P`$.) We define $`\tau _P`$ as the relaxation time associated with the exponential-decay regime; another relaxation time, $`\tau _{\overline{P}}`$, is defined as the time at which $`P(t)`$ equals $`(1+P_P)/2`$, midway to its quasi-stationary value. As we have seen, $`\rho _a(t)`$ exhibits a non-monotonic approach to its stationary value, and does not exhibit a clear exponential regime. Taking advantage of the non-monotonicity, we define $`\tau _m`$ as the time at which $`\rho _a`$ takes its minimum value. Finally, we noted that restricting the sample to trials that survive up to $`t_{max}`$ results in a monotonic, exponential approach to $`\overline{\rho _a}`$ (see Fig. 11 ). A fit to the linear portion of a semi-log plot of the excess density $`\rho _a(t)\overline{\rho _a}`$ yields $`\tau _a`$. Relaxation times in a critical system are expected to diverge with system size as $`\tau (\zeta _c,L)L^{\nu _{||}/\nu _{}}`$. The data for all four relaxation times, plotted in Fig. 12, are consistent with a power law, but due to fluctuations, linear fits to the data (for $`L160`$) yield exponent ratios ranging from $`\nu _{||}/\nu _{}=1.59`$ to 1.74. Since the four data sets do seem to follow a common trend, and since there is no reason to expect different relaxation times to be governed by different exponents, we define $`\overline{\tau }(L)`$ as the geometric mean of all four relaxation times. The behavior of $`\overline{\tau }(L)`$ is quite regular; linear fits to the data for $`L80`$, 160, and 320 yield $`\nu _{||}/\nu _{}=1.671`$, 1.668 and 1.657, respectively, leading to an estimate of 1.665(20) for this ratio.
Another manifestation of scaling is the short-time decay of the order-parameter density in a critical system, starting from a spatially homogeneous initial configuration . In Fig. 13 we show the active-site density for short-times. The data exhibit an imperfect collapse, and there is no clear-cut power-law regime. The roughly linear region for $`L=1280`$ yields a decay exponent $`\theta 0.41`$.
Next we consider the scaling behavior of the active- and critical-site densities away from the critical point. We analyze these data using the finite-size scaling form of Eq. (22), which implies that a plot of $`\stackrel{~}{\rho }L^{\beta /\nu _{}}\overline{\rho _a}`$ versus $`\stackrel{~}{\mathrm{\Delta }}L^{1/\nu _{}}\mathrm{\Delta }`$ will show a data collapse for systems of different sizes. The data analysis is as described above for the Manna FES. The best data collapse (see Fig. 14) for $`L80`$ is obtained with $`\beta /\nu _{}=0.75(2)`$ and $`1/\nu _{}=1.15(2)`$. (This value of $`\beta /\nu _{}`$ is slightly smaller than the value obtained above from the scaling of $`\rho _a`$ at $`\zeta _c`$; note however that the latter value, 0.78(3), is based on systems with $`L320`$.) ¿From this finite-size scaling analysis we therefore obtain the values $`\nu _{}=0.87(2)`$ and $`\beta =0.65(2)`$. One again, though, it is important to check for size dependence of the exponent estimates. Fitting the linear portion of the $`\rho _a`$ data in the scaling plot, we obtain $`\beta =0.62`$, 0.63, 0.66 and 0.69 for $`L=80`$, 160, 320 and 640, respectively.
We can apply a similar analysis to the density of critical sites, but here we must isolate the singular part of $`\rho _c`$ from an analytic background. The latter appears because for $`\zeta <\zeta _c`$, $`\rho _c`$ increases smoothly with $`\zeta `$. Above $`\zeta _c`$, $`\rho _c`$ decreases linearly with $`\rho _a\mathrm{\Delta }^\beta `$, so we expect the singular part $`\rho _{c,sing}=A\mathrm{\Delta }^\beta `$ for $`\mathrm{\Delta }>0`$, with $`A<0`$. The simplest reasonable form for the nonsingular background is $`\rho _{c,reg}=\rho _{c,cr}+B\mathrm{\Delta }`$, where $`\rho _{c,cr}=0.4459`$ is the $`L\mathrm{}`$ critical value as noted above. We expect the singular part of $`\rho _c`$ to follow the same finite-size scaling form as the active-site density. This implies that
$$\rho _c^{}(\stackrel{~}{\mathrm{\Delta }},L)L^{\beta /\nu _{}}(\rho _c\rho _{c,cr})=C(\stackrel{~}{\mathrm{\Delta }})+BL^{(\beta 1)/\nu _{}}\stackrel{~}{\mathrm{\Delta }}.$$
(24)
Thus the singular contributions cancel in $`\rho _c^{}(L)\rho _c^{}(L^{})`$. Using the values for $`\nu _{}`$ and $`\beta /\nu _{}`$ found in the scaling analysis of $`\rho _a`$, we study $`\rho _c^{}(L)\rho _c^{}(L^{})`$ for all pairs of system sizes in the range $`L=80,\mathrm{},640`$, and obtain $`B=0.71(2)`$. We can then construct a scaling plot of the singular part, $`\stackrel{~}{\rho }_{c,sing}L^{\beta /\nu _{}}|\rho _c\rho _{c,cr}B\mathrm{\Delta }|`$, which shows a fair data collapse (see Fig. 14), but with much more scatter than for $`\rho _a`$, presumably because of the uncertainties involved in isolating the singular contribution. As in the case of the active-site density, the $`\beta `$ estimates we obtain from the $`\rho _{c,sing}`$ data increase with $`L`$. Here we find $`\beta =0.65`$, 0.65, 0.67 and 0.70 for $`L=80`$, 160, 320 and 640, respectively. We conclude that $`\beta \stackrel{>}{}0.7`$. Studies of larger lattices will be required to refine this estimate.
We studied the evolution of the interface width $`W(t,L)`$ as defined in Eq (18), at $`\zeta _c`$, in systems of size $`L=20`$ to 640, with sample sizes ranging from $`5\times 10^4`$ for $`L=20`$ to $`10^3`$ for $`L=640`$. As shown in Fig. 15, we obtain a good collapse for $`L40`$ using the values $`\alpha =1.01(1)`$ and $`z=1.63(2)`$. The exponent $`\alpha `$ can be found directly from the data for the saturation value of $`W^2`$ shown in Fig. 16. Fitting the short-time (power-law) data for $`W^2`$ yields an estimate for the growth exponent $`\beta _W`$, which increases systematically with $`L`$, as shown in the inset of Fig. 16. Extrapolating to infinite $`L`$ we obtain $`\beta _W=0.62`$, in agreement with the scaling relation $`\beta _W=\alpha /z`$ . Note also that the value of $`z`$ describing the interface growth crossover time is consistent, as one would expect, with that for $`\nu _{||}/\nu _{}`$, derived from a study of relaxation times.
The size dependence of the critical exponents could be an indication of the failure of the simple scaling hypothesis . A further anomalous aspect of the BTW FES is nonergodicity: in a particular trial, properties such as $`\rho _a`$ typically differ from the mean value computed over a large number of trials. This means that time averages are different from averages over initial configurations, where the latter play the role of “ensemble averages”. It is worth remarking that this nonergodicity is consistent with the existence of toppling invariants . In Fig. 17, for example, we show the evolution of $`\rho _a`$ for five different initial configurations (ICs) in a system with $`L=80`$, at $`\zeta _c`$. Each IC appears to yield a particular active-site density; fluctuations about this value are quite restricted, and typically do not embrace the mean over ICs. Fig. 17 also shows histograms of the stationary mean active-site density (for a given IC), in samples of 10<sup>4</sup> ICs, for $`L=80`$ and 160; the distribution has a single, well-defined maximum, and narrows with increasing $`L`$. The data indicate, however, that the probability distribution for $`\rho _a/\overline{\rho _a}`$ (i.e., the order parameter normalized to its mean value), does not become sharp as $`L\mathrm{}`$, as it would, for example, in directed percolation.
Further evidence of nonergodicity is found in the activity autocorrelation function, defined as
$$C(t)\frac{N_A(t_0+t)N_A(t_0)}{N_A(t_0)^2}1,$$
(25)
where $`N_A(t)`$ is the number of active sites at time $`t`$, and $`\mathrm{}`$ stands for an average over times $`t_0`$ in the stationary state for a given IC, as well as an average over different ICs. The autocorrelation function for the critical BTW FES ($`L=80`$, average over 2000 ICs and $`10^4`$ time units), shown in Fig. 18, exhibits surprisingly little structure. After decaying rapidly to a minimum value at around $`t=34`$, and increasing to a weak local maximum near $`t=62`$, $`C(t)`$ seems to fluctuate randomly about zero. The relaxation occurs on a time scale over an order of magnitude smaller than for $`\rho _a`$ or the survival probability (the relaxation times $`\tau _m`$ and $`\tau _{\overline{P}}`$ $`800`$ for this system size).
The reason for this anomalously rapid decay becomes clear when we examine the autocorrelation function in individual trials ($`C(t)`$ defined as in Eq. (25) but without averaging over ICs). Figures 19 and 20 show some typical results for $`L=80`$. (Here, to get good statistics, we have averaged over $`5\times 10^5`$ to $`10^6`$ time units in the stationary state.) The correlation function in a single trial shows shows considerable structure, including damped oscillations (and in some cases, revivals), which may be superimposed on a more-or-less linear decay. The period (in the range 35 - 70 for $`L=80`$) and other features vary from one IC to another. (Changing the seed for the random choice of toppling sites changes $`C(t)`$ only slightly, if we maintain the same IC .) Evidently, $`C(t)`$ decays rapidly to zero when we average over initial conditions because of dephasing amongst oscillatory signals with varied frequencies. Interestingly, the interface width $`W(t,L)`$ shows much less dependence on the IC than does the active-site density or its autocorrelation.
In summary, the BTW fixed-energy sandpile shows signs of the kind of scaling found at simpler absorbing-state phase transitions, but at the same time exhibits dramatic nonergodic effects. We note unusually strong finite-size effects, which prevent us from determining certain critical exponents precisely. Whether this is a simple finite-size effect or a signature of multiscaling cannot be ascertained definitively with the present data.
### C The Shuffling FES model
The shuffling model has a continuously variable control parameter, since each site has a (non-negative) real-valued energy. Thus we are no longer constrained to vary the energy density $`\zeta `$ in increments of $`1/L^2`$ as we are in discrete models (e.g., the Manna and BTW FES), where the single grain is the smallest energy unit. In the shuffling FES, all sites whose energy exceeds the threshold $`z_{th}=2`$ are considered active. In addition, sites that have received energy from a toppling nearest neighbor can become active if $`z_i<z_{th}`$ with a probability $`p_i=z_i/z_{th}`$. This enlarges considerably the choice of possible initial configurations. In particular, after we have distributed randomly the total amount of energy among the lattice sites, we extract for each site a random number $`\eta _i`$ and we declare active all sites for which $`\eta _iz_i/z_{th}`$. (Obviously, sites with $`z_iz_{th}`$ are active with probability one.) Unlike discrete models, we have the option of generating “flat” initial conditions, in which all sites have the same energy. While stationary properties are not affected by the choice of noisy versus flat initial configurations, we do note differences in the short-time behavior.
Another peculiar characteristic of the shuffling model is the strong non-Abelian character of its dynamics. We implemented the dynamics of the model with parallel updating as in the original definition of Ref. . However, this form of the dynamics contains some non-local effects as described in Sec. II, and does not ensure that parallel and sequential updating generate the same critical behavior. Simulations with sequential updating are in progress.
Simulations of the shuffling model require many calls to the random number generator, and so are extremely time-consuming. Here we present simulations with flat initial conditions and sizes ranging from $`L=32`$ to $`L=384`$. By analyzing the $`L`$-dependence of $`\overline{\rho _a}(\mathrm{\Delta },L)`$ we find the critical point $`\zeta _c=0.20427(5)`$. When $`\zeta =\zeta _c`$ the stationary density has a power-law behavior $`\overline{\rho _a}(0,L)L^{\beta /\nu _{}}`$ that yields $`\beta /\nu _{}=0.76(3)`$. This result is confirmed by the scaling plot of Fig. 21, which, following Eq. (22) shows $`\rho L^{\beta /\nu _{}}\overline{\rho _a}`$ versus $`xL^{1/\nu _{}}\mathrm{\Delta }`$, with $`\beta /\nu _{}=0.76`$ and $`1/\nu _{}=1.266`$. This gives an exponent $`\beta =0.60`$, as confirmed by the straight slope of the upper branch of the scaling plot. An independent measurement of the stationary density versus $`\mathrm{\Delta }`$ for the largest size used ($`L=384`$) gives the estimate $`\beta =0.60(2)`$, where the error bar is due mainly to the uncertainty in $`\zeta _c`$.
We performed a scaling analysis of the temporal behavior by studying the decay of the survival probability $`P(t)\mathrm{exp}(t/\tau _P)`$. At the critical point the $`L`$-dependence of the characteristic time assumes the power-law behavior $`\tau _PL^z`$ with $`z=1.71(5)`$ (see Fig. 22). However, it is worth noting that the scaling behavior with $`L`$ shows a systematic curvature from the smallest to the largest sizes, both below and above the critical point. This could be a signal that the system has not yet reached its asymptotic temporal behavior for the sizes considered ($`L320`$). That the relaxation could be affected by strong finite-size effects is confirmed by the temporal scaling of $`\rho _a(t,L)`$. In Fig. 23 we observe that the active-site density decay does not follow a definite power law before reaching the stationary state. This makes impossible an accurate determination of the exponent $`\theta `$ ($`0.46`$), which is also reflected in the absence of a clear data collapse for the temporal scaling functions.
The roughness analysis is affected by several numerical problems. The short average lifetime of trials at finite size makes it impossible to reach the width-saturation regime. This effect is even more pronounced than in the Manna case. It is therefore impossible to apply a data-collapse analysis, nor is a direct measurement, that would yield $`\alpha `$, feasible. The short-time behavior of the roughness (see Eq. (18)) is governed by the exponent $`\beta _W0.57`$. Applying the scaling relation shown in Sec. IV, and using the dynamical exponent obtained previously, we have $`\alpha 0.96`$. However, in this case the short-time behavior of the roughness appears to have a size dependence, probably due to the lack of complete convergence to the asymptotic scaling behavior, and the numerical values provided here could contain systematic errors that are difficult to estimate.
In summary, the numerical results for the shuffling FES model show also the signature of a continuous phase transition from an absorbing to an active phase. The stationary properties of the model show well defined scaling behavior at the system sizes considered in the present study. The dynamic scaling properties, by contrast, show anomalies and transient effects that could indicate that the system has not yet attained its asymptotic behavior for $`L384`$.
## VI Discussion and open questions
### A Universality classes and critical exponents
Simulations of sandpile models have mainly been performed in the slow driving regime. It is then natural to compare the critical exponents measured in the fixed-energy framework (see Tab. I) with those observed in driven simulations. In driven sandpiles, critical behavior is characterized by the scaling of the number of topplings $`s`$ and the duration $`t`$ following the addition of an energy grain, i.e., an avalanche. The probability distributions of these variables are usually described with the finite-size scaling forms
$$P(s)=s^{\tau _s}𝒢(s/s_c)$$
(26)
$$P(t)=t^{\tau _t}(t/t_c)$$
(27)
where $`s_cL^D`$ and $`t_cL^z`$ are the characteristic avalanche size and time, respectively. Applying the fundamental result (due to conservation), $`<s>L^2`$, we can write the scaling relations $`\tau _s=22/D`$ and $`\tau _t=1+(D2)/z`$. Recently, these simple scaling forms have been questioned in the case of the BTW model. An accurate moment analysis seems to show multiscaling, so that scaling relations obtained from the above finite-size scaling forms do not apply.
While critical exponents governing the deviations from criticality in FES do not have any counterpart in the driven case, which is posed by definition at the critical point, the exponents describing the critical point, including $`z`$ and the fractal dimension $`D`$, can be compared directly. In FES simulations $`D`$ can be calculated by noting that the scaling of an avalanche due to a point seed scales as the total variation of the field $`H(i,t)`$, which represents the total number of topplings. Since the roughness scales with exponent $`\alpha `$, we readily obtain that $`D=d+\alpha `$.
For the Manna model, our simulations yield $`D=2.80(3)`$ and $`z=1.57(4)`$, which should be compared with the most recent analyses of driven sandpiles, which yield $`D=2.76(2)`$ and $`z=1.56(2)`$. By using scaling relations we obtain $`\tau _s1.29`$ and $`\tau _t1.51`$, again in very good agreement with the values obtained in the driven case. For the shuffling model we can compare our results $`z=1.71`$ and $`D=2.96`$ with the simulations of Maslov and Zhang, which give $`z=1.73(5)`$ and $`D=2.92(5)`$. In this case we also see a very good agreement between independent measurements.
More subtle is the case of the BTW model. Here different simulations of the driven sandpile give rather scattered results. A very recent analysis suggesting multiscaling in the (driven) BTW sandpile indicates that neither $`D`$ nor $`z`$ are clearly defined. In particular, the effective value of $`D`$ increases as one studies higher moments, and saturates at $`D3.0`$. This is indeed the result we recover from our analysis ($`D=3.01(1)`$). The possibility of multiscaling is supported by the scaling anomalies and the lack of self-averaging we detected in our simulations of the BTW FES.
We shall attempt, on the basis of our numerical results, to assign the various fixed-energy sandpiles studied to universality classes. This a particularly vexing problem, that has eluded ten years of theoretical and numerical efforts. Soon after the introduction of sandpile models with modified dynamical rules, there were many quests for the precise identification of universality classes. In particular BTW and Manna models, which are prototypes for deterministic and stochastic models, respectively, have been the objects of a longstanding quarrel over their supposed universality classes. The first numerical attempts showed very similar exponents for the avalanche distributions, but the results were afflicted by severe finite-size errors due to the limited sizes attainable using the CPU power available at that time. These results were later questioned by Ben-Hur and Biham, who analyzed the scaling of conditional expectation values of various quantities related to avalanches. These results were, however, biased by the unexpected singular behavior of the distributions, and have been recently reconsidered by applying other numerical methods . From the theoretical standpoint it is very surprising that small modifications of the microscopic dynamics would lead to different universality class. However, no analytical demonstration of distinct universality classes in sandpiles has been presented up to now. On the contrary, many theoretical arguments in favor of a single universality class can be found in the literature.
In Table I we summarize the critical exponents found for each model. The quoted values indicate, beyond numerical uncertainties, that the models discussed here belong to three distinct universality classes. Striking differences appear between the BTW and the Manna model. Beyond the numerical values of critical exponents, we observe for the first time the lack of self-averaging in the BTW FES. This property is related to its deterministic dynamics, and finds consistent analogies in the waves of toppling description. The lack of self-averaging could also be the origin of the multiscaling features recently observed by De Menech et al. in the driven BTW sandpile. From this discussion it appears that the introduction of stochasticity is a relevant modification for the critical behavior. At this point it is worth noting that the Manna model has been considered for a long time as a non-Abelian model. The opposite has been pointed out recently by Dhar, by means of rigorous arguments. The conjecture that Manna and BTW sandpiles belong to different universality classes because the former is non-Abelian has then to be abandoned. Stochasticity per se, however, does not define a unique universality class, as evidenced by the distinct critical properties of the Manna and shuffling FES models. The origin of the different behavior can be traced to the nonlocal nature of the shuffling model dynamics, as we shall make clear later.
In summary, our numerical results are in good agreement with the most recent measurements of driven sandpiles, confirming that the two cases share the same critical behavior. In addition, the FES framework enlarges the set of exponents that can be measured, providing new tools for the characterization of critical behavior and universality classes in different models.
### B Avalanche and spreading exponents
In order to compare the exponents found in fixed-energy simulations with the usual avalanche exponents $`\tau _s`$ and $`\tau _t`$, we relied on scaling relations. However, avalanches can also be studied in the FES case, in simulations of critical “spreading”. Let us first define what is a spreading experiment in a system with an absorbing-state. In such experiments, a small perturbation (a single active site, for instance) is created in an otherwise frozen (absorbing) configuration. In the supercritical regime, the ensuing activity has a finite probability to survive indefinitely, reaching the stationary state deep inside the (growing) active region. In the subcritical regime, activity will decay exponentially. In each spreading sequence, it is customary to measure the spatially integrated activity $`N(t)`$, averaged over all runs, and the survival probability $`P(t)`$ after $`t`$ time steps. At the critical point separating the supercritical and subcritical regimes, these quantities have a singular scaling: $`N(t)t^\eta `$ and $`P(t)t^\delta `$, where $`\eta `$ and $`\delta `$ are called spreading exponents. If we can define the substrate over which the activity spreads uniquely, this spread of activity is the same as an avalanche in a sandpile model.
Sandpile models, however, have infinitely many absorbing configurations. In the infinite-size limit, an infinite number of energy landscapes correspond to the same value $`\zeta `$. (For real-valued energies, as in the shuffling model, this infinite degeneracy already appears for finite systems.) In this case spreading properties at a given value of the control parameter $`\zeta `$ will depend on the initial configuration in which the system is prepared. It is even possible to observe nonuniversality in the spreading exponents, a feature that sandpiles share with the pair contact process (PCP) and other systems with infinitely many absorbing configurations.
In order to have well defined spreading exponents (that can be related to the avalanche exponents of a driven sandpile), we have to define uniquely the properties of the energy landscape for spreading experiments. One possibility is to use the absorbing configurations generated by the fixed-energy sandpile itself for initial configurations. Suppose we use such a configuration for a spreading experiment, by introducing an active site. Repeating this process many times, we obtain the spreading properties for so-called “natural absorbing configurations”. A second option is to use the substrate left by each spreading process as the initial condition for the subsequent one. After a transient time the system will flow to a stationary state with well defined properties, in which each initial configuration is a “natural configuration”. On the other hand, this second definition of a spreading experiment is identical to slow driving, except that energy is strictly conserved (the active site must be generated by a mechanism that does not change the energy).
By performing spreading experiments close to $`\zeta _c`$, it is possible to obtain directly the avalanche and spreading scaling behavior, as well as the divergence of characteristic lengths approaching the critical energy. A preliminary study in this direction for the BTW model confirms the uniqueness of the critical behavior at $`\zeta _c`$. Interesting results have also been obtained for the spreading properties in a FES mean field model. A more complete study of spreading exponents in a variety of sandpile models is a promising path toward the complete characterization of their critical behavior.
### C Comparison with theoretical results
In earlier sections we presented two alternative theoretical descriptions for sandpile models. We compare our numerical results with theoretical predictions in order to assess the validity of these theoretical frameworks, and the eventual improvements needed for a complete description of sandpile models.
In Sec. III we introduced a Langevin description that takes into account the absorbing nature of the phase transition in FES models. Unfortunately, a rigorous derivation of the noise terms has not yet been made. The assumption of RFT-like noise terms leads to the Langevin description of Eq. (8). This differs from the standard DP Langevin description for the presence of a non-Markovian term. Only in the case that this term is irrelevant the theory belongs to the universality class of RFT. From a physical point of view this means that the local coupling between the activity field $`\rho _a(𝐱,t)`$ and the energy field $`\zeta (𝐱,t)`$ is irrelevant on large scales. In other words the activity spreads on an effective average energy substrate whose only role is to tune the spreading probability. This is indeed the same as a DP problem in which the critical parameter is tuned via the average energy $`\zeta `$.
Casting a glance at our numerical results, the only model that has exponents compatible with the DP universality class is the shuffling FES. This is not unexpected; the model was indeed proposed by Maslov and Zhang as a sandpile realization of directed percolation. At the basis of this behavior is nonlocal energy transport. As we emphasized in Sec. II, the shuffling model allows the transfer of the same parcel of energy several times in the same time step. This introduces, on average, a strong mixing effect that makes energy diffusion slower. In this way the spread of activity is effectively decoupled from the local fluctuations that the activity itself generates in the energy field. On the other hand, Maslov and Zhang noted that in $`d=1`$, the nonlocal energy mixing is not capable of destroying correlations and, following a transient, the model exhibits non-DP scaling. While the exponents summarized in Tab. I are compatible with the DP universality class, we note that the dynamic scaling properties of the shuffling model show systematic biases that could signal a nonasymptotic behavior for some observables. We cannot therefore exclude completely that the model is still in a transient regime, that could finally lead to a different critical behavior, as happens in $`d=1`$.
The Manna and BTW FES models, by contrast, exhibit critical exponents different from those of DP. In these models, the energy redistribution during toppling is strictly local, and the spread of activity is always correlated with the energy fluctuations generated during toppling processes. It is then reasonable to expect that a Langevin theory has to take into account fully the non-Markovian term. It may be also possible to derive the pertinent stochastic equations and the noise correlations applying more rigorous treatments, as in Ref .
The moving interface picture is also afflicted by our ignorance of the correlations between the quenched noise terms appearing in the equations (see Sec. IV). By suitable approximations it has been shown that the Manna model could belong to the LIM universality class. Our numerical results show that the stationary critical properties are compatible with this universality class. The dynamic properties, however, show anomalies that are not compatible with LIM. The origin of these anomalies deserves a more accurate analysis, and might be understood if we had a better grasp of the noise terms in the interface representation. It is interesting, in this context, that the BTW model, for which the mapping to the interface representation seems most straightforward, defines a universality class per se, incompatible with linear interface depinning with columnar disorder. This is probably due to the strong nonlinearity introduced by the local velocity constraint implicit in the $`\mathrm{\Theta }`$-function of Eq. (13).
While neither theoretical approach allows an exact characterization of sandpile models, they appear to be conceptually very relevant, because they provide an answer to the basic questions of why driven sandpile models show SOC. The genesis of self-organized criticality in sandpiles is a continuous absorbing-state phase transition. The sandpile exhibiting the latter may be continuous or discrete, deterministic or stochastic. To transform the conventional nonequilibrium phase transition to SOC, we couple the local dynamics of the sandpile to a “drive” (a source with rate $`h`$). The relevant parameter(s) {$`\zeta `$} associated with the phase transition are controlled by the drive, in a way that does not make explicit reference to {$`\zeta `$}. Such a transformation involves slow driving ($`h0`$), in which the interaction with the environment is contingent on the presence or absence of activity in the system (linked to {$`\zeta `$} via the absorbing-state phase transition). Viewed in this light, “self-organized criticality” refers neither to spontaneous or parameter-free criticality, nor to self-tuning. It becomes, rather, a useful concept for describing systems that, in isolation, would manifest a phase transition between active and frozen regimes, and that are in fact driven slowly from outside.
A second class of theoretical questions concern the critical behavior (exponents, scaling functions, power-spectra, etc.) of specific models, and whether these can be grouped into universality classes, as for conventional phase transitions both in and out of equilibrium. In this respect, the theoretical approaches presented here show a very promising path of improvements and modifications that could lead to the solution of many of these questions.
## VII Summary
We studied three fixed-energy sandpile models, whose local dynamics are those of the BTW, Manna, and shuffling sandpiles, studied heretofore under external driving. The former two models are Abelian, the latter two stochastic. The results of extensive simulations, which are in good agreement (via scaling laws), with previous studies of driven sandpiles, place the three models in distinct universality classes. Results for the Manna FES are consistent with the universality class of linear interface depinning, while the shuffling FES appears to follow directed percolation scaling. Both these assignments, however, are somewhat provisional, due to dynamic anomalies and apparent strong finite-size effects. The case of the BTW FES, which appears to define a new universality class, is further complicated by violations of simple scaling and lack of ergodicity. Examining the field-theoretic and interface-height descriptions of sandpiles in light of our simulation results, we find that a more rigorous description of noise correlations will be required, for these approaches to become reliable predictive tools. Our results strongly suggest that there are at least three distinct universality classes for sandpiles. Whether others can be identified, and how the various classes can be accommodated in a unified field-theoretic description, are challenging issues for future study.
ACKNOWLEDGEMENTS We thank M. Alava and R. Pastor-Satorras for the many results on SOC they have discussed and shared with us prior to publication. We are also indebted with P. Grassberger for comments and private communications. We also acknowledge A. Barrat, A. Chessa, D. Dhar, K.B. Lauritsen, E.Marinari, L. Pietronero and A. Stella for very useful discussions and comments. M.A.M., A.V. and S.Z. acknowledge partial support from the European Network contract ERBFMRXCT980183; M.A.M acknowledges also partial support from the Spanish DGESIC project PB97-0842, and Junta de Andalucía project FQM-165. R.D. acknowledges CNPq and CAPES for support of computing facilities.
Appendix: Mean-Field description
We have devised mean-field approximations for fixed-energy sandpiles at the one- and two-site levels. While the present mean-field theory has nothing useful to say about critical behavior, it is nonetheless interesting that a simple analysis can yield reasonable predictions for the order parameter and transition points. Consider first the one-site approximation for the BTW FES. Let $`\rho _z`$ be the density of sites with energy $`z`$. Each site receives a unit of energy at rate $`n_a`$, the number of active neighbors. In the one-site approximation the sites are treated as statistically independent, so that in a homogeneous system, the rate of arrival of particles at any site is $`2d\rho _a`$ where $`\rho _a_{z2d}\rho _z`$ is the density of active sites. In addition to receiving energy, sites with $`zz_{th}=2d`$ make a transition to $`zz_{th}`$ at unit rate. Hence the mean-field equations are
$$\frac{d\rho _z}{dt}=2d\rho _a(\rho _{z1}\rho _z)+\rho _{z+2d}\theta _{z2d}\rho _z,$$
(28)
where $`\theta _j=0`$ for $`j<0`$ and is unity otherwise, and $`\rho _1`$ (in the equation for $`z=0`$) is of course zero.
This set of equations satisfies probability conservation ($`_z\rho _z`$ is constant), and conserves the mean energy $`\zeta _zz\rho _z`$. We try to find a stationary solution by introducing the simplifying assumption that for $`z2d`$, the distribution follows an exponential decay:
$$\rho _z=\alpha ^{z2d}\rho _{2d},z2d.$$
(29)
Under this assumption, the active-site density $`\rho _a=\rho _2/(1\alpha )`$ in one dimension. $`\rho _0`$ can be eliminated using normalization:
$$\rho _0=1\rho _1\frac{\rho _2}{1\alpha }.$$
(30)
Then the mean-field equations become
$$\frac{d\rho _1}{dt}=\rho _2\left[\alpha +\frac{2}{1\alpha }\left(12\rho _1\frac{\rho _2}{1\alpha }\right)\right]$$
(31)
$$\frac{d\rho _2}{dt}=\rho _2\left[\alpha ^21+\frac{2}{1\alpha }(\rho _1\rho _2)\right]$$
(32)
and
$$\frac{d\rho _z}{dt}=\rho _2\left[\alpha ^z\alpha ^{z2}+\frac{2\rho _2}{1\alpha }(\alpha ^{z3}\alpha ^{z2})\right],z3.$$
(33)
The last equation implies that in the stationary state $`\rho _2=\frac{\alpha }{2}(1\alpha ^2)`$, and therefore $`\rho _a=\frac{\alpha }{2}(1+\alpha )`$. From Eq. (32) we then have $`\rho _1=\frac{1}{2}(1\alpha ^2)`$ in the stationary state. Thus $`\rho _2=\alpha \rho _1`$ and the distribution is exponential starting with $`\rho _1`$. One readily verifies that the r.h.s. of Eq. (31) is also zero for the stationary values of $`\rho _1`$ and $`\rho _2`$ given above.
The mean energy is given by:
$$\zeta =\rho _1+\rho _2\underset{n=0}{\overset{\mathrm{}}{}}(n+2)\alpha ^n=\frac{1}{2}\frac{1+\alpha }{1\alpha },$$
(34)
so that
$$\alpha =\frac{2\zeta 1}{2\zeta +1},$$
(35)
which shows that the active stationary state exists only for $`\zeta >\zeta _c=1/2`$. Below this value, $`\rho _a=0`$ and $`\rho _1=\zeta `$. (We have also verified that this solution is stable for $`\zeta >1/2`$.) In the active stationary phase,
$$\rho _a=\frac{4\zeta }{(2\zeta +1)^2}(\zeta \zeta _c),$$
(36)
so the order parameter exponent $`\beta =1`$, the usual mean-field value for systems lacking “up-down” symmetry.
The two-dimensional case is only slightly more complicated; the mft equations may be written so:
$$\frac{d\rho _z}{dt}=4\rho _a(\rho _{z1}\rho _z)\theta _{z4}\rho _z+\rho _{z+4},$$
(37)
where $`\rho _a_{z4}\rho _z`$. We now suppose that in the active stationary state, $`\rho _z=\alpha ^{z4}\rho _4`$ for $`z4`$. Then $`\rho _a=\rho _4/(1\alpha )`$, and proceeding as in the one-dimensional case, one finds the stationary solution:
$$\rho _z=\frac{1}{4}(1\alpha ^{z+1}),z2,$$
(38)
and
$$\rho _z=\frac{\alpha ^{z3}}{4}(1\alpha ^4),z3.$$
(39)
The mean energy is
$$\zeta =\frac{3\alpha }{2(1\alpha )},$$
(40)
so that
$$\alpha =\frac{2\zeta 3}{2\zeta 1},$$
(41)
showing that $`\zeta _c=3/2`$ for the BTW sandpile in 2-d, in this approximation.
It is also possible to derive two-site mean-field equations without much difficulty. Denote the probability for a NN pair of sites to have energies $`i`$ and $`j`$ by $`\rho _{i,j}`$. The gain term, or rate of transitions into the state $`(i,j)`$, due to one of the sites toppling or gaining a unit of energy from a neighbor, is
$$\frac{d\rho _{i,j}^+}{dt}=\rho _{i1,j+2d}+\rho _{j+2d,j1}+(2d1)\left[\frac{\rho _{i,j1}\rho _{j1,a}}{\rho _{j1}}+\frac{\rho _{a,i1}\rho _{i1,j}}{\rho _{i1}}\right],$$
(42)
where $`\rho _{i,a}_{j2d}\rho (i,j)`$, and we have used the fact that in the pair approximation (i) the $`2d`$ neighbors of a given site are mutually independent, and (ii) if $`l`$ and $`n`$ are neighbors of site $`m`$, then the three-site probability $`P(z_l,z_m,z_n)=\rho _{z_l,z_m}\rho _{z_m,z_n}/\rho _{z_m}`$. (The one-site probabilities are given by $`\rho _i=_j\rho _{i,j}`$.) Note that there is no gain term for $`i=j=0`$: we expect no such pairs in the active stationary state. Similarly, the loss term is
$$\frac{d\rho _{i,j}^{}}{dt}=\rho _{i,j}\left[\theta _{i2d}+\theta _{j2d}+(2d1)\left(\frac{\rho _{j,a}}{\rho _j}+\frac{\rho _{a,i}}{\rho _i}\right)\right].$$
(43)
The two-site probabilities are governed by
$$\frac{d\rho _{i,j}}{dt}=(1+\delta _{i,j})\left[\frac{d\rho _{i,j}^+}{dt}\frac{d\rho _{i,j}^{}}{dt}\right].$$
(44)
In the absence of a simple ansatz for the solution of these equations, we analyze them numerically. For this purpose we choose a cutoff and set $`\rho (i,j)0`$ for $`i`$ or $`j>n`$, where $`n`$ is sufficiently large (in the range 20 - 36, depending on $`\zeta `$), that $`\rho _j`$ is completely negligible for $`jn`$. The coupled equations are integrated using a fourth-order Runge-Kutta routine, starting from a product Poisson distribution, $`\rho _{i,j}=\rho _i^0\rho _j^0`$, with $`\rho _j^0=e^\zeta \zeta ^j/j!`$ (We verified that the location of transition points does not depend on the form of the initial distribution.)
The pair mean-field equations for the one-dimensional BTW model predict a first order transition at $`\zeta _c=0.91652`$. The active site density jumps from zero to about 0.052 at this point. Simulations also show a first-order transition, but at $`\zeta _c=1`$, with the order parameter jumping to about 0.14. The energy distribution predicted by pair mft is approximately exponential for $`z3`$ or so.
In two dimensions, the pair approximation yields $`\zeta _c=1.98059`$, to be compared with the exact value of 2.125… The transition is again discontinuous, but the jump in $`\rho _a`$ (from zero to about 0.0061), is very small. (We find no evidence of a discontinuous transition in simulations.) At the critical point, pair mft predicts $`\rho _c=\rho _3=0.3328`$, while simulation yields $`\rho _c=0.434`$. The energy distribution decays exponentially for $`z7`$ or so. Pair approximation and simulation results for the order parameter are compared in Fig. 8.
The pair MFT is readily extended to the Manna FES defined in Sec. II. In one dimension, when site i topples, the two particles are both sent to i-1 with probability 1/4 (similarly for site i+1), and with probability 1/2, one each is sent to i-1 and i+1. In the Manna sandpile some new transitions, not allowed in the BTW model, make their appearance. Enumerating the possibilities as above, one obtains, for the one-dimensional exclusive Manna sandpile, the equations:
$`{\displaystyle \frac{d\rho _{\alpha \beta }}{dt}}=`$ $`{\displaystyle \frac{1}{2}}\left[\rho _{\alpha 1,\beta +2}+\rho _{\alpha +2,\beta 1}+{\displaystyle \frac{\rho _{\alpha ,\beta 1}\rho _{\beta 1,a}}{\rho _{\beta 1}}}+{\displaystyle \frac{\rho _{a,\alpha 1}\rho _{\alpha 1,\beta }}{\rho _{\alpha 1}}}\right]`$ (45)
$`+`$ $`{\displaystyle \frac{1}{4}}\left[\rho _{\alpha +2,\beta }+\rho _{\alpha +2,\beta 2}+\rho _{\alpha ,\beta +2}+\rho _{\alpha 2,\beta +2}+{\displaystyle \frac{\rho _{\alpha ,\beta 2}\rho _{\beta 2,a}}{\rho _{\beta 2}}}+{\displaystyle \frac{\rho _{a,\alpha 2}\rho _{\alpha 2,\beta }}{\rho _{\alpha 2}}}\right]`$ (46)
$``$ $`\rho _{\alpha ,\beta }\left[\theta _{\alpha 2}+\theta _{\beta 2}+{\displaystyle \frac{3}{4}}\left({\displaystyle \frac{\rho _{\beta ,a}}{\rho _\beta }}+{\displaystyle \frac{\rho _{a,\alpha }}{\rho _\alpha }}\right)\right].`$ (47)
These equations predict a continuous transition at $`\zeta _c=0.7500`$, in fair agreement with simulation ($`\zeta _c0.949`$ ). A straightforward generalization to two dimensions yields a continuous transition at $`\zeta _c=0.625`$, about 13% smaller than the value found in simulations ($`\zeta _c=0.7169`$).
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# One Interesting New Sum Rule Extending Bjorken’s to order 𝟏/𝒎_𝑸
## I Introduction
It is well known that quark-hadron duality is valid to a good accuracy in $`b`$-quark decay and particularly in semileptonic decay. A systematic study of the corrections to duality using the powerful tools of Operator Product Expansion (OPE) and Heavy Quark Effective Theory (HQET), in particular Luke’s theorem , has demonstrated that the first corrections to duality only appear at second order, namely $`O(\mathrm{\Lambda }^2/m_Q^2)`$ where $`\mathrm{\Lambda }`$ is for the QCD scale and $`m_Q`$ is one of the heavy quark masses ($`m_b`$ or $`m_c`$). For simplicity we leave aside in this letter the $`O(\alpha _s)`$ radiative corrections notwithstanding their manifest practical relevance.
The OPE based proof is very elegant and circumvents the detailed calculation of the relevant channels. Precisely this feature has generated some doubts or at least some worries. First of all there is the experimental problem of the $`\mathrm{\Lambda }_b`$ life time which has not yet been understood within OPE framework. Second it has been asked if OPE could not miss some subtle kinematical effects related with the delay in the opening of different decay channels . We have shown in a non-relativistic model that the latter effect does not affect the validity of duality.
A numerical calculation of the sum over exclusive channels in the ’t Hooft two dimensional QCD model reported a presence of a duality-violating $`1/m_Q`$ correction in the total width . Later the summation was performed analytically in the case of the massless light quark . Agreement between the OPE and the exact result was found in this case through $`1/m_Q^4`$ order.
The “miraculous” conspiracy of exclusive decay channels to add up to the partonic result and its OPE corrections may be expressed in terms of sum rules which the hadronic matrix elements must satisfy in QCD . OPE was first explicitly used to derive Bjorken sum rule in .
To leading order in $`\mathrm{\Lambda }/m_b`$ Bjorken sum rule straightforwardly implies quark hadron duality for the semileptonic widths (the differential and the total widths). The suppression of the $`O(\mathrm{\Lambda }/m_b)`$ corrections is not so direct. The authors of have done a thorough study of the exclusive contributions of the ground state $`D`$ and $`D^{}`$ mesons up to order $`O(\mathrm{\Lambda }^2/m_b^2)`$. They have chosen the Shifman Voloshin (SV) limit, $`\mathrm{\Lambda }m_bm_cm_b`$, which drastically simplifies the calculation, but did not consider the orbitally excited states, and therefore could not check the matching between the sum of exclusive channels and the OPE prediction to the order $`O(\mathrm{\Lambda }(m_bm_c)/m_b^2)`$.
Our first motivation was precisely to complete this part and add the $`L=1`$ excited states in the sum of exclusive channels. We will discuss in section III why we neglect other excitations.
While performing this task we had a surprise. We found that a new sum rule, eq. (12), was needed beyond Bjorken, Voloshin, and the known tower of higher moment sum rules and we found that this new sum rule could be demonstrated from OPE.
We believe that other new sum rules can be derived along the same line. When the form factors are taken at leading order in $`1/m_b`$, OPE applied to different components of the hadronic tensor, or to different operators, always provides the unique series: Bjorken sum rule, Voloshin sum rule and higher moments. But when the next to leading contribution to the form factors is considered, no such unicity holds anymore. Changing the current operators in the OPE might lead to several other sum rules at order $`1/m_b`$.
In the following we will simplify our task as much as possible. We will neglect radiative corrections. We will also leave aside terms of order $`O(\mathrm{\Lambda }^2/m_b^2)`$, which implies that operators with higher dimension than identity may be neglected in the OPE and consequently that the inclusive results may be computed only via the partonic contribution.
In the next section we will show how the equality of partonic and inclusive widths to the desired order demands for a new sum rule. In section III we will derive the latter sum rule from OPE applied to the T-product of currents. Finally in section IV we show interesting phenomenological consequences of the sum rule. We then conclude.
## II Inclusive semileptonic widths
We work in the SV limit , i.e. we assume the following hierarchy
$$\mathrm{\Lambda }\delta mm_b$$
(1)
where $`\delta mm_bm_c`$ and $`\mathrm{\Lambda }`$ is any energy scale stemming from QCD, for example the hadron-quark mass difference $`\overline{\mathrm{\Lambda }}m_Bm_b=m_Dm_c+O(1/m_b)`$ or the excitation energy.
From OPE one expects quark-hadron duality to be valid up to $`O(\mathrm{\Lambda }^2/m_b^2)`$ corrections, i.e. in terms of the double expansion in $`\delta m/m_b`$ and $`\mathrm{\Lambda }/m_b`$, it should be valid to all orders $`(\delta m/m_b)^n`$ and $`(\delta m/m_b)^n\mathrm{\Lambda }/m_b`$. In fact we will restrict ourselves to check duality up to order $`(\delta m/m_b)^2`$ and $`\delta m\mathrm{\Lambda }/m_b^2`$. The terms of order $`\delta m\mathrm{\Lambda }/m_b^2`$ will turn out to be the trickiest. Of course, in the preceding sentences we mean orders as compared to the leading contribution. For example the inclusive semileptonic width is of order $`(\delta m)^5`$, which implies that we will compute it up to order $`\overline{\mathrm{\Lambda }}(\delta m)^6/m_b^2`$. In this letter the symbol $``$ will always refer to neglecting higher orders than those just mentioned. From OPE the partonic semileptonic decay width should equate the explicit sum of the corresponding exclusive decay widths up to $`O(\mathrm{\Lambda }^2/m_b^2)`$ terms, i.e. :
$$\mathrm{\Gamma }(\overline{B}X_cl\nu )=\mathrm{\Gamma }(bcl\nu )+O(\mathrm{\Lambda }^2/m_b^2)$$
(2)
with the semileptonic partonic width
$$\mathrm{\Gamma }(bcl\nu )=32K(\delta m)^5\left[\frac{2}{5}\frac{3}{5}\frac{\delta m}{m_b}+\frac{9}{35}\frac{(\delta m)^2}{m_b^2}\right]$$
(3)
where
$$K=\frac{G_F^2}{192\pi ^3}|V_{cb}|^2$$
(4)
Using $`M_Bm_b+\overline{\mathrm{\Lambda }}`$ and $`\delta MM_BM_D\delta m`$ we get
$$\mathrm{\Gamma }(\overline{B}X_cl\nu )32K(\delta M)^5\left[\frac{2}{5}\frac{3}{5}\frac{\delta M}{M_B}+\frac{9}{35}\frac{(\delta M)^2}{M_B^2}\frac{21}{35}\frac{\overline{\mathrm{\Lambda }}\delta M}{M_B^2}\right]$$
(5)
The ground state contribution is
$$\mathrm{\Gamma }(\overline{B}(D+D^{})l\nu )32K(\delta M)^5\left[\frac{2}{5}\frac{3}{5}\frac{\delta M}{M_B}+\frac{118\rho ^2}{35}\frac{(\delta M)^2}{M_B^2}\frac{1}{10}\frac{a_+^{(1)}\delta M}{M_B^2}\right]$$
(6)
Strictly speaking nothing compels $`a_+^{(1)}`$ to be real and we must read $`\mathrm{}[a_+^{(1)}]`$ everywhere in this letter instead of $`a_+^{(1)}`$ and $`\mathrm{}[\xi _3]`$ instead of $`\xi _3`$. The contribution of the first orbitally excited states may be computed using results in . We get
$$\mathrm{\Gamma }(\overline{B}(D_1+D_2^{})l\nu )32K|\tau _{3/2}(1)|^2\left[\frac{16}{35}\frac{(\delta M)^2}{M_B^2}\frac{56}{35}\frac{\mathrm{\Delta }_{3/2}\delta M}{M_B^2}\right]$$
(7)
for the states with total angular momentum of the light quanta $`j=3/2`$ and $`\tau _j(w)`$ are the infinite mass limit form factors $`BD^{}`$ as defined in . In all this letter we use for any state $`n`$ the notation
$$\mathrm{\Delta }_n=M_nM_0,$$
(8)
where 0 refers to the ground state.
$$\mathrm{\Gamma }(\overline{B}(D_1^{}+D_0^{})l\nu )32K|\tau _{1/2}(1)|^2\left[\frac{8}{35}\frac{(\delta M)^2}{M_B^2}\frac{49}{35}\frac{\mathrm{\Delta }_{1/2}\delta M}{M_B^2}\right]$$
(9)
for the lowest $`j=1/2`$ states.
To the order considered, quark-hadron duality of the semileptonic decay widths implies the equality of the r.h.s. of eq. (5) with the sum of the r.h.s’s of eqs (6), (7) and (9) to which we need to add the $`L=1`$ radially excited states. Their contributions are identical to eqs. (7) and (9) with the replacement $`\tau _j\tau _j^{(n)}`$ and $`\mathrm{\Delta }_j\mathrm{\Delta }_j^{(n)}`$. The terms proportional to $`(\delta M/M_B)^2`$ match thanks to Bjorken sum rule :
$$\rho ^2\frac{1}{4}=\underset{n}{}\left[|\tau _{1/2}^{(n)}|^2+2|\tau _{3/2}^{(n)}|^2\right]$$
(10)
From now on, unless specified, it is understood that the form factors are taken at $`w=1`$. Taking into account Voloshin sum rule
$$\overline{\mathrm{\Lambda }}=\underset{n}{}\left[2\mathrm{\Delta }_{1/2}^{(n)}|\tau _{1/2}^{(n)}|^2+4\mathrm{\Delta }_{3/2}^{(n)}|\tau _{3/2}^{(n)}|^2\right],$$
(11)
the matching of the terms of order $`\mathrm{\Lambda }\delta M/M_B^2`$ leads to the requirement
$$a_+^{(1)}=4\underset{n}{}\left[\mathrm{\Delta }_{1/2}^{(n)}|\tau _{1/2}^{(n)}|^2\mathrm{\Delta }_{3/2}^{(n)}|\tau _{3/2}^{(n)}|^2\right]$$
(12)
The sum rule (12) is the main result of this paper. The preceding lines can be taken as a derivation of the sum rule, since we simply have made explicit the result from OPE, eq. (5). However, one might feel uncomfortable in view of the peculiarity of the SV kinematics, one might fear that some exception to OPE could happen there. Furthermore, as recalled in the introduction, OPE has been repeatedly submitted to various interrogations. Therefore, we will rederive in the next section the sum rule (12) in a less questionable manner.
Let us note that in the vector current case, we do not need the $`a_+`$ form factor. In that case, matching of the $`(\delta M/M_B)^2`$ and $`\mathrm{\Lambda }\delta M/M_B^2`$ terms occurs thanks to Bjorken and Voloshin sum rule only - or conversely we can invoke duality to demonstrate these sum rules. In particular, it gives a demonstration of Voloshin sum rule just from the same duality requirement invoked by Isgur and Wise to derive Bjorken sum rule: the Voloshin sum rule comes from the matching of $`\mathrm{\Lambda }\delta M/M_B^2`$ terms.
It is in the axial case or in the $`VA`$ case (which corresponds to the sum of vector and axial contribution) that we need the new sum rule. More precisely, we can separate also the contributions with definite helicity of the lepton pair. In the transverse helicity case, there is still matching from just Bjorken and Voloshin sum rule. In fact the need for a new sum rule occurs in the axial current and for longitudinal helicity. We obtain indeed for the $`\lambda =0`$ helicity of the axial current:
$$\mathrm{\Gamma }(bcl\nu )_{A,\lambda =0}4K(\delta M)^5\left[\frac{4}{3}2\frac{\delta M}{M_B}+\frac{4}{5}\frac{(\delta M)^2}{M_B^2}2\frac{\overline{\mathrm{\Lambda }}\delta M}{M_B^2}\right]$$
(13)
$`\mathrm{\Gamma }(\overline{B}D^{}l\nu )_{A,\lambda =0}4K(\delta M)^5`$ $`\left[{\displaystyle \frac{4}{3}}2{\displaystyle \frac{\delta M}{M_B}}+(1{\displaystyle \frac{4}{5}}\rho ^2){\displaystyle \frac{(\delta M)^2}{M_B^2}}{\displaystyle \frac{4}{5}}{\displaystyle \frac{a_+^{(1)}\delta M}{M_B^2}}\right]`$ (14)
$`\mathrm{\Gamma }(\overline{B}D^{}l\nu )_{A,\lambda =0}4K(\delta M)^5`$ $`[{\displaystyle \frac{4}{5}}{\displaystyle \underset{n}{}}[|\tau _{1/2}^{(n)}|^2+2|\tau _{3/2}^{(n)}|^2]{\displaystyle \frac{(\delta M)^2}{M_B^2}}`$ (17)
$`{\displaystyle \frac{28}{5}}{\displaystyle \underset{n}{}}\left[\mathrm{\Delta }_{1/2}^{(n)}|\tau _{1/2}^{(n)}|^2+2\mathrm{\Delta }_{3/2}^{(n)}|\tau _{3/2}^{(n)}|^2\right]{\displaystyle \frac{\delta M}{M_B^2}}`$
$`+{\displaystyle \frac{24}{5}}{\displaystyle \underset{n}{}}[\mathrm{\Delta }_{1/2}^{(n)}|\tau _{1/2}^{(n)}|^2]{\displaystyle \frac{\delta M}{M_B^2}}]`$
whence we get the eq. (12) from the matching of $`\frac{\delta M}{M_B^2}`$ terms.
## III Derivation of the sum rule from OPE
The authors of have derived corrections to Bjorken and Voloshin sum rules and to the resulting inequalities on $`\rho ^2`$. We will follow the same philosophy but including the orbitally excited states in order to derive $`O(\mathrm{\Lambda }/m_b)`$ corrections, within our approximations, to the equalities resulting from the sum rules. We will use the differential semileptonic distributions .
Defining two currents which at present we take arbitrary:
$$J(x)\left(\overline{b}\mathrm{\Gamma }c\right)(x),J^{}(y)\left(\overline{c}\mathrm{\Gamma }^{}b\right)(y).$$
(18)
Their T product is
$$T(q)id^4xe^{iqx}<\overline{B}|T(J(x)J^{}(0))|\overline{B}>$$
(19)
where the states are normalised according to $`<p|p^{}>=(2\pi )^3\delta _3(\stackrel{}{p}^{}\stackrel{}{p})`$.
Neglecting heavy quarks in the “sea”, it is clear that $`x<0`$ receives contributions from intermediate states with one $`c`$ quark and light quanta, usually referred to as the direct channel, while $`x>0`$ receives contributions from intermediate states with $`b\overline{c}b`$ quarks plus light quanta. This will be referred to as the crossed channel, or $`Z`$ diagrams. Expanding the r.h.s of (19) on intermediate states $`X`$ in the $`B`$ rest frame,
$$T=(2\pi )^3\left[\underset{X}{}\delta _3(\stackrel{}{p}_X+\stackrel{}{q})\frac{<\overline{B}|J(0)|X><X|J^{}(0)|\overline{B}>}{M_Bq_0E_X}\underset{X^{}}{}\delta _3(\stackrel{}{p}_X^{}\stackrel{}{q})\frac{<\overline{B}\overline{X}^{}|J(0)|0><0|J^{}(0)|\overline{X}^{}\overline{B}>}{M_B+q_0(E_X^{}+2M_B)}\right]$$
(20)
where $`X,X^{}`$ are charmed states. Let us call $`𝒱`$ the typical virtuality of the direct channels, $`M_Bq_0E_X𝒱`$, we will take $`q_0`$ such that $`\mathrm{\Lambda }𝒱M_B`$. While the direct channels ($`X`$) contribute like $`1/𝒱`$ to (20), the crossed channels ($`X^{}`$) contribute like $`1/(m_D+𝒱)`$. In both cases the denominator is $`\mathrm{\Lambda }`$, which allows to use the leading contribution to OPE:
$$T=id^4xe^{iqx}<\overline{B}|\overline{b}(x)\mathrm{\Gamma }S_c(x,0)\mathrm{\Gamma }^{}b(0)|\overline{B}>+O(1/m_c^2)$$
(21)
where $`S_c(x,0)`$ is the free charmed quark propagator as long as $`O(\alpha _s)`$ corrections are neglected. Assuming as usual that the $`b`$ quark has a momentum $`p_b=m_bv+k`$ with $`k_\mu =O(\mathrm{\Lambda })`$, the charmed quark propagator in (21) has two terms, the positive energy pole with a denominator $`m_bv_0+k_0q_0E_c𝒱`$ and the negative energy one with a denominator $`m_bv_0+k_0q_0+E_cm_c+𝒱`$. Varying $`𝒱`$ independently of $`m_bm_c`$ one can check that the direct channels sum up to the contribution of the positive energy pole of the charmed quark propagator.
As a result, considering now only resonances among the states $`X`$ and fixing $`\stackrel{}{q}`$ in the following, one gets equating the residues
$$\underset{n}{}<\overline{B}|J(0)|n><n|J^{}(0)|\overline{B}>=<\overline{B}|\overline{b}\mathrm{\Gamma }\frac{/v_q^{}+1}{2v_0^{}}\mathrm{\Gamma }^{}b|\overline{B}>$$
(22)
where all the three-momenta are equal to $`\stackrel{}{q}`$ in the $`B`$ rest frame and
$$v_q^{}=\frac{1}{m_c}(\stackrel{}{q},\sqrt{\stackrel{}{q}^{\mathrm{\hspace{0.17em}2}}+m_c^2})$$
(23)
It is well known that to leading order this leads to Bjorken sum rule. Considering successive moments, i.e. multiplying $`T`$ in (19) by $`(q_0E_0)^n`$ ($`E_0`$ being the ground state energy) leads to a tower of sum rules , Voloshin sum rule when $`n=1`$, etc.
In the following we will stick to the $`n=0`$ moment, but include the $`1/m_b`$ correction to the residues. Let us insist on this point. One may discover a tower of sum rules by keeping the form factors to leading order but considering successive moments . One may also discover new sum rules by sticking to the lowest moment but considering the higher orders in the form factors. This is not equivalent and leads to different sum rules, the first moment yields Voloshin sum rule eq (11), the second adds at least one new sum rule, (12), as we shall demonstrate now. The distinction is important since in practice both sum rules apply to the same order in $`1/m_b`$. A significant difference between the two types of subleading sum rules is the following: All the currents provide via OPE the same Voloshin sum rule because the form factors are all related by the heavy quark symmetry. On the contrary, when the form factors are taken at subleading order in $`1/m_b`$, different currents have different corrective terms depending on several independent form factors, and OPE should yield different subleading sum rules. In this letter we only consider eq. (12) for its physical relevance, leaving other sum rules for a forthcoming study.
We now apply eq. (22) with $`J,J^{}`$ substituted by the vector current $`V^\mu `$ and the axial one $`A^\mu `$. One may check that eq. (22) applied to currents projected perpendicularly to the $`v,v^{}`$ plane is trivially satisfied, including the $`O(\mathrm{\Lambda }/m_b)`$ order, by Bjorken sum rule. Let us now consider the vector current projected on the $`B`$ meson four velocity: $`Vv`$. Among the orbitally excited states only the $`J=1`$ states contribute to the wanted order. Dividing both sides of eq (22) by $`(1+w)/(2v_0v_0^{})`$ one gets using the results of and
$$\frac{1+w}{2}|\xi (w)|^2+\underset{n}{}(w1)\left\{2|\tau _{1/2}^{(n)}|^2\left[1+\frac{\mathrm{\Delta }_{1/2}^{(n)}}{m_b}\right]+(w+1)^2|\tau _{3/2}^{(n)}|^2\left[1+\frac{\mathrm{\Delta }_{3/2}^{(n)}}{m_b}\right]\right\}1+(w1)\frac{\overline{\mathrm{\Lambda }}}{m_b}$$
(24)
where we have neglected higher powers of $`(w1)`$ and of $`\mathrm{\Lambda }/m_b`$ than the firstRemember that we take $`\overline{\mathrm{\Lambda }}\mathrm{\Delta }_j\mathrm{\Lambda }`$. The l.h.s is found by a straightforward application of for the ground state and of for the excited ones. The r.h.s yields $`(1+w_q)/(1+w)`$ which has been transformed according to:
$$w_qv.v_q^{}w+\stackrel{}{q}^{\mathrm{\hspace{0.17em}2}}\left[\frac{1}{2m_c^2}\frac{1}{2M_D^2}\right]w+\frac{(w^21)\overline{\mathrm{\Lambda }}}{m_b}.$$
(25)
The leading terms in eq. (24) simply reproduce Bjorken sum rule as expected , while the $`O(\mathrm{\Lambda }/m_b)`$ terms provide Voloshin sum rule. This is another derivation of Voloshin sum rule which does not use higher momenta.
Analogously the axial current projected on the $`D`$ meson velocity $`v^{}`$, $`Av^{}`$ gives, inserted in eq. (22) and after dividing both sides by $`(w1)/(2v_0v_0^{})`$,
$`{\displaystyle \frac{1+w}{2}}|\xi (w)|^2{\displaystyle \frac{4}{m_b}}\xi _3(w)\xi (w)+{\displaystyle \underset{n}{}}\{[2(w1){\displaystyle \frac{6(w+1)\mathrm{\Delta }_{1/2}^{(n)}}{m_b}}]|\tau _{1/2}^{(n)}|^2`$
$$+(w1)(w+1)^2|\tau _{3/2}^{(n)}|^2\}1(w+1)\frac{\overline{\mathrm{\Lambda }}}{m_b}$$
(26)
where $`\xi _3`$ in the notations of is equal to $`a_+^{(1)}/2`$ used in . The matching of the $`1/m_b`$ terms in eq. (26) leads to the sum rule
$$\overline{\mathrm{\Lambda }}+a_+^{(1)}=L_4(1)=+6\underset{n}{}\mathrm{\Delta }_{1/2}^{(n)}|\tau _{1/2}^{(n)}|^2$$
(27)
$`L_4`$ being defined according to . Eliminating $`\overline{\mathrm{\Lambda }}`$ from eqs. (27) and (11) we are left with eq. (12).
We can check this result by using the method for sum rules developed earlier by Bigi and the Minnesota group , which relies on a systematic $`1/m_Q`$ expansion of the moments of the Lorentz invariants of the imaginary part of the hadronic tensor, $`w_i`$. From their equation (131), we read :
$$𝑑q^0w_2^{AA}(q^0,\stackrel{}{q}^{\mathrm{\hspace{0.17em}2}})\frac{m_b}{E_c}$$
(28)
the terms left over being the power corrections due to higher dimension operators. Computing from and the hadronic contribution to the same integral at $`\stackrel{}{q}=0`$ i.e. $`w=1`$, we get the equation (with $`r_0=M_D^{}/M_B`$, $`r_{1/2,3/2}=M_{D_{1/2,3/2}^{}}/M_B`$):
$`{\displaystyle 𝑑q^0w_2^{AA}(q^0,\stackrel{}{q}^{\mathrm{\hspace{0.17em}2}}=0)}={\displaystyle \frac{1}{r_0}}\left\{{\displaystyle \frac{f^2}{4M_B^2r_0^2}}+{\displaystyle \frac{(1r_0)fa_+}{r_0}}\right\}`$ (29)
$`+\left\{{\displaystyle \frac{(1r_{3/2})^2}{r_{3/2}^2}}\left({\displaystyle \frac{k_{A_1}^2}{24}}{\displaystyle \frac{f_A^2}{4}}\right)+{\displaystyle \frac{1}{4r_{1/2}^2}}([(1+r_{1/2})g_+(1r_{1/2})g_{}]^2g_A^2)\right\}`$ (30)
with all form factors taken at $`w=1`$, and with notations for the $`L=1`$ form factors $`g_+,g_{},g_A,f_A,k_{A_1}`$ to be found in . A sum over the $`L=1`$ excitations is unedrstood. If we now work in the SV limit, we see that we need $`g_{},f_A,g_A,k_{A_1}`$ only in the HQET limit, i.e. $`\tau _{1/2,3/2}`$, except for some algebraic factors; as for $`g_+`$, it is subleading, but at $`w=1`$, it is expressible in terms of $`\tau _{1/2}`$ and we do not need to know any of the new subleading form factors. In the $`L=1`$ contributions, only the $`g_+g_{}`$ term remains. We finally end with the equation :
$$\frac{M_B}{M_D}\frac{\delta M}{M_B^2}a_+^{(1)}+6\frac{\delta M}{M_B^2}\underset{n}{}\mathrm{\Delta }_{1/2}^{(n)}|\tau _{1/2}^{(n)}|^2\frac{m_b}{m_c}$$
(31)
which leads directly to eq. (27).
In the preceding calculations we have systematically neglected the contributions from higher orbital excitations or $`L=0`$ radial excitations. This can be justified as follows. The leading $`B`$ transition to radially excited $`L=0`$ final states or to $`L=2`$ final states are suppressed by a factor $`\stackrel{}{q}^{\mathrm{\hspace{0.17em}2}}/m_b^2`$ due to three facts: first, the current operator is proportional at leading order to the identity operator or to $`\stackrel{}{\sigma }_b`$<sup>§</sup><sup>§</sup>§ The heavy quark spin may be factorised out thanks to HQS., second, the orthogonality of the wave functions implies vanishing at $`\stackrel{}{q}=0`$ in the $`B`$ rest frame and, third, parity implies an even power in $`\stackrel{}{q}`$. This suppression leads to the well known fact that these terms appear in the Bjorken sum rule or in the differential widths with a $`(w1)^2`$ factor as compared to the ground state contribution. On the contrary the $`O(\mathrm{\Lambda }/m_b)`$ contributions to the axial form factors for the same type of transitions are not suppressed as compared to the ground state because the current operator is no more proportional to identity neither to $`\stackrel{}{\sigma }_b`$. For example the transition to radially or orbitally excited $`J^P=1^{}`$ states other than the $`D^{}`$ are in principle of the same order of magnitude than the $`a_+^{(1)}`$ terms mentioned above. However, in this letter we have only considered the terms $`a_+^{(1)}`$ via crossed terms, i.e. via cross products of the leading order terms with the $`O(\mathrm{\Lambda }/m_b)`$ ones, because we have neglected all $`O(\mathrm{\Lambda }^2/m_b^2)`$ contributions. Hence we are left with a suppression of a factor $`\stackrel{}{q}^{\mathrm{\hspace{0.17em}2}}/m_b^2`$ in the hadronic tensors or the differential widths, i.e. a factor $`(w1)`$ as compared to the corresponding ground state contribution and we can consequently neglect the $`L=0`$ radial excitations and the $`L=2`$ orbital ones. $`L=3`$ contributions are negligible simply because the total angular momentum $`J2`$ again leads to $`(w1)`$ factors resulting from angular momentum conservation (D-waves). All other operators which are already negligible for the ground state and the $`L=1`$ states are even more so for higher excitations.
Turning now to a comparison of our different demonstrations, we should note that it is not really unexpected that we find consistent results according to three approaches: imposing duality to the widths (section II), imposing duality to the tensors as in eqs. (24) and (26) and finally to the invariant tensors eqs. (28) and (30). Indeed, at fixed $`q^0`$ and $`\stackrel{}{q}`$ there is a linear relation between the tensor components and the invariant tensors. It is as well true that the formula for the decay widths before integrating on the $`q^0`$ variable is, for fixed $`q^0`$ and $`\stackrel{}{q}`$, linear in the tensor components.
We might worry about what happens when we apply duality to the sum of the residues. Integration over $`q^0`$ leads to a sum of residues multiplied by $`\delta `$ functions and the position of the poles is different for each term in the sum and still different for the quark contribution. As a consequence the projector which projects out $`w_2`$ from the tensor residues is different for each term since it depends on $`q^0`$. Still this difference does not lead to a collapse of the sum rule thanks to Voloshin sum rule and the tower of higher momenta sum rules: one can expand the difference between the intervening projectors in powers of $`q^0`$ and the resulting alteration to the sum rule vanishes. Exactly the same happens when one computes the decay widths with the real kinematics on each term.
## IV Phenomenological consequences
Eq. (27) is phenomenologically relevant as it expresses the dominant correction to the zero recoil differential $`BDl\nu `$ decay width as a function of leading form factors and level spacings. Indeed
$$\frac{d\mathrm{\Gamma }(BDl\nu )}{dw}(w^21)^{3/2}\left[12\left(\frac{1}{2m_b}+\frac{1}{2m_c}\right)\frac{M_BM_D}{M_B+M_D}L_4(1)\right].$$
(32)
On the other hand, we may combine our result with an independent estimate of the form factor $`\xi _3`$ from QCD sum rulesThe definitions of $`\xi _3`$ differ by a factor $`\overline{\mathrm{\Lambda }}`$ in and . We use the notations of .:
$$\frac{\xi _3(1)}{\overline{\mathrm{\Lambda }}}=\frac{1}{3}+O(\alpha _s)=0.6\pm 0.2,\frac{a_+^{(1)}}{\overline{\mathrm{\Lambda }}}=\frac{2}{3}O(\alpha _s)=1.2\pm 0.4.$$
(33)
The dispersion formulation of the constituent quark model finds that $`\xi _3(1)`$ is 1/3 the average kinetic energy of the light quark. For a light constituent mass of $`m_u=0.25`$ GeV it gives
$$\xi _3(1)=0.17\mathrm{GeV},\overline{\mathrm{\Lambda }}=0.5\mathrm{GeV}$$
(34)
in perfect agreement with eq. (33) for $`\alpha _s=0`$.
Combining (11), (12) and (33), assuming $`\alpha _s=0`$ since we have neglected radiative corrections all along this letter, we get
$$\frac{\underset{n}{}\mathrm{\Delta }_{1/2}^{(n)}|\tau _{1/2}^{(n)}|^2}{_n\mathrm{\Delta }_{3/2}^{(n)}|\tau _{3/2}^{(n)}|^2}=\frac{1}{4},\mathrm{for}\alpha _s=0$$
(35)
and
$$\underset{n}{}\mathrm{\Delta }_{1/2}^{(n)}|\tau _{1/2}^{(n)}|^2=\frac{1}{18}\overline{\mathrm{\Lambda }},\underset{n}{}\mathrm{\Delta }_{3/2}^{(n)}|\tau _{3/2}^{(n)}|^2=\frac{2}{9}\overline{\mathrm{\Lambda }}$$
(36)
Notice that if we had, somehow inconsistently, taken $`\xi _3(1)/\overline{\mathrm{\Lambda }}=0.6`$ the result would not be qualitatively different.
Since in all spectroscopic models the mass differences between the $`j=1/2`$ and $`j=3/2`$ states turn out to be not so large, we conclude that the $`_n|\tau _{1/2}^{(n)}|^2`$ are significantly smaller than the $`_n|\tau _{3/2}^{(n)}|^2`$ .
Interestingly enough, this hierarchy $`|\tau _{1/2}^{(0)}|^2<|\tau _{3/2}^{(0)}|^2`$ was a clear outcome of a class of covariant quark models . In four different potentials had been used within the Bakamjian-Thomas covariant quark model framework. The potentials labeled ISGW, VD, CCCN, and GI potentials in give respectively for the ratio $`|\tau _{1/2}^{(0)}|^2/|\tau _{3/2}^{(0)}|^2`$ the values 0.33, 0.09, 0.01 and 0.17. As a result these models predict a dominance of the $`BD_{j=3/2}l\nu `$ semileptonic decay widths by one order of magnitude over the $`BD_{j=1/2}l\nu `$. We will comment this prediction later. The same models give for the l.h.s of eq. (35) 0.39, 0.166, 0.151 and 0.247 respectively for the ISGW, VD, CCCN, and GI potentials, in reasonable agreement with 1/4. It might not be mere luck if the GI model, which fits the spectrum in the most elaborate way, yields an almost too good agreement with the expectation (35) We should nevertheless remember that the potentials used in contain a Coulombic part which implies that some part of the $`O(\alpha _s)`$ corrections might be implicit in these models.. From eq. (33) we expect the r.h.s. of eq. (12) divided by that of eq. (11) to be close to $`2/3`$. We have tested this with the numerical calculations of . In all cases we find that the sums in the r.h.s of eqs. (11)-(12) saturate very fast to their symptotic values. At $`n=3`$ they are at less than 3% in all cases. For the ratios $`a_+^{(1)}/\mathrm{\Lambda }`$ computed from the r.h.s of eqs. (11)-(12) one finds -0.51, -0.77, -0.79, -0.67 respectively for the ISGW, VD, CCCN, and GI models. This agreement with (33) is quite striking, and again GI is embarrassingly good.
In more general terms, the prediction that the $`B`$ meson decays dominantly into the narrow resonances $`j=3/2`$ was comforted by a study within a constituent quark-meson model as well as by a semi-relativistic study . A QCD sum rule analysis predicted rather a rough equality between these form factors contrarily to another one which concluded to an overwhelming dominance of the $`j=3/2`$ semileptonic decay over the $`j=1/2`$.
It is fair to say that the general trend of theoretical models is to predict $`3/2`$ dominance and a total semileptonic branching ratio into the orbitally excited states exceeding hardly 1 %. It is well known that the $`j=3/2`$ are expected to be relatively narrow and are identified with the observed narrow resonances $`D_1(2422)`$ and $`D_2^{}(2459)`$. As far as the decay widths into the latter narrow resonances is considered, experimental results are in rough agreement with for the $`BD_1(2422)l\nu `$ and rather below for $`BD_2^{}(2422)l\nu `$. In brief, experiment is rather below the theoretical models for $`BD_{3/2}l\nu `$. The $`j=1/2`$ states are not easy to isolate, being very broad. But thorough studies have been done of the channels $`BD^{()}\pi l\nu `$ and the resulting branching fraction is very large: $`3.4\pm 0.52\pm 0.32\%`$ by DELPHI and $`2.26\pm 0.29\pm 0.33\%`$ by ALEPH.
These experimental results are both welcome and puzzling. Welcome because these $`BD^{()}\pi l\nu `$ fill the gap between the inclusive semileptonic decay branching fraction of 10 - 11 % and the sum $`B(D+D^{})l\nu 7\%`$. They are puzzling when one tries to understand which channels contribute to them. As we have just said, the $`j=3/2`$ channels provide no more than 1 %. The remaining 2 % can come from the $`j=1/2`$, from higher excitations or from a non-resonant continuum. Higher excitations are unlikely to contribute very much, being suppressed both by dynamics and phase space. In the quoted $`bD^{}l\nu `$ branching fractions are very large, exceeding by far what is expected for example in .
The results presented in this letter are doubly relevant in the above discussion. First eq. (35) seems to confirm the models which find a dominance of the $`3/2`$ channels. Of course it is mathematically possible that that eq. (35) is satisfied while $`|\tau _{1/2}^{(0)}|>|\tau _{3/2}^{(0)}|`$, the higher excitations compensating for the sum rule. Admittedly such a situation would look rather queer, and as mentioned above , the models , which agree rather well with the new sum rule eq. (12), also yield $`|\tau _{1/2}^{(0)}|^2<0.35|\tau _{3/2}^{(0)}|^2`$.
It is then hard to understand how the $`bD^{}l\nu `$ branching fractions can be as large as quoted in in view of the smallness of the experimental $`BD_{3/2}l\nu `$ branching fractions. However, the second lesson from our study is that $`1/m_c`$ corrections may play an important role, and a further study of their effect is wanted.
The most serious caveat to our present derivation of a narrow resonance dominance comes from the fact that we have neglected radiative corrections. A priori we expect radiative corrections to provide only corrections and our present estimate to yield the general trend. This is unhappily not always true. As a counterexample see the discussion which follows eq. (7.8) in . It is argued that some radiative corrections to the parameter $`K`$ are parametrically larger than the $`\alpha _s=0`$ estimate. A careful study of radiative corrections to our present sum rule and its consequences would be welcome.
It is not excluded that an important fraction of the $`BD^{()}\pi l\nu `$ decays observed at LEP are non-resonant. Unluckily theoretical works addressing non-resonant decays are rare, find in the soft-pion domain a resonance dominance while Isgur predicts that no more than 5 % of the semileptonic decay is non-resonant. Furthermore, if such a continuum contributes significantly, it should also be included in the sum rules and we might fear that at the end of the day the paradox would still be there.
Finally another experimental result seems to contradict our theoretical expectation: the branching ratio for $`BD_1(j=1/2)\pi ^{}`$ is found to be $`1.5`$ times larger than that of $`BD_1(2420)\pi ^{}`$. Of course the experimental error is still large, and the relation between nonleptonic decays and the semileptonic ones assumes factorisation.
But still there is a puzzle: on one side an increasing amount of theoretical evidence in favor of the narrow resonances dominance, and on the other side an increasing amount of experimental evidence in the opposite direction!
## V Conclusion and outlook
We have explicitly checked quark-hadron duality in the SV limit to order $`\delta m\mathrm{\Lambda }/m_b^2`$ including ground state final hadrons and $`L=1`$ orbitally excited states. We have shown that this duality implied a new sum rule eq. (12) which we have also demonstrated from OPE applied to T-product of axial currents.
We have shown that this sum rule combined with some theoretical estimates of $`\xi _3`$ lead to the conclusion that very probably the $`B`$ decay into narrow $`L=1`$ resonances was dominant over the one into broad resonances. This remark seems to contradict recent experimental claims that the broad resonances dominate. We have discussed this situation which needs urgently further theoretical and experimental work.
Beyond understanding this experimental puzzle, further theoretical work is needed. For example we might wonder if some proof of the new sum rule along the line of is possible. Some progress has been done in this direction. The effect of radiative corrections should also be studied.
Last but not least, other new sum rules derived along the same line with other currents or other components of the currents should be considered.
###### Acknowledgements.
We thank Patrick Roudeau for many stimulating discussions which have triggered this study.
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# Formation of a New Class of Random Fractals in Fragmentation with Mass Loss
## Abstract
We consider the fragmentation process with mass loss and discuss self-similar properties of the arising structure both in time and space focusing on dimensional analysis. This exhibits a spectrum of mass exponents $`\theta `$, whose exact numerical values are given for which $`x^\theta `$ or $`t^{\theta z}`$ has the dimension of particle size distribution function $`c(x,t)`$ where $`z`$ is the kinetic exponent. We also give explicit scaling solution for special case. Finally, we identify a new class of fractals ranging from random to non-random and show that the fractal dimension increases with increasing order and a transition to strictly self-similar pattern occurs when randomness is completely seized.
PACS number(s): 05.20.Dd,02.50.-r,05.40-y
The kinetics of irreversible and sequential breakup of particles occurs in a variety of physical processes and has important applications in science and technology. These include erosion , grinding and crushing of solids, polymer degradation and fiber length reduction , breakup of liquid droplets etc. to name just a few. In recent years there has been an increasing interest in studying fragmentation allowing variations to increase the flexibility of the theory in matching the conditions of real phenomena such as extension to higher dimension , agglomerate erosion , mass loss , volume change , fragmentation-annihilation . The kinetic equation approach of fragmentation is linear in character which makes it analytically tractable for a large class of breakup kernels. This is contrary to the reverse process, describing the kinetics of coagulation whose mean field approach proposed by Smoluchowski is non-linear in character and solved for a limited choice of collision kernels. This may reflect the fact that breaking up of objects follows less constraints than its reverse process. Despite its apparent simplicity and the fact that the first work appeared more than a century ago, the fragmentation process is still producing nontrivial results. For example, only recently it has been observed that when particles are described by more than one dynamical quantity, such as size and shape, the system exhibits multiscaling as it obeys infinitely many conservation laws . Moreover, the resulting fragment distribution was shown to exhibit multifractality on a unique support when describing fragmentation and on one of infinitely many supports when describing stochastic Sierpinski gasket process . It has also been discovered that a shattering transition occurs as the subsequent generation of fragments has a shorter life time than the fragments of previous generation . McGrady and Ziff further showed in that the shattering regime produces a fractal dust with dimension $`0<D_f<1`$ due to mass being lost to the phase of zero sized particles. In one dimension, there is only one phase boundary for shattering transition which is identified by the singularity of kinetic exponent whereas in more than one dimension there are multiple phase boundaries . Shattering transition is also shown to be accompanied by the absence of scaling and self-averaging .
If associating disorder with broken objects is the most natural thing to do, then searching for an order even in this disorder is the next natural thing. This forms part of our motivation of this work. In this letter, we consider the kinetics of fragmentation with continuous mass loss and look more at the geometric and scaling aspects of the process than merely trying to solve the equation. The scaling theory essentially provides solutions in the long-time and short-size limit when the particle size distribution function evolves to a simpler form as well as becomes independent of initial conditions . In reality, the most experimental system evolves to the point where this behaviour is reached. Our aim is to search for an order and quantify the arising geometry of the pattern. In fact, there are many physical processes that provide an intriguing connection between geometry and physics such as Percolation, Diffusion Limited Aggregation (DLA), Self-Organised Criticalitiy (SOC), etc. These processes evolve according to a random process obeying some conservation laws and creating simple geometrical structures that traditional Euclidean geometry cannot describe. Like many statistical physics problems, exact solutions for the distribution of masses of such ramified or stringy objects and their geometry vis-a-vis measuring their fractal dimension by complete analytical means, are still a challenge even in one dimension.
The evolution of particle size distribution function $`c(x,t)`$ for fragmentation with mass loss is
$$\frac{c(x,t)}{t}=c(x,t)_0^{\mathrm{}}F(y,xy)𝑑y+2_x^{\mathrm{}}𝑑yc(y,t)F(x,yx)+\frac{}{x}(m(x)c(x,t))$$
where $`F(x,y)`$ is the breakup kernel describing the rate at which a particle of size $`(x+y)`$ breaks into sizes $`x`$ and $`y`$. Fragmentation is a process whereby cuts are equivalent of seeds being sown on the fragmenting objects, thus producing two new segments. This immediately creates two more new ends belonging to the two different, newly-created fragments; in doing so, fragments start loosing their masses immediately (as if seeds were growing on either sides uniformly) until they encounter another seed or become dust-like thereby stopping loosing their masses. Therefore the model we consider can also mimic nucleation and growth of gap in one dimension which have some relevance in Kolmogorov-Avrami-Johnson-Mehl (KAJM) nucleation and growth processes and space covering by growing rays .
We consider the breakup kernel to be $`F(x,y)=(xy)^\beta (x+y)^{\lambda 1}`$, for which the breakup rate $`a(x)=_0^xF(y,xy)𝑑y=px^{2\beta +\lambda }`$, where $`p=\frac{(\mathrm{\Gamma }(\beta +1))^2}{\mathrm{\Gamma }(2\beta +2)}`$. The first term on the right hand side of the equation $`(1)`$ reveals that $`x^{(2\beta +\lambda )}`$ bears the dimension of time and this put a strong constraint on the mass loss term. So, the dimensional consistency requires $`m(x)=mx^{2\beta +\lambda +1}`$, with $`m`$ a positive real constant. This dimensional consistency has been ignored in all previous studies and $`\gamma <2\beta +\lambda +1`$ was identified as the recession regime and $`\gamma >2\beta +\lambda +1`$ as the fragmentation regime assuming $`m(x)mx^\gamma `$. Since $`x`$ and $`t`$ are inextricably intertwined via the dimensional consistency, any of the two can be taken to be an independent parameter when the other one is expressible in terms of this. If $`x`$ is chosen to be the independent parameter then the spatial scaling ansatz is $`c(x,t)x^\theta \mathrm{\Phi }(t/t_0(x))`$, where $`t_0(x)=x^{(2\beta +\lambda )}`$. On the other hand, if $`t`$ is taken to be the independent parameter then the temporal scaling ansatz is $`c(x,t)t^{\theta z}\varphi (x/x_0(t))`$, with $`x_0(t)=t^{\frac{1}{2\beta +\lambda }}`$ and the kinetic exponent $`z=\frac{1}{2\beta +\lambda }`$. The parameters $`t/t_0(x)=\xi `$ and $`x/x_0(t)=\eta `$ are the dimensionless quantities and so are $`\mathrm{\Phi }(\xi )`$ and $`\varphi (\eta )`$. Consequently, $`\theta `$ takes the value for which $`x^\theta `$ and $`t^{\theta z}`$ have the dimension of $`c(x,t)`$. Note that the spatial and temporal scaling solution are trivially connected via $`\mathrm{\Phi }(\eta ^\gamma )=\eta ^\theta \varphi (\eta )`$, where $`\eta =xt^{\frac{1}{\gamma }}`$. The mass exponent $`\theta `$ can only be found if the system follows some conservation laws. For example, for pure fragmentation ($`m=0`$) the mass or size of the system is a conserved quantity and gives $`\theta =2`$. Defining the $`n^{th}`$ moment $`M_n(t)=_0^{\mathrm{}}x^nc(x,t)𝑑x`$ and combining it with the rate equation $`(1)`$ for the present choice of $`F(x,y)`$ and $`m(x)`$ yields
$$\frac{dM_n(t)}{dt}=[\frac{(\mathrm{\Gamma }(\beta +1))^2}{\mathrm{\Gamma }(2\beta +2)}\frac{2\mathrm{\Gamma }(\beta +1)\mathrm{\Gamma }(n+\beta +1)}{\mathrm{\Gamma }(n+2\beta +2)}+mn]M_{n+2\beta +\lambda }(t).$$
(1)
The interesting feature of the above equation is that for $`m>0`$, there are infinitely many $`n=D_f(\beta ,m)`$ values for which $`M_{D_f(\beta ,m)}(t)`$s are conserved quantities. However, for $`m=0`$, there is only one conserved quantity $`M_1(t)`$, i.e. size or mass of the system, and this does not depend on $`\beta `$. We can find the $`D_f(\beta ,m)`$ value by searching for the positive and real root of the equation
$$\frac{(\mathrm{\Gamma }(\beta +1))^2}{\mathrm{\Gamma }(2\beta +2)}\frac{2\mathrm{\Gamma }(\beta +1)\mathrm{\Gamma }(n+\beta +1)}{\mathrm{\Gamma }(n+2\beta +2)}+mn=0$$
(2)
which is polynomial in $`n`$ of degree determined by the $`\beta `$ value. Substituting the temporal scaling anstaz into the definition of $`M_n(t)`$ gives $`M_n(t)t^{(n(\theta 1))z}_0^{\mathrm{}}\eta ^n\varphi (\eta )𝑑\eta `$ and demanding $`M_{D_f}(\beta ,m)`$ be a conserved quantity immediately gives $`\theta =(1+D_f(\beta ,m))`$, which clearly depends on $`\beta `$ and $`m`$ only if $`m>0`$. Owing to the random nature of the process and due to the presence of mass loss term, it is clear that when the process continues ad infinitum, it creates a distribution of points (dust) along a line at an extreme late stage. This distribution of points will inevitably be different from any known set such as strictly self-similar Cantor set, Julia set, Koch curve , stochastic or random Cantor set . To measure the size of the set created in the long time limit, we define a line segment $`\delta =\frac{M_n(t)}{M_{n1}(t)}t^{\frac{1}{2\beta +\lambda }}`$. We can count the number of such segments needed to cover the set and in the limit $`\delta 0`$ (i.e.$`t\mathrm{}`$), the number $`N(\delta )`$ will simply measure the set and appear to scale as $`N(\delta )\delta ^{D_f(\beta ,m)}`$. The exponent $`D_f(\beta ,m)`$ is known as the Hausdorff-Basicovitch dimension of the set or as the fractal dimension which is simply the real positive root of the equation $`(3)`$.
To get a physical picture of the role played by $`m`$, we set $`\beta =0`$ for the time being for which the equation $`(3)`$ becomes quadratic in $`n`$ and the real positive root is $`D_f(m)=\frac{1}{2}(1+1/m)+\frac{1}{2}\sqrt{(1+1/m)^2+4/m}`$ when the second root is $`D=(D_f(m)+1+1/m)`$. Therefore, the exponent $`\theta `$ is also function of $`m`$. The expression for $`D_f(m)`$ reveals that as $`m`$ value increases, the fractal dimension decreases very sharply and in the limit $`m\mathrm{}`$, $`D_f(m)0`$. This means that as $`m`$ increases the size of the corresponding arising set decreases sharply due to fast disappearance of its member. Whereas, as $`m0`$, $`D_f(m)1`$, that is we recover the full set (pure fragmentation) that describes a line. On the other hand had we kept $`m`$ fixed and let $`p`$ decreases the effect would have been the same as we observed for increasing $`m`$ with $`p=1`$ (i.e. $`\beta =0`$). Thus, it is the ratio between $`m`$ and $`p`$ that matters rather than their individual increases or decreases. To give a physical picture of what these results mean we define mass length relation for the object as $`M_0\delta ^{D_f(m)}`$ and $`M_e\delta ^d`$ for the space where the object is being embedded, here $`d`$ describes the Euclidean space. The density of the property of the object $`\rho `$ then scales as
$$\rho \delta ^{D_f(m)d}.$$
(3)
Note that for $`m>0`$, $`D_f(m)`$ is always less than one. It is thus clear that for a given class of set created by a specific rule, when $`D_f(m)`$ decreases it means that it is increasingly moving away from $`d`$ and hence more and more members from the full set are removed. This in turn creates increasingly ramified or stringy objects since $`D_f(m)=d`$ describes the compact object with uniform density. So, we show that increasing $`m/p`$ ratio means that mass loss process gets stronger than the fragmentation process and vice versa.
We now attempt to find the spatial scaling solution for $`\mathrm{\Phi }(\xi )`$. Note that the dimension of the arising pattern is independent of $`\lambda `$ and consequently independent of how fast or slow the system performs the process. So, we can set $`\lambda =1`$ without fear of missing any physics but it certainly simplifies our calculation. Substituting the spatial scaling ansatz into the rate equation $`(1)`$ for $`F(x,y)=1`$ and $`m(x)=mx^2`$ and differentiating it with respect to $`\xi `$, transforms the partial integro-differential equation into an ordinary differential equation ,
$$\xi (1m\xi )\mathrm{\Phi }^{\prime \prime }(\xi )+[(1\theta )\xi (2m(2\theta )1]\mathrm{\Phi }^{}(\xi )(m(2\theta )(1\theta )(3\theta ))\mathrm{\Phi }(\xi )=0.$$
(4)
For $`m=1`$ this is hypergeometric differential equation whose only physically acceptable linearly independent solutions are $`{}_{2}{}^{}F_{1}^{}(1,(1+2D_f);D_f;\xi )`$ and $`\xi ^{(1+D_f)}{}_{2}{}^{}F_{1}^{}(2+D_f,D_f;2+D_f;\xi )`$, where $`D_f=0.414213`$. From these exact solutions for spatial scaling function we can obtain the asymptotic temporal scaling function $`\varphi (\xi )e^{D_f\xi }`$ that satisfies the condition $`\varphi (\xi )0`$ as $`\xi \mathrm{}`$.
We now attempt to see the role of $`\beta `$ on the system. To judge its role, it is clear from the previous discussion that we ought to give equal weight to all the terms in the equation $`(1)`$ so that each of them can compete on an equal footing. This can be done if only we set $`m=p=\frac{(\mathrm{\Gamma }(\beta +1))^2}{\mathrm{\Gamma }(2\beta +2)}`$ so that the relative strength between fragmentation and mass loss process stays the same as $`\beta `$ value increases. This is a very crucial point to be emphasized. We can obtain the fractal dimension for different values of $`\beta `$, which is simply the real positive root of the equation $`(3)`$. A detailed survey reveals that the fractal dimension increases monotonically with increasing $`\beta `$. To find the fractal dimension in the limit $`\beta \mathrm{}`$, we can use the Stirling’s approximation in $`(3)`$ to obtain $`\mathrm{ln}[n+1]=(1n)\mathrm{ln}[2]`$ when $`n=0.4569997`$ solves this equation. In order to give a physical picture of the role of $`\beta `$ in the limit $`\beta \mathrm{}`$, we consider the following model $`F(x,y)=(x+y)^\gamma \delta (xy)`$. This model describes that cuts are only allowed to be in the middle in order to produce two fragments of equal size at each time event. This makes $`a(x)=\frac{1}{2}x^\gamma `$, so we need to choose $`m(x)=\frac{1}{2}x^{\gamma +1}`$, where $`m=\frac{1}{2}`$ gives the same weight as for the fragmentation process. Then the rate equation for $`M_n(t)`$ becomes
$$\frac{dM_n(t)}{dt}=[\frac{(n+1)}{2}2^n]M_{n+\gamma }(t).$$
(5)
As before we set the numerical factor of the right hand side of this equation equal to zero and then take natural log on both sides to obtain the $`n`$ value for which $`M_n(t)`$ is time independent. In doing so, we arrive at the same functional equation for $`n`$ as we found for $`\beta \mathrm{}`$. This shows that the kernel $`F(x,y)=(xy)^\beta (x+y)^{\lambda 1}`$ behaves exactly in the same fashion as for $`F(x,y)=(x+y)^\gamma \delta (xy)`$. We thus find that in the limit $`\beta \mathrm{}`$, the resulting distribution of points is a set with fractal dimension $`D_f=0.4569997`$ which is a strictly self-similar fractal as randomness is seized by dividing fragments into equal pieces. We are now in a position to give a physical picture of the role played by $`\beta `$. First of all, the process with $`\beta =0`$ describes the frequency curve of placing cuts about the size of the fragmenting particles is Poisson in nature. Consequently, the system enjoys the maximum randomness and the corresponding fractal dimension is $`D_f=0.414213`$. Whereas, for $`\beta >0`$, the frequency curve of placing cuts about the size of the fragmenting particles is Gaussian in nature meaning as $`\beta `$ value increases particles are increasingly more likely to break in the middle than on either end. That is, as $`\beta `$ increases, the variance decreases in such a manner that in the limit $`\beta \mathrm{}`$ the variance of the frequency curve becomes infinitely narrow meaning a delta function distribution for which fragments are broken into two equal pieces. This analysis also specify that the rules determining the location where to place the cut are determined by the details of breakup kernel $`F(x,y)`$ rather than the breakup rate $`a(x)`$. So, there is a spectrum of fractal dimensions between $`\beta 0`$ when $`D_f=0.414213`$ and $`\beta \mathrm{}`$ when $`D_f=0.4569997`$. A detailed numerical survey that we do not present here confirms that fractal dimension increases monotonically with $`\beta `$ and reaches to a constant value when $`\beta \mathrm{}`$ in a similar fashion as the variation of $`q`$ with $`t`$ during charging process in $`RC`$ circuit. According to equation $`(4)`$ increasing $`\beta `$ vis-a-vis increasing order also means that the system losses less and less mass from the system and this happens despite the fact that now $`\frac{m}{p}`$ ratio stays the same. Perhaps it is note worthy to mention that the present model with $`\beta =\lambda =0`$ and $`m=1`$ correspond to Yule-Furry processes for cosmic shower theory with collision loss , though there too dimensional consistency was ignored.
In summary, we have identified a new set with a wide range of subsets produced by tuning the degree of randomness only. The process starts with an initiator of unit interval $`[0.1]`$ and the generator divide the interval into two pieces and deleting some parts from either sides of both the pieces at each time step. The amount of the parts to be deleted is determined by the parameter that control the intensity of randomness. When this operation continues ad infinitum, what remains is an infinite number of dust scattered over the interval. We quantified the size of the arising set by fractal dimension and showed that the fractal dimension increases with increasing order and reaches its maximum value when the pattern described by the set is perfectly ordered, which is contrary to some recently found results . We have also shown that the increase of fractal dimension and the increase of mass exponent $`\theta `$ go hand in hand since they are intimately connected. To the best of our knowledge the exact numerical value of this mass exponent has never been reported. We have given a scaling description of the process both in time and space and obtained explicit scaling function for special case of interest. Finally we argue on the basis of our findings that fractal dimension, degree of order and the extent of ramifications of the arising pattern are interconnected.
The author is grateful to R. M. Ziff for sending valuable comments. The author also acknowledges inspiring correspondence with P. L. Krapivsky and support from the Ministry of Science and Technology of Bangladesh under Grant No. 1/98/112/1(7).
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# Gravitational conformal invariance and coupling constants in Kaluza-Klein theory
## 1 Introduction
The theory of conformal invariance has been playing a particularly important role in the investigation of gravitational models since Weyl, who introduced the notion of conformal rescaling of the metric tensor. Afterwards, it was promoted to the conformal transformations in scalar-tensor theories, in which another transformation on the scalar field was required to represent the conformal invariance in modern gravitational models. There is an open possibility that the gravitational coupling of matter may have its origin in an invariance breaking effect of this conformal invariance. In fact, since the ordinary coupling of matter to gravity is a dimensional coupling (mediated by the gravitational constant), the local conformal transformations which could change the strength of this dimensional coupling, by affecting the local standards of length and time, are expected to play a key role. In a system which includes matter, conformal invariance requires the vanishing of the trace of the stress tensor in the absence of dimensional parameters. However, in the presence of dimensional parameters, the conformal invariance can be also established for a large class of theories if the dimensional parameters are conformally transformed according to their dimensions. One general feature of conformally invariant theories is, therefore, the presence of varying dimensional coupling constants. In particular, one can say that the introduction of a constant dimensional parameter into a conformally-invariant theory breaks the conformal invariance in the sense that a preferred conformal frame is singled out, namely that in which the dimensional parameters have the assumed (constant) configuration. The determination of the corresponding preferred conformal frame depends on the nature of the problem at hand. In a conformally-invariant gravitational model, the symmetry breaking may be considered as a cosmological effect. This means that one breaks the conformal symmetry by defining a preferred conformal frame in terms of the large-scale characteristics of cosmic matter distributed in a universe with finite scale factor $`R_0`$. In this way, the breakdown of conformal symmetry becomes a framework in which one can look for the origin of the gravitational coupling of matter, both classical and quantum , at large cosmological scales.
The purpose of this paper is to show that one may look for the origins of both conformal invariance and its breakdown, leading to gravitational couplings, in a 5-dimensional Kaluza-Klein type gravity theory . In this popular non-compactified approach to Kaluza-Klein gravity, the gravitational field is unified with its source through a new type of 5D manifold in which space and time are augmented by an extra non-compactified dimension which induces 4D matter. This theory basically involves writing the Einstein field equations with matter as a subset of the Kaluza-Klein field equations without matter , a procedure which is guaranteed by an old theorem of differential geometry due to Campbell .
We show that in the context of pure geometry theory, i.e. $`\widehat{R}_{AB}=0`$ in 5D, one may find a generalized conformally-invariant gravitational model. The well-known conformally-invariant model of Deser in 4D is shown to be a special case when we drop the dependence of the 4D metric on the extra dimension. Moreover, we show that the breakdown of conformal invariance which was introduced in by an ad hoc non-conformal invariant term inserted into the action naturally emerges here by $`i)`$ assuming a weak (cosmological) dependence of the 4D metric on the 5th dimension<sup>1</sup><sup>1</sup>1This assumption is reasonable since 4D general relativity is known to be in a very good agreement with present observations. and $`ii)`$ approximating the scalar field with its cosmological background value using the well-known cosmological coincidence usually referred to Mach or Wheeler.
This geometric approach to the subject of conformal invariance and its breakdown in gravitational models accounts properly for coupling of the gravitational field with its source in 5D gravity. It also gives an explanation for the origin of a small cosmological constant emerging from non-compactified extra dimension. This subject is the most recent interest in theories with large extra dimensions .
The paper is organized as follows: In section 2, we briefly review the conformal invariant gravitational model and its breakdown in 4 dimensions due to Deser . In section 3, we introduce a generalized conformal invariant gravitational model in 5 dimensions. In section 4, we study the breakdown of conformal invariance in 5 dimensions and discuss on some relevant interpretations. The paper ends with a conclusion.
## 2 Breakdown of conformal invariance in 4D
In this section we briefly revisit the standard work in 4D conformal invariance due to Deser . Consider the action functional
$$S[\varphi ]=\frac{1}{2}d^4x\sqrt{g}(g^{\alpha \beta }_\alpha \varphi _\beta \varphi +\frac{1}{6}R\varphi ^2),$$
(1)
which describes a system consisting of a real scalar field $`\varphi `$ non-minimally coupled to gravity through the scalar curvature $`R`$. Variations with respect to $`\varphi `$ and $`g_{\alpha \beta }`$ lead to the equations
$$(\mathrm{}\frac{1}{6}R)\varphi =0$$
(2)
$$G_{\alpha \beta }=6\varphi ^2\tau _{\alpha \beta }(\varphi ),$$
(3)
where $`G_{\alpha \beta }=R_{\alpha \beta }\frac{1}{2}g_{\alpha \beta }R`$ is the Einstein tensor and
$$\tau _{\alpha \beta }(\varphi )=[_\alpha \varphi _\beta \varphi \frac{1}{2}g_{\alpha \beta }_\gamma \varphi ^\gamma \varphi ]\frac{1}{6}(g_{\alpha \beta }\mathrm{}_\alpha _\beta )\varphi ^2,$$
(4)
with $`_\alpha `$ denoting the covariant derivative. Taking the trace of (3) gives
$$(\mathrm{}\frac{1}{6}R)\varphi =0,$$
(5)
which is consistent with equation (2). This is a consequence of the conformal symmetry of action (1) under the conformal transformations
$$\varphi \stackrel{~}{\varphi }=\mathrm{\Omega }^1(x)\varphi g_{\alpha \beta }\stackrel{~}{g}_{\alpha \beta }=\mathrm{\Omega }^2(x)g_{\alpha \beta },$$
(6)
where the conformal factor $`\mathrm{\Omega }(x)`$ is an arbitrary, positive and smooth function of space-time. Adding a matter source $`S_m`$ independent of $`\varphi `$ to the action (1) in the form
$$S=S[\varphi ]+S_m,$$
(7)
yields the field equations
$$(\mathrm{}\frac{1}{6}R)\varphi =0$$
(8)
$$G_{\alpha \beta }=6\varphi ^2[\tau _{\alpha \beta }(\varphi )+T_{\alpha \beta }],$$
(9)
where $`T_{\alpha \beta }`$ is the matter energy-momentum tensor. The following algebraic requirement
$$T=0,$$
(10)
then emerges as a consequence of comparing the trace of (9) with (8) which implies that only traceless matter can couple consistently to such gravity models.
We may break the conformal symmetry by adding a dimensional mass term -$`\frac{1}{2}d^4x\sqrt{g}\mu ^2\varphi ^2`$, with $`\mu `$ being a constant mass parameter, to the action (7). This leads to the field equations
$$(\mathrm{}\frac{1}{6}R+\mu ^2)\varphi =0$$
(11)
$$G_{\alpha \beta }+3\mu ^2g_{\alpha \beta }=6\varphi ^2[\tau _{\alpha \beta }(\varphi )+T_{\alpha \beta }].$$
(12)
and we obtain as a result of comparing the trace of (12) with (11)
$$\mu ^2\varphi ^2=T.$$
(13)
Now, the basic input is to consider the invariance breaking as a cosmological effect. This would mean that one may take $`\mu ^1`$ as the length scale characterizing the typical size of the universe $`R_0`$ and $`T`$ as the average density of the large scale distribution of matter $`\overline{T}MR_0^3`$, where $`M`$ is the mass of the universe. This leads, as a consequence of (13) to the estimation of the constant background value of $`\varphi `$
$$\overline{\varphi }^2R_0^2(M/R_0^3)^1R_0/MG_N,$$
(14)
where the well-known empirical cosmological relation $`G_NM/R_01`$ (due to Mach or Wheeler) has been used. In order to well-justify the results we will approximate the correspondence $`\overline{\varphi }^2G_N`$ with $`\overline{\varphi }^2\frac{8\pi }{6}G_N`$. This estimation for the constant background value of the scalar field is usually considered in Brans-Dicke type scalar-tensor gravity theories. Inserting this background value of $`\varphi `$ into the field equations (12) leads to the following set of Einstein equations
$$G_{\alpha \beta }+3\mu ^2g_{\alpha \beta }=6\overline{\varphi }^2T_{\alpha \beta }8\pi G_NT_{\alpha \beta },$$
(15)
with a correct coupling constant $`8\pi G_N`$, and a term $`3\mu ^2`$ which is interpreted as the cosmological constant $`\mathrm{\Lambda }`$ of the order of $`R_0^2`$. The field equation (11) for $`\overline{\varphi }`$ contains no new information. This is because it is not an independent equation, namely it is the trace of Einstein equations (15). One may easily check that using $`\mathrm{}\overline{\varphi }=0`$ and $`\overline{T}=\mu ^2\overline{\varphi }^2`$, equation (11) and the trace of equation (15) result in the same equation as $`\frac{1}{6}R+\mu ^2=0`$.
## 3 5D gravity and generalized conformal invariance
Consider the 5D metric given by
$$dS^2=\widehat{g}_{AB}dx^Adx^B=G\varphi ^2g_{\alpha \beta }dx^\alpha dx^\beta +dl^2$$
(16)
where the 5D line interval is written as the sum of a 4D part relevant to scalar-tensor theory and an extra part due to the 5th dimension. The capital Latin indices $`A,B,\mathrm{}`$ run over 0, 1, 2, 3, 4, Greek indices $`\alpha ,\beta ,\mathrm{}`$ run over 0, 1, 2, 3, and five dimensional quantities are denoted by hats. A constant $`G`$ is also introduced to leave $`G\varphi ^2`$ dimensionless. We proceed keeping $`g_{\alpha \beta }=g_{\alpha \beta }(x^\alpha ,l)`$ and $`\varphi =\varphi (x^\alpha )`$ as in modern Kaluza-Klein theory . The metric is general, since we have only used 4 of the available 5 coordinate degree of freedom to set the electromagnetic potentials, $`g_{4\alpha }`$ to zero.
The corresponding Christoffel symbols are obtained
$$\begin{array}{cc}\widehat{\mathrm{\Gamma }}_{\beta \gamma }^\alpha =\mathrm{\Gamma }_{\beta \gamma }^\alpha +\varphi ^1(\delta _\gamma ^\alpha _\beta \varphi +\delta _\beta ^\alpha _\gamma \varphi g_{\beta \gamma }^\alpha \varphi )\hfill & \\ & \\ \widehat{\mathrm{\Gamma }}_{\beta \alpha }^\alpha =\mathrm{\Gamma }_{\beta \alpha }^\alpha +4\varphi ^1_\beta \varphi \hfill & \\ & \\ \widehat{\mathrm{\Gamma }}_{\beta \gamma }^4=\frac{1}{2}_4\widehat{g}_{\beta \gamma }\hfill & \\ & \\ \widehat{\mathrm{\Gamma }}_{4\alpha }^\alpha =\frac{1}{2}\widehat{g}^{\alpha \beta }_4\widehat{g}_{\alpha \beta }\hfill & \\ & \\ \widehat{\mathrm{\Gamma }}_{\beta 4}^\alpha =\frac{1}{2}\widehat{g}^{\alpha \delta }_4\widehat{g}_{\delta \beta }\hfill & \\ & \\ \widehat{\mathrm{\Gamma }}_{\alpha 4}^4=\widehat{\mathrm{\Gamma }}_{44}^\alpha =\widehat{\mathrm{\Gamma }}_{44}^4=0,\hfill & \end{array}$$
(17)
where $`\widehat{g}_{\alpha \beta }=G\varphi ^2g_{\alpha \beta }`$. The 5D Ricci tensor can be written in terms of the 4D one plus other terms
$$\widehat{R}_{\alpha \beta }=R_{\alpha \beta }2\varphi ^1_\alpha _\beta \varphi +4\varphi ^2_\alpha \varphi _\beta \varphi \varphi ^2[\varphi \mathrm{}\varphi +^\alpha \varphi _\alpha \varphi ]g_{\alpha \beta }$$
$$+\frac{1}{2}G\varphi ^2[g^{\gamma \delta }_4g_{\delta \alpha }_4g_{\beta \gamma }\frac{1}{2}g^{\lambda \delta }_4g_{\alpha \beta }_4g_{\lambda \delta }_4^2g_{\alpha \beta }].$$
(18)
The field equations $`\widehat{R}_{AB}=0`$ then give
$$R_{\alpha \beta }=2\varphi ^1_\alpha _\beta \varphi 4\varphi ^2_\alpha \varphi _\beta \varphi +\varphi ^2[\varphi \mathrm{}\varphi +^\alpha \varphi _\alpha \varphi ]g_{\alpha \beta }\frac{1}{2}G\varphi ^2[g^{\gamma \delta }\dot{g}_{\delta \alpha }\dot{g}_{\beta \gamma }\frac{1}{2}g^{\lambda \delta }\dot{g}_{\alpha \beta }\dot{g}_{\lambda \delta }\ddot{g}_{\alpha \beta }]$$
(19)
$$\widehat{R}_{4\alpha }=_\alpha (k_\beta ^\alpha \delta _\beta ^\alpha k)=0withk_\beta ^\alpha =\frac{1}{2}\widehat{g}^{\alpha \delta }\dot{\widehat{g}}_{\delta \beta }=\frac{1}{2}g^{\alpha \delta }\dot{g}_{\delta \beta }$$
(20)
$$\widehat{R}_{44}=2(\dot{k}4k_\beta ^\alpha k_\alpha ^\beta )=0,$$
(21)
where an overdot denotes differentiation with respect to 5th coordinate $`l`$ ( see ). Equation (19) may lead to a set of 10 Einstein equations. Equation (20) which have the form of conservation law may also lead to a set of 4 Gauss-Codazzi equations for the extrinsic curvature $`k_\beta ^\alpha `$ of a 4D hypersurface $`\mathrm{\Sigma }_l`$ foliating in 5th dimension. Finally, equation (21) is one equation for the scalar combinations of the extrinsic curvature. The Ricci scalar for the space-time part is obtained by contracting equation (19) with the metric $`g_{\alpha \beta }`$
$$R=6\varphi ^1\mathrm{}\varphi \frac{1}{2}G\varphi ^2[g^{\alpha \beta }g^{\gamma \delta }\dot{g}_{\delta \alpha }\dot{g}_{\beta \gamma }\frac{1}{2}g^{\alpha \beta }g^{\lambda \delta }\dot{g}_{\alpha \beta }\dot{g}_{\lambda \delta }g^{\alpha \beta }\ddot{g}_{\alpha \beta }].$$
(22)
Combining equations (19) and (22) we obtain the Einstein-like equations with Einstein tensor $`G_{\alpha \beta }`$ in the left hand side and some terms of scalar field together with 4D metric and their covariant derivatives in the right hand side as follows
$$G_{\alpha \beta }=6\varphi ^2\tau _{\alpha \beta }(\varphi )+\frac{1}{2}G\varphi ^2[𝒯_{\alpha \beta }\frac{1}{2}𝒯g_{\alpha \beta }]$$
(23)
where
$$\tau _{\alpha \beta }(\varphi )=\frac{2}{3}_\alpha \varphi _\beta \varphi +\frac{1}{6}g_{\alpha \beta }^\gamma \varphi _\gamma \varphi \frac{1}{3}\varphi \mathrm{}\varphi g_{\alpha \beta }+\frac{1}{3}\varphi _\alpha _\beta \varphi $$
(24)
and
$$𝒯_{\alpha \beta }=g^{\gamma \delta }\dot{g}_{\delta \alpha }\dot{g}_{\beta \gamma }\frac{1}{2}g^{\lambda \delta }\dot{g}_{\alpha \beta }\dot{g}_{\lambda \delta }\ddot{g}_{\alpha \beta }.$$
(25)
It is easy to show that the tensor $`\tau _{\alpha \beta }`$ in equation (24) is exactly the same one in equation (4). The field equation for the scalar field may be obtained by contracting equation (23) with $`g_{\alpha \beta }`$ or $`\widehat{g}_{\alpha \beta }`$ as
$$\left(\mathrm{}\frac{1}{6}R+\frac{1}{12}G\varphi ^2𝒯\right)\varphi =0.$$
(26)
We notice that equation (26) has a dynamical mass term $`\frac{1}{12}G\varphi ^2𝒯`$ with the dimension of $`(mass)^2`$. In the presence of dimensional parameters, the conformal invariance can be established for a large class of theories if the dimensional parameters are conformally transformed according to their dimensions. In this regard, equation (26), although modified by the mass term compared to (5), but is still invariant under the generalized conformal transformations
$$\varphi \stackrel{~}{\varphi }=\mathrm{\Omega }^1(x,l)\varphi g_{\alpha \beta }\stackrel{~}{g}_{\alpha \beta }=\mathrm{\Omega }^2(x,l)g_{\alpha \beta }.$$
(27)
This is simply because the 5D metric (16) is invariant under the above conformal transformations. Obviously, the following combination
$$\widehat{G}_{\alpha \beta }\widehat{R}_{\alpha \beta }\frac{1}{2}\widehat{g}^{\gamma \lambda }\widehat{R}_{\gamma \lambda }\widehat{g}_{\alpha \beta }=\widehat{R}_{\alpha \beta }\frac{1}{2}g^{\gamma \lambda }\widehat{R}_{\gamma \lambda }g_{\alpha \beta }$$
is invariant under (27) due to the invariance of the metric $`\widehat{g}_{\alpha \beta }`$. Therefore, equation (23) which arises as a result of $`\widehat{G}_{\alpha \beta }=\widehat{R}_{\alpha \beta }\frac{1}{2}g^{\gamma \lambda }\widehat{R}_{\gamma \lambda }g_{\alpha \beta }=0`$ is invariant under (27). And equation (26) as a consequence of $`\widehat{g}^{\alpha \beta }\widehat{G}_{\alpha \beta }=0`$ or $`g^{\alpha \beta }\widehat{G}_{\alpha \beta }=0`$ is invariant under (27) as well, regardless of which metric is used to contraction since the right hand side is zero. Note that although the initial $`l`$-independent scalar field $`\varphi `$ transforms to an $`l`$-dependent one $`\stackrel{~}{\varphi }`$, but the $`l`$-dependent function $`\mathrm{\Omega }^1(x,l)`$ will not appear in the transformed scalar field equation because the metric also transforms in such a way that the function $`\mathrm{\Omega }^1(x,l)`$ is factored out throughout the transformed equation rendering the initial scalar field equation. Therefore, by pure 5D approach we are able to introduce a generalized conformal invariant gravitational model defined by equations (23), (26) and (27) subject to the subsidiary equations (20) and (21).
## 4 Breakdown of conformal invariance in 5D
Now, we are in a position to compare equations (26), (23) with the corresponding equations (11), (12). By this comparison it turns out that we are able to revisit the breakdown of conformal invariance in 4D by a 5D approach since we have derived the field equations (26), (23) which can be identified with (11), (12) in the broken phase of the conformal invariance in 4D.
To this end, we take a dimensional analysis. The dimension of $`𝒯_{\alpha \beta }`$ or $`𝒯`$ will no doubt be $`(length)^2`$. Now, we assume the cosmological effect $`\dot{g}_{\alpha \beta }\frac{1}{R_0}`$ which fixes a very slow variation of $`g_{\alpha \beta }`$ over the absolute cosmological scale $`R_0`$. This assumption leads to the breakdown of the conformal invariance since it means that we have fixed our standard of length by the scale of the universe and that (comparing equation (26) with (11) and using $`G\varphi ^21`$) $`𝒯`$ may be identified with 12$`\mu ^2`$ which is a constant mass term breaking the conformal invariance. Now, we put the above identification into the Einstein-like equation (23). We then have
$$G_{\alpha \beta }+3\mu ^2g_{\alpha \beta }=6\varphi ^2[\tau _{\alpha \beta }(\varphi )+\frac{1}{12}\varphi ^2𝒯_{\alpha \beta }],$$
(28)
which, comparing with equation (13), leads to the identification
$$T_{\alpha \beta }=\frac{1}{12}\varphi ^2𝒯_{\alpha \beta },$$
(29)
which is the desired result in the context of induced matter theory since the matter energy-momentum tensor $`T_{\alpha \beta }`$ is dynamically induced by the scalar field $`\varphi `$ and higher dimension, namely $`𝒯_{\alpha \beta }`$. Taking the trace of (29) we find
$$T=\frac{1}{12}\varphi ^2𝒯,$$
(30)
and by taking $`𝒯=12\mu ^2`$ we obtain the equation (13). Now, according to (30) we may discuss on the background value $`\overline{\varphi }`$ corresponding to the absolute cosmological scale $`R_0`$. We have already fixed $`𝒯`$ (or $`\mu ^2`$) by cosmological considerations, namely $`𝒯R_0^2`$. This was achieved by the 5th coordinate degree of freedom through $`\dot{g}_{\alpha \beta }R_0^1`$. The 5th coordinate degree of freedom accounts for the scalar field in the general metric (16). Thus, (see Eq.(14) and the following discussion) we may take a background value $`\overline{\varphi }`$, using this coordinate degree of freedom, as
$$\overline{\varphi }^2\frac{8\pi }{6}G_N,$$
(31)
which identifies $`G`$ with $`\frac{8\pi }{6}G_N`$ such that $`G\overline{\varphi }^21`$. This condition reduces the general metric (16) to the canonical one . If we now insert this constant background value $`\overline{\varphi }^2`$ into equation (28) and use (29) we find
$$G_{\alpha \beta }+3\mu ^2g_{\alpha \beta }8\pi G_NT_{\alpha \beta },$$
(32)
in which
$$T_{\alpha \beta }=\frac{1}{16\pi }G_N^1𝒯_{\alpha \beta }.$$
(33)
Equation (32) is the well-known Einstein equation in the broken phase of the conformal invariance with a cosmological constant $`\mathrm{\Lambda }=3\mu ^2`$ and a coupling of matter to gravity, $`G_N`$. The scalar field equation (26) is the trace of Einstein equations, so its information is already included in them (see the discussion in section 2).
Now, the relevance of 5D approach manifests. It relates the current upper bound value of the cosmological constant $`\mathrm{\Lambda }R_0^2`$ to a geometric phenomenon in which the cosmological constant is generated by the very slow variation of 4D metric with respect to 5th dimension<sup>2</sup><sup>2</sup>2In a recent work of Arkani-Hamed et al , a small effective cosmological constant is emerged from a large extra dimension in a non-compactified approach to 5D Kaluza-Klein gravity. Also, in a compactified model of Kaluza-Klein cosmology , smallness of the cosmological constant is related to smallness of the compactified dimension. Therefore, it seems that the subject of cosmological constant in higher dimensional (at least in 5D) models is inevitably involved with extra dimension.. Moreover, it unifies the origins of the matter and the cosmological constant in that they appear as “two manifestations of higher dimensional geometry”.
The traditional Einstein equation (32) may alternatively be written in its pure geometric form
$$G_{\alpha \beta }+3\mu ^2g_{\alpha \beta }\frac{1}{2}𝒯_{\alpha \beta },$$
(34)
in which the coupling constant $`G_N`$ is removed from theory. To say, although the Einstein tensor $`G_{\alpha \beta }`$ couples to the matter $`T_{\alpha \beta }`$ by $`G_N`$ but the matter itself couples by $`G_N^1`$ to the geometry $`𝒯_{\alpha \beta }`$ (33) and so the coupling $`G_N`$ is removed. In this level, the appearance of $`G_N`$ in the traditional Einstein equation seems to be a mathematical tool only for dimensional consistency. However, in the physical level equations (32) and (33) exhibit an interesting phenomenon, with varying $`G_N`$, in that if $`G_N`$ decreases with time leading to a weakly coupling of gravity $`G_{\alpha \beta }`$ to the matter $`T_{\alpha \beta }`$ (32), the matter itself will then be coupled strongly to the hidden geometry $`𝒯_{\alpha \beta }`$ (33). Regarding the present small value of $`G_N`$ we find an strong coupling of matter $`T_{\alpha \beta }`$ to the higher dimensional geometry $`𝒯_{\alpha \beta }`$. This strong coupling may account for non-observablity of the 5th dimension. In other words, the effects of the 5th dimension may be hidden behind this strong coupling and what we observe as the matter may be the manifestation of a weak effect of 5th dimension which is strengthened by a strong coupling $`G_N^1`$. This means that going back in time in $`G_N`$ varying scenarios we will encounter with an era $`G_N1`$ in which $`T_{\alpha \beta }`$ may decouple from $`𝒯_{\alpha \beta }`$ leading to a naked geometry of 5th dimension without the concept of matter, as indicated in Eq.(34). In conclusion it may be said that two equations (32) and (33) define dual weak-strong regimes, in 5D approach to coupling constants, and that equation (33) defines a dual-Einstein equation coupling matter to higher dimension.
It is worth noting that the conformal invariance in 4D may be easily recovered in this 5D approach by restricting the 4D metric $`g_{\alpha \beta }`$ to be independent of 5th dimension (simply by assuming Kaluza-Klein compactification condition for higher dimension). The relevant field equations in this choice are
$$R_{\alpha \beta }=2\varphi ^1_\alpha _\beta \varphi 4\varphi ^2_\alpha \varphi _\beta \varphi +\varphi ^2[\varphi \mathrm{}\varphi +^\alpha \varphi _\alpha \varphi ]g_{\alpha \beta },$$
(35)
where by taking the trace of (35) and combining it with (35) we obtain the conformal invariant equations (2) and (3). The origin of this conformal invariance in 4D turns out to be the invariance of the 4D part of 5D metric
$$\widehat{g}_{AB}(x^A)=\left(\begin{array}{cc}G\varphi ^2(x^\alpha )g_{\alpha \beta }(x^\alpha )& 0\\ & \\ 0& 1\end{array}\right)_,$$
(36)
under the conformal transformations (6).
## Conclusion
A key feature of any fundamental theory consistent with a given symmetry is that its breakdown would lead to effects which can have various manifestations of physical importance. Therefore, in the case of conformal symmetry in gravitational models, one would expect that the corresponding cosmological invariance breaking would have important effects generating the gravitational coupling and the cosmological constant. In this paper we have introduced a generalized conformally-invariant gravitational model of 5D gravity theory $`\widehat{R}_{AB}=0`$, with 4D part that is dependent on the extra dimension. The conformal invariance in 4D then becomes a special case when we take the 4D metric to be independent of the extra dimension. Moreover, we have shown that the cosmological breakdown of conformal symmetry in a conformally-invariant gravitational model in 4D may be naturally derived in this context if we assume a weak (cosmological) dependence of the 4D metric on the higher dimension and use the cosmological coincidence due to Mach or Wheeler to approximate the scalar field by its cosmological background value. This approach to the issue of couplings and parameters in gravity leads to a geometric interpretation for the small cosmological constant $`\mathrm{\Lambda }`$. Moreover, a dual coupling $`G_N^1`$ is introduced by which the matter couples strongly to the geometric effects of higher dimension through a dual Einstein equation, and non-observability of higher dimension is then justified.
We also mention to the generality of the 5D conformal invariance. In Deser’s model the conformal symmetry is broken once a constant mass term is introduced. However, in 5D approach a dynamical mass term is appeared without breaking the conformal symmetry. This generalized symmetry is broken when we take a preferred conformal frame by introducing an absolute length scale $`R_0`$ through $`\dot{g}_{\alpha \beta }R_0^1`$. In other words, what we call the conformal invariance in Deser’s model is not really a conformal invariance; it is just scale invariance which is a special case of conformal invariance. This is because the dimensional constant mass term could not transform conformally<sup>3</sup><sup>3</sup>3The conformal invariance is more general than scale invariance which is used in Deser’s model. If scale invariance is characterized by vanishing of the trace of the energy-momentum tensor, conformal invariance implies scale invariance in the absence of dimensional parameters in the theory..
There is a natural question in the context of induced matter theory about its possible connection to quantum theory. This is because we can induce the matter geometrically from the 5th dimension whereas we know the matter has a underlying quantum structure. Therefore, it deserves to pay attention to this issue. First, it is well-known that the existence of a dimensional gravitational constant $`G_N`$ is the main source of non-renormalizablity of quantum gravity. On the other hand, the quantum theory approach to the traditional Einstein equation suffers from the problem that the left hand side is geometry and the right hand side is the matter. Equation (34), however, as a pure geometric Einstein equation is free of $`G_N`$. Moreover, both sides of this equation has the geometric structure. Perhaps, it is helpful to study the 4D quantum gravity in this pure geometric 5D approach. Second, in a study of 4D conformal invariance in QFT in , the following equation like our scalar field equation (26) is obtained
$$\left(\mathrm{}\frac{1}{6}+\varphi ^2𝒮_\alpha ^\alpha \{\omega \}\right)\varphi =0$$
in which $`𝒮_\alpha ^\alpha \{\omega \}`$ is the trace of the tensor $`𝒮_{\alpha \beta }\{\omega \}`$ describing the distribution of matter due to local quantum effects. It is therefore very appealing to think about how the higher dimensional effects in 5D may play the role of quantum effects in 4D.
## Acknowledgment
F. Darabi would like to thanks B. Mashhoon, W. N. Sajko and L. de Menezes for useful comments.
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# A Class of Einstein-Maxwell Fields Generalizing the Equilibrium Solutions
## 1 Introduction: Einstein-Maxwell fields of rotating stationary sources
A stationary Einstein-Maxwell field is characterized by the two complex Ernst potentials $``$ and $`\mathrm{\Phi }`$ of the gravitational and electromagnetic fields, respectively. For the investigation of the internal symmetries of the system, it is more advantageous to use a pair of new potentials $`\xi `$ and $`q`$, defined by the relations
$$=\frac{\xi 1}{\xi +1}\mathrm{\Phi }=\frac{q}{\xi +1}.$$
(1)
In terms of the $`\xi `$ and $`q`$ potentials, the Lagrangian of the Einstein-Maxwell system takes the form
$$L=\theta ^2[_\mu \xi ^\mu \overline{\xi }_\mu q^\mu \overline{q}+(\xi _\mu qq_\mu \xi )(\overline{\xi }^\mu \overline{q}\overline{q}^\mu \overline{\xi })]$$
(2)
with
$$\theta =\xi \overline{\xi }+q\overline{q}1.$$
Here the metric $`g_{\mu \nu }`$ is that of the Euclidean 3-space in any suitable coordinate system. Hence the Euler-Lagrange equations follow:
$`\theta \mathrm{\Delta }\xi 2(\overline{\xi }\xi +\overline{q}q)\xi `$ $`=`$ $`0`$ (3)
$`\theta \mathrm{\Delta }q2(\overline{\xi }\xi +\overline{q}q)q`$ $`=`$ $`0.`$
The global symmetries of this system form the $`SU(2,1)`$ group. In the next section, we shall briefly review the basic theory of conserved currents and their relation to global symmetries. We then compute the conserved currents of the Einstein-Maxwell system and find a highly symmetrical form of the field equations using the currents. We introduce the concept of the ’swirl vector’ $`\stackrel{}{C}`$ which is a group invariant. In the subsequent sections the fields satisfying the condition of vanishing swirl are investigated.
## 2 Symmetries and Currents
The action as a functional of the fields $`\varphi ^i`$ and $`_\mu \varphi ^i`$ is
$$S=d^nx(\varphi ^i,_\mu \varphi ^i).$$
Associated with a symmetry transformation
$$\varphi ^i\varphi ^i+\delta \varphi ^i,$$
such that the Lagrangian density is invariant, $`\delta =0,`$ and
$$\delta \varphi ^i=\epsilon _k^i\varphi ^k,$$
(4)
there exists a conserved current
$$J^\mu =\frac{}{\left(_\mu \varphi ^i\right)}\epsilon _k^i\varphi ^k$$
(5)
(Noether current) satisfying
$$_\mu J^\mu =0.$$
We now proceed with the application of this general theory to the Einstein-Maxwell system. The Lagrangian $`\left(\text{2}\right)`$ possesses an $`SU(2,1)`$ global symmetry group . The potentials $`\alpha ,`$ $`\beta `$ and $`\gamma `$ belonging to the fundamental representation of this group are given by
$$\xi =\frac{\alpha }{\beta }q=\frac{\gamma }{\beta }.$$
(6)
The global invariance transformations are
$$\left(\begin{array}{c}\alpha \hfill \\ \beta \hfill \\ \gamma \hfill \end{array}\right)𝐔\left(\begin{array}{c}\alpha \hfill \\ \beta \hfill \\ \gamma \hfill \end{array}\right)$$
(7)
where $`𝐔SU(2,1)`$ is a constant matrix.
The defining representation of the global symmetry is given by the spinor potential $`\mathrm{\Psi }_A=(\alpha ,\beta ,\gamma ),`$ and its group adjoint $`\mathrm{\Psi }^A=(\overline{\alpha },\overline{\beta },\overline{\gamma })`$. Hence also the Ernst potentials are written
$$=\frac{\mathrm{\Psi }_1\mathrm{\Psi }_2}{\mathrm{\Psi }_1+\mathrm{\Psi }_2}\mathrm{\Phi }=\frac{\mathrm{\Psi }_3}{\mathrm{\Psi }_1+\mathrm{\Psi }_2}.$$
These spinors are determined up to an overall complex multiplying function; the equivalence classes are given by the relation
$$\mathrm{\Psi }_A^{}=\mathrm{\Omega }\mathrm{\Psi }_A$$
(8)
where $`\mathrm{\Omega }`$ is an arbitrary complex scalar.
The field equations $`\left(\text{3}\right)`$ will take the $`SU(2,1)`$ invariant form
$$[\mathrm{\Theta }^1(\mathrm{\Psi }_A\mathrm{\Psi }_B\mathrm{\Psi }_B\mathrm{\Psi }_A)]=2\mathrm{\Theta }^2\stackrel{}{C}(\mathrm{\Psi }_A\mathrm{\Psi }_B\mathrm{\Psi }_B\mathrm{\Psi }_A)$$
(9)
where $`\mathrm{\Theta }=\mathrm{\Psi }^{}\mathrm{\Psi }`$ is the $`SU(2,1)`$ norm of $`\mathrm{\Psi }`$.
We now turn our attention to the group invariant vector
$$\stackrel{}{C}=2i\mathrm{I}m\{\mathrm{\Psi }^{}\mathrm{\Psi }\}.$$
For static fields, all potentials are real, hence we have $`\stackrel{}{C}=0`$. We thus infer that the vector $`\stackrel{}{C}`$ is a concomittant of the rotation of the source, so that we will call $`\stackrel{}{C}`$ the swirl of the field. The form of the vector $`\stackrel{}{C}`$ depends on the gauge (8). In this paper, we consider Einstein-Maxwell fields for which it is possible to choose a gauge in which the swirl vanishes:
$$\stackrel{}{C}=0.$$
(10)
In addition to static metrics, this condition characterizes the equilibrium ($`|e|=m`$) class of fields.
The currents $`\left(\text{5}\right)`$ of the SU(2,1) symmetry can be expressed by use of the swirl $`\stackrel{}{C}`$ as follows:
$`J^{(1)}`$ $`=`$ $`\mathrm{\Theta }^2(\alpha \overline{\alpha }+\beta \overline{\beta })\stackrel{}{C}+\mathrm{\Theta }^1(\overline{\alpha }\alpha \alpha \overline{\alpha }+\overline{\beta }\beta \beta \overline{\beta })`$
$`J^{(2)}`$ $`=`$ $`\mathrm{\Theta }^2\alpha \overline{\beta }\stackrel{}{C}+\mathrm{\Theta }^1(\overline{\beta }\alpha \alpha \overline{\beta })`$
$`J^{(3)}`$ $`=`$ $`\mathrm{\Theta }^2\beta \overline{\alpha }\stackrel{}{C}+\mathrm{\Theta }^1(\overline{\alpha }\beta \beta \overline{\alpha })`$
$`J^{(4)}`$ $`=`$ $`\mathrm{\Theta }^2\alpha \overline{\gamma }\stackrel{}{C}+\mathrm{\Theta }^1(\overline{\gamma }\alpha \alpha \overline{\gamma })`$
$`J^{(5)}`$ $`=`$ $`\mathrm{\Theta }^2\gamma \overline{\alpha }\stackrel{}{C}+\mathrm{\Theta }^1(\overline{\alpha }\gamma \gamma \overline{\alpha })`$ (11)
$`J^{(6)}`$ $`=`$ $`\mathrm{\Theta }^2\beta \overline{\gamma }\stackrel{}{C}+\mathrm{\Theta }^1(\overline{\gamma }\beta \beta \overline{\gamma })`$
$`J^{(7)}`$ $`=`$ $`\mathrm{\Theta }^2\gamma \overline{\beta }\stackrel{}{C}+\mathrm{\Theta }^1(\overline{\beta }\gamma \gamma \overline{\beta })`$
$`J^{(8)}`$ $`=`$ $`\mathrm{\Theta }^2\gamma \overline{\gamma }\stackrel{}{C}+\mathrm{\Theta }^1(\overline{\gamma }\gamma \gamma \overline{\gamma }).`$
Here
$$\mathrm{\Theta }=\alpha \overline{\alpha }\beta \overline{\beta }+\gamma \overline{\gamma }$$
and the swirl has the detailed form
$$\stackrel{}{C}=\alpha \overline{\alpha }\beta \overline{\beta }+\gamma \overline{\gamma }\overline{\alpha }\alpha +\overline{\beta }\beta \overline{\gamma }\gamma .$$
From $`J^{(1)}`$ and $`J^{(8)}`$ we get
$`J^{(1a)}`$ $`=`$ $`\mathrm{\Theta }^2\alpha \overline{\alpha }\stackrel{}{C}\mathrm{\Theta }^1(\alpha \overline{\alpha }\overline{\alpha }\alpha )`$
$`J^{(1b)}`$ $`=`$ $`\mathrm{\Theta }^2\beta \overline{\beta }\stackrel{}{C}\mathrm{\Theta }^1(\beta \overline{\beta }\overline{\beta }\beta ).`$
For fields satisfying the condition of vanishing swirl, $`\stackrel{}{C}=0,`$ the Einstein-Maxwell Eqs. (9) can be written in the simple form
$$[\mathrm{\Theta }^1(\mathrm{\Psi }_A\mathrm{\Psi }_B\mathrm{\Psi }_B\mathrm{\Psi }_A)]=0$$
(12)
From the equations (11) of current conservation we get the symmetrical set
$$[\mathrm{\Theta }^1(\mathrm{\Psi }^A\mathrm{\Psi }_B\mathrm{\Psi }_B\mathrm{\Psi }^A)]=0.$$
(13)
Equations (12) and (13) govern the class of Einstein-Maxwell fields characterized by a vanishing swirl.
## 3 Formulation using the vectors $`G`$ and $`H`$:
An Einstein-Maxwell system with one Killing vector may be fully characterized by the complex 3-vectors :
$$𝐆=\frac{+2\overline{\mathrm{\Phi }}\mathrm{\Phi }}{2(\mathrm{R}e+\overline{\mathrm{\Phi }}\mathrm{\Phi })},𝐇=\frac{\mathrm{\Phi }}{(\mathrm{R}e+\overline{\mathrm{\Phi }}\mathrm{\Phi })^{1/2}}.$$
In the notation referring to the metric of the three-space, the field equations can be written
$$R_{\mu \nu }=G_\mu \overline{G}_\nu \overline{G}_\mu G_\nu +H_\mu \overline{H}_\nu +\overline{H}_\mu H_\nu $$
(14)
$`(𝐆)𝐆`$ $`=`$ $`\overline{𝐇}𝐇\overline{𝐆}𝐆`$ (15)
$`(𝐆)\times 𝐆`$ $`=`$ $`\overline{𝐇}\times 𝐇\overline{𝐆}\times 𝐆`$ (16)
$`(𝐆)𝐇`$ $`=`$ $`{\displaystyle \frac{1}{2}}(𝐆\overline{𝐆})𝐇`$ (17)
$`\times 𝐇`$ $`=`$ $`{\displaystyle \frac{1}{2}}(𝐆+\overline{𝐆})\times 𝐇.`$ (18)
The vectors $`𝐆`$and $`𝐇`$ can be expressed in terms of the gradients of the complex potentials as follows,
$`\mathrm{\Theta }𝐆`$ $`=`$ $`\left(\overline{\xi }+1q\overline{q}\right){\displaystyle \frac{\xi }{\xi +1}}+\overline{q}q`$ (19)
$`\mathrm{\Theta }^{1/2}𝐇`$ $`=`$ $`\left({\displaystyle \frac{\overline{\xi }+1}{\xi +1}}\right)^{1/2}\left(qq{\displaystyle \frac{\xi }{\xi +1}}\right).`$ (20)
Solving for the gradients, we have
$`\xi `$ $`=`$ $`\mathrm{\Theta }{\displaystyle \frac{\xi +1}{\overline{\xi }+1}}𝐆\mathrm{\Theta }^{1/2}\overline{q}\left({\displaystyle \frac{\xi +1}{\overline{\xi }+1}}\right)^{3/2}𝐇`$ (21)
$`q`$ $`=`$ $`\left(\overline{\xi }+1\right)^1q\mathrm{\Theta }𝐆+\left(\overline{\xi }+1q\overline{q}\right)\mathrm{\Theta }^{1/2}q\left(\xi +1\right)^{1/2}\left(\overline{\xi }+1\right)^{3/2}𝐇.`$
By use of the definitions $`\left(\text{6}\right)`$ and the vanishing of the vector $`\stackrel{}{C},`$ we get the relation
$$2i\mathrm{I}m\left(\frac{\xi \overline{\xi }+q\overline{q}}{\xi \overline{\xi }+q\overline{q}1}\right)=\left(\mathrm{ln}\frac{\overline{\beta }}{\beta }\right).$$
Inserting here the gradients $`\left(\text{21}\right),`$ we obtain the integrability conditions of $`\overline{\beta }/\beta `$ in the form:
$$𝐆\times \overline{G}=H\times \overline{H}.$$
(22)
These constraints are apparently milder than the equilibrium ($`R_{ij}=0`$) condition $`Re(𝐆\overline{G}H\overline{H})=0`$ characterizing the Einstein-Maxwell fields with balanced electromagnetic and gravitational forces.
With the help of the integrability condition (22), Eq. (16) takes the simple form
$$\times 𝐆=0.$$
(23)
Hence there exists a complex potential $`\psi `$ such that
$$𝐆=\mathrm{ln}\psi .$$
Using this in Eq. (18),
$$\times 𝐇=\frac{1}{2}(\mathrm{ln}\psi \overline{\psi })\times 𝐇.$$
Hence
$$𝐇=(\psi \overline{\psi })^{1/2}\chi $$
(24)
for some complex function $`\chi `$. Eq. (14) takes the form
$$R_{\mu \nu }=2\frac{\psi _{_{(,\mu }}\overline{\psi }_{,\nu )}\chi _{_{(,\mu }}\overline{\chi }_{,\nu )}}{\psi \overline{\psi }}.$$
The essential field equations can now be written down in terms of the complex potentials $`\psi `$ and $`\chi `$ as follows,
$`[\psi ^2(\overline{\psi }\psi \overline{\chi }\chi )]`$ $`=`$ $`0`$ (25)
$`(\psi ^2\chi )`$ $`=`$ $`0`$ (26)
$`\psi \times \overline{\psi }\chi \times \overline{\chi }`$ $`=`$ $`0.`$ (27)
These are two complex and one real equations for the two complex functions $`\chi `$ and $`\psi `$. Despite this apparent overdetermination of the system, at least two large classes of solutions exist. One of these consists of stationary fields that are $`SU(2,1)`$ rotations of a static state. This is due to the fact that the potentials are real for static fields and that the vector $`\stackrel{}{C}`$ is $`SU(2,1)`$ invariant. The other subset of solutions of these equations has the form $`\chi =\psi .`$ Field Eqs. (25)-(27) then reduce to the simple condition that $`\psi ^1`$ is a harmonic function. This is the equilibrium class characterized by a vanishing Ricci tensor.
## 4 Nonstatic fields
An example of the $`SU(2,1)`$ global transformations generating a nonstatic state from a given static electrovacuum is the Kramer-Neugebauer transformation , . This has the form
$$\alpha ^{}=(1z\overline{z})\alpha ,\beta ^{}=(1+z\overline{z})\beta 2\overline{z}\gamma ,\gamma ^{}=\left(\frac{\overline{z}}{z}\right)^{1/2}\left[(1+z\overline{z})\gamma 2z\beta \right]$$
(28)
where $`z`$ is a complex constant and $`(\alpha ,\beta ,\gamma )`$ is the real triplet of potentials for a static solution. Choosing $`(\alpha ,\beta ,\gamma )`$ to be the potentials of an asymptotically flat electrovacuum, the metric so obtained remains asymptotically flat.
There is a wide range of static electrovacuum space-times which is suitable to be chosen as the seed metric of global symmetry transformations. All vacuum Weyl metrics and their electrovacuum counterparts, characterized by the condition $`q=const.,`$ belong to these. A further simple case of an explicit static electrovacuum is the Bonnorized Kerr metric in oblate spheroidal coordinates $`(x,y)`$:
$$\alpha =x^2+p^2q^2y^2,\beta =2px,\gamma =2pqy$$
(29)
where $`p^2+q^2=1.`$ This solution describes the gravitational and magnetic field of a dipole source. An infinite sequence of static solutions is given by the Bonnorized Tomimatsu-Sato metrics. A wealth of nonstationary solutions is readily available by the $`SU(2,1)`$ rotations of these metrics. Among these, the rotating solutions obtained by the Kramer-Neugebauer transformation (28) will possess controllable asymptotic properties.
## 5 Conclusions
The Einstein-Maxwell fields with no swirl comprise a large number of space-times, and among those ones with astrophysical significance. In addition to the equilibrium ($`|e|=m`$) electrovacua, all static fields as well as their stationary counterparts resulting from the global $`SU(2,1)`$ symmetry belong to this class. To be able to better assess the extent of generality of this class, it would be clearly be of much help to find the solution of the following open research problem:
Find an example which is neither an equilibrium solution nor an $`SU(2,1)`$ rotation of a static field.
Likewise, it would be of considerable interest to see if the Kinnersley group $`K^{}`$ has a subgroup leaving the condition of vanishing swirl invariant.
## 6 Acknowledgments
I thank Mária Süveges and Mátyás Vasúth for computing the $`SU(2,1)`$ currents and for discussions. This work has been supported by the OTKA grant T031724.
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# Universal Similarity Factorization Equalities over Complex Clifford Algebras
## 1 Introduction
Let $`𝒞_n`$ be the complex Clifford algebra, with the identity 1, defined on $`n`$ generators $`e_1,e_2,\mathrm{},e_n`$ subject to the multiplication laws:
$`e_i^2`$ $`=1,i=1,2,\mathrm{},n,`$ (1)
$`e_ie_j+e_je_i`$ $`=0,ij,i,j=1,2,\mathrm{},n,`$ (2)
and $`e_1e_2\mathrm{}e_n\pm 1.`$ In that case $`𝒞_n`$ is spanned as a $`2^n`$-dimensional vector space with $`2^n`$ basis $`\{e_A\},`$ where the multi index $`A`$ ranges all naturally ordered subsets of the first positive integer set $`\{1,2,\mathrm{},n\};`$ the basis element $`e_A,`$ where $`A=(i_1,i_2,\mathrm{},i_k)`$ with $`1i_1<i_2<\mathrm{}<i_kn,`$ is defined as the product
$$e_Ae_{(i_1,i_2,\mathrm{},i_k)}=e_{i_1}e_{i_2}\mathrm{}e_{i_k},e_{A=\mathrm{}}=1.$$
For simplicity, the volume element $`e_{12\mathrm{}n}=e_1e_2\mathrm{}e_n`$ of $`𝒞_n`$ will be denoted by $`e_{[n]}`$ in the sequel. The square of the volume element is
$$e_{[n]}^2=(1)^{\frac{1}{2}n(n1)}e_1^2e_2^2\mathrm{}e_n^2.$$
In that case, all $`a𝒞_n`$ can be expressed as
$$a=\underset{A}{}a_Ae_A,a_A𝒞,$$
where $`A`$ ranges all naturally ordered subsets of $`\{1,2,\mathrm{},n\}.`$ We shall adopt the following notation from now on: $`𝒞_n:=𝒞\{e_1,\mathrm{},e_n\}.`$
Clifford algebras (real or complex) have been studied for many years and their algebraic properties are well known. In particular, all Clifford algebras are classified as matrix algebras, or as direct sums of matrix algebras over the fields of real or complex numbers, or the quaternion ring (see, e.g., ). For complex Clifford algebras $`𝒞_n,`$ it is well known that they can faithfully be realized as certain matrix algebras over $`𝒞,`$ and the general algebraic isomorphism is
$$𝒞_n\{\begin{array}{cc}𝒞(2^{\frac{n}{2}})\hfill & \text{if }n\text{ is even,}\hfill \\ {}_{}{}^{2}𝒞(2^{\frac{n1}{2}})\hfill & \text{if }n\text{ is odd,}\hfill \end{array}$$
(3)
where $`𝒞(s)`$ stands for the $`s\times s`$ total complex matrix algebra, and $`{}_{}{}^{2}𝒞(s)`$ stands for the complex matrix algebra
$${}_{}{}^{2}𝒞(s)=\left\{\left[\begin{array}{cc}A& 0\\ 0& B\end{array}\right]\right|A,B𝒞(s)\}.$$
In this article we improve this relationship to a new level by establishing a set of valuable universal similarity factorization equalities between elements of $`𝒞_n`$ and complex matrices over $`𝒞(2^{\frac{n}{2}})`$ or $`𝒞(2^{\frac{n1}{2}})`$ for all $`n.`$ Through these universal factorization equalities the complex matrix representations of elements in $`𝒞_n`$ can explicitly be established. Based on them various results in the complex matrix theory can directly be extended to complex Clifford numbers. Moreover, one can also easily develop matrix analysis over complex Clifford algebras.
We first establish some basic universal similarity factorization equalities for elements in Clifford algebras with dimensions $`2`$ and $`4,`$ which will be directly applied to establish some more general results for $`𝒞_n.`$
Lemma 1. Let $``$ be an algebraically closed field and $`_1=\{e|e^2=u\}`$ be a Clifford algebra defined on a generator $`e`$ with $`e^2=u`$ and $`u0.`$ Then all $`a_1`$ can be written as $`a=a_0+a_1e,`$ where $`a_0,a_1.`$ Moreover, define $`\overline{a}=a_0a_1e.`$ In that case, $`a`$ and $`\overline{a}`$ satisfy the following universal similarity factorization equality
$$P\left[\begin{array}{cc}\hfill a& \hfill 0\\ \hfill 0& \hfill \overline{a}\end{array}\right]P^1=\left[\begin{array}{cc}a_0+\sqrt{u}a_1& 0\\ 0& a_0\sqrt{u}a_1\end{array}\right],$$
(4)
where $`P`$ and $`P^1`$ have the following universal forms (no relation with $`a):`$
$`P`$ $`={\displaystyle \frac{1}{2}}\left[\begin{array}{cc}1+\frac{1}{\sqrt{u}}e& (\sqrt{u}e)\\ \frac{1}{u}(\sqrt{u}e)& 1+\frac{1}{\sqrt{u}}e\end{array}\right],`$ (7)
$`P^1`$ $`={\displaystyle \frac{1}{2}}\left[\begin{array}{cc}1+\frac{1}{\sqrt{u}}e& \sqrt{u}e\\ \frac{1}{u}(\sqrt{u}e)& 1+\frac{1}{\sqrt{u}}e\end{array}\right].`$ (10)
Proof. Note that $`e^2=u.`$ It is easy to verify that
$$\left[\begin{array}{cc}1& e\\ e^1& 1\end{array}\right]\left[\begin{array}{cc}a& \hfill 0\\ 0& \hfill \overline{a}\end{array}\right]\left[\begin{array}{cc}1& e\\ e^1& 1\end{array}\right]=\left[\begin{array}{cc}a_0& ua_1\\ a_1& a_0\overline{a}\end{array}\right].$$
On the other hand, it is also easy to verify that
$$\left[\begin{array}{cc}1& \sqrt{u}\\ \frac{1}{\sqrt{u}}& 1\end{array}\right]\left[\begin{array}{cc}a_0+ua_1& 0\\ 0& a_0ua_1\end{array}\right]\left[\begin{array}{cc}1& \sqrt{u}\\ \frac{1}{\sqrt{u}}& 1\end{array}\right]=\left[\begin{array}{cc}a_0& ua_1\\ a_1& a_0\overline{a}\end{array}\right].$$
Combining the above two equalities yields (4) and (10). $`\mathrm{}`$
Lemma 2. Let $`M_2()`$ be the $`2\times 2`$ total matrix algebra over an arbitrary field $``$ with its basis satisfying the multiplication rules
$$\tau _{pq}\tau _{st}=\{\begin{array}{cc}\tau _{pt},\hfill & q=s,\hfill \\ 0,\hfill & qs,\hfill \end{array}$$
(11)
for $`p,q,s,t=1,2.`$ Then all $`a=a_{11}\tau _{11}+a_{12}\tau _{12}+a_{21}\tau _{21}+a_{22}\tau _{22}M_2(),`$ where $`a_{pq},`$ satisfy the following universal similarity factorization equality
$$Q\left[\begin{array}{cc}a& 0\\ 0& a\end{array}\right]Q^1=\left[\begin{array}{cc}a_{11}& a_{12}\\ a_{21}& a_{22}\end{array}\right],$$
(12)
where $`Q`$ has the following universal form
$$Q=Q^1=\left[\begin{array}{cccc}\tau _{11}& \tau _{21}& & \\ \tau _{12}& \tau _{22}& & \end{array}\right].$$
(13)
Proof. Follows directly from a verification. $`\mathrm{}`$
Lemma 3. Let $``$ be an algebraically closed field of characteristic not two, and $`_2=\{e_1,e_2\}`$ be a Clifford algebra defined on two generators $`e_1,e_2`$ with $`e_1^2=u,e_2^2=v`$ and $`u0,v0.`$ Then all $`a_2`$ can be written as
$$a=a_0+a_1e_1+a_2e_2+a_3e_{12},$$
(14)
where $`a_0,\mathrm{},a_3.`$ In that case, $`aI_2`$ satisfies the following universal similarity factorization equalities
$$R\left[\begin{array}{cc}\hfill a& \hfill 0\\ \hfill 0& \hfill a\end{array}\right]R^1=\left[\begin{array}{cc}a_0+\sqrt{u}a_1& v(a_2+\sqrt{u}a_3)\\ a_2\sqrt{u}a_3& a_0\sqrt{u}a_1\end{array}\right],$$
(15)
and
$$T\left[\begin{array}{cc}\hfill a& \hfill 0\\ \hfill 0& \hfill a\end{array}\right]T^1=\left[\begin{array}{cc}a_0+\sqrt{v}a_1& u(a_1\sqrt{v}a_3)\\ a_1+\sqrt{v}a_3& a_0\sqrt{v}a_1\end{array}\right],$$
(16)
where
$$R=R^1=\frac{1}{2}\left[\begin{array}{cc}1+\frac{1}{\sqrt{u}}e_1& e_2\frac{1}{\sqrt{u}}e_{12}\\ \frac{1}{v}(e_2+\frac{1}{\sqrt{u}}e_{12})& 1\frac{1}{\sqrt{u}}e_1\end{array}\right],$$
(17)
and
$$T=T^1=\frac{1}{2}\left[\begin{array}{cc}1+\frac{1}{\sqrt{v}}e_2& e_1+\frac{1}{\sqrt{v}}e_{12}\\ \frac{1}{u}(e_1\frac{1}{\sqrt{v}}e_{12})& 1\frac{1}{\sqrt{v}}e_2\end{array}\right].$$
(18)
Proof. By Lemma 2, we take the change of basis of $`_2`$ as follows
$`\tau _{11}`$ $`={\displaystyle \frac{1}{2}}\left(1+{\displaystyle \frac{1}{\sqrt{u}}}e_1\right),`$ $`\tau _{12}`$ $`={\displaystyle \frac{1}{2v}}\left(e_2+{\displaystyle \frac{1}{\sqrt{u}}}e_{12}\right),`$ (19a)
$`\tau _{21}`$ $`={\displaystyle \frac{1}{2}}\left(e_2{\displaystyle \frac{1}{\sqrt{u}}}e_{12}\right),`$ $`\tau _{22}`$ $`={\displaystyle \frac{1}{2}}\left(1{\displaystyle \frac{1}{\sqrt{u}}}e_1\right).`$ (19b)
Then it is not hard to verify that this new basis satisfies the multiplication rules in (11). In this new basis, any $`a=a_0+a_1e_1+a_2e_2+a_3e_{12}_2`$ can be rewritten as
$$a=(a_0+\sqrt{u}a_1)\tau _{11}+(va_2+\sqrt{u}a_3)\tau _{12}+(a_2\sqrt{u}a_3)\tau _{21}+(a_0\sqrt{u}a_1)\tau _{22}.$$
Substituting this and (19) into (12) and (13), we can obtain (15) and (17). By the similar approach, we can get (16) and (18). $`\mathrm{}`$
## 2 Main results
Notice that $`𝒞`$ is an algebraically closed field. We then can establish a set of universal similarity factorization equalities for elements in $`𝒞_n`$ on the basis of Lemmas 1 and 3.
Theorem 4. Let $`a𝒞_1=𝒞\{e\}`$ be given. Then $`a`$ can be written as $`a=a_0+a_1e,`$ where $`a_0,a_1𝒞=\{x+iy|x,y\}.`$ Moreover denote by $`\overline{a}=a_0a_1e,`$ called the conjugate of $`a.`$ Then $`a`$ and $`\overline{a}`$ satisfy the following universal similarity factorization equality
$$P_1\left[\begin{array}{cc}\hfill a& \hfill 0\\ \hfill 0& \hfill \overline{a}\end{array}\right]P_1^1=\left[\begin{array}{cc}a_0+a_1i& 0\\ 0& a_0a_1i\end{array}\right],$$
(20)
where $`P_1`$ and $`P_1^1`$ have the following universal forms (no relation with $`a):`$
$$P_1=\frac{1}{2}\left[\begin{array}{cc}1ie& (ie)\\ (ie)& 1ie\end{array}\right],P_1^1=\frac{1}{2}\left[\begin{array}{cc}1ie& ie\\ ie& 1ie\end{array}\right].$$
(21)
Proof. Let $`=𝒞`$ and $`u=1`$ as in Lemma 1. Then (20) and (21) follow directly from (4) and (10). $`\mathrm{}`$
It is easy to verify that
$$\overline{\overline{a}}=a,\overline{a+b}=\overline{a}+\overline{b},\overline{ab}=\overline{a}\overline{b},\overline{\lambda a}=\overline{a\lambda }=\lambda \overline{a}.$$
(22)
hold for all $`a,b𝒞_1`$ and $`\lambda 𝒞.`$ According to (20), we define a map from $`𝒞_1`$ to the double field $`𝒞𝒞`$ by
$$\varphi _1:a=a_0+a_1e𝒞_1\left[\begin{array}{cc}a_0+a_1i& 0\\ 0& a_0a_1i\end{array}\right]𝒞𝒞.$$
(23)
Then it is easy to derive from (20) and (22) the following properties.
Corollary 5. Let $`a=a_0+a_1e,b=b_0+b_1e𝒞_1,\lambda 𝒞`$ be given, and $`\varphi _1`$ be defined by (23). Then
* $`a=b\varphi _1(a)=\varphi _1(b).`$
* $`\varphi _1(a+b)=\varphi _1(a)+\varphi _1(b),\varphi _1(ab)=\varphi _1(a)\varphi _1(b),\varphi _1(\lambda a)=\lambda \varphi _1(a),`$ $`\varphi _1(1)=I_2.`$
* $`\varphi _1(\overline{a})=\left[\begin{array}{cc}0& 1\\ 1& 0\end{array}\right]\varphi _1(a)\left[\begin{array}{cc}0& 1\\ 1& 0\end{array}\right].`$
* Denote $`a^\mathrm{\#}=\overline{a_0}\overline{a_1}e,`$ then $`\varphi _1(a^\mathrm{\#})=\varphi _1^{}(a),`$ the conjugate transpose of the complex matrix $`\varphi _1(a).`$
* $`a=\frac{1}{4}[1ie,ie]\varphi _1(a)[1ie,ei]^T.`$
* $`det[\varphi _1(a)]=a_0^2+a_1^2.`$
* $`a`$ is invertible $``$ $`\varphi _1(a)`$ is invertible, in which case $`\varphi _1(a^1)=\varphi _1^1(a).`$
The properties in Corollary 5(a) and (b) show that through the bijective map (23), the Clifford algebra $`𝒞_1`$ is algebraically isomorphic to the double field $`𝒞𝒞`$, and $`\varphi _1(a)`$ is a faithful matrix representation of $`a`$ in $`𝒞𝒞.`$
Notice that $`P_1`$ and $`P_1^1`$ in (20) have no relation with $`a.`$ Thus the equality in Theorem 4 can also be extended to all matrices over the complex Clifford algebra $`𝒞_1.`$
Theorem 6. Let $`A=A_0+A_1e𝒞_1^{m\times n}`$ be given, where $`A_0,A_1𝒞^{m\times n}.`$ Then $`A`$ and its conjugate $`\overline{A}=A_0A_1e`$ satisfy the following universal factorization equality
$`J_{2m}\left[\begin{array}{cc}A& 0\\ 0& \overline{A}\end{array}\right]J_{2n}^1`$ $`=\left[\begin{array}{cc}A_0+A_1i& 0\\ 0& A_0A_1i\end{array}\right],`$ (28)
where
$`J_{2m}`$ $`={\displaystyle \frac{1}{2}}\left[\begin{array}{cc}(1ie)I_m& (ie)I_m\\ (ie)I_m& (1ie_1)I_m\end{array}\right],`$ (31)
$`J_{2n}^1`$ $`={\displaystyle \frac{1}{2}}\left[\begin{array}{cc}(1ie_1)I_n& (ie)I_n\\ (ie)I_n& (1ie)I_n\end{array}\right].`$ (34)
In particular, if $`m=n,`$ then (28) becomes a universal similarity factorization equality over $`𝒞_1`$.
It is easy to verify that for any matrices $`A,B𝒞_1^{m\times n},C𝒞_1^{n\times p},`$ and $`\lambda 𝒞`$
$$\overline{\overline{A}}=A,\overline{A+B}=\overline{A}+\overline{B},\overline{AC}=\overline{A}\overline{C},\overline{\lambda A}=\overline{A\lambda }=\lambda \overline{A}.$$
(35)
Now according to (28), we define the complex matrix representation of a matrix $`A=A_0+A_1e𝒞_1^{m\times n}`$ by $`\mathrm{\Phi }_1(A)=\left[\begin{array}{cc}A_0+A_1i& 0\\ 0& A_0A_1i\end{array}\right].`$ Then the following properties can be easily derived from (28) and (35).
Corollary 7. Let $`A,B𝒞_1^{m\times n},C𝒞_1^{n\times p}`$ and $`\lambda 𝒞`$ be given. Then
* $`A=B\mathrm{\Phi }_1(A)=\mathrm{\Phi }_1(B).`$
* $`\mathrm{\Phi }_1(A+B)=\mathrm{\Phi }_1(A)+\mathrm{\Phi }_1(B).`$
* $`\mathrm{\Phi }_1(AC)=\mathrm{\Phi }_1(A)\mathrm{\Phi }_1(C),\mathrm{\Phi }_1(\lambda A)=\lambda \mathrm{\Phi }_1(A),\mathrm{\Phi }_1(I_m)=I_{2m}.`$
* $`\mathrm{\Phi }_1(\overline{A})=\left[\begin{array}{cc}0& I_m\\ I_m& 0\end{array}\right]\mathrm{\Phi }_1(A)\left[\begin{array}{cc}0& I_n\\ I_n& 0\end{array}\right].`$
* Let $`A=A_0+A_1e`$ and denote $`A^\mathrm{\#}=A_0^{}A_1^{}e,`$ where $`A_0^{}`$ and $`A_1^{}`$ are the conjugate transposes of the complex matrices $`A_0`$ and $`A_1.`$ Then $`\mathrm{\Phi }_1(A^\mathrm{\#})=\mathrm{\Phi }_1^{}(A),`$ the conjugate transpose of the complex matrix $`\mathrm{\Phi }_1(A).`$
* $`A=\frac{1}{4}[(1ie)I_m,(ie)I_m]\mathrm{\Phi }_1(A)[(1ie)I_n,(ei)I_n]^T.`$
* $`A`$ is invertible $``$ $`\mathrm{\Phi }_1(A)`$ is invertible, in which case $`\mathrm{\Phi }_1(A^1)=\mathrm{\Phi }_1^1(A).`$
* $`p_A(A)=0,`$ where $`p_A(x)=det[xI_{2m}\mathrm{\Phi }_1(A)].`$
Theorem 8. Let $`a𝒞_2=𝒞\{e_1,e_2\},`$ the complex quaternion algebra. Then $`a`$ can be written as
$$a=a_0+a_1e_1+a_2e_2+a_3e_{12},$$
where $`a_0,\mathrm{},a_3𝒞.`$ In that case, $`aI_2`$ satisfies the following universal similarity factorization equality
$$P_2\left[\begin{array}{cc}a& 0\\ 0& a\end{array}\right]P_2^1=\left[\begin{array}{cc}a_0+a_1i& (a_2+a_3i)\\ a_2a_3i& a_0a_1i\end{array}\right],$$
(36)
where $`P_2`$ has the universal form
$$P_2=P_2^1=\frac{1}{2}\left[\begin{array}{cc}1ie_1& e_2+ie_{12}\\ e_2+ie_{12}& 1+ie_1\end{array}\right].$$
(37)
Proof. Let $`=𝒞`$ and $`u=v=1`$ in Lemma 3. Then (36) and (37) follow from (15) and (17). $`\mathrm{}`$
According to (36), define a map from $`𝒞_2`$ to the $`2\times 2`$ complex matrix algebra $`𝒞^{2\times 2}`$ by
$$\varphi _2:a=a_0+a_1e_1+a_2e_2+a_3e_{12}𝒞_2\left[\begin{array}{cc}a_0+a_1i& (a_2+a_3i)\\ a_2a_3i& a_0a_1i\end{array}\right]𝒞^{2\times 2}.$$
(38)
Then we can easily derive from (36) the following properties.
Corollary 9. Let $`a,b𝒞_2=𝒞\{e_1,e_2\}`$ and $`\lambda 𝒞`$ be given. Then
* $`a=b\varphi _2(a)=\varphi _2(b).`$
* $`\varphi _2(a+b)=\varphi _2(a)+\varphi _2(b),\varphi _2(ab)=\varphi _2(a)\varphi _2(b),\varphi _2(\lambda a)=\lambda \varphi _2(a),`$ $`\varphi _2(1)=I_2.`$
* Let $`a=a_0+a_1e_1+a_2e_2+a_3e_{12}𝒞_2,`$ and denote $`a^\mathrm{\#}=\overline{a_0}\overline{a_1}e_1\overline{a_2}e_2\overline{a_3}e_{12}.`$ Then $`\varphi _2(a^\mathrm{\#})=\varphi _2^{}(a),`$ the conjugate transpose of the complex matrix $`\varphi _2(a).`$
* $`a=\frac{1}{4}[1ie,e_2+ie_{12}]\varphi _2(a)[1ie,e_2+ie_{12}]^T.`$
* $`det[\varphi _2(a)]=a_0^2+a_1^2+a_2^2+a_3^2.`$
* $`a`$ is invertible $`\varphi _2(a)`$ is invertible, in which case $`\varphi _2(a^1)=\varphi _2^1(a).`$
* The two elements $`a`$ and $`b`$ in $`𝒞_2`$ are similar, i.e., there is an invertible $`x𝒞_2`$ such that $`ax=xb`$ if and only if the two complex matrices $`\varphi _2(a)`$ and $`\varphi _2(b)`$ are similar over $`𝒞.`$
The properties in Corollary 9(a) and (b) clearly show that through the bijective map (38) the Clifford algebra $`𝒞_2,`$ i.e., the complex quaternion algebra, is algebraically isomorphic to the complex matrix algebra $`𝒞^{2\times 2},`$ and $`\varphi _2(a)`$ is a faithful matrix representation of $`a`$ in $`𝒞^{2\times 2}.`$
Notice that $`P_2`$ and $`P_2^1`$ in (37) have no relation to $`a.`$ Thus the equality in Theorem 8 can also be extended to all matrices over $`𝒞_2.`$
Theorem 10. Let $`A=A_0+A_1e_1+A_2e_2+A_3e_{12}𝒞_2^{m\times n}=𝒞^{m\times n}\{e_1,e_2\}`$ be given where $`A_0,\mathrm{},A_3𝒞^{m\times n}.`$ Then $`A`$ satisfies the following universal factorization equality
$`K_{2m}\left[\begin{array}{cc}A& 0\\ 0& A\end{array}\right]K_{2n}^1`$ $`=\left[\begin{array}{cc}A_0+A_1i& (A_2+A_3i)\\ A_2A_3i& A_0A_1i\end{array}\right],`$ (43)
where
$$K_{2t}=K_{2t}^1=\frac{1}{2}\left[\begin{array}{cc}(1ie_1)I_t& (e_2+ie_{12})I_t\\ (e_2+ie_{12})I_t& (1+ie_1)I_t\end{array}\right].$$
(44)
In particular, if $`m=n,`$ then (43) becomes a universal similarity factorization equality over $`𝒞_2.`$
According to (43), we define the complex representation of a matrix $`A=A_0+A_1e𝒞_2^{m\times n}`$ by $`\mathrm{\Phi }_2(A)=\left[\begin{array}{cc}A_0+A_1i& (A_2+A_3i)\\ A_2A_3i& A_0A_1i\end{array}\right].`$ Then the following properties are easy to verify by (43).
Corollary 11. Let $`A,B𝒞_2^{m\times n},C𝒞_2^{n\times p},`$ and $`\lambda 𝒞`$ be given. Then
* $`A=B\mathrm{\Phi }_2(A)=\mathrm{\Phi }_2(B).`$
* $`\mathrm{\Phi }_2(A+B)=\mathrm{\Phi }_2(A)+\mathrm{\Phi }_2(B).`$
* $`\mathrm{\Phi }_2(AC)=\mathrm{\Phi }_2(A)\mathrm{\Phi }_2(C),\mathrm{\Phi }_2(\lambda A)=\lambda \mathrm{\Phi }_2(A),\mathrm{\Phi }_2(I_m)=I_{2m}.`$
* Let $`A=A_0+A_1e_1+A_2e_2+A_3e_{12}`$ and denote $`A^\mathrm{\#}=A_0^{}A_1^{}e_1A_2^{}e_2A_3^{}e_{12},`$ then $`\mathrm{\Phi }_2(A^\mathrm{\#})=\mathrm{\Phi }_2^{}(A),`$ the conjugate transpose of the complex matrix $`\mathrm{\Phi }_2(A).`$
* $`A=\frac{1}{4}[(1ie_1)I_m,(e_2+ie_{12})I_m]\mathrm{\Phi }_2(A)[(1ie_1)I_n,(e_2+ie_{12})I_n]^T.`$
* $`A`$ is invertible $`\mathrm{\Phi }_2(A)`$ is invertible, in which case $`\mathrm{\Phi }_2(A^1)=\mathrm{\Phi }_2^1(A).`$
* $`p_A(A)=0,`$ where $`p_A(\lambda )=det[\lambda I_{2m}\mathrm{\Phi }_2(A)],`$ the characteristic polynomial of $`\mathrm{\Phi }_2(A).`$
* Two square matrices $`A`$ and $`B`$ are similar over $`𝒞_2,`$ i.e., there is an invertible matrix $`X`$ over $`𝒞_2`$ such that $`AX=XB`$ if and only if $`\mathrm{\Phi }_2(A)`$ and $`\mathrm{\Phi }_2(B)`$ are similar over $`𝒞.`$
In the next several results we only present the basic universal similarity factorization equalities without listing their operation properties and their extensions to matrices over $`𝒞_n.`$
Theorem 12. Let $`a𝒞_3=𝒞\{e_1,e_2,e_3\}`$ be given. Then $`a`$ can factor as
$$a=a_0+a_1e_{[3]},$$
(45)
where
$$a_0,a_1𝒞_2=𝒞\{e_1,e_2\},e_{[3]}^2=1.$$
Moreover, define $`\overline{a}=a_0a_1e_{[3]}.`$ In that case, the diagonal matrix $`D_a=\text{diag}(aI_2,\overline{a}I_2)`$ satisfies the following universal similarity factorization equality
$`P_3D_aP_3^1`$ $`=\left[\begin{array}{cc}\varphi _2(a_0)+\varphi _2(a_1)& 0\\ 0& \varphi _2(a_0)\varphi _2(a_1)\end{array}\right]`$ (48)
$`:=\varphi _3(a)^2𝒞^{2\times 2},`$ (49)
where $`\varphi _2(a_t),t=0,1,`$ is the matrix representation of $`a_t`$ in $`𝒞^{2\times 2}`$ defined in (38) and
$`P_3`$ $`={\displaystyle \frac{1}{2}}\left[\begin{array}{cc}(1+e_{[3]})P_2& (1e_{[3]})P_2\\ (1e_{[3]})P_2& (1+e_{[3]})P_2\end{array}\right],`$ (52)
$`P_3^1`$ $`={\displaystyle \frac{1}{2}}\left[\begin{array}{cc}P_2^1(1+e_{[3]})& P_2^1(1e_{[3]})\\ P_2^1(1e_{[3]})& P_2^1(1+e_{[3]})\end{array}\right],`$ (55)
where $`P_2`$ and $`P_2^1`$ are given by (37).
Proof. Notice that $`be_{[3]}=e_{[3]}b`$ holds for all $`b𝒞_2=𝒞\{e_1,e_2\}.`$ We have by applying (36) to $`a`$ in (45) that
$$P_2(aI_2)P_2^1=P_2(a_0I_2)P_2^1+P_2(a_1I_2)P_2^1e_{[3]}=\varphi _2(a_0)+\varphi _2(a_1)e_{[3]}:=\psi (a),$$
and
$$P_2(\overline{a}I_2)P_2^1=\varphi _2(a_0)\varphi _2(a_1)e_{[3]}:=\psi (\overline{a}).$$
Next we build, according to Lemma 1, a matrix and its inverse as follows
$`V`$ $`={\displaystyle \frac{1}{2}}\left[\begin{array}{cc}(1+e_{[3]})I_2& (1e_{[3]})I_2\\ (1e_{[3]})I_2& (1+e_{[3]})I_2\end{array}\right],`$ (58)
$`V^1`$ $`={\displaystyle \frac{1}{2}}\left[\begin{array}{cc}(1+e_{[3]})I_2& (1e_{[3]})I_2\\ (1e_{[3]})I_2& (1+e_{[3]})I_2\end{array}\right].`$ (61)
and then calculate to get
$$V\left[\begin{array}{cc}\psi (a)& 0\\ 0& \psi (\overline{a})\end{array}\right]V^1=\left[\begin{array}{cc}\varphi _2(a_0)+\varphi _2(a_1)& 0\\ 0& \varphi _2(a_0)\varphi _2(a_1)\end{array}\right].$$
Finally substituting $`\psi (a)=P_2(aI_2)P_2^1`$ and $`\psi (\overline{a})=P_2(\overline{a}I_2)P_2^1`$ into the left-hand side of the above equality yields (49), (52), and (55). $`\mathrm{}`$
Theorem 13. Let $`a𝒞_4=𝒞\{e_1,e_2,e_3,e_4\}`$ be given. Then $`a`$ can factor as
$$a=a_0+a_1e_{123}+a_2e_{124}+a_3e_{43}=a_0+e_{123}a_1+e_{124}a_2+e_{43}a_3,$$
(62)
where
$$a_0,a_1,a_2,a_3𝒞_2=𝒞\{e_1,e_2\},$$
$$e_{123}^2=1,e_{124}^2=1,e_{43}=e_{123}e_{124}=e_{124}e_{123}.$$
In that case, $`aI_4`$ satisfies the following universal similarity factorization equality
$`P_4(aI_4)P_4^1`$ $`=\left[\begin{array}{cc}\varphi _2(a_0)+\varphi _2(a_1)& \varphi _2(a_2)+\varphi _2(a_3)\\ \varphi _2(a_2)\varphi _2(a_3)& \varphi _2(a_0)\varphi _2(a_1)\end{array}\right]`$ (65)
$`:=\varphi _4(a)𝒞^{4\times 4},`$ (66)
where $`\varphi _2(a_t),t=0,\mathrm{},3,`$ is the matrix representation of $`a_t`$ in $`𝒞^{2\times 2}`$ defined in (38) and
$$P_4=P_4^1=\frac{1}{2}\left[\begin{array}{cc}(1+e_{[3]})P_2& (e_{124}e_{43})P_2\\ (e_{124}+e_{43})P_2& (1e_{[3]})P_2\end{array}\right],$$
(67)
where $`P_2`$ is given in (37).
Proof. Note that the commutative rules $`be_{123}=e_{123}b,be_{124}=e_{124}b,be_{43}=e_{43}b`$ hold for all $`b𝒞_2=𝒞\{e_1,e_2\}.`$ Thus, it follows from (36) that
$`P_2(aI_2)P_2^1`$
$`=P_2(a_0I_2)P_2^1+P_2(a_1I_2)P_2^1e_{123}+P_2(a_2I_2)P_2^1e_{124}+P_2(a_3I_2)P_2^1e_{43}`$
$`=\varphi _2(a_0)+\varphi _2(a_1)e_{123}+\varphi _2(a_2)e_{124}+\varphi _2(a_3)e_{43}`$
$`:=\psi (a).`$
Next building, according to Lemma 3, a matrix and its inverse as follows
$$V=V^1=\frac{1}{2}\left[\begin{array}{cc}(1+e_{123})I_2& (e_{124}e_{43})I_2\\ (e_{124}+e_{43})I_2& (1e_{123})I_2\end{array}\right],$$
and applying them to $`\psi (a)`$ given above, we obtain
$$V\left[\begin{array}{cc}\psi (a)& 0\\ 0& \psi (a)\end{array}\right]V^1=\left[\begin{array}{cc}\varphi _2(a_0)+\varphi _2(a_1)& \varphi _2(a_2)+\varphi _2(a_3)\\ \varphi _2(a_2)\varphi _2(a_3)& \varphi _2(a_0)\varphi _2(a_1)\end{array}\right].$$
Finally substituting $`\psi (a)=P_2(aI_2)P_2^1`$ into its left-hand side yields (66) and (67). $`\mathrm{}`$
By induction, we have the following two general results.
Theorem 14. Suppose that there is an independent invertible matrix $`P_n`$ over $`𝒞_n=𝒞\{e_1,\mathrm{},e_n\}`$ with $`n`$ even such that
$$P_n(aI_{2^{\frac{n}{2}}})P_n^1:=\varphi _n(a)𝒞^{2^{\frac{n}{2}}\times 2^{\frac{n}{2}}}=𝒞(2^{\frac{n}{2}})$$
(68)
holds for all $`a𝒞_n.`$ Now let $`a𝒞_{n+1}=𝒞\{e_1,\mathrm{},e_{n+1}\}.`$ Then $`a`$ can factor as
$$a=a_0+a_1e_{[n+1]}=a_0+e_{[n+1]}a_1,$$
(69)
where
$$a_0,a_1𝒞_n=𝒞\{e_1,\mathrm{},e_n\},e_{[n+1]}^2=(1)^{\frac{1}{2}(n+1)(n+2)}:=r.$$
Moreover define $`\overline{a}=a_0a_1e_{[n+1]}.`$ In that case, $`D_a=\text{diag}(aI_{2^{\frac{n}{2}}},\overline{a}I_{2^{\frac{n}{2}}})`$ satisfies the following universal similarity factorization equality
$`P_{n+1}D_aP_{n+1}^1`$ $`=\left[\begin{array}{cc}\varphi _n(a_0)+\sqrt{r}\varphi _n(a_1)& 0\\ 0& \varphi _n(a_0)\sqrt{r}\varphi _n(a_1)\end{array}\right]`$ (72)
$`:=\varphi _{n+1}(a)^2𝒞(2^{\frac{n}{2}}),`$ (73)
where
$$P_{n+1}=\frac{1}{2}\left[\begin{array}{cc}(1+\frac{1}{\sqrt{r}}e_{[n+1]})P_n& (\sqrt{r}e_{[n+1]})P_n\\ \frac{1}{r}(\sqrt{r}e_{[n+1]})P_n& (1+\frac{1}{\sqrt{r}}e_{[n+1]})P_n\end{array}\right],$$
(74)
and
$$P_{n+1}^1=\frac{1}{2}\left[\begin{array}{cc}P_n^1(1+\frac{1}{\sqrt{r}}e_{[n+1]})& P_n^1(\sqrt{r}e_{[n+1]})\\ P_n^1\frac{1}{r}(\sqrt{r}e_{[n+1]})& P_n^1(1+\frac{1}{\sqrt{r}}e_{[n+1]})\end{array}\right].$$
(75)
Proof. Note that the commutative rule $`be_{[n+1]}=e_{[n+1]}b`$ holds for all $`b𝒞_n=𝒞\{e_1,\mathrm{},e_n\}.`$ By applying (68) to (69), we obtain
$`P_n(aI_{2^{\frac{n}{2}}})P_n^1`$ $`=P_n(aI_{2^{\frac{n}{2}}})P_n^1+P_n(a_1I_{2^{\frac{n}{2}}})P_n^1e_{[n+1]}`$
$`=\varphi _n(a_0)+\varphi _n(a_1)e_{[n+1]}:=\psi (a).`$
and
$$P_n(\overline{a}I_{2^{\frac{n}{2}}})P_n^1=\varphi _n(a_0)\varphi _n(a_1)e_{[n+1]}:=\psi (\overline{a}).$$
Next, setting
$$V=\frac{1}{2}\left[\begin{array}{cc}(1+\frac{1}{\sqrt{r}}e_{[n+1]})I_{2^{\frac{n}{2}}}& (\sqrt{r}e_{[n+1]})I_{2^{\frac{n}{2}}}\\ \frac{1}{\sqrt{r}}(\sqrt{r}e_{[n+1]})I_{2^{\frac{n}{2}}}& (1+\frac{1}{\sqrt{r}}e_{[n+1]})I_{2^{\frac{n}{2}}}\end{array}\right],$$
$$V^1=\frac{1}{2}\left[\begin{array}{cc}(1+\frac{1}{\sqrt{r}}e_{[n+1]})I_{2^{\frac{n}{2}}}& (\sqrt{r}e_{[n+1]})I_{2^{\frac{n}{2}}}\\ \frac{1}{r}(\sqrt{r}e_{[n+1]})I_{2^{\frac{n}{2}}}& (1+\frac{1}{\sqrt{r}}e_{[n+1]})I_{2^{\frac{n}{2}}}\end{array}\right],$$
and applying them to $`D_a=\text{diag}(\psi (a),\psi (\overline{a}))`$ we get
$$V\left[\begin{array}{cc}\psi (a)& 0\\ 0& \psi (\overline{a})\end{array}\right]=\left[\begin{array}{cc}\varphi _n(a_0)+\sqrt{r}\varphi _n(a_1)& 0\\ 0& \varphi _n(a_0)\sqrt{r}\varphi _n(a_1)\end{array}\right].$$
Finally, substituting $`\psi (a)=P_n(aI_{2^{\frac{n}{2}}})P_n^1`$ and $`\psi (\overline{a})=P_n(\overline{a}I_{2^{\frac{n}{2}}})P_n^1`$ into its left-hand side yields the desired results. $`\mathrm{}`$
Theorem 15. Suppose that there is an independent invertible matrix $`P_n`$ over $`𝒞_n=𝒞\{e_1,\mathrm{},e_n\}`$ with $`n`$ even such that
$$P_n(aI_{2^{\frac{n}{2}}})P_n^1:=\varphi _n(a)𝒞(2^{\frac{n}{2}})$$
(76)
holds for all $`a𝒞_n.`$ Now, let $`a𝒞_{n+2}=𝒞\{e_1,\mathrm{},e_{n+2}\}.`$ Then it can factor as
$$a=a_0+a_1e_{[n]}e_{n+1}+a_2e_{[n]}e_{n+2}+a_3\mu _{n+2},$$
(77)
where
$$a_0,a_1,a_2,a_3𝒞_n=𝒞\{e_1,\mathrm{},e_n\},$$
$$(e_{[n]}e_{n+1})^2=(e_{[n]}e_{n+2})^2=(1)^{\frac{1}{2}(n+1)(n+2)}=r,$$
$$\mu _{n+2}=(e_{[n]}e_{n+1})(e_{[n]}e_{n+2})=(e_{[n]}e_{n+2})(e_{[n]}e_{n+1}).$$
In that case, $`aI_{2^{\frac{n+1}{2}}}`$ satisfies the following universal similarity factorization equality
$`P_{n+2}(aI_{2^{\frac{n+1}{2}}})P_{n+2}^1`$ $`=\left[\begin{array}{cc}\varphi _n(a_0)+\sqrt{r}\varphi _n(a_1)& r[\varphi _n(a_2)+\sqrt{r}\varphi _n(a_3)]\\ \varphi _n(a_2)\sqrt{r}\varphi _n(a_3)& \varphi _n(a_0)\sqrt{r}\varphi _n(a_1)\end{array}\right]`$
$`:=\varphi _{n+2}(a)𝒞(2^{\frac{n+1}{2}}),`$
where
$$P_{n+2}=\frac{1}{2}\left[\begin{array}{cc}(1+\frac{1}{\sqrt{r}}e_{[n+1]})P_n& (e_{[n]}e_{n+2}\frac{1}{\sqrt{r}}\mu _{n+2})P_n\\ \frac{1}{r}(e_{[n]}e_{n+2}+\frac{1}{\sqrt{r}}\mu _{n+2})P_n& (1\frac{1}{\sqrt{r}}e_{[n+1]})P_n\end{array}\right],$$
$$P_{n+2}^1=\frac{1}{2}\left[\begin{array}{cc}P_n^1(1+\frac{1}{\sqrt{r}}e_{[n+1]})& P_n^1(e_{[n]}e_{n+2}\frac{1}{\sqrt{r}}\mu _{n+2})\\ P_n^1\frac{1}{r}(e_{[n]}e_{n+2}+\frac{1}{\sqrt{r}}\mu _{n+2})& P_n^1(1\frac{1}{\sqrt{r}}e_{[n+1]})\end{array}\right].$$
The proof of this result is analogous to that of Theorem 13, and is therefore omitted here.
## 3 Conclusions
In this article we have established a set of universal similarity factorization equalities for elements over the complex Clifford algebra $`𝒞_n.`$ These equalities reveal two basic facts about $`𝒞_n:`$
* Each element $`a`$ in $`𝒞_n`$ has a complex matrix representation $`\varphi _n(a).`$ Moreover, a diagonal matrix constructed by the element $`a`$ is uniformly similar to its complex matrix representation $`\varphi _n(a).`$
* Conversely, each element $`a`$ in $`𝒞_n`$ could be regarded as an eigenvalue of its complex representation matrix $`\varphi _n(a).`$ In other words, all complex matrices with the form $`\varphi _n(a)`$ can uniformly be diagonalized over $`𝒞_n.`$
Based on the above two facts, one easily see that almost all known results in complex matrix theory can be extended to complex Clifford algebras. On the other hand, some problems related to complex matrices can also transform to the problems related to Clifford numbers. One such a problem (see ) is concerning the exponential $`e^A`$ of a complex matrix $`A.`$ In fact, we can see from Theorems 14 and 15 that for any $`a𝒞_n,`$ there is
$$P_n(e^aI)P_n^1=e^{\varphi _n(a)},$$
or
$$P_n\text{diag}(e^aI,e^{\overline{a}}I)P_n^1=e^{\varphi _n(a)},$$
which implies that the exponential of a complex matrix $`\varphi _n(a)`$ can be determined by the exponential of its corresponding Clifford numbers.
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# Resonant tunnelling and quenching of tunnel splitting in Wess-Zumino nanospin systems
Recently, the oscillations of the level splitting of the nanospin systems with magnetic field along the hard axis has been studied by many authors theoretically and experimentally. This oscillation phenomenon is semiclassically understood as a result of the interference between instanton and anti-instanton in the Euclidean space. However, although this instanton method gives exact evaluation of the oscillation of the ground level splitting, it does not give a satisfactory explanation of how the quenching disappears at higher magnetic field and how higher level splitting behaves with magnetic field.
On the other hand, there have been many works about complex periodic orbit theory for the tunnelling systems. This theory is based on the extension of well-known periodic orbit theory by including tunneling paths going through energy barrier. The tunnelling path have a minimal Euclidean action with a given energy, so they can be regarded as a ’bounce’, ’instanton’, or ’periodic instanton’. The advantage of complex periodic orbit theory is that for tunnelling systems it gives energy levels and their splittings throughout whole energy range. It is worth while to see how the geometrical phase can be incorperated with the complex periodic orbit theory, because, to our knowledge, for the systems with geometrical phase the theory has not been applied so far.
The purpose of this Letter is to apply the complex periodic orbit theory to the nanospin system showing geometrical phase effect, and to give semiclassical explanation for the existence of the level-splitting oscillations at low magnetic field and the absence of the oscillations at higher magnetic field for excited states as well as ground state. As will be shown clearly, we find out that the mechanism of the disappearance of quenching is the exact cancellation of quenching with the divergence of resonant tunnelling, which is very novel feature not reported so far. This strong correlation between quenching and resonant tunnelling seems to be inherent property in Wess-Zumino tunnelling systems.
Let us start with the Hamiltonian of system with magnetic field along the hard axis,
$$=K_1S_z^2+K_2S_y^2g\mu _BHS_z,$$
(1)
where $`K_1`$ and $`K_2`$ are anisotropy coefficients with $`K_1>K_2`$, and $`g`$ and $`\mu _B`$ are the gyromagnetic ratio and the Bohr magneton, respectively. The effctive Euclidean action for this Hamiltonian is obtained by expressing it in $`(\theta ,\varphi )`$\- representation and integrating over $`\theta `$,
$$S_E=iS_W+S\sqrt{\lambda }𝑑\tau \left[\frac{1}{2}m(\varphi )\dot{\varphi }^2+V(\varphi )\right]$$
(2)
where $`S_W`$ denotes the Wess-Zumino term which is expressed as
$$S_W=S𝑑\tau \dot{\varphi }(1\frac{h}{1\lambda \mathrm{sin}^2\varphi }),$$
(3)
and the position dependent mass and effective potential are
$`m(\varphi )`$ $`=`$ $`{\displaystyle \frac{1}{1\lambda \mathrm{sin}^2\varphi }},`$ (4)
$`V(\varphi )`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{sin}^2\varphi (1{\displaystyle \frac{h^2}{1\lambda \mathrm{sin}^2\varphi }}).`$ (5)
Here we put $`\lambda \frac{K_2}{K_1}`$ and $`h=\frac{H}{H_c}`$ with $`H_c\frac{2K_1S}{g\mu _B}`$. Note that the geometrical phase factor $`S_W`$ does not change under analytic continuation of Euclidean time into real time, i.e., $`\tau it`$. This means that one should consider the phase effect of $`S_W`$ for the path in real space in the same way as for the tunnelling path in the Euclidean space. This point of view is somewhat different from that of Ref. where only the phase effect of tunnelling path is taken. The lack of the contribution of real path to the phase effect has created unreasonable oscillation of level splitting with temperature.
Now, we apply complex periodic orbit theory to this effective Hamiltonian system with Wess-Zumino term. In the regime of $`h<h_11\lambda `$, as depicted in Fig. 1 (a), the potential shows a simple barrier between $`\varphi =0`$ and $`\pi `$. Let us first consider the energy range of $`E<E_b`$, $`E_b`$ being the energy of the barrier top. A periodic orbit starting from $`\varphi =\varphi _1`$ can be considered as a combination of four elements; (tunnelling segment + tunnelling segment ), (real segment + real segment), (tunnelling segment + real segment), and (real segment + tunnelling segment). With a prescription for Maslov index for tunnelling segment in Ref., we can write semiclassical trace of Green’s function in terms of a periodic orbit sum as
$`g(E)`$ $`=`$ $`{\displaystyle \frac{T}{i}}{\displaystyle \underset{m=0}{\overset{1}{}}}{\displaystyle \underset{N=0}{\overset{\mathrm{}}{}}}[(e^{\theta +i\frac{\pi }{2}})^2+(e^{iWi\frac{\pi }{2}})^2`$ (6)
$`+`$ $`(e^{\theta +i\frac{\pi }{2}})(e^{iWi\frac{\pi }{2}})e^{i\alpha }e^{im\pi }`$ (7)
$`+`$ $`(e^{iWi\frac{\pi }{2}})(e^{\theta +i\frac{\pi }{2}})e^{i\alpha }e^{im\pi }]^N`$ (8)
where $`T`$ is the period of the real orbit in the well, $`W`$ is the action for the real segment, i.e., $`W(E)=_{\varphi _1}^{\varphi _1}𝑑\varphi m(\varphi )\frac{d\varphi }{dt}`$, $`\theta `$ is the action for the tunnelling segment, i.e., $`\theta (E)=_{\varphi _1}^{\pi \varphi _1}𝑑\varphi m(\varphi )\frac{d\varphi }{d\tau }`$, and $`\alpha =S\pi (1\frac{h}{\sqrt{1\lambda }})`$ is the geometrical phase for an element corresponding to $`\pi `$ increment in $`\varphi `$ axis. The first two terms correspond to the elements (tunnelling segment + tunnelling segment) and (real segment + real segment), respectively. These elements represent oscillatory motions confined in the wells and barriers, which do not enclose $`z`$-axis and thus have no geometrical phase factor. While, the next two terms, (tunelling segment + real segment) and (real segment + tunnelling segment) represent half rotations in the direction of increasing $`\varphi `$ and decreasing $`\varphi `$, respectively, so that they have the geometrical phase factor $`e^{i\alpha }`$ or $`e^{i\alpha }`$ according to the rotation direction. The contributions of nonperiodic orbits ending at $`\varphi _1+\pi `$ in the combinations are removed by introducing $`m`$ sum. Using the same approximation as Ref., we can write the sum as
$$g(E)\underset{k}{}\frac{1}{E(E_k\pm \frac{2}{T}\mathrm{cos}\alpha e^\theta )}$$
(9)
where $`E_k`$ is the EBK(Einstein-Brillouin-Keller) energies determined by $`2W(E_k)=2\pi (k+1/2)`$, $`k`$ being integer, and $`T=\frac{d(2W)}{dE}`$. Since the poles of $`g(E)`$ indicate the eigenvalues of the system, we can see the oscillation of the energy splitting and the quenching fields are determined by $`\alpha (h)=\pi (n+1/2)`$, $`n`$ being integer. It is noted that these quenching fields do not shift in excited energy level splittings.
Now, consider the regime of $`h>h_1`$. In this case the barrier top becomes local mininum and then equivalent two barriers and new meta-stable well appear as shown in Fig. 1 (b). Eq.(9) is still valid in the energy range of $`E<E_m`$. Consider the energy range of $`E_m<E<E_b`$, $`E_m`$ being the local minimum energy. In this structure one can expect that the resonant tunnelling may happen due to the states in the meta-stable well. As will be shown below, however, the resonant tunnelling does not take place, since it compensates the quenching of the level splittings. As a result, there is no more oscillations of level splittings. The periodic orbit sum can be done simply if one attaches a small correction
$$\left(1\frac{e^{2\theta _m+i\pi }}{1e^{2iW_mi\pi }}\right)^1$$
(10)
to each segment and replace the contribution of tunnelling segment by
$$\frac{e^{2\theta _m+i\pi }e^{iW_mi\pi /2}}{1e^{2iW_mi\pi }}$$
(11)
where $`W_m(E)=_{\varphi _2}^{\pi \varphi _2}𝑑\varphi m(\varphi )\frac{d\varphi }{dt}`$ and $`\theta _m(E)=_{\varphi _1}^{\varphi _2}𝑑\varphi m(\varphi )\frac{d\varphi }{d\tau }`$. The small correction, Eq.(10), comes from the consideration of possible paths, leaving from $`\varphi _1`$ or $`\varphi _1+\pi `$ and returning back after some oscillatory motion in the meta-stable well, and their repetition between two segments. In fact, this small correction does not play any important role under the approximation used in getting the final $`g(E)`$. The contribution of the new tunnelling segment, Eq.(11), shows all possible oscillatory motions in the meta-stable well between two equivalent tunnelling subsegments. Then, the trace of Green’s function can be written approximately as
$$g(E)\underset{k}{}\frac{1}{E(E_k\pm \frac{1}{T}\frac{\mathrm{cos}\alpha }{\mathrm{cos}W_m}e^{2\theta _m})}$$
(12)
Here, it is noted that the geometrical phase effect $`\mathrm{cos}\alpha `$ is modulated by the action of meta-stable state. This effect is the main point in this Letter and will be discussed in detail soon. Note that in the limit $`EE_m`$, Eq.(12) gives smaller energy splittings than those given by Eq.(9), i.e., shows a discrepancy of factor 2. This fact leads to a discontinuity at $`E=E_m`$ in energy splittings (see Fig.3). To remove this discrepancy at $`E=E_m`$, one should use an uniform approximation which may also give a smooth behavior near $`E_b`$ as well as near $`E_m`$. However, we do not use the uniform approximation because the simple complex periodic orbit theory used here is sufficient for our purpose of this Letter.
If the energy is larger than $`E_b`$, the particle motion is simple rotation and the poles of the trace of Green’s function appear when
$$S_{rot}(E)\pm 2\alpha =2\pi k,k=integer,$$
(13)
where the action $`S_{rot}(E)=_0^{2\pi }𝑑\varphi m(\varphi )\frac{d\varphi }{dt}`$ and this condition gives EBK energies $`E_k`$ for $`E>E_b`$. In this energy range the energy splitting does not come from tunnelling, but rather from the Zeeman effect. Indeed, this simple semicalssical evalution of EBK energies gives exact eigenvalues for uniaxial spin system of $`=S_z^2HS_z`$ which is not a tunnelling system but an example of simple rotational orbit. In this uniaxial spin system the mass is $`1/2`$, the potential is $`H^2/4`$, and $`\alpha =\pi (SH/2)`$, and then from Eq.(13) the EBK energies is given as $`E_k=k^2kH`$ where $`k`$ is integer or half integer according to the spinvalue $`S`$.
All actions, $`W(E)`$, $`\theta (E)`$, $`W_m(E)`$, $`\theta _m(E)`$, and $`S_{rot}(E)`$ are analytically calculated by the complete elliptic integrals of the first and third kind (see Table 1).
In Fig.2 (a) the whole energy spectrum with $`h`$ are shown for $`S=10`$ case. The solid line and the dashed line denote the barrier energy $`E_b`$ and the local minimum energy $`E_m`$, respectively, with $`h`$. In order to avoid complicated view, we draw only EBK energies for $`E<E_b`$ without showing the splitting of the EBK energies. For the comparison we draw the result of direct diagonalization for the same system. In this direct numerical calculation we take the energy such that $`E=0`$ at the well minimum. In fact, the effective potential in Eq.(5) has been derived after the same energy translation. We can see that the entire behaviour of energy spectrum with $`h`$ is very similar to the result of direct diagonalization except for the region of $`EE_b`$. As mentioned before, an uniform approximation would give a consistent explanation for this transition region. In our semiclassical evaluation of the energy spectrum the physics in the energy range of $`E>E_b`$ is very different from that in the range of $`E<E_b`$; the former is the Zeeman splitting and the latter is tunnelling splitting. However, in quantum mechanical result these seem to be mixed around $`EE_b`$ and no abrupt change at $`E=E_b`$ is shown.
Fig.3 shows the lowest three level splittings, $`\mathrm{\Delta }E_0`$, $`\mathrm{\Delta }E_1`$, and $`\mathrm{\Delta }E_2`$. Here, the solid lines describe the results of complex periodic orbit theory($`\mathrm{\Delta }E_i^s`$), while the dashed lines show the result of direct diagonalization($`\mathrm{\Delta }E_i^q`$). It should be emphasized that in the evaluation of the actions related to tunnelling segment, i.e., $`\theta (E)`$, $`W_m(E)`$, and $`\theta _m(E)`$, the energy value used is not the EBK energies $`E_k`$, but $`E_k^{}`$ determined from the condition $`2W(E_k^{})=2\pi k`$. This choice of energy is consistent with usual instanton method where the ground state splitting is represented in terms of Euclidean action at the well minimum energy $`E_{min}`$(action of instanton) instead of that at the ground state energy $`E_0`$. The discrepancies between the quantal results ($`\mathrm{\Delta }E_i^q`$) and the semiclassical results ($`\mathrm{\Delta }E_i^s`$) would be resolved provided that the quantum fluctuations around the tunnelling path are taken into account. It is very interesting to see that the quenching of the level splitting does not exist at higher $`h`$, and for $`\mathrm{\Delta }E_0`$ the number of quenching is $`S`$ and $`\mathrm{\Delta }E_i`$ has $`(Si)`$ quenchings in the range of $`0<h<1`$. This can be understood if one knows the role of $`\mathrm{cos}W_m`$ in Eq.(12). The factor $`1/\mathrm{cos}W_m`$ diverges at $`2W_m=2\pi (n+1/2)`$, $`n`$ being integer, showing the resonant tunnelling, which means that if $`E_k^{}`$ coincides with the EBK energy $`E_j^{meta}`$ in the meta-stable well the level splitting becomes divergent. The fact that in the range of $`E_m<E_i<E_b`$ the splittings $`\mathrm{\Delta }E_i`$ $`(i=0,1,2,3)`$ do not show any quenching and any resonant tunnelling implies that the quenching fields are same with those of resonant tunnelling. This exact coincidences are shown in Fig. 4. The dashed lines denote the EBK energies in the meta-stable well, $`E_j^{meta}`$. Note that the cross points with $`E_0^{}(=0)`$, $`E_1^{}`$, and $`E_2^{}`$ (the fields of resonant tunnelling) are located exactly on quenching fields given by the condition $`\mathrm{cos}\alpha (h)=0`$, which is the mechanism for the disappearance of the level-splitting oscillation.
In conclusion we examine the energy spectrum of the quantum tunnelling Wess-Zumino spin system by means of the complex periodic orbit theory. All important features in the spectrum can be understood within this semiclassical scheme. Especially, we find out that the absence of the oscillation of the level splittings at higher magnetic field has its physical origin in the exact cancellation of quenching with the resonant tunnelling phenomena. It is more reasonable to conclude that this cancellation is an inherent property of Wess-Zumino tunnelling system.
We thank Stephen Creagh and C.-H. Kim for useful discussions.
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# EXAFS indication of double-well potential for oxygen vibration in Ba1-xKxBiO3
## I Introduction
Though superconductivity in BaPb<sub>1-x</sub>Bi<sub>x</sub>O<sub>3</sub> (BPBO) was discovered significantly earlier than in cuprates, the question of the nature of superconducting state in this oxide as well as in cognate system Ba<sub>1-x</sub>K<sub>x</sub>BiO<sub>3</sub> (BKBO) is still unsolved.
The structures of crystal lattice of copper oxide high temperature superconductors (HTSC’s) and bismuth-based oxides have some important common characteristics. The both oxide classes have perovskite-like lattice with CuO<sub>n</sub> ($`n`$=4, 5, 6) or Bi(Pb)O<sub>6</sub> complexes joined by the common oxygen ions. In bismuthates, the intersection of octahedral complexes in the three crystallographic directions determines their three-dimensional cubic structure. The CuO<sub>n</sub> complexes are joined in CuO<sub>2</sub> planes, which makes the two-dimensional structure of copper-oxides.
Because of strong hybridisation of covalent Bi(Pb)$`6s`$, Cu$`3d`$ – O$`2p_\sigma `$ bonds, the above mentioned complexes are the most tightly bound items of the perovskite-like structure. Therefore such important peculiarities of perovskite structure as lattice instability in respect to soft tilting mode of CuO<sub>n</sub> or BiO<sub>6</sub> complexes (see for review ) and highly anisotropic thermal factors of oxygen ions vibration , which point out the large amplitude of rotation oscillations, are inherent to the both classes of superconducting oxides and cause anharmonic vibrations of oxygen atoms that may be described by movement in a double-well potential . These structural instabilities of perovskite-like lattice can be related with the transition to superconducting state .
The layered structure of copper oxide compounds, presence of several non-equivalent copper positions, and a number of different Cu-O bonds complicate the local structure analysis. At the same time, the simplicity of cubic three-dimensional structure of BPBO-BKBO systems makes the interpretation of experimental data easier to a great extent. Relatively low temperatures of superconducting transition $`T_c13`$ K in BaPb<sub>0.75</sub>Bi<sub>0.25</sub>O<sub>3</sub> and $`T_c30`$ K in Ba<sub>0.6</sub>K<sub>0.4</sub>BiO<sub>3</sub> , the values of superconducting gap $`2\mathrm{\Delta }(0)/kT_c=3.6\pm 0.1`$ for BPBO and $`2\mathrm{\Delta }(0)/kT_c=3.5\pm 0.5`$ for BKBO , and a sizeable oxygen isotope effect allow one to rely on standard BCS-theory of superconductivity, not excluding a possible realisation of other mechanisms. The simpler electronic structure of $`sp`$ valence band of BPBO-BKBO systems in comparison with $`dp`$ band of cuprates favours the establishment of relationship of crystal and electronic structures in these compounds.
However, even for these relatively simple systems there is no agreement so far on a number of crucial aspects: crystal structure symmetry and lattice dynamics , electronic structure , bismuth valence state . The most part of unusual properties of BaBiO<sub>3</sub>-family compounds mentioned in early review are still unexplained.
In addition, the average structural data contradicts the local ones. Integral crystallographic methods point to the simple cubic lattice in BKBO ($`x>0.37`$) . In contrast to this, the EXAFS-analysis of four nearest spheres of bismuth environment reveals a local tilting of octahedra by 4–5, and Raman spectra evidence for a local lowering of symmetry from simple cubic .
EXAFS-analysis of the nearest oxygen octahedral environment of Bi in metallic BKBO compound with $`x=0.4`$ previously was made in harmonic approach and pointed out at least one unresolved problem. Temperature dependence of Debye-Waller factor $`\sigma ^2(T)`$ found by Heald et al. can be described in Einstein approximation, which is suitable for harmonic systems, only if one takes into account the temperature independent factor of 0.0025 Å<sup>2</sup>. The attempt to explain this fact due to some static disorder with doping was denied since $`\sigma ^2(T)`$ of the next Bi-Ba and Bi-Bi shells show the absence of any disorder. Similar weakly varying $`\sigma ^2(T)`$ dependence was found in . Thus, it was observed that amplitude of low temperature Bi-O vibrations in metallic BKBO compound with $`x=0.4`$ is too high (more than twice larger than in metallic BaPbO<sub>3</sub> compound) .
Besides, our careful measurements of temperature dependent EXAFS-spectra of undoped BaBiO<sub>3</sub> which were also treated in harmonic approach showed absolutely abnormal increasing of $`\sigma ^2`$ values of both Bi-O bonds with temperature decreasing from 90 K. We explained these anomalies by the influence of anharmonic rotational vibrations of oxygen octahedra which become pronounced at low temperatures , the model of connection of superconductivity with rotation mode asymmetry was proposed .
The above contradictions in description for local structure lead us to necessity of EXAFS analysis of the nearest Bi-O shell in anharmonic approach.
In the present paper we report the results of temperature dependent EXAFS-investigation of BPBO-BKBO systems firstly treated in anharmonic approach based on the idea of such an analysis for apical oxygen atoms in YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-δ</sub> proposed by J. Mustre de Leon et al. We have analysed EXAFS function $`\chi (k)`$, using the new program “VIPER for Windows” , by construction of the model potentials of atomic vibrations, subsequent calculation of the pair radial distribution function, and calculation of the model $`\chi (k)`$. This new approach for the EXAFS-analysis of Bi-based oxides gives us an opportunity to investigate the character of oxygen atom vibrations and, in combination with our previous local electronic structure studies by XANES (x-ray absorption near edge structure) spectroscopy , to understand the nature of the structural phase transitions in these systems and to explain practically all above contradictions.
In Sec. II we describe the experimental details and present the procedure of data treatment. The general results are given in Sec. III. In Sec. IV we relate the local crystal structure with local electronic structure, which is the base of our processing of experimental data. In Sec. V we discuss the possible superconductive mechanism and demonstrate how the oscillations of oxygen atoms in double-well potential contributes to superconductivity in BKBO.
## II Experimental and data analysis
In this work we have investigated ceramic samples of BaPbO<sub>3</sub> and BKBO with $`x=`$ 0, 0.4, 0.5 synthesised as described in . The materials were examined by x-ray powder diffraction for phase purity. The samples were controlled by transport and susceptibility measurements. BaBiO<sub>3</sub> showed semiconductor-like behaviour, the compositions BKBO with $`x=0.4`$ and 0.5 showed typical metallic $`\rho (T)`$ dependence and superconducting properties with $`T_c30`$ K and $`T_c16`$ K, correspondingly. BaPbO<sub>3</sub> samples were metallic but not superconducting at all temperatures.
For the XAS measurements, a crushed fine powder was precipitated onto a micropore substrate. The thickness of samples was about two absorption lengths at the chosen absorption edge.
The x-ray absorption spectra were collected at D-21 line (XAS-13) of DCI (LURE,Orsay, France) synchrotron operated at energy 1.85 GeV and the average current $`250`$ mA of positron beam at the $`L_3`$ edges of Bi (13040.6 eV) and Pb (13426 eV). Energy resolution of the double-crystal Si monochromator (detuned to reject 50% of the incident signal in order to minimise harmonic contamination) with a 0.4 mm slit at 13 keV was about 2–3 eV. The low temperature measurements were carried out using a liquid helium circulation type cryostat with a temperature control of $`\pm 1`$ K.
The background in the experimental spectra was removed as described in , taking care to remove the low frequency oscillations. EXAFS-functions $`\chi (k)k^2`$ obtained from absorption spectra were Fourier transformed in the wave number range $`k`$ from 1.5 to 16.5 Å<sup>-1</sup>, using Kaiser-Bessel windowing function. Back Fourier transform was done using a Hanning window from $`1`$ to $`2`$ Å corresponding to the first Bi-O near-neighbour shell. In such a case, the number of independent experimental points was $`N_{exp}=2\mathrm{\Delta }k\mathrm{\Delta }r/\pi +211`$.
The model EXAFS-function for pair atomic absorber-scatterer oscillations is constructed as follows. Suppose we know the potential of these oscillations as a parametric function of interatomic distance. Solving the stationary Srödinger equation for the particle with reduced mass of the atomic pair , one obtains a pair radial distribution function (PRDF) of atoms in $`i`$-th sphere:
$$g_i(r)=N_i\underset{n}{}|\mathrm{\Psi }_n(r)|^2e^{E_n/kT}/\underset{n}{}e^{E_n/kT},$$
(1)
where $`N_i`$ is the coordination number, $`E_n`$ and $`\mathrm{\Psi }_n`$ are $`n`$-th energy level and its corresponding wave function. Given PRDF’s, the model EXAFS function calculated as
$$\chi (k)=\frac{1}{k}\underset{i}{}F_i(k)\underset{r_{min}}{\overset{r_{max}}{}}g_i(r)\mathrm{sin}[2kr+\varphi _i(k)]/r^2𝑑r,$$
(2)
where $`k=\sqrt{2m_e/\mathrm{}^2(EE_{th})}`$ is the photoelectron wave number referenced to the ionisation threshold $`E_{th}`$, $`r_{min}`$ and $`r_{max}`$ are determined by the windowing function of back Fourier transform. The phase shift $`\varphi _i(k)`$ and scattering amplitude $`F_i(k)`$ were calculated using FEFF-6 code for six-shell cluster with crystallographic data from neutron diffraction study and using default set of the FEFF-6 parameters. The potential parameters were extracted from the model-to-experimental EXAFS-function fits.
## III General results
In Fig. 1 (left) we show the experimental $`\chi (k)k^2`$ for BKBO x=0, 0.4, 0.5 and for BaPbO<sub>3</sub> measured at $`L_3`$ Bi(Pb) absorption edge at 7 K. Good signal-to-noise ratio seen even for maximal wave number values $`k16`$ Å<sup>-1</sup> indicates the high spectra quality. The absence of signal on the Fourier transform in the low-$`r`$ range Fig. 1 (right) testifies for correct background removal procedure.
### A BaBiO<sub>3</sub>
EXAFS-researches of BaBiO<sub>3</sub> confirm the results of crystallographic works . According to them, there exist two inequivalent bismuth positions characterised by two Bi-O bond lengths. The equality of coordination numbers of the two BiO<sub>6</sub> spheres points out that the BaBiO<sub>3</sub> structure represents the ordered alternation of small and large BiO<sub>6</sub> octahedra in barium lattice. Such an alternation together with static rotation distortion around axis produce the monoclinic distortion of cubic lattice (Fig. 2). As will be shown in Sec. IV, to the larger soft octahedra corresponds the configuration BiO<sub>6</sub>, and to the smaller rigid octahedra corresponds BiL<sup>2</sup>O<sub>6</sub>. Here, L<sup>2</sup> denotes the hole pair in antibonding Bi$`6s`$O$`2p_\sigma ^{}`$ orbital of the octahedral complex.
In Fig. 3(a) the experimental $`\chi (k)k^2`$ EXAFS for BaBiO<sub>3</sub> at 7 K for the first Bi-O near-neighbour shell is shown. The pronounced beating near 8 Å<sup>-1</sup> evidences for existence of at least two different lengths of Bi-O bonds. In all previous EXAFS-researches the EXAFS-function was quite successfully fitted in harmonic approximation as a sum of two independent harmonic functions with different bond lengths and Debye-Waller factors. However, temperature dependencies of Debye-Waller factors contradict to the harmonic Einstein model . This argues against the independence of oxygen atom vibrations in two neighbouring octahedra.
In this work, we model the oscillatory Bi-O potential as follows. Inequivalence of the two types of BiO<sub>6</sub> octahedra is due to presence or absence of a hole pair in the hybridised molecular orbitals Bi$`6s`$O$`2p_\sigma ^{}`$ (see Sec. IV). Suppose, the movement of oxygen atoms may transfer a hole pair from one octahedron to another. Such a movement exchanges the roles of two inequivalent octahedra and requires a double-well form of oscillatory Bi-O potential. Here, we take a parabolic form of each well $`U_1=\kappa _1(rr_1)^2/2`$ and $`U_2=\kappa _2(rr_2)^2/2`$ which are joined continuously. Given the calculated $`\chi (k)`$, defined by Eqs. (1) and (2), we perform a least-squares fit between the model and experimental $`\chi (k)`$ over the range $`k=`$ 2–16 Å<sup>-1</sup> (see Fig. 3(a)). The six parameters determined by the fit were $`E_{th}`$, $`r_1`$, $`r_2`$, $`\kappa _1`$, $`\kappa _2`$, and $`N`$. The number of parameters here is the same as for fitting in the harmonic approximation.
The analysis of parameters of double-well potential allows one to draw the following conclusions on the oscillations of oxygen atoms in BaBiO<sub>3</sub>. (i) At low temperatures there exist two explicit peaks on PRDF, that is why x-ray and neutron diffraction detect the static distortions in BaBiO<sub>3</sub>. (ii) Even at low temperatures, the tunnelling probability between the two wells is non-zero. As it will be discussed in Sec. IV, such a tunnelling is equivalent to the dynamic exchange BiL<sup>2</sup>O$`{}_{6}{}^{}`$ BiO<sub>6</sub> and explains observed activation conductivity in BaBiO<sub>3</sub>. (iii) The temperature rise leads to gradual increase of probability of interwell tunnelling and to the structural transition to cubic phase found at 750–800 K . (iv) Non-equidistance of energy levels implies that oscillations of breathing- or stretching types with several frequencies exist: $`\omega _0=E_1E_0`$, $`\omega _1=E_2E_0`$, $`\omega _2=E_3E_0`$, etc. The latter are gradually manifested at high temperatures. (v) The positions of minima of double-well potential in our model are the average positions. They are spaced widely at maximal octahedra tilting and spaced closely at minimal octahedra tilting. In the latter case the probability of interwell tunnelling is maximal. In this sense the tunnelling frequency $`\omega _0`$ is bounded by soft rotation mode frequency.
### B BaPbO<sub>3</sub>
The experimental $`\chi (k)k^2`$ EXAFS for BaPbO<sub>3</sub> for the first Pb-O near-neighbour shell at all temperatures represents a sinusoid with no beatings and phase breaks (Fig. 4) and fitted well, using a single parabolic potential. This means that all Pb-O bonds in BaPbO<sub>3</sub> are equivalent and the breathing-type distortion is absent.
### C Ba<sub>1-x</sub>K<sub>x</sub>BiO<sub>3</sub>
The experimental $`\chi (k)k^2`$ EXAFS function for BKBO ($`x=0.4`$) at 7 K for the first Bi-O near-neighbour shell is shown in Fig. 5(a), solid curve. Also, the model calculated in harmonic approach (as in ) is shown as a dash-dot curve. It is seen that in the range $`k12`$ Å<sup>-1</sup> the harmonic approach fails. This argues for the anharmonic Bi-O vibration behaviour similar to BaBiO<sub>3</sub> case also in superconducting compositions BKBO. To a possible tendency in the cubic superconducting materials towards the same type of distortions as in BaBiO<sub>3</sub> was pointed earlier .
The K doping of BaBiO<sub>3</sub> leads to partial replacement of the larger soft octahedra BiO<sub>6</sub> by the smaller rigid octahedra BiL<sup>2</sup>O<sub>6</sub> (see details in Sec. IV). This causes a decrease and disappearance of static breathing and tilting distortions, but keeps the different rigidities of Bi-O bonds. Hence, the movement of an oxygen atom depends on that to which neighbouring octahedra this atom belongs. If the neighbouring octahedra are different, the oxygen atom oscillates in double-well potential as in BaBiO<sub>3</sub>. If the octahedra are equal, the oxygen atom oscillates in single-parabolic potential as in BaPbO<sub>3</sub> (see Fig. 6). The statistical weights of these two cases depend on the potassium content $`x`$, and are $`(1x)`$ and $`x`$, correspondingly. Here, the force constants of the two parabolas are assumed to be equal since their independent varying does not improve the fits, and the number of fitting parameters equals 7. The resulting model curve is presented in Fig. 5(a), dotted. It is seen from Fig. 5(b) that the total PRDF is unsplit, so it is not surprising that crystallographic measurements reveal the cubic structure (e.g. at $`T=10`$ K $`a=4.2742(1)`$). Meanwhile, the first momentum of the radial distribution function derived from our model is in a good agreement with the diffraction data , which validates our calculated amplitudes and phases.
In this paragraph we address to the important question of statistical grounds for the choice among several possible models. Consider, first, the statistical chi-square function (labelled just like EXAFS-function, but this is the different value)
$$\chi ^2=\frac{N_{exp}}{M}\underset{i}{\overset{M}{}}\left(\frac{\chi _{exp}(k_i)\chi _{mod}(k_i)}{\epsilon _i}\right)^2,$$
(3)
where M is the number of data points in the fit, $`N_{exp}`$ was introduced above, $`\epsilon _i`$ are the individual errors in the experimental data points. The latter were calculated as average of $`\sqrt{1/I_0+1/I_t}`$, $`I_0`$ and $`I_t`$ being the intensities measured in the transmission EXAFS experiment. The $`\chi ^2`$ value must follow the $`\chi ^2`$ distribution law with degrees of freedom $`\nu =N_{exp}P`$, where $`P`$ is the number of parameters varied during the fit. That is $`\chi ^2`$ must be less than the critical value $`X_\nu ^c`$ of the $`\chi ^2`$ distribution with $`\nu `$ degrees of freedom and the confidence level $`c`$. For our model in Fig. 5(a) $`\nu _2=117=4`$ and $`\chi _2^2=5.3<X_4^{0.95}=9.5`$, but for single-Gaussian model $`\nu _1=114=7`$ and $`\chi _1^2=16.8>X_7^{0.95}=14.1`$. Thus, our model meets the $`\chi ^2`$-test while the simple harmonic model does not. It is quite natural that having increased the number of parameters we got decreased $`\chi ^2`$ value. But what the gain should be? The comparison between the two models can be performed on the basis of $`F`$-test. If the difference $`\chi _1^2\chi _2^2`$ is physically meaningful, not due to presence of the noise, that is if the simpler model cannot describe some features in principle, this difference must not follow the $`\chi ^2`$ distribution law with $`\nu _1\nu _2`$ degrees of freedom. Provided that $`\chi _2^2`$ follows the $`\chi ^2`$ distribution with $`\nu _2`$ degrees of freedom, the value $`F(\nu _1\nu _2,\nu _2)=(\chi _1^2/\chi _2^21)\nu _2/(\nu _1\nu _2)`$ must not follow Fisher’s $`F`$-distribution with $`\nu _1\nu _2`$ and $`\nu _2`$ degrees of freedom. That is $`F`$ must be greater than the critical value $`F_{\nu _1\nu _2,\nu _2}^c`$ of the $`F`$ distribution with $`\nu _1\nu _2`$ and $`\nu _2`$ degrees of freedom and the confidence level $`c`$. Comparing the two models in Fig. 5(a), we find $`F(3,4)=2.89`$, which is equal to $`F_{3,4}^{0.84}`$. Therefore, we can claim with 84% probability that we propose a better model than the simplest single-Gaussian.
The complete results for Bi(Pb) octahedral ($`N=6`$) oxygen environment in BaPbO<sub>3</sub> and Ba<sub>1-x</sub>K<sub>x</sub>BiO<sub>3</sub> with $`x=`$ 0, 0.4, 0.5 at various temperatures are listed in Table I. The uncertainties in the EXAFS distances and force constants are less than $`\pm 0.4\%`$ and $`\pm 50\%`$, respectively. (These increments cause an increase of less than $`10\%`$ in the value of misfit.) It is seen from Figs. 3 and 5 that, in general, the positions of potential minima are not equal to the positions of PRDF’s maxima. Because of that, among other parameters, we give the values of PRDF’s centers of gravity, i.e. values that can be defined crystallographically. For BaBiO<sub>3</sub>, the centers of gravity were calculated for the two peaks of PRDF separately.
The parameters of potentials were obtained through the fit and, of course, depend on the form of the model potentials. To elucidate this influence, we constructed the model EXAFS function, using the polynomial of degree four: $`U=\kappa (rr_0)^2/2+\xi (rr_0)^4+\kappa ^2/(16\xi )`$, where the last term is introduced to zero the minimum values (Fig. 7). The number of fitting parameters here is the same as for double-parabolic potential. It turned out, that for the both potential forms the mean frequencies of Bi-O oscillations are practically identical for all temperatures. Because of this, in the present work we use double-parabolic potentials with parameters of clear meaning (potential minima positions and force constants) rather than a polynomial potential with abstract coefficients. Since the model EXAFS function weakly depends on the shape and value of interwell energy barrier $`U_0`$, we can not determine the $`U_0`$ value from EXAFS measurements exactly.
It worth to notice that in the case of BaBiO<sub>3</sub> the above model of oxygen vibrations has some limitation due to the existence of the static rotation (tilting) distortion $`11^{}`$. Thus, the distances observed from EXAFS spectra treatment are measured along Bi-O-Bi zigzag curve, but not along the Bi-Bi () direction as in BKBO. This may result in some difference between $`\kappa 1`$ and $`\kappa 2`$ values (see Fig. 3(b) and Table I) and some deviation of calculated frequency from real ones in BaBiO<sub>3</sub>.
## IV Relationship between the local crystal and local electronic structures
### A BaBiO<sub>3</sub>
The co-existence in BaBiO<sub>3</sub> of two different types of octahedra with two different Bi-O bond lengths and strengths reflects the different electronic structures of BiO<sub>6</sub> complexes.
Octahedral complexes represent the most tightly bound items of the perovskite-like structures because of strong covalence of Bi$`6s`$-O$`2p_\sigma `$ bonds. The valence band structure of BaBiO<sub>3</sub> is determined by overlapping of Bi$`6s`$ and O$`2p`$ orbitals , and, owing to strong Bi$`6s`$-O$`2p_\sigma `$ hybridisation, the octahedra can be considered as quasi-molecular complexes . Each complex has ten electron levels consisting of a pair of bonding Bi$`6s`$-O$`2p_\sigma `$, six nonbonding O$`2p_\pi `$, and a pair of antibonding Bi$`6s`$O$`2p_\sigma ^{}`$ orbitals. A unit cell, which includes two octahedra, has 38 valence electrons (10 from two bismuth ions, 4 from two barium ions, and 24 from six oxygen atoms). However, the numbers of occupied states in the two octahedtral complexes are different: octahedron BiL<sup>2</sup>O<sub>6</sub> carries 18 electrons and has one free level or a hole pair L<sup>2</sup> in the upper antibonding Bi$`6s`$O$`2p_\sigma ^{}`$ orbital, in octahedron BiO<sub>6</sub> with 20 electrons both antibonding orbitals are filled (Fig. 8). It is quite natural that BiL<sup>2</sup>O<sub>6</sub> octahedra have stiff (quasi-molecular) Bi-O bonds and the smaller radius, and BiO<sub>6</sub> octahedra represent non-stable molecules with filled antibonding orbitals and the larger radius. Because the sum of two nearest octahedra radii overcomes the lattice parameter $`a`$, the octahedral system must tilt around axis, producing a monoclinic distortion in BaBiO<sub>3</sub> (see Fig. 2).
The assumption of equal electron filling for nearest octahedra (BiL<sup>1</sup>O$`{}_{6}{}^{}+`$BiL<sup>1</sup>O<sub>6</sub>) contradicts to experiments, since in this case equal Bi-O bond lengths and local magnetic ordering should be observed.
Therefore our new scheme of bismuth disproportionation 2BiL<sup>1</sup>O$`{}_{6}{}^{}`$ BiL<sup>2</sup>O<sub>6</sub>+BiO<sub>6</sub> is in full agreement with charge balance, presence of two types of octahedron complexes and absence of any local magnetic moment .
Because Ba<sup>2+</sup> in perovskite-type lattice is bound by pure ionic bond, the electron density is concentrated mainly in octahedra volume $`V_0`$ . Hence, the energy of the highest occupied level $`E_f`$ is related to the octahedron radius $`R`$ and the number of valence electrons $`N`$ in a unit cell as
$$E_fh^2N^{2/3}/m_eV_0^{2/3}h^2N^{2/3}/m_eR^2,$$
(4)
where $`h`$ is the Planck’s constant, $`m_e`$ is the electron mass. This qualitative relation leaves out of account the deviation of Fermi surface from sphere . Nevertheless, the value of $`E_f`$, which transforms to the Fermi level $`E_F`$ at spatial overlapping of equal octahedra in the crystal, is in close connection with the octahedron radius and with the number of valence electrons.
According to expression (4), the local electronic structure is connected with the local crystal structure. In the left of Fig. 9(a), the case of hypothetical simple cubic structure of BaBiO<sub>3</sub> is shown. Real monoclinic structure arises from combined breathing (alternating octahedra with different radii $`R`$) and rotation distortions, which, according to (4), leads to the lowering of the energy $`E_f`$ of the highest occupied Bi$`6s`$O$`2p_\sigma ^{}`$ orbital in BiO<sub>6</sub> octahedron in comparison with energy $`E_h`$ of unoccupied Bi$`6s`$O$`2p_\sigma ^{}`$ orbital in BiL<sup>2</sup>O<sub>6</sub> octahedron. Hence, the repetitive sequence of empty level $`E_h`$ on the background of fully filled valence band represents the electronic structure of the ground state of BaBiO<sub>3</sub> (Fig. 9(a), right). In such a system there are no free carriers and conductivity occurs at hopping of carrier pair from one octahedron to another with the activation energy $`E_a=E_hE_f`$. In this process the movement of hole pair in real space leads to change of large octahedra to small ones and vice versa or to a dynamic exchange BiL<sup>2</sup>O$`{}_{6}{}^{}`$ BiO<sub>6</sub>, in full accordance with vibrations in double-well potential showed in Fig. 3, and causes the hole-type conductivity. In this case, the activation energy $`E_a`$ means the pair localisation energy.
This picture agrees with results of investigation of temperature dependencies of conductivity and Hall coefficient: $`n(T)=n_0\mathrm{exp}(E_a^{}/k_BT)`$ , where the value of hole concentration $`n_0=1.1\times 10^{22}`$ cm<sup>-3</sup> coincides with the concentration of unit cells, and the activation energy $`E_a^{}=0.24`$ eV, if one takes into account that in the case of two-particle activation conductivity, the number of hole pairs is equal to the number of BiL<sup>2</sup>O<sub>6</sub> complexes (concentration of which is half $`n_0`$), and the activation energy $`E_a=2E_a^{}=0.48`$ eV. It should be mentioned that the possibility of the two-particle (bipolaron) conductivity as a most probable mechanism of conductivity in BaBiO<sub>3</sub> was pointed to earlier .
Pair splitting and hopping of a single electron from one octahedron to another cost energy and lead to electronic structure reconstruction of the both octahedral complexes. Such a splitting can be achieved under optical excitation, which was observed experimentally as a photoconductivity peak at the photon energy $`h\nu =1.9`$ eV . The optical gap $`E_g`$ is the energy difference between the excited and ground states. In the excited state, BaBiO<sub>3</sub> has the local lattice of equivalent octahedra BiL<sup>1</sup>O<sub>6</sub> and possesses the non-compensated spin. Admittedly, it must have anti-ferromagnetic ordering, as the ground state in undoped cuprates La<sub>2</sub>CuO<sub>4</sub> and Nd<sub>2</sub>CuO<sub>4</sub>. Such optical excitation leads to the local dynamical lattice deformation observed in Raman spectra as a breathing mode $``$570 cm<sup>-1</sup> of giant amplitude under resonant coinciding of photon energy of Ar<sup>+</sup> laser with $`E_g`$ . Using lasers with other quantum energies destroys this resonant effect and leads to abrupt decrease of breathing mode amplitude .
Thus, the scheme proposed accounts for the nature of the two energy gaps in BaBiO<sub>3</sub>. The activation energy $`E_a`$ appears in transport measurements and is connected with coherent delocalisation of hole pairs owing to dynamic exchange BiL<sup>2</sup>O$`{}_{6}{}^{}`$ BiO<sub>6</sub>. The optical gap is radically differs from a traditional gap in semiconductors. It corresponds to energy difference between the excited and ground states.
### B BaPbO<sub>3</sub>
In BaPbO<sub>3</sub>, each octahedron also has ten molecular orbitals, nine of which are filled and one, antibonding Pb$`6s`$O$`2p_\sigma ^{}`$, is unoccupied as in octahedron BiL<sup>2</sup>O<sub>6</sub>, since lead has one electron less than bismuth. Thus, all the octahedral complexes represent the bound molecules PbL<sup>2</sup>O<sub>6</sub> with equal radii. At spatial overlapping of PbL<sup>2</sup>O<sub>6</sub> in crystal BaPbO<sub>3</sub>, wave functions of Pb$`6s`$O$`2p_\sigma ^{}`$ states become delocalised, unoccupied levels $`E_h`$ of neighbouring octahedra split into a free carrier band, and, merging with the top of valence band consisting of initially filled Pb$`6s`$O$`2p_\sigma ^{}`$ orbitals, make the half filled conduction band (Fig. 9(b)). It is important to notice that unoccupied Bi(Pb)$`6s`$O$`2p_\sigma ^{}`$ orbitals in the case of absence of spatial overlapping in BaBiO<sub>3</sub> behave as localised hole pairs, but transform to the conduction band at the spatial overlapping in BaPbO<sub>3</sub>. Because of this, BaBiO<sub>3</sub> appears to be a semiconductor of $`p`$-type, and BaPbO<sub>3</sub> — a semimetal of $`n`$-type. Since the radius of PbL<sup>2</sup>O<sub>6</sub> complex is slightly greater than that of BiL<sup>2</sup>O<sub>6</sub> complex, the Fermi level $`E_F`$ in BaPbO<sub>3</sub>, according to expression (4), lies lower in comparison with $`E_f`$ in BaBiO<sub>3</sub>, which agrees with results of photoelectron spectroscopy .
Thus, the lead valence state PbL<sup>2</sup>O<sub>6</sub> in BaPbO<sub>3</sub> is similar to the bismuth state in BiL<sup>2</sup>O<sub>6</sub> complexes in BaBiO<sub>3</sub>. According to Fig. 4, oxygen atoms in BaPbO<sub>3</sub> oscillate in the parabolic potential without any exchange between equal PbL<sup>2</sup>O<sub>6</sub> octahedra. Therefore, there is no charge pair transfer in semimetallic BaPbO<sub>3</sub> and the itinerant electrons behave as usual Fermi liquid.
### C Ba<sub>1-x</sub>K<sub>x</sub>BiO<sub>3</sub>
The doping of BaBiO<sub>3</sub> by lead or potassium leads to decrease of integral structural lattice distortions and causes essential changes in both local crystal and electronic systems. Firstly, we discuss the changes in electronic structure of BKBO as a simpler case, and then will extend the obtained conclusions to the BPBO system.
Since K<sup>+</sup> ion has one valence electron instead of two of Ba<sup>2+</sup> ion, the substitution of each two Ba<sup>2+</sup> for two K<sup>+</sup> ions produces an additional unoccupied level or hole pair in a Bi$`6s`$O$`2p_\sigma ^{}`$ orbital and modifies the BiO<sub>6</sub> complex to the BiL<sup>2</sup>O<sub>6</sub>. As a result, the ratio of the numbers of BiL<sup>2</sup>O<sub>6</sub> and BiO<sub>6</sub> complexes changes from $`1:1`$ in BaBiO<sub>3</sub> to $`(1+x):(1x)`$ and equals $`7:3`$ in Ba<sub>0.6</sub>K<sub>0.4</sub>BiO<sub>3</sub> and $`3:1`$ in Ba<sub>0.5</sub>K<sub>0.5</sub>BiO<sub>3</sub>. The spatial overlapping of BiL<sup>2</sup>O<sub>6</sub> complexes appear, which, taking into account their small radii and rigid bonds, contracts the lattice. This leads to decrease and disappearance (at $`x>0.37`$) of both static rotation- and breathing-type distortions. The lattice is forced to contract, despite the ionic radii of Ba<sup>2+</sup> and K<sup>+</sup> are practically equal.
Structural changes are accompanied by essential changes in physical properties of BKBO: at $`x0.37`$ the phase transition insulator-metal occurs and the superconductivity arises that remains up to the dopant concentration $`x=0.5`$ corresponding to the solubility limit of potassium in BaBiO<sub>3</sub>. The type of the temperature dependence of conductivity changes from semiconducting to metallic one, Hall coefficient changes its sign, and, in the normal state, BKBO compound with $`x>0.37`$ behaves as a metal with $`n`$-type conductivity .
These changes are well described in the framework of the above scheme (Fig. 9 for BKBO). At the low dopant amount ($`x<0.37`$), the contraction of the larger (soft) and the stretching of the smaller (rigid) octahedra bring, according to expression (4), $`E_h`$ and $`E_f`$ energies close together and the activation energy $`E_a=E_hE_f`$ decreases. As dopant concentration rises, the number of BiL<sup>2</sup>O<sub>6</sub> octahedra increases, and at $`x>0.37`$, in the lattice arise the three-dimensional chains of spatially overlapped BiL<sup>2</sup>O<sub>6</sub> octahedral complexes , their unoccupied levels split and form the conductivity band, which is equivalent to the percolation threshold reaching that determines the insulator-metal phase transition. Here, the itinerant electrons show Fermi liquid behaviour with Fermi level $`E_F`$ as in BaPbO<sub>3</sub>. However, in contrast to the picture in BaPbO<sub>3</sub>, in the BKBO structure there exist complexes BiO<sub>6</sub> through which itinerant electrons can not move because all the levels in these complexes are occupied. The movement of electrons occupying the upper Bi$`6s`$O$`2p_\sigma ^{}`$ orbital is possible only in pairing state at the dynamic exchange BiL<sup>2</sup>O$`{}_{6}{}^{}`$ BiO<sub>6</sub>, since the unpairing costs energy consumption. But, in contrast to the case of BaBiO<sub>3</sub>, the static rotation- and breathing-type distortions of the octahedra disappear, their mean radii become approximately equal (and equal to half the lattice parameter), and the pair localisation energy approaches zero: $`E_a=E_hE_f0`$ (Fig. 9(d)). In this case, the delocalised carrier pairs can freely move through the octahedral system that explains the appearance of superconductivity in BKBO ($`x>0.37`$). This scheme of the local electronic structure is in full agreement with the picture of local oxygen vibration in double-well potential (see Figs.5,6,7).
The unpairing of electrons, as in BaBiO<sub>3</sub>, is possible under optical excitation and is manifested by a pseudo-gap observed in reflectivity spectra even in the metallic phase of BPBO-BKBO systems. For instance, in BKBO with $`x=0.4`$ the pseudo-gap is about 0.5 eV and about 0.6 eV in BPBO with $`x=0.25`$.
Our model suggests that in the metallic phase two carrier types are present: itinerant electrons and pairs of initially (in BaBiO<sub>3</sub>) localised carriers. The co-existence of the two carrier types was confirmed by experiments on investigation of conductivity, Hall effect, and thermoelectric power , as well as by photoemission spectra and Raman spectra investigations . These results are in a good agreement with observed zero-bias conductance . The change in carriers behaviour from localised to itinerant at doping of BaBiO<sub>3</sub> by potassium was pointed to in x-ray absorption spectra analysis at the O $`K`$ edge . The analysis of EPR spectra showed the presence of the localised carrier pairs, which also was confirmed by the observation of two-particle tunnelling in the normal-state of BKBO .
### D BaPb<sub>1-x</sub>Bi<sub>x</sub>O<sub>3</sub>
Practically the same changes in electronic structure arise at the doping of BaBiO<sub>3</sub> by lead. The electronic structure of octahedral PbL<sup>2</sup>O<sub>6</sub> complex is entirely equivalent to that of BiL<sup>2</sup>O<sub>6</sub> complex. Unoccupied levels of overlapping octahedra in BaPbO<sub>3</sub> form the conduction band, as at overlapping of BiL<sup>2</sup>O<sub>6</sub> in BKBO. At the doping of BaPbO<sub>3</sub> by bismuth ($`0<x<0.37`$), the case is similar to BKBO ($`x>0.37`$). However, in BPBO the following combinations of neighbouring octahedra are possible: PbL<sup>2</sup>O<sub>6</sub>-PbL<sup>2</sup>O<sub>6</sub>, PbL<sup>2</sup>O<sub>6</sub>-BiL<sup>2</sup>O<sub>6</sub>, PbL<sup>2</sup>O<sub>6</sub>-BiO<sub>6</sub>, BiL<sup>2</sup>O<sub>6</sub>-BiO<sub>6</sub>. This leads to different local shifts of $`E_h`$ and $`E_f`$ energies, depending on different octahedra neighbouring pairs. At the further increase of bismuth content, the number of filled BiO<sub>6</sub> octahedra rises, the tilting and breathing distortions enlarge, the pair localisation energy raises, the spatial overlapping of unoccupied levels in PbL<sup>2</sup>O<sub>6</sub> octahedra disappears, which destroys the metallic type of conductivity and the system becomes a semiconductor with $`p`$-type conductivity.
Thus, superconductive properties of BPBO in our model are connected with dynamic exchange PbL<sup>2</sup>O$`{}_{6}{}^{}`$ BiO<sub>6</sub>. Unfortunately, it is practically impossible to observe the double-well-potential oxygen vibrations for superconducting BPBO compositions ($`0<x<0.35`$) by EXAFS study because of overlapping of Pb $`L_3`$ and Bi $`L_3`$ edges. Though an analysis of EXAFS spectra for these compounds is possible under elaborate treatment procedure , the precise enough values can be obtained only for interatomic distances, but not for amplitude factors. For this reason, we do not present the EXAFS data of BPBO ($`0<x<1`$) in this paper.
### E Photoemission spectra anomalies
The analysis of photoemission spectra has revealed a serious contradiction concerning the experimental observation of the splitting of the Bi$`4f`$(5/2, 7/2) spectral lines in superconducting compounds BKBO at the absence of any peculiarities of those lines in parent BaBiO<sub>3</sub> . The contradiction consists in that that though the analysis of x-ray diffraction and EXAFS data points to existence of two different Bi-O bonds in BaBiO<sub>3</sub>, no splitting of Bi$`4f`$ doublet lines was observed, which indicates the absence of considerable difference in bismuth valence states. The doping of BaBiO<sub>3</sub> by potassium relieves the monoclinic distortion and equalises the Bi-O bond lengths up to $`a/2`$. However, the Bi$`4f`$ spectral lines become broaden and, at measurements on a high resolution spectrometer, even split , which, in contrast with simple cubic lattice in BKBO ($`x>0.37`$), points to existence of two different bismuth valence states.
The above scheme of local electronic structure completely resolves this issue. Indeed, in BaBiO<sub>3</sub> the binding energies of the Bi$`4f`$ core levels in different BiL<sup>2</sup>O<sub>6</sub> and BiO<sub>6</sub> octahedra are almost equal and hence, the Bi$`4f`$ doublet lines in photoemission spectra are unsplit (Fig. 10(a)). Contrary to that, in Ba<sub>0.6</sub>K<sub>0.4</sub>BiO<sub>3</sub>, the binding energies of the Bi$`4f`$ core levels in the different octahedra $`E_1^{}`$, $`E_2^{}`$ and $`E_1`$, $`E_2`$ differ on $`\mathrm{\Delta }E=E_hE_FE_a`$, which causes the broadening (see curve $`x=0.5`$ in Fig. 10(a) from and splitting (see curve $`x=0.4`$ in Fig. 10(b) from of the Bi$`4f`$ doublet lines. The value of the splitting coincides well with our estimation of $`E_a=0.48`$ eV, and the ratio of intensities $`I`$ of the core-level peaks corresponding to the binding energies in different octahedral complexes BiO<sub>6</sub> and BiL<sup>2</sup>O<sub>6</sub> is in accordance with the relative content of these complexes in Ba<sub>0.6</sub>K<sub>0.4</sub>BiO<sub>3</sub>: $`I(E_1^{})/I(E_1)I(E_2^{})/I(E_2)3/7`$.
## V On possible nature of superconducting state
Essentially anharmonic character of oxygen ion oscillations in soft rotation mode is the base of anharmonic model of the high temperature superconductivity . It was shown that anharmonic coupling constant $`\lambda _s`$ exceeds the constant in harmonic approach $`\lambda _{ph}`$:
$$\lambda _s/\lambda _{ph}J_s^2d^2\overline{\omega }/J_{ph}^2u^2\omega _s1,$$
(5)
which explains high critical temperatures in cuprate HTSC compounds as well as in BPBO-BKBO systems. Here, $`d`$ and $`\omega _s`$ are the amplitude and frequency of oscillation in a double-well potential, $`u^2`$ is the mean-square displacement of the ions with the mean frequency $`\overline{\omega }`$ in harmonic approach, $`J_{ph}^2`$ and $`J_s^2`$ are, respectively, the deformation potentials of harmonic and anharmonic oscillations averaged on the Fermi surface.
In that model, the oscillatory movement of oxygen ions in rotation mode with a large amplitude in the direction perpendicular to the Bi(Pb)-O bonds was considered. A considerable gain in value of the coupling constant in expression (5) is due to low frequency of soft rotation mode ($`\omega _s<\overline{\omega }`$) and the excess of its amplitude over amplitude of harmonic oscillations ($`d^2u^2`$).
Meanwhile, previous as well as recent ( and references therein) electronic structure calculations showed that such a movement hardly affects the Bi-O bond lengths, keeping $`sp(\sigma )`$ nearest-neighbour interaction nearly constant. Thus, despite the existence of double-well potential in rotation (tilting) mode, Meregalli and Savrasov obtained very small anharmonic contribution to $`\lambda `$. Also, they assumed the breathing-type oxygen vibrations to be harmonic with small amplitude and high frequency and found too small $`\lambda 0.3`$ to explain high $`T_c`$ in BKBO.
Our investigation firstly showed that in BKBO the double-well vibrations with strong deformation potential also exist along the Bi-O-Bi direction. Such vibrations have the breathing-like character and low frequency bounded by soft rotation mode. Thus, our results resolve the above issue and explain the reason for strong electron-phonon coupling in Bi-based oxides. Also, to these vibrations we can assign the low-frequency part, below 40 meV, of phonon density of states observed in neutron scattering for superconducting BKBO and absent in calculations . Besides, the anomalous phonon softening along direction becomes clear due to our observation of coherent breathing-like vibrations along axis.
Hardy and Flocken evaluated values of $`\lambda `$, using abstract model double-well potentials:
$$\lambda (T)=N(0)\underset{kk^{}}{\overset{\mathrm{FS}}{}}\underset{n^{}>n}{}\frac{|n|M_{kk^{}}|n^{}|^2}{E_n^{}E_n}(f_nf_n^{}),$$
(6)
where $`N(0)`$ is the density of electron states at the Fermi level, $`M_{kk^{}}`$ is the e-ph matrix element between electronic states $`|k`$ and $`|k^{}`$ on the Fermi surface, $`|n`$ and $`|n^{}`$ are the oscillatory states with energies $`E_n`$ and $`E_n^{}`$, $`f_n`$ and $`f_n^{}`$ are the thermal weighting factors; $`f_n=\mathrm{exp}(E_n/kT)/_n^{}\mathrm{exp}(E_n^{}/kT)`$.
Given the oscillatory energies from our EXAFS-experiment, we have calculated $`\lambda `$ in terms of arbitrary multiplicative constant which is proportional to $`N(0)`$ and $`M_{kk^{}}`$ (Fig. 11). Although “phonon” part of $`\lambda `$ is the strongest for BaBiO<sub>3</sub>, this composition is not superconducting. On the one hand, because of pair localisation energy $`E_a`$. On the other, the rotation oscillations of rigid BiL<sup>2</sup>O<sub>6</sub> octahedra and ordinary phonons (of stretching and bending types) differ in collective character of motion in BiO<sub>2</sub> planes. In BaBiO<sub>3</sub>, rigid BiL<sup>2</sup>O<sub>6</sub> octahedra are separated by soft BiO<sub>6</sub> octahedra and are not spatially overlapped. The partial replacement of the larger soft octahedra by the smaller rigid ones at the potassium doping leads to decrease and disappearance of localisation energy $`E_a`$ (Fig. 9(c,d)) and to spatial coherence of rotation oscillations with length of several lattice parameters, depending on the doping level $`x`$. For superconducting compositions, $`T_c`$ means the temperature up to which $`\lambda `$ is considerably decreased (Fig. 11) and/or the coherence in oscillations of neighbouring octahedra is thermally destroyed by ordinary phonons.
The ground state of BaBiO<sub>3</sub> can be considered as a bipolaronic state as well, and the movement of carrier pairs correlated with oxygen atoms oscillations at the dynamic exchange BiL<sup>2</sup>O$`{}_{6}{}^{}`$ BiO<sub>6</sub> evidences for possible application of bipolaron theory , which is supported by the small size of the pair (octahedron size) and by the existence of the pair state above $`T_c`$.
## VI Conclusion
From the EXAFS investigation of BPBO-BKBO systems, the local crystal structure peculiarities connected with non-equivalence of BiL<sup>2</sup>O<sub>6</sub>-BiO<sub>6</sub> octahedra were observed. It was found that oxygen vibrations are well described using double-well potentials, which leads to the strong electron-phonon coupling due to coherent modulation of Bi-O bond lengths with low frequency and causes correlated carrier pair movement with oxygen ions oscillations. The underlying relationship between the local crystal and local electronic structures was established that explains the full list of unusual properties of BPBO-BKBO systems: the existence of two energy gaps (transport $`E_a`$ and optical $`E_g`$) in BaBiO<sub>3</sub>; the mechanism of two-particle conductivity in BaBiO<sub>3</sub>; the co-existence of two different carrier types; pseudo-gap observation in metallic compositions of BKBO and BPBO; the observation of localised pairs from EPR spectra; the observation of low-frequency ($`<40`$ meV) phonons; the nature of concentration and temperature phase transitions; the contradictions between local (Raman, and EXAFS) and integral structural methods for metallic BKBO; the anomalies of XPS spectra in BaBiO<sub>3</sub> and BKBO. We observed the correlated movement of carrier pairs and low-frequency breathing-like oxygen octahedra vibrations, which corresponds to superconducting state. The model proposed combines some principal features of the real-space pairing , anharmonic models and bipolaron theory of high-$`T_c`$ superconductivity.
The likeness of rotation mode peculiarities of BiO<sub>6</sub> octahedra and CuO<sub>n</sub> complexes and the anomalies in temperature dependencies of Debye-Waller factors in cuprates allow one to hope that similar model approach can be applied to the Cu-based superconductors.
###### Acknowledgements.
We acknowledge LURE Program Committee for beamtime providing, Professors S. Benazeth and J. Purans for help during x-ray absorption measurements. We are also grateful to Professor A. P. Rusakov for high-quality BKBO samples and to Dr. A. V. Kuznetsov and Dr. A. A. Ivanov for helpful discussions. The work is supported by RFBR (Grant No. 99-02-17343) and Program “Superconductivity” (Grant No. 99010).
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# Entropy driven intermitency
## Abstract
This note reviews some physical aspects of the chaotic hypothesis in nonequilibrium statistical mechanics and attempts at the physical interpretation of the fluctuation theorem as a quantitative intermittency property.
The main assumption, for the foundation of a nonequilibrium statistical mechanics or a theory of developed turbulence, will be:
Chaotic hypothesis: Asymptotic motions of a chaotic system, be it a system of $`N`$ particles or a viscous fluid, can be regarded as motions of a mixing Anosov system, for the purposes of computing time averages.
The point of view goes back to Ruelle, , , and in this specific form has been proposed in . I shall not define here “mixing Anosov system”, see . It will suffice to say that Anosov systems are very well understood dynamical systems in spite of being in a sense the most chaotic: they are so well understood that they can be regarded as the paradigm of chaotic systems much as the harmonic oscillators are the paradigm of regular and orderly motions.
This immediately implies the existence and uniqueness of an invariant probability distribution $`\mu `$ which gives the statistics if the motions of the system in the sense
$$\underset{T\mathrm{}}{lim}T^1_0^TF(S_tx)𝑑t=\mu (dy)F(y)$$
(1)
for almost all initial data, i.e. outside a set of $`0`$ volume in phase space (volume being measured in the ordinary sense of the Lebesgue measure). Here $`S_t`$ denotes the time evolution map, solution of the differential equations of motion, and $`\mu `$ is called the SRB distribution.
Particular interest will be reserved to time reversible systems: i.e. systems for which there is an isometry $`I`$ of phase space such that
$$I^2=1,andIS_t=S_tI$$
(2)
Denoting $`\dot{x}=f(x)`$ the equations of motion and $`\sigma (x)`$ the divergence $`\mathrm{div}f(x)=_j_{x_j}f_j(x)`$ a key quantity to study will be
$$\sigma _+=\mathrm{time}\mathrm{average}\mathrm{of}\sigma $$
(3)
that shall be called average entropy production rate or average phase space contraction rate; it will be assumed that $`\sigma _+>0`$ (this quantity is, in general, non negative, ).
The main result on $`\sigma (x)`$ concerns the probability distribution in the stationary state $`\mu `$ of the quantity (an observable, i.e. a function of the phase space point $`x`$)
$$p=T^1_{T/2}^{T/2}\frac{\sigma (S_tx)}{\sigma _+}𝑑t$$
(4)
which will be called the $`T`$–average dimensionless entropy creation rate. The result is a symmetry property of its probability distribution $`\pi _T(p)`$: which can be written in the form $`\pi _T(p)\stackrel{def}{=}conste^{\zeta (p)T+o(T)}`$, defining implicitly $`\zeta (p)`$, as it follows, on general grounds, from the theory of mixing Anosov systems, cf. , see for a more mathematical statement. The function $`\zeta (p)`$ is called the large deviation rate for the observable $`p`$.
Then (see : the delicate continuous time extension of a similar result for discrete time)
Fluctuation theorem: The “rate function” $`\zeta (p)`$ verifies
$$\zeta (p)=\zeta (p)p\sigma _+$$
(5)
which is a parameterless relation, valid under the above hypotheses of chaoticity and of time reversibility. It is a “mechanical” identity valid for systems with arbitrarily many particles (i.e. $`N=1,2,\mathrm{}10^{23},\mathrm{}`$).
Mathematically the above result holds for mixing Anosov flows which are time reversible. The chaotic hypothesis extends the result to rather general systems: of course for such systems it is no longer a theorem much in the same way as the consequences of the ergodic hypothesis in equilibrium statistical mechanics are not theorems for most systems to which they are applied.
Therefore we can say that the chaotic hypothesis implies in equilibrium (when the equations of motion are Hamiltonian and therefore $`\sigma (x)\mathrm{\hspace{0.33em}0}`$) the ergodic hypothesis (quite clearly): hence it implies classical statistical mechanics, starting with Boltzmann’s heat theorem which says that $`(dU+pdV)/T`$ is an exact differential (with $`U,p,V,T`$ defined as time averages of suitable mechanical quantities, see ). Likewise, out of equilibrium and in reversible systems, the chaotic hypothesis implies a general relation that is given by the fluctuation theorem, valid for systems with arbitrarily large numbers of particles. Both the heat theorem and the fluctuation theorem are “universal”, i.e. parameterless, system independent relations. They could perhaps be considered a curiosity for $`N`$ small, but certainly the first, at least, is an important property for $`N=10^{23}`$.
It becomes, at this point, clear that one should attempt at an interpretation of the fluctuation theorem (5). It is however convenient to analyze it first in more detail to get some familiarity with the physical questions that it is necessary to address to grasp its meaning.
We begin with an attempt at justifying the name “entropy creation rate” for $`\sigma _+`$. Without too many comments we quote here the simple but relevant remark, , that the so called “Gibbs entropy” of an evolving probability distribution which starts, at $`t=0`$, as absolutely continuous with density $`\rho (x)`$ on phase space has the following property
$$\frac{d}{dt}\rho _t(x)\mathrm{log}\rho _t(x)𝑑x=\sigma (x)\rho _t(x)𝑑x$$
(6)
where $`\rho _t(x)=\rho (S_tx)det\frac{S_tx}{x}`$ is the evolving phase space density (here $`S_tx/x`$ is the matrix of the derivatives of the time evolution map $`S_t`$). Since the r.h.s. of (6) formally tends as $`t\mathrm{}`$ to the average of $`\sigma `$ with respect to the distribution $`\mu `$, i.e. to $`\sigma _+`$, we realize that (6) is a possible justification of the name used for $`\sigma _+`$.
Before proceeding it is also necessary to understand the role of the reversibility assumption in order to see whether it is a serious limitation in view of a possible physical interpretation and physical interest of the fluctuation theorem. Here one can argue that in many cases an irreversible system is “equivalent” to a reversible one: in fact several reversibility conjectures have been proposed, see , , (Sect. 2 and 5), (Sec. 8), , (Sec. 9.11), .
Rather than trying to be general I shall consider an example, and analyze a Navier–Stokes fluid, in a periodic container of side $`L`$, subject to a force with intensity $`F`$ and with viscosity $`\nu `$. If $`R=FL^3\nu ^2`$ is the “Reynolds number”, $`p`$ is the pressure field, the density is $`\rho =1`$, and $`\underset{¯}{g}`$ is a force field of intensity $`1`$ the equations are
$`\underset{¯}{\overset{\dot{}}{u}}=Ru_{\stackrel{~}{}\text{ }}_{\stackrel{~}{}\text{ }}\underset{¯}{u}+\mathrm{\Delta }\underset{¯}{u}+\underset{¯}{g}\underset{¯}{}p`$ (7)
$`\underset{¯}{}\underset{¯}{u}=0`$ (8)
which are irreversible equations. According to the K41 theory (Kolmogorov theory, see ) the above equations can be truncated retaining only a few harmonics of the field $`\underset{¯}{u}`$: i.e. replacing $`\underset{¯}{u}(\underset{¯}{x})=_{\underset{¯}{k}}\underset{¯}{u}_{\underset{¯}{k}}e^{i\underset{¯}{k}\underset{¯}{x}}`$ by $`\underset{¯}{u}(\underset{¯}{x})=_{|\underset{¯}{k}|<K(R)}\underset{¯}{u}_{\underset{¯}{k}}e^{i\underset{¯}{k}\underset{¯}{x}}`$ with $`K(R)=R^{3/4}`$ so that “effectively” (and accepting the K41 theory) the fluid has $`N(R)=O(R^{9/4})`$ degrees of freedom.
If $`R`$ is large (hence the motion is turbulent) the quantity $`𝒟=(_{\stackrel{~}{}\text{ }}\underset{¯}{u})^2𝑑\underset{¯}{x}`$ fluctuates in time with some average $`𝒟_{\mu _R}\stackrel{def}{=}𝒟(R)`$, where $`\mu _R`$ is the probability distribution giving the statistics of the time averages.
We can also consider the following equations, introduced in and called GNS equations,
$`\underset{¯}{\overset{\dot{}}{u}}=Ru_{\stackrel{~}{}\text{ }}_{\stackrel{~}{}\text{ }}\underset{¯}{u}+\nu (\underset{¯}{u})\mathrm{\Delta }\underset{¯}{u}+\underset{¯}{g}\underset{¯}{}p`$ (9)
$`\underset{¯}{}\underset{¯}{u}=0`$ (10)
where the “multiplier” $`\nu (\underset{¯}{u})`$ is so defined that $`𝒟(\underset{¯}{u})`$ is a constant of motion. An elementary calculation yields
$$\nu (\underset{¯}{u})=\frac{_V\left(\underset{¯}{\phi }\mathrm{\Delta }\underset{¯}{u}R\mathrm{\Delta }\underset{¯}{u}(u_{\stackrel{~}{}\text{ }}_{\stackrel{~}{}\text{ }}\underset{¯}{u})\right)𝑑\underset{¯}{x}}{_V(\mathrm{\Delta }\underset{¯}{u})^2𝑑\underset{¯}{x}}$$
(11)
which shows that $`\nu (\underset{¯}{u})`$ is odd in $`\underset{¯}{u}`$ so that (10) is reversible if time reversal is defined as $`(I\underset{¯}{u})(\underset{¯}{x})=\underset{¯}{u}(\underset{¯}{x})`$.
Call “local observable” any function $`\underset{¯}{u}F(\underset{¯}{u})`$ of the velocity field which depends on $`\underset{¯}{u}`$ only via fionitely many of its Fourier compopnents $`\underset{¯}{u}_{\underset{¯}{k}}`$; and let $`\stackrel{~}{\mu }_{𝒟,R}`$ be the SRB distribution for (10). Then the following conjecture was proposed in :
Equivalence conjecture: If $`F(\underset{¯}{u})`$ is a local observable with non zero average then
$$\frac{\mu _R(F)}{\mu _{𝒟(R),R}(F)}\begin{array}{c}R\mathrm{}\hfill \end{array}\mathrm{\hspace{0.17em}1}$$
(12)
i.e. at large Reynolds number the irreversible NS and the reversible GNS equations are equivalent.
Note the analogy between the conjecture and the equivalence property of statistical ensembles in equilibrium statistical mechanics: here the coefficient of the Laplacian in (8) (which is $`1`$ by our definitions) plays the role of temperature in a canonical ensemble while the quantity $`𝒟`$ plays the role of the energy in a microcanonical ensemble: if the energy is suitably tuned the distributions $`\mu _R,\mu _{𝒟(R),R}`$ are statistically equivalent as $`R\mathrm{}`$ and the Reynolds number $`R`$ plays the role of the volume and $`R\mathrm{}`$ the role of thermodynamic limit.
It is interesting to remark that in nonequilibrium physics the statistical ensembles may be defined not only by the parameters that one regards naturally as control parameters (like energy in the microcanonical ensemble and temperature in the canonical) but also by the equations of motion that are used: this is perhaps not so strange because in equilibrium systems no friction is necessary and the equations of motion do not suffer from the ambiguity due to the arbitrariness of the thermostatting mechanisms that remove heat from the system (making possible the evolution towards a statistically stationary state).
The above conjecture is beginning to be tested with results that are at least encouraging, see . Progress in the theory is needed as experiments on fluctuations of entropy production are already available and one would like to interpret them theoretically, .
The second question that one has to clarify preliminarly is that the “rate function” $`\zeta (p)`$ in (5) should be expected to be proportional to some macroscopic parameter measuring the size (like volume or number of degrees of freedom) so that the probability of observing the value $`p`$ in a stationary state is $`\pi _T(p)=conste^{\zeta (p)T}`$ hence it is not observable if $`p1`$: any attempt at measuring $`p`$ will inexorably lead to $`p=1`$ (given that by our normalizations the average value of $`p`$ is $`1`$).
Therefore one should investigate if, or when, a “local version” of the above fluctuation theorem holds telling us some properties, relative to a small volume or to a few degrees of freedom, which could have fluctuations that are frequent enough to be observable.
To get some inspiration we consider an analogous problem: suppose that a high temperature low density gas has density $`\rho `$ and that it occupies the whole space. Given a volume $`V`$ we can consider the observable $`p=N_V/\rho V`$. Then its probability distribution will have the form
$$\pi _V(p)=conste^{\overline{\zeta }(p)V}$$
(13)
where the exponent is affected by an error $`O(V)`$ of the size of the boundary area of $`V`$ and $`\overline{\zeta }(p)`$ is $`V`$–independent.
This shows that density fluctuations that are not observable in volumes $`V`$ of macroscopic size do become observable in small enough volumes when the quantity $`(\overline{\zeta }(p)\overline{\zeta }(1))V`$ becomes reasonably small. Furthermore since $`\overline{\zeta }(p)`$ is essentially $`V`$–independent we can infer that the probability of density fluctuations in a large volume is also measurable, being trivially related to $`\overline{\zeta }(p)`$ which is visible via the fluctuations in small volumes.
Coming back to our dynamical questions we can ask whether there is at least one model for which one can establish a “local fluctuation theorem” in a sense analogous to the above result on density fluctuations. Indeed there is a class of models that is very suitable for illustration purposes: these are the chains of coupled maps. Although their importance is mainly illustrative, they clearly show the possibility of local fluctuation theorems: here it will be enough to refer to the literature, , .
The conclusion is that if a local fluctuation relation holds then it becomes possible to perform tests of the fluctuation theorem, hence of the chaotic hypothesis.
Therefore we can try to attack the main question of interest here, namely “which is the physical interpretation” of the fluctuation theorem. The key is the following theorem, which is a simple extension of it and holds under the same hypotheses (i.e. chaotic hypothesis and time reversibility). It can be regarded as an extension of the Onsager–Machlup theory of fluctuation patterns, . Let $`F,G`$ be time reversal odd observables (for simplicity and to fix the ideas): $`F(Ix)=F(x),G(IX)=G(x)`$; and let $`h,k:[T/2,T/2]R^1`$ be two real valued functions or “patterns”. We call $`h^{}(t)=h(t)`$, $`k^{}(t)=k(t)`$ the “time–reversed patterns” or “antipatterns” of the patterns $`h,k`$. If $`F(S_tx)=h(t)`$ for $`t[T/2,T/2]`$ we say that $`F`$ follows the pattern $`h`$ around the reference point $`x`$ in the time interval $`[T/2,T/2]\stackrel{def}{=}W_T`$. Then, see ,
Theorem (extension of Onsager–Machlup theory): The probabilities of the patterns $`h,k`$ conditioned to a $`T`$–average dimensionless entropy production $`p`$, see (4), denoted $`\pi \left(F(S_t)\right)=h(t),tW_T\left|p\right)`$ and $`\pi \left(G(S_t)\right)=k(t),tW_T\left|p\right)`$ respectively verify
$$\frac{\pi (F(S_t)=h(t),tW_T|p)}{\pi (F(S_t)=h(t),tW_T|p)}=e^{p\sigma _+T}$$
(14)
and (consequently)
$`{\displaystyle \frac{\pi (F(S_t)=h(t),tW_T|p)}{\pi (G(S_t)=k(t),tW_T|p)}}=`$ (15)
$`={\displaystyle \frac{\pi (F(S_t)=h(t),tW_T|p)}{\pi (G(S_t)=k(t),tW_T|p)}}`$ (16)
Hence relative probabilities of patterns in presence of $`T`$–average entropy production $`p`$ are the same as those of the corresponding antipatterns in presence of the opposite $`T`$–average entropy production rate.
In other words it suffices to change the sign of the entropy production to reverse the arrow of time. In a reversible system the quantity $`\zeta (p)`$ measures the degree of irreversibility of a motion observed to have the value $`p`$ of dimensionless entropy creation rate during an observation time of size $`T`$: if we observe patterns over time intervals of size $`T`$ then the fraction of such intervals in which we shall see an entropy production $`p`$ rather than $`1`$ (which is the most probable value) will be
$$e^{(\zeta (p)\zeta (1))T}$$
(17)
More generally, in a situation in which a local fluctuation theorem holds and $`\zeta (p)=V\overline{\zeta }(p)`$ we can divide the volume occupied by the system into small boxes of size $`V_0`$, small enough so that one can observe entropy production fluctuations within them, and we can divide the time axis into time intervals of size $`T`$. Then
Proposition (intermittency of fluctuations): The fraction of time intervals in which we shall observe $`p`$ in a given box $`V_0`$ will be $`e^{(\overline{\zeta }(p)\overline{\zeta }(1))V_0T}`$ and this same quantity will be the fraction of boxes $`V_0`$ where we shall observe, within a given time interval of size $`T`$, entropy production $`p`$.
Normally we shall see $`p=1`$ in a fixed box $`V_0`$ but “seldom” we shall see $`p=1`$ and then, by the above extension of the Onsager–Machlup theory, everything will look wrong: every improbable pattern will appear as frequently as we would expect its (probable) antipattern to appear. This will last only for a moment and then things will return normal for a very long time (as the fraction of times in which this can happen in a given bix is $`e^{\overline{\sigma }_+V_0T}`$). This is a kind of intermittency phenomenon.
In fact we see that when a local fluctuation theorem holds we shall see intermittency, in the form of a reversed time arrow, happening in a small volume $`V_0`$ somewhere in the volume $`V`$ of the system, provided
$$VV_{0}^{}{}_{}{}^{1}e^{\overline{\sigma }_+V_0T}1$$
(18)
furthermore there is a simple relation between fraction of volumes and fraction of times where time reversal occurs: namely they are equal and directly measured by $`e^{\overline{\sigma }_+V_0T}`$, i.e. by the average entropy creation rate.
We conclude by noting that the above remarks set up a possible way of measuring $`\sigma _+`$ and of attempting to measure $`\zeta (p)`$ in systems for which it is hard develop a reasonably good numerical simulation and/or an expression for $`\sigma _+`$. The quantity $`\overline{\sigma }_+`$ can be measured by considering some “current” $`J`$ associated with the system; i.e. an observable $`J(x)`$ which, among other properties, is odd under time reversal. One then looks at the observable obtained by averaging this current over a time $`T`$ and averaged over a time $`T`$, namely $`J_T(x)\stackrel{def}{=}T^1_{T/2}^{T/2}J(S_tx)𝑑t`$ and one measures how often $`J_T`$ takes a value close to the opposite of its (infinite time) average value $`J_+`$, assuming that the latter is $`>0`$: this should happen with a frequency $`e^{\overline{\sigma }_+V_0T}`$ giving us access to $`\overline{\sigma }_+`$.
Given the special role that entropy generation plays it is very tempting to think that there might be many currents $`J`$ associated with the system: for each of them one could define $`p=J_T/J_+`$; then the new quantity $`p`$ has the same probability distribution as the variable with the same name that we have associated with the entropy production. This is true at least for the special case $`p=1`$ as just noted: if true in general then we could have easily access to the function $`\zeta (p)`$ for several values of $`p`$. Hence analysizing this “universality” property in special models seems to be an interesting problem.
For a general review on recent developments in nonequilibrium statistical mechanics see .
Acknowledgements: This is a contribution to the Proceedings of “Inhomogeneous random systems” January 25-26, 2000, (Université de Cergy-Pontoise, Paris); supported partially by “Cofinanziamento 1999”.
REVTEX
Preprints at: http://ipparco.roma1.infn.it
e-mail: giovanni.gallavotti@roma1.infn.it
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# Brane worlds: the gravity of escaping matter
## I Introduction and summary
It has been found recently that four-dimensional gravity may emerge as a low energy effective theory in models with non-compact extra dimensions (see also Refs ). One ingredient of this scenario is a relatively old , and recently revived idea that ordinary matter may reside on a 3-brane embedded in higher-dimensional space. Another key point is the existence in a class of models of a bound state of a multi-dimensional graviton which is localized near the brane . Even though gravity at short distances is higher-dimensional, this bound state dominates gravitational interactions in the brane world at large distances and gives rise to four-dimensional behavior of the gravity force experienced by matter residing on the brane. Modifications of this scenario include models with a metastable graviton bound state or more than one very light four-dimensional graviton ; these models predict violations of the four-dimensional Newton’s law at ultra-large distances as well.
If the extra dimensions are non-compact, it is conceivable that energy may leak from the brane into the bulk. As an example, field theoretic models of 3-branes can be viewed as defects in higher dimensions (domain walls in $`(4+1)`$-, cosmic strings in $`(4+2)`$-dimensions etc.). In these models, highly energetic ordinary particles are able to leave the brane and propagate into the bulk. Another possibility is that besides ordinary matter, there exist particle species that are not trapped to the brane at all; pair creation of these particles would also lead to the transfer of energy from the brane to the bulk. The same role may be played by Kaluza–Klein gravitons whose interactions with the brane matter are weak but non-vanishing.
At first glance, the possibility that energy is carried away from the brane to arbitrarily large distances in the bulk seems, from the four-dimensional point of view, to be in conflict with locality, causality, and the four-dimensional Newton’s law (“Gauss’ law of general relativity”): changing mass in a finite region of three-dimensional space would seem to result in an instantaneous changes of the gravitational potential everywhere in this space. By this argument, an observer living on the brane would immediately realize that energy had been emitted from the brane to the bulk, no matter how far away this event occurred.
Off hand, one may think of several possible resolutions of this apparent paradox. Three of them are as follows:
(i) One may recall that in $`(1+1)`$ dimensions, energy non-conservation is in fact consistent with locality, causality and the long-range character of forces analogous to gravity . Even though the reasons for this property are peculiar to $`(1+1)`$ dimensions, one may wonder whether a similar consistency may be possible in $`(3+1)`$ dimensions.
(ii) The bulk geometry in the Randall–Sundrum (RS) model and its analogs is anti-de Sitter, so that particles leaving the brane get accelerated away from the brane. Their energy increases as they move away, and one may wonder whether this effect could compensate for the decrease of the strength of gravitational interactions of these particles with matter on the brane. In that case, the mass measured through gravitational forces by an observer residing on the brane would remain constant, and the particles continuously accelerating in the bulk would behave, from the four-dimensional point of view, as dark matter particles of fixed mass which participate in gravitational interactions in a standard way.
(iii) Finally, gravity in the RS model is guaranteed to be effectively four-dimensional only as far as interactions of matter residing on the brane are concerned. No argument implies that the effective four-dimensional description is valid for gravity induced by bulk matter. In other words, the gravitational field induced on the brane by particles moving in the bulk need not obey the four-dimensional Einstein equations, so the above discussion of the conflict between causality and Newton’s law may not apply.
In this paper we decide between these possibilities by calculating, in the linear theory about the RS background, the gravitational field of a pair of particles escaping from the brane and propagating in the bulk. We show that the most exotic option (i) has nothing to do with the RS model. The possibility (ii) would be realized if the graviton bound state were the only relevant mode, i.e., if the zero mode approximation were reliable. This may be viewed as a consistency check: in the zero-mode approximation gravity is effectively four-dimensional irrespective of the position of its source in extra dimensions, so the four-dimensional gravitating mass must be conserved.
However, the zero-mode approximation is in fact not adequate to the problem in question, and the Kaluza–Klein (KK) tower of gravitons plays an important role. When KK states are included in the analysis, the outcome is option (iii). The final picture is that the particles moving away from the brane produce a spherical gravity wave on the brane, which expands in three-dimensional space with the speed of light (or almost the speed of light). The four-dimensional space-time left behind this spherical wave is flat, whereas in front of this wave the four-dimensional metric is of the usual Schwarzschild asymptotic form, in accord with causality. The spherical gravitational wave itself does not obey the four-dimensional Einstein equations, so the gravitational effects of particles moving in the bulk are intrinsically higher-dimensional, even if these effects are measured by an observer residing on the brane.
## II The physical set up
The Randall–Sundrum model contains a 3-brane with tension $`\sigma >0`$ embedded in five-dimensional space-time. The bulk cosmological constant between the branes, $`\mathrm{\Lambda }`$, is negative and tuned in such a way that the intrinsic geometry on the brane is that of flat four-dimensional Minkowski spacetime. The solution to the five-dimensional Einstein equations is
$$ds^2=a^2(z)\eta _{\mu \nu }dx^\mu dx^\nu dz^2$$
(1)
where
$$a(z)=e^{k|z|}$$
(2)
Here $`\mu ,\nu =0,1,2,3`$ and $`\eta _{\mu \nu }`$ is the Minkowski metric. The constant $`k`$ is related to $`\sigma `$ and $`\mathrm{\Lambda }`$ as follows: $`\sigma =\frac{3k}{4\pi G_5}`$, $`\mathrm{\Lambda }=\sigma k`$, where $`G_5`$ is the five-dimensional Newton constant. The four-dimensional hypersurfaces $`z=const.`$ are flat, the five-dimensional space-time to the right of the brane is a part of anti-de Sitter (adS) space.
It is sometimes convenient to introduce another coordinate,
$$\zeta =\frac{1}{k}e^{kz}$$
(3)
in terms of which the background metric is conformally flat,
$$ds^2=\frac{1}{k^2\zeta ^2}\left(\eta _{\mu \nu }dx^\mu dx^\nu d\zeta ^2\right)$$
(4)
Consider now a process in which some matter is emitted from the brane, and then freely moves in the bulk adS background. Let us assume for simplicity that this matter is symmetric with respect to $`zz`$; this will enable us to take symmetric metric perturbations and effectively consider only the space to the right of the brane. An example which we discuss throughout this paper is two particles of mass $`m`$ which are emitted from the brane at time $`t=0`$ along the line $`𝐱=0`$ in opposite directions in $`z`$ with zero initial velocities. The coordinates of these particles obey the geodesic equations in adS. It is straightforward to see that the world line of a particle right to the brane is described as follows,
$$𝐱_c(t)=0,z_c(t)=\frac{1}{2k}\mathrm{ln}(1+k^2t^2)$$
(5)
In terms of the coordinate $`\zeta `$, this means
$$\zeta _c(t)=\sqrt{k^2+t^2}$$
(6)
The particle accelerates towards $`z\mathrm{}`$, and at $`tk^1`$ its world line approaches the light cone, $`\zeta =t`$. Similar observations have been made independently in Ref..
The energy-momentum tensor of this particle,
$$\widehat{T}^{ab}=\frac{m}{\sqrt{g}}\frac{dx^a}{ds}\frac{dx^b}{dt}\delta (𝐱𝐱_c(t))\delta (zz_c(t))$$
(7)
has the following non-vanishing components,
$`T_{zz}`$ $`=`$ $`{\displaystyle \frac{m}{a^3}}{\displaystyle \frac{v^2}{\sqrt{1v^2}}}\delta (zz_c(t))\delta (𝐱)`$ (9)
$`T_{z0}`$ $`=`$ $`{\displaystyle \frac{m}{a^2}}{\displaystyle \frac{v}{\sqrt{1v^2}}}\delta (zz_c(t))\delta (𝐱)`$ (10)
$`T_{00}`$ $`=`$ $`{\displaystyle \frac{m}{a}}{\displaystyle \frac{1}{\sqrt{1v^2}}}\delta (zz_c(t))\delta (𝐱)`$ (11)
where
$$v=\frac{\dot{z_c}}{a(z_c)}=\frac{kt}{\sqrt{1+k^2t^2}}$$
(12)
Hereafter we consider tensors with all lower indices as basic ones; quantities with and without hats have upper indices raised by full adS and Minkowski metrics respectively. As an example, $`T_\nu ^\mu =\eta ^{\mu \lambda }T_{\lambda \nu }`$, $`\widehat{T}_\nu ^\mu =a^2\eta ^{\mu \lambda }T_{\lambda \nu }`$.
Equations (7) and (7) are valid at $`t>0`$; before that the energy-momentum tensor is concentrated on the brane. We will assume for definiteness that at $`t<0`$, the source on the brane is a point-like particle with mass $`2m`$. The physical picture is that this brane particle decays at $`t=0`$ into two particles of equal mass, one of which moves according to eq. (5) and another which is its mirror image.
The linearized five-dimensional Einstein equations in RS background have been considered in a number of papers, see, e.g., Refs. . In particular, the advantages and disadvantages of the Gaussian Normal (GN) gauge have been emphasized . As we are interested in the gravitational effects on the brane of a particle moving outside the brane, we find it convenient to work entirely in GN coordinates, however, this will mean that there will not be a global coordinate system which is GN once a matter source has been introduced. In a GN frame one has
$$g_{zz}=1,g_{z\mu }=0$$
(13)
and the linearized theory is described by the metric
$$ds^2=a^2(z)\eta _{\mu \nu }dx^\mu dx^\nu +h_{\mu \nu }(x,z)dx^\mu dx^\nu dz^2$$
(14)
The linearized Einstein equations have the following form,
$`\delta R_{zz}`$ $`=`$ $`8\pi G_5\theta _{zz}`$ (16)
$`\delta R_{z\mu }`$ $`=`$ $`8\pi G_5\theta _{z\mu }`$ (17)
$`\delta R_{\mu \nu }4k^2h_{\mu \nu }`$ $`=`$ $`8\pi G_5\theta _{\mu \nu }`$ (18)
where
$`\theta _{zz}`$ $`=`$ $`\left({\displaystyle \frac{2}{3}}T_{zz}+{\displaystyle \frac{1}{3a^2}}T_\lambda ^\lambda \right)`$ (20)
$`\theta _{z\mu }`$ $`=`$ $`T_{z\mu }`$ (21)
$`\theta _{\mu \nu }`$ $`=`$ $`\left(T_{\mu \nu }{\displaystyle \frac{1}{3}}\eta _{\mu \nu }T_\lambda ^\lambda +{\displaystyle \frac{a^2}{3}}\eta _{\mu \nu }T_{zz}\right)`$ (22)
and
$`\delta R_{zz}`$ $`=`$ $`\left({\displaystyle \frac{h^{}}{2a^2}}\right)^{}`$ (24)
$`\delta R_{z\mu }`$ $`=`$ $`\left[{\displaystyle \frac{1}{2a^2}}(h_{\mu ,\nu }^\nu h_{,\mu })\right]^{}`$ (25)
$`\delta R_{\mu \nu }`$ $`=`$ $`{\displaystyle \frac{1}{2}}h_{\mu \nu }^{\prime \prime }+2k^2h_{\mu \nu }\left(k^2h+{\displaystyle \frac{k}{2}}h^{}\right)\eta _{\mu \nu }`$ (27)
$`+{\displaystyle \frac{1}{2a^2}}(2h_{(\mu ,\nu )\lambda }^\lambda h_{\mu \nu ,\lambda }^{,\lambda }h_{,\mu \nu })`$
Hereafter $`h=h_\mu ^\mu \eta ^{\mu \nu }h_{\mu \nu }`$.
We will explicitly consider times $`t>0`$ when the source is in the bulk. Since there is no matter on the brane, the metric should obey the Israel junction condition on the brane, which, due to the $`Z_2`$ symmetry reads
$$h_{\mu \nu }^{}+2kh_{\mu \nu }=0$$
(28)
Let us first solve eq. (16). Its general solution is
$$h(z,x)=\frac{8\pi G_5}{k}\left[a^2(z)_z^{\mathrm{}}𝑑z^{}\theta _{zz}(z^{})_z^{\mathrm{}}𝑑z^{}a^2(z^{})\theta _{zz}(z^{})\right]+a^2(z)C(x)+D(x)$$
(29)
From the junction condition, $`h^{}+2kh=0`$, we find
$$D(x)=\frac{8\pi G_5}{k}_0^{\mathrm{}}𝑑z^{}a^2(z^{})\theta _{zz}(z^{})$$
(30)
Consider now eq. (18). Let us introduce
$$\xi _\mu =h_{\mu ,\lambda }^\lambda \frac{1}{2}h_{,\mu }$$
(31)
This function can be found from eq. (17), but its explicit form will be irrelevant. With the definition (31), eq. (18) becomes
$$\frac{1}{2}h_{\mu \nu }^{\prime \prime }2k^2h_{\mu \nu }\frac{1}{2a^2}\mathrm{}h_{\mu \nu }=8\pi G_5\theta _{\mu \nu }+\left(k^2h+\frac{k}{2}h^{}\right)\eta _{\mu \nu }\frac{1}{2a^2}(\xi _{\mu ,\nu }+\xi _{\nu ,\mu })$$
(32)
Hence, one can write
$$h_{\mu \nu }=\overline{h}_{\mu \nu }+(u_{\mu ,\nu }+u_{\nu ,\mu })$$
(33)
where $`u_\mu `$ accounts for $`\xi _\mu `$ in eq. (32), and $`\overline{h}_{\mu \nu }`$ obeys the following equation,
$$\frac{1}{2}\overline{h}_{\mu \nu }^{\prime \prime }2k^2\overline{h}_{\mu \nu }\frac{1}{2a^2}\mathrm{}\overline{h}_{\mu \nu }=8\pi G_5\theta _{\mu \nu }^{eff}$$
(34)
with
$$8\pi G_5\theta _{\mu \nu }^{eff}=8\pi G_5\theta _{\mu \nu }+\left(k^2h+\frac{k}{2}h^{}\right)\eta _{\mu \nu }$$
(35)
The last term in this equation is known explicitly from eqs. (29) and (30):
$$\left(k^2h+\frac{k}{2}h^{}\right)=8\pi G_5k_0^z𝑑z^{}a^2(z^{})\theta _{zz}(z^{})$$
(36)
We are interested in the induced metric on the brane,
$$h_{\mu \nu }(z=0)=\overline{h}_{\mu \nu }(z=0)+(u_{\mu ,\nu }+u_{\nu ,\mu })(z=0)$$
(37)
The last term in this equation is a pure gauge in the four-dimensional brane world and has no effect on the motion of matter confined to the brane, so we may omit it. In effect, we have to solve eq. (34) and then find $`\overline{h}_{\mu \nu }`$ on the brane, i.e., at $`z=0`$.
It is worth looking at the additional term (36) in the case of the point-like particle moving in the bulk. We have
$$\left(k^2h+\frac{k}{2}h^{}\right)=8\pi G_5k\theta (zz_c(t))a^2(z_c(t))\mathrm{\Phi }_{zz}(x)$$
(38)
where $`\mathrm{\Phi }_{ab}`$ is defined by
$$\theta _{ab}(z,x)=\delta (zz_c(t))\mathrm{\Phi }_{ab}(x)$$
(39)
Explicitly
$`\mathrm{\Phi }_{zz}(x)`$ $`=`$ $`{\displaystyle \frac{m}{a^3(z_c)\sqrt{1v^2}}}\left({\displaystyle \frac{2}{3}}v^2+{\displaystyle \frac{1}{3}}\right)\delta (𝐱)`$ (41)
$`\mathrm{\Phi }_{00}(x)`$ $`=`$ $`{\displaystyle \frac{m}{a(z_c)\sqrt{1v^2}}}\left({\displaystyle \frac{1}{3}}v^2+{\displaystyle \frac{2}{3}}\right)\delta (𝐱)`$ (42)
$`\mathrm{\Phi }_{ij}(x)`$ $`=`$ $`{\displaystyle \frac{m}{a(z_c)\sqrt{1v^2}}}\left({\displaystyle \frac{1}{3}}v^2+{\displaystyle \frac{1}{3}}\right)\delta (𝐱)\delta _{ij}`$ (43)
A special feature of the additional term (38) is that this is a non-local expression with a “string” extending from $`z=z_c`$ to $`z=\mathrm{}`$. This string is of course a gauge artifact due to the breakdown of the brane GN gauge as we pass the particle in the $`z`$-direction, and represents a caustic of the normal geodesic congruence used to define the affine parameter $`z`$. In fact, a similar string is also characteristic to the asymptotic four-dimensional Schwarzchild solution transformed into a gauge which is GN with respect to an arbitrarily chosen 2-plane. This artifact is easily removable (as noted in ) via a construction reminiscent of the Wu-Yang gauge patching for a Dirac monopole . One introduces an additional GN gauge patch to the right of the accelerating particle, valid for $`z>z_c(t)ϵ\mathrm{tan}^1|𝐱|`$, with the brane GN patch being valid for $`z<z_c(t)+ϵ\mathrm{tan}^1|𝐱|`$. The gauge transformation on the overlap is readily seen to be a bulk analog of the Garriga-Tanaka transformation:
$$ϵ_\mu =\frac{ϵ_z(x)}{2k};^2ϵ_z=8\pi G_5a^2\left(z_c(t)\right)\mathrm{\Phi }_{zz}(x)$$
(44)
One could always choose a harmonic bulk gauge ($`^ah_{ab}\frac{1}{2}h_{,b}=0`$), however, the computation of the Green’s function in this gauge is cumbersome and not particularly illuminating, therefore we simply use a GN gauge. Indeed, since we are primarily interested in the metric induced on the brane, we stick with the brane GN system, the string singularity outside the brane not leading to any inconvenience.
## III Solution of the linearized problem
### A First trial: zero mode approximation
We are going to treat the parameter $`k`$ of the RS model as microscopic, and are interested in the induced metric on the brane, $`\overline{h}_{\mu \nu }(x,z=0)`$, at distance scales large compared to $`k^1`$. In particular, we consider the region of the four-dimensional space-time such that $`r|𝐱|k^1`$, $`tk^1`$. The solution to eq. (34) involves the retarded Green’s function of the operator
$$\frac{1}{2}\frac{d^2}{dz^2}2k^2\frac{1}{2a^2}\mathrm{}$$
(45)
with boundary conditions enforcing the junction property (28). This Green’s function is expressed in terms of the eigenfunctions of the corresponding one-dimensional problem ,
$$G_R(xx^{},z,z^{})=4ka^2(z)a^2(z^{})D_0(xx^{})+2_0^{\mathrm{}}𝑑mu_m(z)u_m(z^{})D_m(xx^{})$$
(46)
where $`D_0`$ and $`D_m`$ are retarded Green’s functions of massless and massive scalar fields in four flat dimensions, and
$$u_m(z)=\sqrt{\frac{m}{k}}\frac{J_1(m/k)Y_2(m\zeta )Y_1(m/k)J_2(m\zeta )}{\sqrt{J_1(m/k)^2+Y_1(m/k)^2}}$$
(47)
Here $`\zeta (z)`$ is defined by eq. (3), $`J_n`$ and $`Y_n`$ are the Bessel functions. The first term in eq. (46) is the contribution of the bound state of the five-dimensional graviton (the zero mode), whereas the second term comes from the continuous KK spectrum. Note that our expression for the Green’s function, eq. (46), differs from that of Refs. by an overall factor of 4. One factor of 2 has to do with our form of the operator (45), and the other is due to symmetry $`zz`$ and effectively accounts for the two particles moving left and right from the brane.
If the source in eq. (34) were on the brane, the long-distance behaviour of the induced metric would be governed by the zero mode contribution,
$$G_R^{zm}(xx^{},z,z^{})=4ka^2(z)a^2(z^{})D_0(xx^{})$$
(48)
As our first trial, let us boldly use the zero mode approximation (48) in the problem at hand. As mentioned in the Introduction, this approximation is not adequate in our case, but rather provides a non-trivial consistency check.
In the zero mode approximation, the metric induced on the brane is effectively determined by the following equation,
$$\mathrm{}\overline{h}_{\mu \nu }(x)=8\pi G_N\tau _{\mu \nu }(x)$$
(49)
where $`G_N=kG_5`$ is the four-dimensional Newton constant and
$$\tau _{\mu \nu }(x)=4_0^{\mathrm{}}𝑑za^2(z)\theta _{\mu \nu }^{eff}(z)$$
(50)
In the case of a point particle, the right hand side of this equation is straightforward to evaluate,
$`\tau _{00}`$ $`=`$ $`2m{\displaystyle \frac{a(z_c(t))}{\sqrt{1v^2}}}\delta (𝐱)`$ (52)
$`\tau _{ij}`$ $`=`$ $`2m{\displaystyle \frac{a(z_c(t))}{\sqrt{1v^2}}}\delta (𝐱)\delta _{ij}`$ (53)
Making use of eq. (5) one finds
$$\frac{a(z_c(t))}{\sqrt{1v^2}}=1$$
(54)
Therefore, eq. (49) coincides with the linearized Einstein equation in four dimensions with a static source of mass $`2m`$. As discussed in the Introduction, this is consistent with the general property that gravity in the zero mode approximation is always effectively four-dimensional, irrespective of the position of its source in the fifth dimension.
### B Full treatment
To obtain the correct expression for the induced metric on the brane, we have to consider the complete Green’s function (46). Is is convenient to work in the coordinate representation, and make use of the explicit form of the retarded scalar Green’s function in four flat dimensions ,
$$D_m(x)=\frac{1}{2\pi }\theta (t)\delta (\lambda ^2)+\frac{m}{4\pi \lambda }\theta (t|𝐱|)J_1(m\lambda ),\lambda =\sqrt{t^2|𝐱|^2}$$
(55)
At $`m=0`$ (i.e., in the zero mode case), only the first term in this expression survives.
The first term in eq. (55) is independent of $`m`$. Now, one recalls that the functions $`u_m(z)`$, with the zero mode included, form a complete set, so that
$$2ka^2(z)a^2(z^{})+u_m(z)u_m(z^{})𝑑m=a^2(z)\delta (zz^{})$$
(56)
We see that the contribution of the zero mode to the five-dimensional Green’s function (46) is cancelled out at $`zz^{}`$ by KK states. Since we are interested in the metric on the brane induced by the particle outside the brane, we set $`zz^{}`$ and obtain
$$G_R(x;z,z^{})=_0^{\mathrm{}}𝑑m\frac{m}{2\pi \lambda }\theta (t|𝐱|)J_1(m\lambda )u_m(z)u_m(z^{}),zz^{}$$
(57)
To analyze the long-distance physics, we consider the regime $`t,|𝐱|,\zeta ^{}k^1`$ (recall that $`\zeta t`$ at the position of the source, see eq. (6)). Then the main contribution to the Green’s function comes from modes with small $`m`$, that is $`mk`$. Thus, we retain the terms in eq. (47) which are leading order in $`m/k`$ (not assuming that $`m\zeta ^{}`$ is small) and obtain for the Green’s function with one argument on the brane and another off the brane
$$G_R(x;0,z^{})=\frac{\theta (t|𝐱|)}{2\pi k\lambda }_0^{\mathrm{}}𝑑mm^2J_1(m\lambda )J_2(m\zeta ^{})$$
(58)
The integration here is performed by making use of the following relations,
$`{\displaystyle _0^{\mathrm{}}}𝑑xx^2J_1(\alpha x)J_2(\beta x)`$ $`=`$ $`\left({\displaystyle \frac{^2}{\beta ^2}}{\displaystyle \frac{1}{\beta }}{\displaystyle \frac{}{\beta }}\right){\displaystyle _0^{\mathrm{}}}𝑑xJ_1(\alpha x)J_0(\beta x)`$ (59)
$`=`$ $`\left({\displaystyle \frac{^2}{\beta ^2}}{\displaystyle \frac{1}{\beta }}{\displaystyle \frac{}{\beta }}\right){\displaystyle \frac{1}{\alpha }}\theta (\alpha \beta )`$ (60)
In this way we find
$$G_R(x;0,z^{})=\frac{\theta (t|𝐱|)}{2\pi k\lambda ^2}\left(\delta ^{}(\lambda \zeta ^{})+\frac{1}{\zeta ^{}}\delta (\lambda \zeta ^{})\right)$$
(61)
The Green’s function $`G_R(xx^{};0,z^{})`$ is concentrated on the five-dimensional light cone, $`\lambda =\zeta ^{}`$, i.e.,
$$(tt^{})^2(𝐱𝐱^{})^2\zeta ^2=0$$
(62)
This feature may (or may not) be an artifact of our coarse-graining: if we were able to resolve short distances, the Green’s function might spread over the region of size $`k^1`$ inside the light cone.
To obtain the metric induced on the brane, we have to integrate the Green’s function (61) with $`\theta _{\mu \nu }^{eff}`$,
$`\overline{h}_{\mu \nu }(x)\overline{h}_{\mu \nu }(x,z=0)`$ $`=`$ $`8\pi G_5{\displaystyle d^4x^{}_0^{\mathrm{}}𝑑z^{}G_R(xx^{};0,z^{})\theta _{\mu \nu }^{eff}(x^{},z^{})}`$ (63)
$`=`$ $`8\pi G_5{\displaystyle d^4x^{}_{1/k}^{\mathrm{}}\frac{d\zeta ^{}}{k\zeta ^{}}G_R(xx^{};0,\zeta ^{})\theta _{\mu \nu }^{eff}(x^{},\zeta ^{})}`$ (64)
The source here is given by eqs. (35), (38), (39),
$$\theta _{\mu \nu }^{eff}(x,\zeta )=k\zeta \delta (\zeta \zeta _c(t))\mathrm{\Phi }_{\mu \nu }(x)\frac{1}{k\zeta _c^2(t)}\eta _{\mu \nu }\theta (\zeta \zeta _c(t))\mathrm{\Phi }_{zz}(x)$$
(65)
Even though this source contains the non-local “string”, the integrand in eq. (64) is in effect local. This is due to the particular structure of the Green’s function (61),
$$G_R(x;0,z^{})=\frac{\theta (t|𝐱|)}{2\pi k\lambda ^2}\zeta ^{}\frac{}{\zeta ^{}}\left(\frac{1}{\zeta ^{}}\delta (\lambda \zeta ^{})\right)$$
(66)
Upon integrating by parts in eq. (64), the second term in eq. (65) becomes a delta-function of $`(\zeta \zeta _c(t))`$, i.e., only the region $`\zeta \zeta _c(t)`$ in fact contributes to this integral. This of course had to be the case, as the “string” is a gauge artifact.
In the regime considered, $`tk^1`$, one has $`\zeta _c(t)=t`$, and the leading terms in $`\mathrm{\Phi }_{ab}`$ are (see eqs. (41,42,43))
$`\mathrm{\Phi }_{zz}`$ $`=`$ $`m(kt)^4\delta (𝐱)`$ (68)
$`\mathrm{\Phi }_{00}`$ $`=`$ $`m(kt)^2\delta (𝐱)`$ (69)
$`\mathrm{\Phi }_{ij}`$ $`=`$ $`{\displaystyle \frac{m}{3}}\delta _{ij}\delta (𝐱)`$ (70)
Substituting these expressions into eq. (65) we get
$`\theta _{00}^{eff}(x,\zeta )`$ $`=`$ $`m[(kt)^3\delta (\zeta t)k^3t^2\theta (\zeta t)]\delta (𝐱)`$ (72)
$`\theta _{ij}^{eff}(x,\zeta )`$ $`=`$ $`mk^3t^2\theta (\zeta t)\delta _{ij}\delta (𝐱)`$ (73)
The integration in eq. (64) is now straightforward. We end up with
$`\overline{h}_{00}(x)`$ $`=`$ $`4G_Nm{\displaystyle \frac{2t^2r^2}{t^3}},`$ (74)
$`\overline{h}_{ij}(x)`$ $`=`$ $`4G_Nm{\displaystyle \frac{1}{t}}\delta _{ij},tr>0`$ (75)
Here we stress the fact that these expressions are valid inside the four-dimensional future light cone, $`tr>0`$, as is clear from the explicit $`\theta `$-function in eq. (57). Outside this light cone, the induced metric is determined by the source existing at $`t<0`$. With our model for this source (particle of mass $`2m`$ on the brane at rest at $`𝐱=0`$), the induced metric outside the light cone is the asymtotic Schwarzschild solution (see, e.g. Ref. for explicit derivation in RS model), which in an appropriate gauge reads
$`\overline{h}_{00}(x)`$ $`=`$ $`4G_Nm{\displaystyle \frac{1}{r}},`$ (76)
$`\overline{h}_{ij}(x)`$ $`=`$ $`4G_Nm{\displaystyle \frac{1}{r}}\delta _{ij},tr<0`$ (77)
We see that the induced metric (75), (77) is continuous on the four-dimensional light cone, but its derivatives are not. Because of the latter property, the induced metric does not obey the (linearized) four-dimensional Einstein equations. It describes a spherical gravity wave propagating in four dimensions with the speed of light. Again, within our approximation we do not resolve distances of order $`k^1`$, so the wave may actually spread over the region of this size.
The four-dimensional space-time on the brane left behind the spherical wave is in fact flat. Indeed, a coordinate transformation on the brane
$$\delta t=4G_Nm\left(\frac{r^2}{4t^2}+\mathrm{ln}t\right),\delta x^i=4G_Nm\frac{x^i}{2t}$$
(78)
reduces the four-dimensional metric perturbation (75) to $`\overline{h}_{\mu \nu }=0`$ inside the light cone $`tr>0`$. The gravitational field induced on the brane by matter escaping into the bulk finally disappears, this process occuring in a causal way.
## IV Discussion
We have shown that within linearized perturbation theory, the metric on the brane does indeed react to the ‘loss’ of the sources in the bulk in an intrinsically five-dimensional fashion, a spherical shock wave expanding outwards from the moment of emission leaving behind flat space. It is tempting to ask if something similar could perhaps be obtained by an appropriate use of the far field Schwarzschild-adS solution
$$ds^2=\left(1+k^2\rho ^2\frac{\mu }{k^2\rho ^2}\right)d\tau ^2\left(1+k^2\rho ^2\frac{\mu }{k^2\rho ^2}\right)^1d\rho ^2\rho ^2d\mathrm{\Omega }_{III}^2$$
(79)
In the absence of the mass term, $`\mu `$, the transformation between the brane coordinates and the spherical coordinates is
$`\rho `$ $`=`$ $`{\displaystyle \frac{1}{2k\zeta }}\left[(\zeta ^2k^2)^2+k^2(𝐱^2t^2)^2+2k^2\zeta ^2(𝐱^2t^2)+2(t^2+𝐱^2)\right]^{1/2}`$ (81)
$`\mathrm{tan}k\tau `$ $`=`$ $`{\displaystyle \frac{2t}{k}}\left[k^2+\zeta ^2+𝐱^2t^2\right]^1`$ (82)
$`\mathrm{tan}\chi `$ $`=`$ $`{\displaystyle \frac{2|𝐱|}{k}}\left[\zeta ^2+𝐱^2t^2k^2\right]^1`$ (83)
where the parametrisation of S<sup>3</sup> in the spherical system is the direct generalisation of the S<sup>2</sup> one, rather than Euler angles, and the angular $`\theta ,\varphi `$ variables coincide with the angular variables on the brane-world. Notice the periodic relation of the $`\tau `$-variable to the brane coordinates, this signifies that the spherical coordinates represent in fact the universal covering space of adS, whereas the brane exists in a single patch. In a sense, there are an infinite family of branes existing in the maximally extended spacetime, as discussed in Chamblin, Hawking and Reall . These branes start off ‘planar’, become parabolic, then return to planar again, oscillating indefinitely in the maximally extended spherical coordinates. From the perspective of the brane spacetime, $`\rho =0`$ corresponds to the geodesic trajectory (5). It is tempting therefore to ask whether we can, by modifying our spacetime to Schwarzschild-adS, come up with a metric corresponding to the particle accelerating away from the brane, at least at large spherical $`\rho `$-coordinate. From the perspective of the metric (79), this would correspond to placing a brane at some angular coordinate $`\chi =\chi (t,\rho )`$.
In Ref. , it was shown that it was not possible to find a static trajectory, $`\chi (\rho )`$, which might correspond to a particle sitting on the brane, however, it was suspected that this was due to the non-accelerating nature of the black hole – a more appropriate exact metric, such as some sort of C-metric would be a better candidate. These expectations were partly backed up by the lower-dimensional calculation of . In our case however, the accelerating particle in the bulk translates into a non-accelerating $`\rho =0`$ geodesic in the spherical spacetime (79) so we might hope that a time dependent brane trajectory will work. Unfortunately, as we show in the appendix, it is not possible even to find a time-dependent trajectory corresponding to a particle accelerating in the bulk. The reason why this approach fails (as well as that in ) becomes apparent once we think of the difference in the causal structure of the adS and Schwarzschild-adS spacetimes. In the former, we have an infinite family of oscillating branes, whereas the latter has simply one copy of the brane. The trajectory of the brane therefore cannot be a simple perturbation of the pure adS trajectory, which oscillates periodically in $`\tau `$. Indeed, if one were to try this approach, one would be trying to consider the cumulative effect of the attractive central potential generated by the Schwarzschild source over an infinite number of oscillations! It becomes clear that (79) is an inappropriate metric to use in this context when one tries to use it as a far-field approximation, i.e. assuming that the matter source is extended in such a fashion as to avoid an event horizon. A brane-world observer sitting outside the domain of influence of this accelerating particle (the past light cone of $`t=k^1`$ for example) should see no effect, however, the metric (79) still induces a nonzero perturbation in this region, which is of course due to the infinitely extended nature of the maximal adS spacetime. This behaviour should be contrasted with the causal behaviour of the metric perturbation derived above.
Acknowledgments
We would like to thank Victor Berezin, Christos Charmousis, Ed Copeland, Savas Dimopoulos, Sergei Dubovsky, Gia Dvali, Dmitry Gorbunov, Maxim Libanov, John March-Russell, Lisa Randall, Simon Ross and Sergei Troitsky for useful discussions. R.G. and V.R. acknowledge the hospitality of the Isaac Newton Institute for Mathematical Sciences, where this work was begun. R.G. is supported by the Royal Society, and V.R. and S.S by the Russian Foundation for Basic Research, grant 990218410.
## Appendix
In this appendix, we show that it is not possible to use the Schwarzschild-adS solution (79) to construct a spacetime corresponding to a brane with an accelerating particle in the bulk. To do this, we need only two of the four seperate Israel conditions appropriate to the hypersurface defined by the brane:
$$X^a=(\tau ,\rho ,\chi (\tau ,\rho ),\theta ,\varphi )$$
(84)
with normal
$$n_a=(\dot{\chi },\chi ^{},1,0,0)/n$$
(85)
where $`n^2==1/\rho ^2+A^2\chi ^2\dot{\chi }^2A^2`$, writing $`A^2=g_{tt}`$ for convenience. The $`\theta `$ and $`\varphi `$ Israel conditions for the brane are identical, and give
$$\frac{(\mathrm{cos}\chi \rho A^2\chi ^{}\mathrm{sin}\chi )}{n\rho ^2\mathrm{sin}\chi }=k$$
(86)
In order to get the remaining Israel conditions, rather than working with the fundamental forms of the brane hypersurface, it is easier to generalise the technique used in Ipser and Sikivie, , and use the normal jump in the parallel derivatives of the unit vectors corresponding to the remaining time and space-like directions on the brane:
$$u^a=\dot{X}^a/|\dot{X}|,v^a=X^a/|X^{}|$$
(87)
which satisfy
$$n_a_uu^a=k,n_a_vv^a=k,n_a_vu^a=n_a_uv^a=ku^av_a$$
(88)
This last relation giving
$$\dot{\chi }^{}+\dot{\chi }\left(\frac{A^{}}{A}A^2\rho \chi ^2\frac{1}{\rho }\right)=k\rho ^2\dot{\chi }\chi ^{}n$$
(89)
Combining (86) and (89) requires
$$\dot{\chi }\mathrm{sin}\chi =\frac{A}{\rho }f(k\tau )\mathrm{cos}\chi =\frac{A}{\rho }F(k\tau )+C(\rho )$$
(90)
where $`F=𝑑\tau f`$. But then (86) implies
$$f^2+F^22\frac{AF}{k\rho }\left(C+\rho C^{}\right)\left(1+\frac{2\mu }{k^4\rho ^4}\right)+A^2\left(C+\rho C^{}\right)^2+\frac{\mu }{k^4\rho ^4}\left(3F^2C^2\rho ^2A^2C^2\right)=1$$
(91)
It is not difficult to see that since $`f,F`$ are functions of $`\tau `$ this cannot be satisfied once $`\mu 0`$. For $`\mu =0`$ the solution is of course the adS brane trajectory: $`f=\mathrm{sin}k\tau `$, $`C=1/k\rho `$.
We should also note that this approach is distinct from that of Kraus and Ida , who use the Schwarzschild-adS spacetime to obtain a cosmological brane world, i.e. a homogeneous universe with an evolving scale factor, which we would not expect to obtain from a localised particle accelerating in the bulk.
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# The European Large Area ISO Survey II: mid-infrared extragalactic source counts
## 1 Introduction
The IRAS mission enjoyed huge successes, including the sensational discoveries of ultra- and hyper-luminous galaxies and of an enormous population of evolving starbursts. However, the survey had several drawbacks. For example, the bright limiting flux densities restricted the samples to low redshifts ($`z\stackrel{<}{_{}}0.3`$) for all but a few ultraluminous objects. Also, only $`1000`$ galaxies were detected at $`12\mu `$m over the whole sky. These deficiencies restricted the study of IR-luminous galaxies at all redshifts.
The Infrared Space Observatory (ISO) offered $`\times 1000`$ improvements in sensitivity in the mid-IR over IRAS, and the large allocations of guaranteed and discretionary time for deep surveys on ISO will greatly improve on the IRAS surveys in the mid-IR. For instance, ISO observations of the northern Hubble Deep Field (Serjeant et al. 1997, 1999, Goldschimdt et al. 1997, Oliver et al. 1997, Aussel et al. 1999, Désert et al. 1999) reached the $`15\mu `$m confusion limit ($`0.1`$ mJy) over $`17`$ square arcminutes, while the CAM-Deep and CAM-Shallow surveys (Elbaz et al. 1998a,b) were slightly less sensitive but had wider areal coverage ($`0.5`$ mJy over $`0.3`$ square degrees and $`0.8`$ mJy over $`0.41`$ square degrees). These have also been complemented by deep ISO photometry of selected high-$`z`$ galaxies (e.g. Flores et al. 1999).
The European Large Area ISO Survey (ELAIS, Oliver et al. 1999 (paper I), Rowan-Robinson et al. 1998) was the largest open time project on ISO, complementing the deep ISO samples by surveying $`12`$ square degrees to a depth of $`2`$ mJy at $`15\mu `$m and $`\stackrel{<}{_{}}100`$ mJy at $`90\mu `$m. Around half the area was also mapped at $`6.7\mu `$m to $`1`$ mJy. Three fields in the Northern hemisphere (N1, N2, N3) collectively comprised around two-thirds of the $`15\mu `$m areal coverage, with the remaining area taken by the Southern S1 field and several small areas in both hemispheres. The ambitious cosmological aims include tracing the extinguished star formation history of the Universe to $`z12`$, orientation-independent selection of dust-shrouded quasars, and the potential discovery of hyperluminous galaxies (with comparable intrinsic luminosities to IRAS FSC 10214+4724) out to redshifts $`z\stackrel{<}{_{}}5`$. A more detailed discussion of the diverse scientific aims of ELAIS, the selection of areas and observational parameters can be found in the ELAIS survey paper (Oliver et al. 1999); in summary, the survey areas were selected to have low galactic cirrus emission, high visibility by ISO, high ecliptic latitude and avoiding $`12\mu `$m IRAS sources brighter than $`0.6`$ Jy. In another companion paper, Efstathiou et al. 1999, we discuss the $`90\mu `$m source counts from the Preliminary Analysis of the ELAIS ISOPHOT data, and in Crockett et al. (1999) we discuss the stellar mid-infrared source counts. The ELAIS areas have also been the subject of intensive multi-wavelength follow up, summarised to date in Oliver et al. (1999) and presented in detail in other papers (e.g. Ciliegi et al. 1999, Gruppioni et al. 1999). Here we present the completeness, reliability and extragalactic source counts from our initial Preliminary $`6.7\mu `$m and $`15\mu `$m ISOCAM catalogues. A future paper will present the Final Analysis products from the ISOCAM ELAIS data, which is expected to improve on the Preliminary Analysis presented here.
This paper is structured as follows. In section 2 we describe the Preliminary Analysis CAM pipeline, explaining the artefacts in the data (section 2.1), and the pipeline algorithm (section 2.2). The results from the Preliminary Analysis catalogue are presented in section 3. Our various completeness and reliability estimates are discussed in section 3.2, and the segregation of extragalactic from stellar sources in section 3.3. Section 3.4 presents the source counts in both wavebands. These results are compared with source count models and previous results in section 4, where we also discuss the implications for the evolution of star forming galaxies and on the star formation history of the Universe.
## 2 ELAIS CAM Preliminary Analysis
### 2.1 Data quality
The ELAIS CAM survey proper was conducted in raster mode (astronomical template CAM01), with the LW-2 ($`6.7\mu `$m) and LW-3 ($`15\mu `$m) filters. Details of the CAM01 Astronomical Observation Template (AOT) can be found in paper I. The CAM detector is stepped across the sky in a grid pattern, with roughly half-detector-width steps in one direction and roughly whole detector widths in the other, covering approximately half a square degree per raster. This pattern leads to a redundancy of at least $`\times 2`$ over most of the area surveyed. At each raster pointing (i.e. each grid position of the raster) the $`32\times 32`$ CAM detector is read out several times.
Like the ISO-HDF North data (Serjeant et al. 1997, 1999, Aussel et al. 1999, Désert et al. 1999) the ELAIS CAM data contains many problematic artefacts. Because of the frequent and complicated glitches, we do not take the approach of reconstructing a sky-map and searching for sources in these maps. Rather, we look for the characteristic signatures of sources and glitches in the time histories of individual pixels. See Starck et al. (1998) or Aussel et al. (1999) for more details.
The CAM detector also exhibits hysteresis. Source fluxes are initially around a factor of $`2`$ fainter in instrumental units than the stabilised (i.e. asymptotic) value. Our survey strategy ensures that sources almost always have corroborating sightings in separate pixels. This permits a filter to remove glitch events from candidate sources.
### 2.2 Preliminary source extraction pipeline
#### 2.2.1 Preprocessing and deglitching
The available CAM data reduction software underwent several substantial improvements over the lifetime of ISO, as the knowledge of the detector characteristics improved. However, from the outset we needed a method of preliminary data analysis, to feed for example immediate follow-up projects. Such a Preliminary Analysis pipeline may not of course represent best practice at the end of the ISO mission, but should at least provide reasonably complete and reliable preliminary source list. It was decided that the CAM Preliminary Analysis data reduction should be as uniform as possible, which required that the Preliminary Analysis pipeline be fixed at an early stage. Accordingly our adopted pipeline could not incorporate the accurate field distortion corrections (Abergel et al. 1998), which were not established at the start of the mission, nor the analytic models for the cosmic ray transients and detector hysteresis which were developed in the course of the ISO lifetime (e.g. Lari 1999, Abergel et al. 1998). Nevertheless, such improvements will be incorporated in future ELAIS Final Analysis products.
The data reduction was performed using the Interactive Data Language (IDL) software. The initial steps of the CAM Preliminary Analysis pipeline are straightforward. The edited raw data supplied by ESA were converted to IDL structures using the CAM Interactive Analysis (CIA, April 1996 version), and converted from ADUs (analogue-to-digital units) to ADUs per second. The default dark frame was subtracted from each exposure. To estimated the noise level in each pixel, we performed an iterative Gaussian fit to the histogram of readout values for each pixel. Cosmic ray spikes were then identified by $`>4\sigma `$ rises followed (one or two readouts later) by $`>4\sigma `$ falls. A similar algorithm was used to identify occasional readout troughs. The readout histograms were then re-fit and an empirical sky flat field was obtained; the default ESA-supplied flat field was found to give very unsatisfactory results. An attempt was made to model the initial detector stabilisation using the IAS model within the CAM Interactive Analysis (CIA) software package.
#### 2.2.2 Background estimation and Source Detection
Unlike the detectors on the IRAS satellite (e.g. Neugebauer et al. 1984), the transients in the ISOCAM detector make the background levels in each pixel vary strongly and discontinuously with time. The approach adopted in the Preliminary Analysis was to estimate the background level in a given pixel and pointing from linear fits (in time) to the readouts in previous and subsequent pointings, then identify candidate sources from $`>3\sigma `$ features above the background in the pixel readout timelines. The source extraction pipeline is therefore spatially $`1`$ dimensional. This simple approach avoids any explicit parameterisations of the transient profiles, which were not available at the time, and allows a local error estimate of the background level. This error was taken to be the formal error on the fit, but did not incorporate the instrumental noise at each data point. Since cosmic rays were frequently observed to cause discontinuities, the range for the linear fits extended no further than the nearest (in time) cosmic ray, and in any case not longer than $`3`$ pointings. If an acceptable $`\chi ^2`$ was not obtained in the linear fit, the range for the fit was decreased; if a good fit was still not obtainable, an average of the $`10`$ nearest readouts was taken. Sources could then be identified from their excess above the extrapolated backgrounds. Since sources also create discontinuities, the data stream was iteratively re-fit treating $`>3.5\sigma `$ sources in the same way as cosmic rays above. An initial list of brightest sources was obtained before the iterative fitting, by searching for their discontinuous rise at the start of a pointing, and discontinuous fall at the end. Where sources were found (whether initially or in the iterative fitting) the background level in later iterations was extrapolated from only the previous pointings, ignoring the subsequent pointings, since the hysteresis after a source would otherwise lead to an overestimation of the background level. Sources were not extracted from the first or final pointings of the rasters, owing to limitations in the background fitting routines at the time the pipeline was frozen. Note that this iterative source extraction does not distinguish genuine sources from “fader” transient events, so a further filter for source candidates is still required.
#### 2.2.3 Source Corroboration
The source detections in the pixel readout histories were spatially merged in each pointing using a connected pixel algorithm. (Note that there was no minimum number of pixel detections, unlike e.g. connected pixel source detection on an oversampled CCD image.) We then used the $`\times 2`$ redundancy in the CAM rasters to search for corroborating observations of each source candidate. Genuine sources should be present in both observations, but glitch events should not be confirmed except by chance. The adopted search radius of $`2`$ pixels, while large enough to safely encompass the (then uncertain) field distortion, nevertheless led to a large number of spurious detections, with the majority of source candidates at $`15\mu `$m due to glitch corroboration. Each candidate corroborated source was therefore examined by eye independently by at least two observers, who assigned quality flags of $`14`$, where $`1`$ refers to a “definite source,” $`2`$ is a “probable source,” $`3`$ “probably not” and $`4`$ “definitely not.” Each raster in the LW-2 ($`6.7\mu `$m) filter yielded typically $`100`$ events in total of which around one half were classified as probable sources by at least one observer. In LW-3 ($`15\mu `$m) both the number and fraction of spurious events was much higher: there were typically $`30`$ strong source candidates, and a further $`3050`$ sources where the classification was ambiguous or debatable, with typically around $`200500`$ spurious events. Note that the Preliminary Analysis algorithm will necessarily miss sources with (a) only one observation or (b) corroborating observations in only the first or last pointings, so the effective area is slightly less than the nominal $`12`$ square degrees.
#### 2.2.4 Astrometric corrections
After the Preliminary Analysis reduction was complete, we improved the astrometry by incorporating the latest field distortion correction into the $`15\mu `$m source catalogues. Several sources with strong transient events nearby had their centroids strongly affected by the glitches. We therefore adopted a simple strategy for our Preliminary Analysis astrometry and flux calibration: the flux and (distortion-corrected) position of a source are taken from those of the brightest single-pixel detection of that source, excluding transients. We found this to be superior to (eg) masking nearby transients by eye then recalculating the centroids of the eyeball-accepted sources, particularly if the PSF wings lie on the detector but the source itself is just outside. Our adopted algorithm should yield astrometry accurate to $`\pm 3^{\prime \prime }`$ in the absence of any other systematic errors. Two such systematics were expected in our data: firstly, errors in the position of the lens introduce a random astrometric offset to each raster of order one pixel; secondly, any errors in the calculation of the pointing position by CIA would offset any sources in that pointing. By examining the offsets with the likelihood-ratio identifications we can determine the lens offset empirically (section 3.3). However, several rasters were found to have bimodal distributions of ISO–optical offsets, due to some unknown error in the CIA-derived astrometry in at least part of the raster. We therefore rederived the pointing astrometry using the ESA-supplied IIPH.FIT astrometry file, using the median coordinate positions in the duration of the pointing. This was found to remove the bimodality.
#### 2.2.5 Aperture Corrections to Photometry
As discussed above, our source extraction method involves looking for characteristic time signatures in individual pixels. Without (at the time) a reliable and exact model for the glitch events, nor a reliable glitch event identification, we found that aperture photometry around our source positions was often seriously affected by nearby glitches. Instead of aperture photometry, we simply took the brightest flux of the pixels detecting the source, excluding (by eye) those pixels affected by glitches. Clearly, some aperture correction is needed to correct these peak pixel fluxes to total fluxes.
We can quantify these aperture corrections using a PSF model. In figure 1 we show the predicted flux in the brightest pixel of two randomly positioned observations of a point source (recall that at least two observations of a source are required for it to pass the Preliminary Analysis selection). At $`15\mu `$m, the brightest pixel has a flux of $`0.4\pm 0.1`$ times the total flux of that source. At $`6.7\mu `$m the histogram is more sharply peaked, since the PSF is undersampled; the peak flux is always less than $`0.69`$ times the total at $`6.7\mu `$m, but is greater than $`0.5`$ ($`0.6`$) of the total in $`>90\%`$ ($`>75\%`$) of occasions. We therefore applied global aperture corrections of $`2.36`$ at $`15\mu `$m and $`1.54`$ at $`6.7\mu `$m to our peak fluxes.
## 3 Results
### 3.1 Eyeballing results
Our eyeballing results imply our catalogue is highly reliable to at least $`3`$ mJy at $`15\mu `$m. In the ELAIS areas considered in this paper, there were $`715`$ $`15\mu `$m sources accepted by two observers, of which $`510`$ had APM identifications (section 3.3); of the $`816`$ singly-accepted sources at $`15\mu `$m, $`212`$ had APM identifications. The singly-accepted sources are heavily skewed to faint fluxes: $`90\%`$ are fainter than $`2.2`$ mJy. If we choose to regard all sources accepted by only one observer as false positives, and all sources accepted by two as true positives, then we obtain $`80\%`$ reliability for sources accepted by any observer at $`3`$ mJy and $`95\%`$ at $`5`$ mJy. These are one of the most pessimistic assumptions we could make; if alternatively we merge the doubly-accepted sources with the optically identified singly-accepted sources, the reliability (fraction doubly-accepted) of the sources accepted by any observer rises to $`95\%`$ at $`3`$ mJy and increases at brighter fluxes. The blank-field singly-accepted sources are heavily skewed to faint fluxes, so we can neglect their effect on the counts above $`3`$ mJy and combine the doubly-accepted sources with the optically identified singly-accepted for the purposes of the source counts. However this is not to say that the fainter singly-accepted sources necessarily have lower reliability, since the genuine fainter sources may have fainter optical counterparts. The situation is similar at $`6.7\mu `$m. Of the singly-accepted sources $`95\%`$ have fluxes less than $`1.5`$ mJy.
### 3.2 Completeness, reliability and flux calibration
#### 3.2.1 Repeated observations
There are several potential approaches to estimating the Preliminary Analysis completeness and reliability, the most obvious of which is repeated observations over small areas. Oliver et al. (1999) present $`10`$ repeat observations of a small ELAIS raster, covering at least six known sources by design. The source extractions in the individual reductions confirmed our result that the source extraction is highly complete and reliable at flux densities above $`2.5`$mJy at $`15\mu `$m.
#### 3.2.2 PSF-wing test
One robust test of our temporal source extraction is to examine the point spread function wings of bright sources. Like the cross-correlation with bright stars, this relies on the knowledge that a given pixel is illuminated by a known source flux, but has the advantage that the comparison can be taken to arbitrarily faint flux levels. We selected the $`30`$ brightest $`15\mu `$m sources with stellar optical identifications, and selected a model theoretical Point Spread Function (PSF; Aussel, priv. comm.) with the smallest RMS difference in the central $`3\times 3`$ PSF pixels. We normalised the PSF to the source flux using the mean observed flux in the central pixel. Using this model we predicted the mean flux in each pixel of the PSF wings, and hence determined the single-pixel detection efficiency $`f`$ of the temporal source extraction. The results are shown in figure 1. We can also compare the predicted PSF fluxes with the fluxes extracted by the Preliminary Analysis algorithm, shown in figure 1. The relation is encouragingly linear, though the scatter is larger than expected (by $`\times 2`$ at $`15\mu `$m) based on the quoted errors in the Preliminary Analysis sky background fits. This is perhaps not surprising, since the background fitting algorithm does not make use of the detector noise in the fit, so will tend to underestimate the background level uncertainty; also, non-white noise features may prevent the noise scaling as $`\sqrt{N}`$. Oddly, the discrepancy in the scatter is largest at brightest fluxes. Plausible explanations include slight errors in the source centroids or in the theoretical PSF shape, both of which sensitively affect the brightest flux predictions. Analogous results for the $`30`$ brightest $`6.7\mu `$m stars are also shown in figure 1. The undersampling of the PSF at $`6.7\mu `$m makes the scatter harder to interpret: more of the flux is contained in the central pixel making it less sensitive to errors in the assumed PSF shape, but is more affected by the much larger uncertainty in the centroid.
There are several caveats which apply to the high apparent completeness in figure 1. The chance detections of “faders” will tend to increase this estimate, but since both the $`6.7\mu `$m and $`15\mu `$m completeness appear to asymptote to $`0.015`$ (i.e. probability of detection of nearby spurious events is $`1.5\%`$) this appears not to be a serious problem. This is also only applicable to single-pixel detections, whereas in fact the CAM detector Nyquist samples the PSF at $`15\mu `$m. Indeed at $`15\mu `$m the fraction of the flux in the brightest pixel rarely exceeds $`0.5`$ though the same does not apply at $`6.7\mu `$m. The completeness also is not the detected source fraction $`F`$ but rather $`F^2`$, since we require at least $`2`$ independent detections for a source to be accepted by our algorithm. There is also a slight bias in that the brightest detected sources are not typically observed when the detector is suffering the strongest transients, because the sources would not otherwise be detected. Finally, any incompleteness caused by the eyeballing, or by any areal coverage lost to e.g. cosmic ray impacts, would not be included in these estimates as they stand.
We can therefore estimate the Preliminary Analysis completeness from the PSF-wing test in the following way. We can convolve the single-pixel detection rate in figure 1 with the peak flux distribution expected from the PSF models to obtain the source sensitivity $`F`$ for single observations. The Preliminary Analysis completeness, before eyeballing, will then be proportional to $`F^2`$, assuming the PSF wings themselves are representative of the data as a whole.
We fit the detected fractions in figure 1 with $`\mathrm{tanh}`$ functions, i.e.
$$f(S)=\frac{(f_{\mathrm{max}}f_{\mathrm{min}})}{2}\mathrm{tanh}(\alpha \mathrm{log}_{10}(S/S_0)+1)+f_{\mathrm{min}}$$
(1)
where $`f_{\mathrm{min}}`$ and $`f_{\mathrm{max}}`$ represent the asymptotic limits at faint and bright fluxes respectively. We define the single-pixel source detection fraction to be $`f^{}(S)=f(S)f_{\mathrm{min}}`$, and use a grid of PSF models (each normalised to 1) spanning the possible range of centroid positions to estimate the single-pointing source detection:
$$F(S)=\frac{1}{N}\underset{i=1}{\overset{N}{}}f^{}(S\times S_{\mathrm{peak},\mathrm{i}})$$
(2)
where $`S_{\mathrm{peak},\mathrm{i}}`$ is the peak flux in the $`i`$th PSF. This assumes that if a source is not detected in the peak pixel, it will not be detected in any pixel. The maximum areal coverage of the Preliminary Analysis, ie excluding the first and last pointings of each raster, and regions with no redundancy, is $`\mathrm{\Omega }_{\mathrm{max}}=10.0`$ square degrees at $`15\mu `$m, and $`\mathrm{\Omega }_{\mathrm{max}}=6.51`$ square degrees at $`6.7\mu `$m. From this we obtain the Preliminary Analysis areal coverage as a function of flux:
$$\mathrm{\Omega }_{\mathrm{PreliminaryAnalysis}}(S)=\mathrm{\Omega }_{\mathrm{max}}F^2(S)$$
(3)
The final areal coverage from the PSF-wing test is plotted in figures 2 and 3.
#### 3.2.3 Comparison with IAS and CEA pipelines
As a final check of the completeness of our Preliminary Analysis catalogue, we compared our source extraction in the repeat observation regions (c.f. Oliver et al. 1999 and section 3.2.1) with extractions made with the CAM-Deep pipeline developed at the Commisariat á L’Energie Atomique, Saclay (CEA; Baker, priv. comm.) and a pipeline based on the “Triple Beamswitch” method developed at Institut d’Astrophysique Spatiale (IAS; Clements, priv. comm; Désert et al. 1998). Of our six robust sources (section 3.2.1), CEA and IAS both identify three (the same three), with fluxes brighter than $`3`$mJy. This appears to be mainly because the Preliminary Analysis pipeline extracts lower signal-to-noise sources, but supplements with greater manual eyeballing. Nevertheless, this confirms that we are not significantly underestimating the surface density of sources brighter than $`3`$mJy.
#### 3.2.4 Flux calibration and Cross-correlation with bright stars
Another useful test of the completeness is to search for detections at the locations of bright stars (see section 3.3 for details of the optical identification algorithm). If the spectral types of the stars are known, one can predict their mid-infrared fluxes. An exception would be dust-shell stars, but these are expected to be rare in the survey. In Crockett et al (1999) and Crockett (1999) the sources are cross-correlated with stars from the Simbad and Hipparcos databases. All the $`22`$ stars with predicted fluxes brighter than $`3`$mJy at $`15\mu `$m were detected in the Preliminary Analysis, and all but two of the $`20`$ stars with predicted $`6.7\mu `$m fluxes above $`1`$mJy appeared in the Preliminary Analysis. To assess the level of random associations, we randomised the stellar positions within the ELAIS fields and repeated the cross-association, and found none. The $`15\mu `$m completeness is shown in figure 4. Note that this is an extremely robust estimator of the completeness, since the stellar flux prediction algorithm reproduces IRAS fluxes well, and the stars cannot be argued to lie on atypical regions of the detector.
However, the flux calibration implied by these associations is around a factor of $`2`$ discrepant with the expectation that ADU/gain/second $``$ mJy at $`15\mu `$m (see figure 5). Across the entire range in flux, the $`\pm 1\sigma `$ limits on the log of the calibration ratio are $`0.246\pm 0.050`$, ie the ratio is $`1.75_{0.23}^{+0.26}`$. As shown in figure 5, there are hints that this calibration is a function of flux, with a calibration ratio of $`2`$ preferred at faint fluxes (where indeed most of our sources lie). (The ISOCAM Observer’s Manual recommends a conversion of approximately $`2`$ ADU/gain/second to $`1`$ mJy at $`15\mu `$m for fully stabilised sources, and after correcting for the loss of flux due to lack of stabilisation (e.g. section 2.1) becomes around a $`1:1`$ conversion.) To predict the fluxes we followed the procedure of Crockett et al. (1999) and Misoulis et al. (1998), incorporating the passband profiles (for more details we refer the reader to these papers; the algorithm accurately reproduces stellar mid-IR IRAS fluxes so there are unlikely to be significant systematics in the CAM flux predictions). It is not immediately clear what might cause such a discrepancy, though there are several possibilities, such as the uncertainties in the PSF for the aperture correction, and the assumed $`\times 2`$ loss in flux (section 2.1) from the lack of an upward source stabilisation correction. (Note that differences in the PSF due to the differing spectral slopes of the stars are too small to account for the discrepancy: e.g. Aussel et al. (1999) find it only to be a $`10\%`$ effect.) For the purposes of the source counts we will adopt the mJy : ADU/g/s = $`1:2`$ stellar calibration implied in figure 5, where the ADUs are not corrected for losses due to lack of stabilisation.
At $`6.7\mu `$m a lower bound of $`95\%`$ can be made on the completeness at fluxes $`>10`$ mJy, but the uncertain aperture corrections make applying the stellar flux calibration more difficult. As a preliminary measure we therefore simply take the pre-flight ISOCAM Observer’s Manual calibration at $`6.7\mu `$m, corrected by a factor $`2`$ (section 2.1) to account for the loss in flux from lack of stabilisation.
In summary, our various completeness estimates yield a $`1`$ mJy limit to the completeness at $`15\mu `$m, and $`0.5`$mJy at $`6.7\mu `$m (figures 2 and 3) using the ISOCAM observer’s manual flux calibration corrected by a factor of $`2`$ to account for stabilisation loss (section 2.1). However our stellar cross-correlation suggests we have underestimated the fluxes by a factor of $`2`$ at $`15\mu `$m (figure 5) so the $`15\mu `$m completeness quoted should by revised upwards to $`2`$ mJy. In all subsequent discussion, the $`15\mu `$m ELAIS fluxes are assumed to obey this stellar flux calibration.
### 3.3 Optical identifications
In this section we summarise the optical identification algorithm used for the Preliminary Analysis catalogue, for APM stars and galaxies (McMahon & Irwin 1992). We adopted the likelihood ratio procedure of Sutherland & Saunders (1992) to associate our Preliminary Analysis sources with known objects. The surface density of catalogue objects as a function of magnitude is incorporated into the likelihood of each identification of the Preliminary Analysis sources with the catalogue objects. Following Mann et al. (1997), we define the likelihood ratio to be the ratio of the probability of detecting a genuine counterpart to the source with the position and flux of the catalogue object, to the probability of such an association occurring by chance given the positional errors. For a catalogue surface density $`n(f)`$ (where $`f`$ is the flux), a positional uncertainty $`ϵ(x,y)`$ (where $`x`$ and $`y`$ are e.g. Cartesian coordinates) and an a priori flux distribution of IDs given by $`q(f)`$, the likelihood ratio is given by
$$LR(f,x,y)\frac{q(f)ϵ(x,y)}{n(f)}$$
(4)
Using this expression and, assuming $`q(f)`$ to be constant as a function of optical magnitude, we found, for each Preliminary Analysis source, the object in the APM catalogue giving the highest likelihood ratio. We normalised the likelihood ratios by finding the maximum likelihood-ratio associations around random source positions, to yield probabilities for each Preliminary Analysis identification, and accepted identifications with random probabilities of less than $`0.3`$. Slight errors in the lens positioning lead to systematic astrometric shifts in each raster, of order a few arcseconds. To correct for this, the high-likelihood identifications in each raster were used to determine any systematic astrometric offset. The identifications were subsequently rederived. There were not enough reliable optical identifications to obtain a robust estimate of the lens offset in the smaller ELAIS areas, so the analysis was restricted to the main ELAIS areas of N1, N2, N3 and S1. Further discussion of the identifications is deferred to later papers in this series.
### 3.4 Source counts
Using the completeness and reliability from section 3.2, we can extract the extragalactic source counts at $`6.7\mu `$m and $`15\mu `$m. For the purposes of the counts we included all sources accepted by (at least) two observers. From this list, we exclude $`15\mu `$m stellar identifications brighter than approximately $`B=16.5`$ (assuming a monotonic stellar plate saturation correction converting $`B=18`$ galaxy magnitudes to $`B\stackrel{<}{_{}}16.5`$ stars), but include fainter stellar IDs since these are expected to be predominantly AGN. All stellar IDs at $`6.7\mu `$m were eliminated from the extragalactic source counts. It is highly unlikely for stars with $`B>16.5`$ to be detected at $`15\mu `$m to these limits ($`B=16.5`$ is equivalent to $`S_\mathrm{B}=1`$ mJy; see also Crockett et al. 1999 and Crockett 1999), with the exception of rare dust-shell stars. Note also that all the $`15\mu `$m ($`6.7\mu `$m) sources accepted by only one observer are fainter than $`2.2`$mJy ($`2.5`$mJy), so we can be highly confident of the reliability of sources brighter than these limits.
The segregation of extragalactic from stellar counts is robust at $`15\mu `$m, but the large fraction of stellar IDs at $`6.7\mu `$m make it possible that some faint stars have been included in the extragalactic counts at this wavelength: of the $`794`$ doubly-accepted sources at $`6.7\mu `$m, only $`79`$ have stellar APM classifications faint enough to be accepted in the extragalactic counts. We eyeballed the DSS images of every optical identification of the $`6.7\mu `$m and $`15\mu `$m sources. We found the by-eye classifications to agree with the APM in nearly all cases. The resultant extragalactic counts are virtually indistinguishable from the automated segregations at both $`6.7\mu `$m and $`15\mu `$m.
The extragalactic source counts are plotted in figure 6, using the by-eye star-galaxy separation and the stellar flux calibration. The $`15\mu `$m Lockman Hole ISOCAM survey data will be discussed in Elbaz et al. 1999. The counts from this survey at around the $`5`$mJy level, which overlap with the ELAIS counts, appear significantly lower than those of ELAIS. However this is entirely attributable to the $`\times 2`$ differences in assumed flux calibration. The counts are in excellent agreement in instrumental units or in mJy with the same flux calibration assumption; alternatively, a $`30\%`$ reduction of the ELAIS Preliminary Analysis $`15\mu `$m flux calibration factor would also bring the counts into formal $`1\sigma `$ agreement while remaining consistent with figure 1. However this may require a commensurate reduction in the IRAS counts. Such a reduction has been argued to be necessary by Elbaz et al. (1999) in order to account for large-scale-structure effects in the Rush et al. sample.
Also plotted are the $`12\mu `$m source counts calculated by Oliver et al. (1997) from the Rush et al. (1993) sample, using the QMW IRAS Galaxy Mask (Rowan-Robinson et al. 1991), and transposed to $`15\mu `$m assuming a population mix matching the Pearson & Rowan-Robinson (1996) predictions. The faint source counts from the ISO-HDF North survey (Oliver et al. 1997) are also shown in the counts figures. The ELAIS extragalactic Preliminary Analysis counts at $`15\mu `$m are consistent with an interpolation between the ISO-HDF North and Rush et al. (1993) data sets, and at $`6.7\mu `$m with reasonable extrapolations from the ISO-HDF North. The source density at $`15\mu `$m is also in good agreement with that obtained at $`12\mu `$m by Clements et al. (1999), though the differing K-corrections make it not immediately clear that the counts must necessarily agree (e.g. Xu et al. 1998, Serjeant 1999). Note that the $`6.7\mu `$m counts have significant photometric errors independent of flux, due to the undersampling of the PSF (figure 1). For a constant source count slope the shape of the counts is unaffected by flux-independent errors, so we can regard the $`6.7\mu `$m counts as subject to a potential systematic error in the form of a horizontal shift. Such a systematic is less than or of order a factor $`2`$ in flux. We overplot the model predictions from Pearson & Rowan-Robinson (1996) in figures 6 and 7, as well as the model predictions of Franceschini et al. (1994, 1997). Also overplotted are the Guiderdoni et al. (1997) $`15\mu `$m model counts and the evolving models from Xu et al. (1998), with and without mid-IR spectral features. Note that all the Xu et al. models have been renormalised to match the IRAS counts. Figure 6 also compares the observations to a variety of non-evolving models. Apart from the renormalisation to the IRAS Rush et al. (1993) $`\times `$ QIGC counts, these are the same non-evolving models as discussed in Elbaz et al. 1998.
## 4 Discussion
The experimental agreement with the evolving models of Pearson & Rowan-Robinson (1996), Franceschini et al. (1994), Guiderdoni et al. (1997) and Xu et al. (1998) at $`15\mu `$m over seven orders of magnitude in flux density is striking. The starbursts in the Pearson & Rowan-Robinson (1996) models are normalised to the $`60\mu `$m IRAS counts, but not explicitly to the $`12\mu `$m counts. The slight overprediction of the IRAS counts is also present in the ELAIS counts, until a slight upturn at around $`10`$mJy (which is not an effect of incompleteness or low reliability) departing from the Euclidean slope brings the data into closer agreement with the model. Of the four Xu et al. (1998) models, the $`(1+z)^3`$ luminosity evolution models have the stronger upturn, reproducing the observed counts slightly more well than the alternative $`(1+z)^4`$ density evolution. This is clearer still in the $`15\mu `$m differential counts (figure 8), where luminosity evolution is shown to make a much better fit to the ELAIS counts, though an even larger excess is suggested by the ISO-HDF counts. Note that we renormalised the Xu et al. (1998) predictions by $`\times 1.8`$ to match the IRAS counts.
The source count models present rather different predictions at $`6.7\mu `$m. The Pearson & Rowan-Robinson (1996) model has only a slight overestimate, accountable for instance to the flux calibration uncertainties. However, the underprediction in the Franceschini et al. (1997) model is over an order of magnitude, probably too large to be an artefact of our albeit uncertain flux calibration at this wavelength. This discrepancy is most likely to be due to deficiencies in the assumed spectral energy distributions, which are not well-constrained in this wavelength range (e.g. Serjeant 1999), rather than due to incorrect assumptions about the evolution or population mixes.
The loss of the Wide Field Infrared Explorer (WIRE) satellite was a serious blow to infrared extragalactic astronomy. In the hope or expectation of a new WIRE mission, we can compare our source counts with the expectations of the WIRE team. WIRE was to survey at least $`170`$ square degrees to a limiting $`12\mu `$m sensitivity of $`1.9`$mJy, and smaller areas to deeper limits. Our source counts imply a surface density of approximately $`100`$ galaxies per square degree to the projected WIRE $`12\mu `$m bright survey limit, using the Pearson & Rowan-Robinson or Franceschini et al. $`15\mu `$m source count models to extrapolate. This is larger than projected source density of the WIRE team (q.v. figure 6, bottom-right panel, and 8 left panel, where the WIRE predictions were renormalised upwards by a factor of $`2`$ to match e.g. the observed IRAS and ELAIS counts). Although encouraging for the numbers of sources, it suggests that surveys much deeper than this will very rapidly be confusion limited. This is in good agreement with the observations of Oliver et al. (1997). The composite counts from the various $`15\mu `$m surveys also imply that confusion will be irrelevant e.g. for NGST, but will be significant for the MIPS imager on SIRTF even in short ($`2000`$ second) exposures. The gains in detector size and sensitivity in future missions will offer very large improvements in wide-field survey efficiency: for example, a survey of a similar depth and areal coverage to ELAIS at $`15\mu `$m will be possible with MIPS on SIRTF in only $`10`$ hours.
The upturn in the counts at faint fluxes continues in the Hubble Deep Field ISO counts, and is too large to be attributable to starburst K-corrections (e.g. Xu et al. 1998, Elbaz et al. 1998a,b). This excess is above any reasonable no-evolution predictions (figure 6), which all have shallower slopes than the evolving models independent of K-correction and world models. This implies that the faintest sources in ELAIS are sampling a significantly cosmologically evolving mid-IR population. Given the strong evolution in both the comoving star formation rate claimed from optical-UV observations and in the quasar comoving number density, as well as the large fraction of AGN and starbursts expected in ELAIS from the counts models, we can reasonably expect that optical spectroscopy of ELAIS and fainter samples will have a major impact on the study of dust-shrouded star formation and quasar activity and their evolution. The prediction from the Pearson & Rowan-Robinson (1996) source count models is for of order a few objects as (intrinsically) luminous as IRAS F10214+4724 out to $`z=4`$ over the entire ELAIS areas. The MIR luminosity function and luminosity density can be used to derive a comoving star formation rate without the difficult problems of extinction corrections that affect optical-UV estimates, and (for the ISOCAM LW3 filter) is reasonably free of K-correction uncertainties (Xu et al. 1998, Serjeant 1999). The ELAIS survey is also highly sensitive to obscured quasars and will be an exceptional resource for active galaxy unification models. The ELAIS limits can also provide important constraints for sources detected at other wavelengths: for example, the $`15\mu `$m, $`90\mu `$m and $`175\mu `$m limits at $`850\mu `$m source positions can provide robust redshift constraints (e.g. Hughes et al. 1998).
## 5 Conclusions
The extragalactic source counts agree extremely well with all evolving model predictions (Franceschini et al. 1994, Pearson & Rowan-Robinson 1996, Guiderdoni et al. 1997, Xu et al. 1998) over seven orders of magnitude in $`15\mu `$m flux density. The Pearson & Rowan-Robinson (1996) models can broadly reproduce the $`6.7\mu `$m extragalactic counts, but the observations are in excess of the Franceschini et al. (1997) predicted counts at this wavelength using our preliminary $`6.7\mu `$m flux calibration. All no-evolution models are clearly excluded, and imply a cosmologically evolving population of obscured starbursts and/or active galaxies dominates below $`10`$ mJy at $`15\mu `$m, independent of K-correction assumptions. Source confusion appears to have been underestimated in the WIRE and SIRTF missions, but will not impact significantly on the NGST.
## Acknowledgements
We would like to thank Dave Alexander and Ruth Carballo for helpful comments and proofreading of this paper. This paper is based on observations with ISO, an ESA project with instruments funded by ESA member states (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) and with participation of ISAS and NASA. The ISOCAM data presented in this paper was analysed using “CIA”, a joint development by the ESA Astrophysics Division and the ISOCAM Consortium. The ISOCAM Consortium is led by the ISOCAM PI, C. Cesarsky, Direction des Sciences de la Matiere, C.E.A., France. This work was supported by PPARC (grant number GR/K98728) and by the EC TMR Network programme (FMRX-CT96-0068).
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# 1 Introduction
## 1 Introduction
In recent years, attention has been paid to properties of supersymmetric gauge theories with matter multiplets in various dimensions. In particular, parity odd interactions obtained from one-loop fermionic determinant or equivalently hexagon diagrams have received much interest till now in their interpretation as interactions among branes . They also give distinct contribution to the imaginary part of the effective action and make sense beyond renormalizability. These interactions are often called anomalous interactions (or Wess-Zumino type terms albeit the fact that can be represented locally in $`D>4`$ dimensions).
In this paper, we will provide another example of anomalous interactions. We consider the supersymmetric $`USp(2k)`$ gauge theory in five dimensions with one matter hypermultiplet in the antisymmetric representation and $`n_f`$ matter hypermultiplets in the fundamental representation. We consider this theory on the new phase where the $`vev`$s of the scalars belonging to the antisymmetric hypermultiplet are also nonvanishing. This theory is a many-probe generalization of $`SU(2)`$ gauge theory with $`n_f`$ fundamental matters and is also related to the $`USp(2k)`$ matrix model via matrix T duality operation. In the former case, our result exhibits a magnetic interaction among D4-branes in nontrivial gauge backgrounds. On the Coulomb phase, the anomalous interaction has been computed in .
In section 2, we exhibit the five dimensional SYM lagrangian dealt with in this paper and find the background configurations of our model. In section 3, we present our calculation and our final result is eq. (3. 15)
## 2 Set up
### 2.1 Lagrangian of Five Dimensional $`USp(2k)`$ Gauge Theory
We discuss the five dimensional worldvolume gauge theory associated with $`USp(2k)`$ matrix model. The lagrangian of this five dimensional theory is given by
$$=_{adj}+_{asym}+_{fund},$$
(2. 1)
where
$$_{adj}=\frac{1}{g^2}\mathrm{Tr}\left(\frac{1}{4}v_{\mu \nu }v^{\mu \nu }+\frac{1}{2}[𝒟_\mu ,v_7][𝒟^\mu ,v_7]+\frac{i}{2}\overline{\mathrm{\Psi }}_{(adj)}\mathrm{\Gamma }^\mu [𝒟_\mu ,\mathrm{\Psi }_{(adj)}]\frac{1}{2}\overline{\mathrm{\Psi }}_{(adj)}\mathrm{\Gamma }^7[v_7,\mathrm{\Psi }_{(adj)}]\right),$$
(2. 2)
$`_{asym}`$ $`=`$ $`{\displaystyle \frac{1}{g^2}}\mathrm{Tr}\left({\displaystyle \underset{M_+=5,6,8,9}{}}{\displaystyle \frac{1}{2}}[𝒟_\mu ,v_{M_+}][𝒟^\mu ,v_{M_+}]+{\displaystyle \underset{M_+=5,6,8,9}{}}{\displaystyle \frac{1}{2}}[v_7,v_{M_+}][v_7,v_{M_+}]\right)`$ (2. 3)
$`+`$ $`{\displaystyle \frac{1}{g^2}}\mathrm{Tr}\left({\displaystyle \frac{i}{2}}\overline{\mathrm{\Psi }}_{(asym)}\mathrm{\Gamma }^\mu [𝒟_\mu ,\mathrm{\Psi }_{(asym)}]{\displaystyle \frac{1}{2}}\overline{\mathrm{\Psi }}_{(asym)}\mathrm{\Gamma }^7[v_7,\mathrm{\Psi }_{(asym)}]\right)`$
$``$ $`{\displaystyle \frac{1}{g^2}}\mathrm{Tr}\left({\displaystyle \underset{M_+=5,6,8,9}{}}{\displaystyle \frac{1}{2}}\overline{\mathrm{\Psi }}_{(asym)}\mathrm{\Gamma }^{M_+}[v_{M_+},\mathrm{\Psi }_{(adj)}]+{\displaystyle \underset{M_+=5,6,8,9}{}}{\displaystyle \frac{1}{2}}\overline{\mathrm{\Psi }}_{(adj)}\mathrm{\Gamma }^{M_+}[v_{M_+},\mathrm{\Psi }_{(asym)}]\right)`$
$`+`$ $`{\displaystyle \frac{1}{g^2}}\mathrm{Tr}{\displaystyle \underset{M_+,N_+=5,6,8,9}{}}{\displaystyle \frac{1}{4}}[v_{M_+},v_{N_+}][v_{M_+},v_{N_+}],`$
$`_{fund}`$ $`=`$ $`{\displaystyle \frac{1}{g^2}}{\displaystyle \underset{f=1}{\overset{n_f}{}}}({\displaystyle \underset{M_+=5,6,8,9}{}}{\displaystyle \frac{1}{2}}𝒟_\mu v_{(f)}{}_{M_+}{}^{}𝒟^\mu v_{(f)}{}_{M_+}{}^{}+{\displaystyle \underset{M_+=5,6,8,9}{}}{\displaystyle \frac{1}{2}}v_7v_{(f)}{}_{M_+}{}^{}v_7v_{(f)}{}_{M_+}{}^{})`$ (2. 4)
$`+`$ $`{\displaystyle \frac{1}{g^2}}{\displaystyle \underset{f=1}{\overset{n_f}{}}}\left({\displaystyle \frac{i}{2}}\overline{\mathrm{\Psi }}_{(f)}\mathrm{\Gamma }^\mu 𝒟_\mu \mathrm{\Psi }_{(f)}{\displaystyle \frac{1}{2}}\overline{\mathrm{\Psi }}_{(f)}\mathrm{\Gamma }^7v_7\mathrm{\Psi }_{(f)}\right)`$
$``$ $`{\displaystyle \frac{1}{g^2}}{\displaystyle \underset{f=1}{\overset{n_f}{}}}\left({\displaystyle \underset{M_+=5,6,8,9}{}}{\displaystyle \frac{1}{2}}m_f{}_{}{}^{2}v_{(f)}^{}{}_{M_+}{}^{}v_{(f)}{}_{M_+}{}^{}+{\displaystyle \frac{1}{2}}m_f\overline{\mathrm{\Psi }}_{(f)}\mathrm{\Gamma }^7\mathrm{\Psi }_{(f)}\right)`$
$``$ $`{\displaystyle \frac{1}{g^2}}{\displaystyle \underset{f=1}{\overset{n_f}{}}}\left({\displaystyle \underset{M_+=5,6,8,9}{}}{\displaystyle \frac{1}{2}}\overline{\mathrm{\Psi }}_{(f)}\mathrm{\Gamma }^{M_+}\mathrm{\Psi }_{(adj)}v_{(f)}{}_{M_+}{}^{}+{\displaystyle \underset{M_+=5,6,8,9}{}}{\displaystyle \frac{1}{2}}v_{(f)}{}_{M_+}{}^{}\overline{\mathrm{\Psi }}_{(adj)}\mathrm{\Gamma }^{M_+}\mathrm{\Psi }_{(f)}\right)`$
$`+`$ $`{\displaystyle \frac{1}{g^2}}{\displaystyle \underset{f=1}{\overset{n_f}{}}}{\displaystyle \frac{1}{4}}({\displaystyle \underset{M_+=5,6,8,9}{}}v_{(f)}^2{}_{M_+}{}^{})^2.`$
Here $`v_\mu `$ is the five dimensional gauge field, and $`v_7`$, $`v_{M_+}`$, and $`v_{(f)}`$ are respectively $`USp(2k)`$ adjoint, antisymmetric, and fundamental scalars. $`\mathrm{\Psi }_{(adj)}`$, $`\mathrm{\Psi }_{(asym)}`$, and $`\mathrm{\Psi }_{(f)}`$ are respectively $`USp(2k)`$ adjoint, antisymmetric, and fundamental fermions. These fermions can be represented, using thirty-two component Majorana-Weyl spinors in 9+1 dimensions, which satisfy $`C\overline{\mathrm{\Psi }}^t=\mathrm{\Psi }`$, $`\mathrm{\Gamma }_{11}\mathrm{\Psi }=\mathrm{\Psi }`$,
$$\gamma \mathrm{\Psi }_{(adj)}=\mathrm{\Psi }_{(adj)},\gamma \mathrm{\Psi }_{(asym)}=\mathrm{\Psi }_{(asym)},\gamma \mathrm{\Psi }_{(f)}=\mathrm{\Psi }_{(f)},$$
(2. 5)
where $`\gamma \mathrm{\Gamma }^5\mathrm{\Gamma }^6\mathrm{\Gamma }^8\mathrm{\Gamma }^9`$.
Let us pause for a moment to discuss that this five dimensional lagrangian can be understood from the action of type IIB matrix model , followed by the $`USp`$ projections .
$$S(\underset{¯}{v}_M,\underset{¯}{\mathrm{\Psi }})=\frac{1}{g^2}\mathrm{Tr}\left(\frac{1}{4}[\underset{¯}{v}_M,\underset{¯}{v}_N][\underset{¯}{v}^M,\underset{¯}{v}^N]\frac{1}{2}\overline{\underset{¯}{\mathrm{\Psi }}}\mathrm{\Gamma }^M[\underset{¯}{v}_M,\underset{¯}{\mathrm{\Psi }}]\right),$$
(2. 6)
where $`\underset{¯}{v}_M`$ are bosonic coordinates, and $`\underset{¯}{\mathrm{\Psi }}`$ is a thirty-two component Majorana-Weyl spinor, which satisfies $`C\overline{\underset{¯}{\mathrm{\Psi }}}^t=\underset{¯}{\mathrm{\Psi }}`$, $`\mathrm{\Gamma }_{11}\underset{¯}{\mathrm{\Psi }}=\underset{¯}{\mathrm{\Psi }}`$. All underlined fields are in $`u(2k)`$-valued. We can obtain the action of $`USp(2k)`$ matrix model by introducing the projectors $`\widehat{\rho }_b`$, $`\widehat{\rho }_f`$ which act on underlined fields,
$$S(\widehat{\rho }_b\underset{¯}{v}_M,\widehat{\rho }_f\underset{¯}{\mathrm{\Psi }})+\mathrm{\Delta }S,$$
(2. 7)
where
$`\widehat{\rho }_b\underset{¯}{v}_M_{}`$ $`=`$ $`v_M_{}={\displaystyle \underset{a}{}}v_M_{}^aT_a,`$ (2. 8)
$`\widehat{\rho }_b\underset{¯}{v}_{M_+}`$ $`=`$ $`v_{M_+}={\displaystyle \underset{a}{}}v_{M_+}^aX_a,`$ (2. 9)
$`\widehat{\rho }_f\underset{¯}{\mathrm{\Psi }}`$ $`=`$ $`{\displaystyle \underset{a}{}}{\displaystyle \frac{1}{2}}(1+\gamma )\mathrm{\Psi }^aT_a+{\displaystyle \underset{a}{}}{\displaystyle \frac{1}{2}}(1\gamma )\mathrm{\Psi }^aX_a`$ (2. 10)
$`=`$ $`{\displaystyle \underset{a}{}}\mathrm{\Psi }_{(adj)}^aT_a+{\displaystyle \underset{a}{}}\mathrm{\Psi }_{(asym)}^aX_a.`$
Here $`T_a`$ and $`X_a`$ are, respectively, adjoint and antisymmetric representation matrices of $`USp(2k)`$. $`M_{}=0,1,2,3,4,7`$, $`M_+=5,6,8,9`$. We find that $`S(\widehat{\rho }_b\underset{¯}{v}_M,\widehat{\rho }_f\underset{¯}{\mathrm{\Psi }})=S_{adj}+S_{asym}S_{adj+asym}`$ is the reduced action of $`d=4`$, $`𝒩=2`$ super Yang-Mills with one antisymmetric matter. $`\mathrm{\Delta }S`$ contains $`S_{fund}`$, which is the zero dimensional reduced action of $`n_f`$ $`𝒩=2`$ fundamental matters in $`d=4`$. The part in $`\mathrm{\Delta }S`$ which is not contained in $`S_{fund}`$ is irrelevant to the rest of our discussion. For more detail, see ref. . We obtain the lagrangian of the five dimensional gauge theory via matrix T-dual transformation with respect to $`x^0`$, $`x^1`$, $`\mathrm{}`$, $`x^4`$ directions, or replacement $`iv_\mu `$ with the covariant derivative $`𝒟_\mu =_\mu +iv_\mu `$ for $`\mu =0,1,\mathrm{},4`$.
For the purpose of our calculation, we would like to regard the present five dimensional lagrangian as the reduction of higher dimensional one, when we compute the anomalous interaction . On the Coulomb phase, we can consider $`S_{adj}`$, $`S_{asym}`$ and $`S_{fund}`$ as the reductions of $`d=6`$, $`𝒩=1`$ supersymmetric theories. However, on the new phase where the $`vev`$s of the scalars belonging to the antisymmetric hypermultiplet are also nonvanishing, it is easier to regard $`S_{adj+asym}`$ as the reduction, with projections, of $`d=10`$, $`𝒩=1`$ supersymmetric theory.
### 2.2 Vacuum Solutions
We will compute the anomalous interaction on the new phase in the next section. Here we set all fermionic backgrounds to zero and we find the background configurations of our model. From the equations of motion for bosonic fields,
$`[v_\mu ,v_7]`$ $`=`$ $`0,`$
$`[v_M_{},v_{N_+}]`$ $`=`$ $`0,`$
$`v_{(f)}_{M_+}`$ $`=`$ $`0.`$ (2. 11)
We find
$`v_7`$ $`=`$ $`diag(v_7^{(1)},v_7^{(2)},\mathrm{},v_7^{(k)},v_7^{(1)},v_7^{(2)},\mathrm{},v_7^{(k)}),`$ (2. 12)
$`v_{M_+}`$ $`=`$ $`diag(v_{M_+}^{(1)},v_{M_+}^{(2)},\mathrm{},v_{M_+}^{(k)},v_{M_+}^{(1)},v_{M_+}^{(2)},\mathrm{},v_{M_+}^{(k)}),`$ (2. 13)
and all the fundamental bosonic fields $`v_{(f)}`$ vanish. The gauge field is in Cartan subalgebra of $`USp(2k)`$.
### 2.3 Adjoint and Antisymmetric Representation Matrices
We present all the elements $`T`$ of $`usp(2k)`$ Lie algebra, which satisfy $`T^tF+FT=0`$ and $`T^{}=T`$,
$`T_{0i}`$ $`=`$ $`\sigma _ze_{ii}(i=1,\mathrm{},k),`$ (2. 14)
$`T_{1ij}`$ $`=`$ $`\sigma _x{\displaystyle \frac{1}{\sqrt{2}}}e_{\{ij\}}(1i<jk),`$ (2. 15)
$`T_{2ij}`$ $`=`$ $`\sigma _y{\displaystyle \frac{1}{\sqrt{2}}}e_{\{ij\}}(1i<jk),`$ (2. 16)
$`T_{3ij}`$ $`=`$ $`\sigma _z{\displaystyle \frac{1}{\sqrt{2}}}e_{\{ij\}}(1i<jk),`$ (2. 17)
$`T_{4ij}`$ $`=`$ $`\mathrm{𝟏}_2{\displaystyle \frac{i}{\sqrt{2}}}e_{[ij]}(1i<jk),`$ (2. 18)
$`T_{5ij}`$ $`=`$ $`\sigma _xe_{ii}(i=1,\mathrm{},k),`$ (2. 19)
$`T_{6ij}`$ $`=`$ $`\sigma _ye_{ii}(i=1,\mathrm{},k),`$ (2. 20)
where $`\sigma _x`$, $`\sigma _y`$, and $`\sigma _z`$ are Pauli matrices and $`\mathrm{𝟏}_2`$ is the unit matrix of size 2. $`e_{ij}`$ is a $`k\times k`$ matrix such that the element $`(e_{ij})_{kl}=\delta _{ik}\delta _{jl}`$, and $`e_{\{ij\}}e_{ij}+e_{ji}`$, $`e_{[ij]}e_{ij}e_{ji}`$. $`T_{0i}`$ is Cartan subalgebra of $`usp(2k)`$. We define
$`H_{e^i}`$ $`=`$ $`T_{0i}(i=1,\mathrm{},k),`$ (2. 21)
$`T_{\pm (e_i+e_j)}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(T_{1ij}\pm iT_{2ij})(1i<jk),`$ (2. 22)
$`T_{\pm (e_ie_j)}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(T_{3ij}\pm iT_{4ij})(1i<jk),`$ (2. 23)
$`T_{\pm 2e_i}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(T_{5ij}\pm iT_{6ij})(i=1,\mathrm{},k),`$ (2. 24)
where $`e^i`$ are $`k`$ dimensional basis vectors and $`e_j`$ are dual basis vectors, and $`e^ie_j=\delta _j^i`$. The commutation relation of Cartan subalgebra $`H_{e^i}`$ and $`T_𝐰`$ is
$$[H_{e^i},T_𝐰]=e^i𝐰T_𝐰,$$
(2. 25)
where $`𝐰𝐖_{adj}\{\{\pm (e_i+e_j),\pm (e_ie_j),\pm 2e_i\}\}`$ is the root vector of $`USp(2k)`$.
Next, we present all the elements of antisymmetric representation matrices. The element $`X`$ is expressed as
$$X=\left(\begin{array}{cc}A+iC& BiD\\ BiD& AiC\end{array}\right),$$
(2. 26)
where $`A`$ is a real symmetric matrix, and $`B`$, $`C`$, $`D`$ are real skew-symmetric matrices. All the elements of antisymmetric representation matrices are
$`X_{0i}`$ $`=`$ $`\mathrm{𝟏}_2e_{ii}(i=1,\mathrm{},k),`$ (2. 27)
$`X_{1ij}`$ $`=`$ $`\sigma _x{\displaystyle \frac{i}{\sqrt{2}}}e_{[ij]}(1i<jk),`$ (2. 28)
$`X_{2ij}`$ $`=`$ $`\sigma _y{\displaystyle \frac{i}{\sqrt{2}}}e_{[ij]}(1i<jk),`$ (2. 29)
$`X_{3ij}`$ $`=`$ $`\mathrm{𝟏}_2{\displaystyle \frac{1}{\sqrt{2}}}e_{\{ij\}}(1i<jk),`$ (2. 30)
$`X_{4ij}`$ $`=`$ $`\sigma _z{\displaystyle \frac{i}{\sqrt{2}}}e_{[ij]}(1i<jk).`$ (2. 31)
We define
$`\stackrel{~}{H}_{e^i}`$ $`=`$ $`X_{0i}(i=1,\mathrm{},k),`$ (2. 32)
$`X_{\pm (e_i+e_j)}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(X_{1ij}\pm iX_{2ij})(1i<jk),`$ (2. 33)
$`X_{\pm (e_ie_j)}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(X_{3ij}\pm iX_{4ij})(1i<jk).`$ (2. 34)
The commutation relation of Cartan subalgebra $`H_{e^i}`$ and $`X_𝐰`$ is
$$[H_{e^i},X_𝐰]=e^i𝐰X_𝐰,$$
(2. 35)
where $`𝐰𝐖_{asym}\{\{\pm (e_i+e_j),\pm (e_ie_j)\}\}`$ is the weight vector of antisymmetric representation. Diagonal matrices $`\stackrel{~}{H}_{e^i}`$ commute with $`H_{e^i}`$.
The commutation relation of a diagonal matrix $`\stackrel{~}{H}_{e_i}`$ and $`T_𝐰`$ is
$$[\stackrel{~}{H}_{e^i},T_𝐰]=e^i\stackrel{~}{𝐰}X_𝐰,$$
(2. 36)
where
$`\stackrel{~}{𝐰}=\pm (e_ie_j)`$ $`\mathrm{for}𝐰=\pm (e_ie_j),`$
$`\stackrel{~}{𝐰}=\pm (e_ie_j)`$ $`\mathrm{for}𝐰=\pm (e_i+e_j),`$
$`\stackrel{~}{𝐰}=0`$ $`\mathrm{for}𝐰=\pm 2e_i.`$ (2. 37)
Similarly, the commutation relation of a diagonal matrix $`\stackrel{~}{H}_{e_i}`$ and $`X_𝐰`$ is
$$[\stackrel{~}{H}_{e^i},X_𝐰]=e^i\stackrel{~}{𝐰}T_𝐰.$$
(2. 38)
In terms of $`H_{e_i}`$ and $`\stackrel{~}{H}_{e_i}`$, we can express the vacuum solutions as
$$v_7=\underset{i=1}{\overset{k}{}}v_7^{(i)}H_{e_i}𝐯_7𝐇,$$
(2. 39)
$$v_{M_+}=\underset{i=1}{\overset{k}{}}v_{M_+}^{(i)}\stackrel{~}{H}_{e_i}𝐯_{M_+}\stackrel{~}{𝐇}.$$
(2. 40)
## 3 Anomalous Interactions of Five Dimensional $`USp(2k)`$ <br>Gauge Theory
### 3.1 Computation of the Anomalous Interactions
We compute the anomalous interactions on the new phase where the $`vev`$s of the scalars belonging to the antisymmetric hypermultiplet are also nonvanishing.
Firstly, we compute the contribution to the one-loop effective action by the adjoint and antisymmetric fermions,
$$\mathrm{\Gamma }_{1\mathrm{loop}}^{adj+asym}=\frac{i}{2}\mathrm{Tr}\frac{1+\mathrm{\Gamma }_{11}}{2}\widehat{\rho }_f\mathrm{ln}D\text{/},$$
(3. 1)
where $`D\text{/}\mathrm{\Gamma }^\mu (_\mu +iv_\mu )+\mathrm{\Gamma }^7iv_7+\mathrm{\Gamma }^{M_+}iv_{M_+}`$.
Under the variation of the gauge field, eq. (3. 1) is
$`\delta \mathrm{\Gamma }_{1\mathrm{loop}}^{adj+asym}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle d^5x\mathrm{Tr}\frac{1+\mathrm{\Gamma }_{11}}{2}\widehat{\rho }_f\mathrm{\Gamma }^\mu 𝐰\delta 𝐯_\mu x|\frac{1}{D\text{/}}|x}`$ (3. 2)
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle d^5x\mathrm{Tr}\frac{1+\mathrm{\Gamma }_{11}}{2}\widehat{\rho }_f\mathrm{\Gamma }^\mu 𝐰\delta 𝐯_\mu x|\frac{D\text{/}}{D\text{/}^2}|x}.`$
From $`\mathrm{\Gamma }^M\mathrm{\Gamma }^N=\frac{1}{2}\{\mathrm{\Gamma }^M,\mathrm{\Gamma }^N\}+\frac{1}{2}[\mathrm{\Gamma }^M,\mathrm{\Gamma }^N]=\eta ^{MN}+\frac{1}{2}[\mathrm{\Gamma }^M,\mathrm{\Gamma }^N]`$, we obtain
$`D\text{/}^2`$ $`=`$ $`D^\mu D_\mu +(𝐰𝐯_7)^2+{\displaystyle \underset{M_+}{}}(\stackrel{~}{𝐰}𝐯_{M_+})^2+i\mathrm{\Gamma }^\mu \mathrm{\Gamma }^\nu 𝐰_\mu 𝐯_\nu `$ (3. 3)
$`+`$ $`i\mathrm{\Gamma }^\mu \mathrm{\Gamma }^7𝐰_\mu 𝐯_7+i\mathrm{\Gamma }^\mu \mathrm{\Gamma }^{M_+}\stackrel{~}{𝐰}_\mu 𝐯_{M_+}`$
We note that all $`v_M`$ are diagonal matrices, and $`[v_M,v_N]`$ vanish.
We want the terms which is proportional to the epsilon symbol and does not involve the metric, so we keep the contribution to the imaginary part of the one-loop effective action in the Euclidean formalism.
$`\delta \mathrm{\Gamma }_{1\mathrm{loop}}^{adj+asym}`$
$``$ $`{\displaystyle \frac{i}{2}}{\displaystyle d^5x\mathrm{Tr}\frac{1+\mathrm{\Gamma }_{11}}{2}\widehat{\rho }_f\mathrm{\Gamma }^\mu 𝐰\delta 𝐯_\mu (\mathrm{\Gamma }^7𝐰𝐯_7+\mathrm{\Gamma }^{M_+}\stackrel{~}{𝐰}𝐯_{M_+})x|\frac{1}{D\text{/}^2}|x}`$
$``$ $`{\displaystyle \frac{i}{2}}{\displaystyle d^5x\mathrm{Tr}\frac{1+\mathrm{\Gamma }_{11}}{2}\widehat{\rho }_f\mathrm{\Gamma }^\mu 𝐰\delta 𝐯_\mu (\mathrm{\Gamma }^7𝐰𝐯_7+\mathrm{\Gamma }^{M_+}\stackrel{~}{𝐰}𝐯_{M_+})}`$
$`\times x|{\displaystyle \frac{1}{_\varphi ^2+i\mathrm{\Gamma }^\mu \mathrm{\Gamma }^\nu 𝐰_\mu 𝐯_\nu +i\mathrm{\Gamma }^\mu \mathrm{\Gamma }^7𝐰_\mu 𝐯_7+i\mathrm{\Gamma }^\mu \mathrm{\Gamma }^{M_+}\stackrel{~}{𝐰}_\mu 𝐯_{M_+}}}|x`$
$``$ $`{\displaystyle \frac{i}{2}}{\displaystyle \underset{𝐫=adj,asym}{}}{\displaystyle \underset{𝐰𝐖_𝐫}{}}{\displaystyle d^5x\mathrm{Tr}\frac{1+\mathrm{\Gamma }_{11}}{2}\frac{1+()^{|𝐫|}\gamma }{2}\mathrm{\Gamma }^\mu 𝐰\delta 𝐯_\mu (\mathrm{\Gamma }^7𝐰𝐯_7+\mathrm{\Gamma }^{M_+}\stackrel{~}{𝐰}𝐯_{M_+})}`$
$`\times {\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}x|{\displaystyle \frac{(i)^n(\mathrm{\Gamma }^\nu \mathrm{\Gamma }^\lambda 𝐰_\nu 𝐯_\lambda +\mathrm{\Gamma }^\nu \mathrm{\Gamma }^7𝐰_\nu 𝐯_7+\mathrm{\Gamma }^\nu \mathrm{\Gamma }^{M_+}\stackrel{~}{𝐰}_\nu 𝐯_{M_+})^n}{(_\varphi ^2)^{n+1}}}|x`$
$``$ $`{\displaystyle \frac{i}{8}}{\displaystyle \underset{𝐫=adj,asym}{}}()^{|𝐫|}{\displaystyle \underset{𝐰𝐖_𝐫}{}}{\displaystyle d^5x\mathrm{Tr}\mathrm{\Gamma }_{11}\gamma \mathrm{\Gamma }^\mu 𝐰\delta 𝐯_\mu \mathrm{\Gamma }^7𝐰𝐯_7x|\frac{(i)^2(\mathrm{\Gamma }^\nu \mathrm{\Gamma }^\lambda 𝐰_\nu 𝐯_\lambda )^2}{(_\varphi ^2)^3}|x}`$
$`+`$ $`{\displaystyle \frac{i}{8}}{\displaystyle \underset{𝐫=adj,asym}{}}{\displaystyle \underset{𝐰𝐖_𝐫}{}}{\displaystyle d^5x\mathrm{Tr}\mathrm{\Gamma }_{11}\mathrm{\Gamma }^\mu 𝐰\delta 𝐯_\mu \mathrm{\Gamma }^7𝐰𝐯_7x|\frac{(i)^4(\mathrm{\Gamma }^\nu \mathrm{\Gamma }^{M_+}\stackrel{~}{𝐰}_\nu 𝐯_{M_+})^4}{(_\varphi ^2)^5}|x}`$
$`+`$ $`{\displaystyle \frac{i}{8}}{\displaystyle \underset{𝐫=adj,asym}{}}{\displaystyle \underset{𝐰𝐖_𝐫}{}}{\displaystyle d^5x\mathrm{Tr}\mathrm{\Gamma }_{11}\mathrm{\Gamma }^\mu 𝐰\delta 𝐯_\mu \mathrm{\Gamma }^{M_+}\stackrel{~}{𝐰}𝐯_{M_+}}`$
$`\times x|{\displaystyle \frac{4(i)^4(\mathrm{\Gamma }^\nu \mathrm{\Gamma }^7𝐰_\nu 𝐯_7)(\mathrm{\Gamma }^\lambda \mathrm{\Gamma }^{M_+}\stackrel{~}{𝐰}_\lambda 𝐯_{M_+})^3}{(_\varphi ^2)^5}}|x`$
$`=`$ $`4i{\displaystyle \underset{𝐫=adj,asym}{}}()^{|𝐫|}{\displaystyle \underset{𝐰𝐖_𝐫}{}}{\displaystyle d^5x𝐰𝐯_7ϵ^{\mu \nu \lambda \rho \sigma }𝐰\delta 𝐯_\mu (𝐰_\nu )(𝐯_\lambda 𝐰_\rho 𝐯_\sigma )x|\frac{1}{(_\varphi ^2)^3}|x}`$
$`+`$ $`4i{\displaystyle \underset{𝐫=adj,asym}{}}{\displaystyle \underset{𝐰𝐖_𝐫}{}}{\displaystyle d^5x𝐰𝐯_7ϵ^{\mu \nu \lambda \rho \sigma }𝐰\delta 𝐯_\mu (\stackrel{~}{𝐰}_\nu 𝐯_{M_+})(\stackrel{~}{𝐰}_\lambda 𝐯_{N_+})(\stackrel{~}{𝐰}_\rho 𝐯_{P_+})(\stackrel{~}{𝐰}_\sigma 𝐯_{Q_+})}`$
$`\times ϵ^{M_+N_+P_+Q_+}x|{\displaystyle \frac{1}{(_\varphi ^2)^5}}|x`$
$`+`$ $`16i{\displaystyle \underset{𝐫=adj,asym}{}}{\displaystyle \underset{𝐰𝐖_𝐫}{}}{\displaystyle d^5xϵ^{\mu \nu \lambda \rho \sigma }𝐰\delta 𝐯_\mu 𝐰_\nu 𝐯_7(\stackrel{~}{𝐰}_\lambda 𝐯_{N_+})(\stackrel{~}{𝐰}_\rho 𝐯_{P_+})(\stackrel{~}{𝐰}_\sigma 𝐯_{Q_+})(\stackrel{~}{𝐰}𝐯_{M_+})}`$
$`\times ϵ^{M_+N_+P_+Q_+}x|{\displaystyle \frac{1}{(_\varphi ^2)^5}}|x,`$
where
$$_\varphi ^2_\mu ^\mu +(𝐰𝐯_7)^2+\underset{M_+}{}(\stackrel{~}{𝐰}𝐯_{M_+})^2,$$
(3. 5)
and $`()^{|𝐫|}=1`$ for $`𝐫=adj`$, $`()^{|𝐫|}=1`$ for $`𝐫=asym`$. The value of $`x|(_\varphi ^2)^n|x`$ is given by
$$x|\frac{1}{(_\varphi ^2)^n}|x=\frac{i}{(2\sqrt{\pi })^5}\frac{\mathrm{\Gamma }(n\frac{5}{2})}{\mathrm{\Gamma }(n)}\left[(𝐰𝐯_7)^2+\underset{M_+}{}(\stackrel{~}{𝐰}𝐯_{M_+})^2\right]^{\frac{5}{2}n}.$$
(3. 6)
We substitute this equation into eq. (LABEL:eqn:a5),
$`\delta \mathrm{\Gamma }_{1\mathrm{loop}}^{adj+asym}`$
$`=`$ $`{\displaystyle \frac{1}{16\pi ^2}}{\displaystyle \underset{𝐫=adj,asym}{}}()^{|𝐫|}{\displaystyle \underset{𝐰𝐖_𝐫}{}}{\displaystyle d^5x\left[(𝐰𝐯_7)^2+\underset{M_+}{}(\stackrel{~}{𝐰}𝐯_{M_+})^2\right]^{\frac{1}{2}}𝐰𝐯_7}`$
$`\times ϵ^{\mu \nu \lambda \rho \sigma }𝐰\delta 𝐯_\mu (𝐰_\nu 𝐯_\lambda )(𝐰_\rho 𝐯_\sigma )`$
$`+`$ $`{\displaystyle \frac{1}{256\pi ^2}}{\displaystyle \underset{𝐫=adj,asym}{}}{\displaystyle \underset{𝐰𝐖_𝐫}{}}{\displaystyle d^5x\left[(𝐰𝐯_7)^2+\underset{M_+}{}(\stackrel{~}{𝐰}𝐯_{M_+})^2\right]^{\frac{5}{2}}𝐰𝐯_7}`$
$`\times ϵ^{\mu \nu \lambda \rho \sigma }𝐰\delta 𝐯_\mu \stackrel{~}{𝐰}_\nu 𝐯_{M_+}(\stackrel{~}{𝐰}_\lambda 𝐯_{N_+})(\stackrel{~}{𝐰}_\rho 𝐯_{P_+})(\stackrel{~}{𝐰}_\sigma 𝐯_{Q_+})ϵ^{M_+N_+P_+Q_+}`$
$`+`$ $`{\displaystyle \frac{1}{64\pi ^2}}{\displaystyle \underset{𝐫=adj,asym}{}}{\displaystyle \underset{𝐰𝐖_𝐫}{}}{\displaystyle d^5x\left[(𝐰𝐯_7)^2+\underset{M_+}{}(\stackrel{~}{𝐰}𝐯_{M_+})^2\right]^{\frac{5}{2}}\stackrel{~}{𝐰}𝐯_{M_+}}`$
$`\times ϵ^{\mu \nu \lambda \rho \sigma }𝐰\delta 𝐯_\mu (𝐰_\nu 𝐯_7)(\stackrel{~}{𝐰}_\lambda 𝐯_{N_+})(\stackrel{~}{𝐰}_\rho 𝐯_{P_+})(\stackrel{~}{𝐰}_\sigma 𝐯_{Q_+})ϵ^{M_+N_+P_+Q_+}.`$
We can simplify the first term in eq. (LABEL:eqn:a3),
$`\delta \mathrm{\Gamma }_{1\mathrm{loop}}^{adj+asym}`$
$`=`$ $`{\displaystyle \frac{1}{16\pi ^2}}{\displaystyle \underset{𝐰\{\{\pm 2e_i\}\}}{}}{\displaystyle d^5x\left[(𝐰𝐯_7)^2\right]^{\frac{1}{2}}𝐰𝐯_7ϵ^{\mu \nu \lambda \rho \sigma }𝐰\delta 𝐯_\mu (𝐰_\nu 𝐯_\lambda )(𝐰_\rho 𝐯_\sigma )}+\mathrm{}`$
$`=`$ $`{\displaystyle \frac{1}{16\pi ^2}}{\displaystyle \underset{𝐰\{\{\pm 2e_i\}\}}{}}{\displaystyle d^5xsgn(𝐰𝐯_7)ϵ^{\mu \nu \lambda \rho \sigma }𝐰\delta 𝐯_\mu (𝐰_\nu 𝐯_\lambda )(𝐰_\rho 𝐯_\sigma )}+\mathrm{}.`$
Finally, we obtain
$`\mathrm{\Gamma }_{1\mathrm{loop}}^{adj+asym}`$
$`=`$ $`{\displaystyle \frac{1}{48\pi ^2}}{\displaystyle \underset{𝐰\{\{\pm 2e_i\}\}}{}}{\displaystyle d^5xsgn(𝐰𝐯_7)ϵ^{\mu \nu \lambda \rho \sigma }𝐰𝐯_\mu (𝐰_\nu 𝐯_\lambda )(𝐰_\rho 𝐯_\sigma )}`$
$`+`$ $`{\displaystyle \frac{1}{256\pi ^2}}{\displaystyle \underset{𝐫=adj,asym}{}}{\displaystyle \underset{𝐰𝐖_𝐫}{}}{\displaystyle d^5x\left[(𝐰𝐯_7)^2+\underset{M_+}{}(\stackrel{~}{𝐰}𝐯_{M_+})^2\right]^{\frac{5}{2}}𝐰𝐯_7}`$
$`\times ϵ^{\mu \nu \lambda \rho \sigma }𝐰𝐯_\mu (\stackrel{~}{𝐰}_\nu 𝐯_{M_+})(\stackrel{~}{𝐰}_\lambda 𝐯_{N_+})(\stackrel{~}{𝐰}_\rho 𝐯_{P_+})(\stackrel{~}{𝐰}_\sigma 𝐯_{Q_+})ϵ^{M_+N_+P_+Q_+}`$
$`+`$ $`{\displaystyle \frac{1}{64\pi ^2}}{\displaystyle \underset{𝐫=adj,asym}{}}{\displaystyle \underset{𝐰𝐖_𝐫}{}}{\displaystyle d^5x\left[(𝐰𝐯_7)^2+\underset{M_+}{}(\stackrel{~}{𝐰}𝐯_{M_+})^2\right]^{\frac{5}{2}}\stackrel{~}{𝐰}𝐯_{M_+}}`$
$`\times ϵ^{\mu \nu \lambda \rho \sigma }𝐰𝐯_\mu (𝐰_\nu 𝐯_7)(\stackrel{~}{𝐰}_\lambda 𝐯_{N_+})(\stackrel{~}{𝐰}_\rho 𝐯_{P_+})(\stackrel{~}{𝐰}_\sigma 𝐯_{Q_+})ϵ^{M_+N_+P_+Q_+}.`$
Similarly, We compute the contribution to the one-loop effective action by the fundamental fermions,
$$\mathrm{\Gamma }_{1\mathrm{loop}}^{fund}=\frac{i}{2}\underset{f=1}{\overset{n_f}{}}\frac{1+\mathrm{\Gamma }_{11}}{2}\frac{1+()^{|fund|}\gamma }{2}\mathrm{ln}D\text{/}_{(f)},$$
(3. 10)
where $`D\text{/}_{(f)}=\mathrm{\Gamma }^\mu (_\mu +iv_\mu )+\mathrm{\Gamma }^7i(v_7+m_f)+\mathrm{\Gamma }^{M_+}iv_{(f)}_{M_+}`$ and $`()^{|fund|}=1`$. We take the variation with respect to the gauge field and pick up the relevant terms,
$`\delta \mathrm{\Gamma }_{1\mathrm{loop}}^{fund}`$
$`=`$ $`{\displaystyle \frac{1}{16\pi ^2}}()^{|fund|}{\displaystyle \underset{f=1}{\overset{n_f}{}}}{\displaystyle \underset{𝐰𝐖_{fund}}{}}{\displaystyle d^5xsgn(𝐰𝐯_7+m_f)ϵ^{\mu \nu \lambda \rho \sigma }𝐰\delta 𝐯_\mu (𝐰_\nu 𝐯_\lambda )(𝐰_\rho 𝐯_\sigma )},`$
where $`𝐖_{fund}\{\{\pm e_i\}\}`$.
Finally, we obtain
$`\mathrm{\Gamma }_{1\mathrm{loop}}^{fund}`$
$`=`$ $`{\displaystyle \frac{1}{48\pi ^2}}()^{|fund|}{\displaystyle \underset{f=1}{\overset{n_f}{}}}{\displaystyle \underset{𝐰𝐖_{fund}}{}}{\displaystyle d^5xsgn(𝐰𝐯_7+m_f)ϵ^{\mu \nu \lambda \rho \sigma }𝐰𝐯_\mu (𝐰_\nu 𝐯_\lambda )(𝐰_\rho 𝐯_\sigma )}.`$
### 3.2 Summary of Our Results
We have exhibited the anomalous interaction on the new phase,
$`\mathrm{\Gamma }_{1\mathrm{loop}}`$ $`=`$ $`\mathrm{\Gamma }_{1\mathrm{loop}}^{adj+asym}+\mathrm{\Gamma }_{1\mathrm{loop}}^{fund},`$
$`\mathrm{\Gamma }_{1\mathrm{loop}}^{adj+asym}`$ $`=`$ $`{\displaystyle \frac{i}{2}}\mathrm{Tr}{\displaystyle \frac{1+\mathrm{\Gamma }_{11}}{2}}\widehat{\rho }_f\mathrm{ln}D\text{/},`$
$`\mathrm{\Gamma }_{1\mathrm{loop}}^{fund}`$ $`=`$ $`{\displaystyle \frac{i}{2}}{\displaystyle \underset{f=1}{\overset{n_f}{}}}{\displaystyle \frac{1+\mathrm{\Gamma }_{11}}{2}}{\displaystyle \frac{1\gamma }{2}}\mathrm{ln}D\text{/}_{(f)},`$ (3. 13)
where
$`D\text{/}`$ $`=`$ $`\mathrm{\Gamma }^\mu (_\mu +iv_\mu )+\mathrm{\Gamma }^7iv_7+\mathrm{\Gamma }^{M_+}iv_{M_+},`$
$`D\text{/}_{(f)}`$ $`=`$ $`\mathrm{\Gamma }^\mu (_\mu +iv_\mu )+\mathrm{\Gamma }^7i(v_7+m_f).`$ (3. 14)
We have obtained
$`\mathrm{\Gamma }_{1\mathrm{loop}}`$
$`=`$ $`{\displaystyle \frac{1}{48\pi ^2}}{\displaystyle \underset{𝐰\{\{\pm 2e_i\}\}}{}}{\displaystyle d^5xsgn(𝐰𝐯_7)ϵ^{\mu \nu \lambda \rho \sigma }𝐰𝐯_\mu (𝐰_\nu 𝐯_\lambda )(𝐰_\rho 𝐯_\sigma )}`$
$`+`$ $`{\displaystyle \frac{1}{48\pi ^2}}{\displaystyle \underset{f=1}{\overset{n_f}{}}}{\displaystyle \underset{𝐰𝐖_{fund}}{}}{\displaystyle d^5xsgn(𝐰𝐯_7+m_f)ϵ^{\mu \nu \lambda \rho \sigma }𝐰𝐯_\mu (𝐰_\nu 𝐯_\lambda )(𝐰_\rho 𝐯_\sigma )}`$
$`+`$ $`{\displaystyle \frac{1}{256\pi ^2}}{\displaystyle \underset{𝐫=adj,asym}{}}{\displaystyle \underset{𝐰𝐖_𝐫}{}}{\displaystyle d^5x\left[(𝐰𝐯_7)^2+\underset{M_+}{}(\stackrel{~}{𝐰}𝐯_{M_+})^2\right]^{\frac{5}{2}}𝐰𝐯_7}`$
$`\times ϵ^{\mu \nu \lambda \rho \sigma }𝐰𝐯_\mu (\stackrel{~}{𝐰}_\nu 𝐯_{M_+})(\stackrel{~}{𝐰}_\lambda 𝐯_{N_+})(\stackrel{~}{𝐰}_\rho 𝐯_{P_+})(\stackrel{~}{𝐰}_\sigma 𝐯_{Q_+})ϵ^{M_+N_+P_+Q_+}`$
$`+`$ $`{\displaystyle \frac{1}{64\pi ^2}}{\displaystyle \underset{𝐫=adj,asym}{}}{\displaystyle \underset{𝐰𝐖_𝐫}{}}{\displaystyle d^5x\left[(𝐰𝐯_7)^2+\underset{M_+}{}(\stackrel{~}{𝐰}𝐯_{M_+})^2\right]^{\frac{5}{2}}\stackrel{~}{𝐰}𝐯_{M_+}}`$
$`\times ϵ^{\mu \nu \lambda \rho \sigma }𝐰𝐯_\mu (𝐰_\nu 𝐯_7)(\stackrel{~}{𝐰}_\lambda 𝐯_{N_+})(\stackrel{~}{𝐰}_\rho 𝐯_{P_+})(\stackrel{~}{𝐰}_\sigma 𝐯_{Q_+})ϵ^{M_+N_+P_+Q_+},`$
(3. 15)
where $`𝐖_{fund}\{\{\pm e_i\}\}`$.
The first and the second terms have been computed in . The third and the fourth terms are the anomalous interactions we have found. These interactions represent a generalized Lorentz force among D4-branes in the multiprobe picture .
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# 1 Introduction
## 1 Introduction
Temperley-Lieb (TL) algebra has been widely used in the construction of solutions of the Yang-Baxter equation , which is a sufficient condition for integrability of lattice models through the quantum inverse scattering method. On the other hand, quantum group in some cases is a symmetry of the integrable model . Thus, quantum group together with Temperley-Lieb algebra provide a powerful algebraic framework to build and study various kinds of lattice models in two-dimensional statistical mechanics. One particular way of building models which are quantum group invariant uses the TL algebra. It is a algebra generated by the Hamiltonian density $`U_k`$ , $`k=1,2,\mathrm{},N1`$ subject to the following constraints
$$\begin{array}{ccc}U_k^2=(Q+Q^1)U_k\hfill & ,\hfill & U_kU_{k\pm 1}U_k=U_k\hfill \\ & & \\ U_kU_j=U_jU_k\hfill & & |kj|>1.\hfill \end{array}$$
(1.1)
with $`QC`$ a given number. This algebra appears in a large class of quantum lattice models and leads, at level of free and ground state energies, to some equivalence among the models .
Taking into account usual boundary conditions, the TL Hamiltonians take the form
$$=\underset{k=1}{\overset{N1}{}}U_{kk+1}+\mathrm{𝐛𝐭}$$
(1.2)
where $`U_{kk+1}U_k`$ operates in a direct product of complex spaces at positions $`k`$ and $`k+1`$. In general, they are not invariant with respect to quantum groups since the boundary terms, $`\mathrm{𝐛𝐭}`$, break translational invariance, reflecting the non-cocommutativity of the co-product. Indeed, we know from that very special boundary terms must be considered when we seek these quantum group invariant spin chains. In particular, one possibility is to consider free boundary conditions, i.e., $`\mathrm{𝐛𝐭}=0`$ . For The XXZ-Hamiltonian with free boundary conditions one has to apply the Bethe ansatz techniques introduced first by Alcaraz et al and after by Sklyanin using Cherednik’s reflection matrices . By this method the XXZ-Heisenberg model , the $`spl_q(2,1)`$ invariant supersymmetric t-J model , the $`U_q[sl(n)]`$ invariant generalization of the XXZ-chain and the $`SU_q(n|m)`$ spin chains have been solved for free boundary conditions.
Recently, by means of a generalized algebraic nested Bethe ansatz, Karowski and Zapletal presented a class of quantum group invariant $`n`$-state vertex models with periodic boundary conditions. Also an extension of this method to the case of graded vertex models was analyzed in , where a $`spl_q(2|1)`$ invariant susy $`t`$-$`J`$ model with closed boundary conditions was presented.
In fact, these kind of models were first discussed by Martin from the representations of the Hecke algebra. In this case, the boundary term is a non-local operator, $`\mathrm{𝐛𝐭}=𝒰_0`$ (see Section $`4`$). The study of closed quantum group invariant closed spin chains in the framework of the coordinate Bethe ansatz was presented by Grosse et al for the fundamental representation of $`SU_q(2)`$ and generalized by . In this context it would be interesting to discuss other quantum group invariant closed spin chains.
In the present paper we find the spectrum of Hamiltonians based on representations of the TL algebra associated with quantum groups. The Bethe ansatz equations for different types of boundary conditions (periodic, closed and free) are obtained through a modified version of the coordinate Bethe ansatz. Therefore we generalize the results of ref..Our method was recently applied to solve graded T-L Hamiltonian and anisotropic correlated electron model associated with this algebra .
The paper is organized as follows. In Section $`2`$, we describe the representations of the TL algebra, constructed as projectors on total spin zero of two neighboring spins. In Section $`3`$, we introduce the modified coordinate Bethe Ansatz for the TL Hamiltonians with periodic boundary conditions. In Section $`4`$ their Bethe Ansatz solution is presented with non-local boundary conditions. In section $`5`$, contains the solution for free boundary conditions. Finally the conclusions are reserved for section $`6`$.
## 2 Representations of the TL algebra
Representations of the TL algebra, commuting with quantum groups, can be constructed in the following way . Suppose $`𝒰_q(X_n)`$ is the universal enveloping algebra of a finite dimensional Lie algebra $`X_n`$, equipped with the coproduct $`\mathrm{\Delta }:𝒰_q𝒰_q𝒰_q`$ . If now $`\pi :𝒰_q\mathrm{En}𝖽V_\mathrm{\Lambda }`$ is a finite dimensional irreducible representation with highest weight $`\mathrm{\Lambda }`$ and we assume that the decomposition $`V_\mathrm{\Lambda }V_\mathrm{\Lambda }`$ is multiplicity free and includes one trivial representation on $`V_0`$, then the projector $`𝒫_0`$ from $`V_\mathrm{\Lambda }V_\mathrm{\Lambda }`$ onto $`V_0`$ is a representation of the TL algebra. The deformation parameter $`q`$, which plays the role of a coupling constant in the Hamiltonian, is related to $`Q`$ as:
$$Q+Q^1=\mathrm{Tr}_{V_\mathrm{\Lambda }}(q^{2\rho }),$$
(2.1)
where $`\rho `$ is half the sum of the positive roots.
We will consider the following specific cases, $`(V_\mathrm{\Lambda },𝒰_q(X_n))=(V_{2s\mathrm{\Lambda }_1},𝒰_q(A_1))`$ for spin $`s`$, $`(V_{\mathrm{\Lambda }_1},𝒰_q(B_n)`$,$`(V_{\mathrm{\Lambda }_1},𝒰_q(C_n)`$ and $`(V_{\mathrm{\Lambda }_1},𝒰_q(D_n)`$. Namely we treat the $`q`$-deformations of the spin-$`s`$ representation of $`sl(2)`$ and the vector representations of $`so(2n+1),sp(2n)`$ and $`so(2n)`$. $`V_\mathrm{\Lambda }`$ denotes the $`𝒰_q(X_n)`$ module with highest weight $`\mathrm{\Lambda }`$ and $`\mathrm{\Lambda }_1`$ is a highest weight of $`X_n`$.
Since we are not going to use any group-theoretical machinery, we will just lift the relevant formulas off Batchelor and Kuniba in order to display explicitly the Hamiltonians to be diagonalized.
We introduce orthonormal vectors $`e_\mu ,`$ $`<e_\mu ,e_\nu >=\delta _{\mu \nu }`$, to express $`\rho `$ and the set $`𝒜`$ of weights appearing in the representation $`\pi `$ of $`𝒰_q(X_n)`$ as follows:
* : Spin-$`s`$ representation of the $`sl(2)`$ algebra
$`𝒜=\{se_{12},(s1)e_{12},\mathrm{},se_{12}\},\rho ={\displaystyle \frac{1}{2}}e_{12},(e_{12}=e_1e_2),`$
$`J=\{s,s1,\mathrm{},s\},ϵ(\mu )=(1)^{\stackrel{~}{\mu }},Q+Q^1=[2s+1].`$ (2.2)
* : Vector representation of the $`so(2n+1)`$ algebra $`(n2)`$
$`𝒜=\{0,\pm e_1,\mathrm{},\pm e_n\},\rho ={\displaystyle \underset{\alpha =1}{\overset{n}{}}}(n(\alpha {\displaystyle \frac{1}{2}}))e_\alpha ,J=\{0,\pm 1,\mathrm{},\pm n\},`$
$`ϵ(\mu )=(1)^{\stackrel{~}{\mu }},Q+Q^1={\displaystyle \frac{[2n1][n+1/2]}{[n1/2]}}.`$ (2.3)
* : Vector representation of the $`sp(2n)`$ algebra $`(n1)`$
$`𝒜=\{\pm e_1,\mathrm{},\pm e_n\},\rho ={\displaystyle \underset{\alpha =1}{\overset{n}{}}}(n\alpha +1)e_\alpha ,J=\{\pm 1,\mathrm{},\pm n\},`$
$`ϵ(\mu )=\mathrm{sign}(\mu ),Q+Q^1={\displaystyle \frac{[n][2n+2]}{[n+1]}}.`$ (2.4)
* : Vector representations of the $`so(2n)`$ algebra $`(n3)`$
$`𝒜=\{\pm e_1,\mathrm{},\pm e_n\},\rho ={\displaystyle \underset{\alpha =1}{\overset{n1}{}}}(n\alpha )e_\alpha ,J=\{\pm 1,\mathrm{},\pm n\},`$
$`ϵ(\mu )=1,Q+Q^1={\displaystyle \frac{[2n2][n]}{[n1]}}.`$ (2.5)
For $`\mu J`$ the symbol $`\stackrel{~}{\mu }`$ is defined as $`\stackrel{~}{\mu }=\mu +1/2`$ for $`A_1`$ with $`s`$ semi-integer and $`\stackrel{~}{\mu }=\mu `$ for $`A_1`$ with $`s`$ integer. For $`B_n`$, $`C_n`$ and $`D_n`$, $`\stackrel{~}{\mu }=0`$ with the exception of $`\stackrel{~}{0}=1`$ for $`B_n`$. The $`q`$-number notation is $`[x]=(q^xq^x)/(qq^1)`$.
For $`X_n=B_n,C_n`$ and $`D_n`$, we extend the suffix of $`e_\mu `$ to $`n\mu n`$ by setting $`e_\mu =e_\mu `$ (hence $`e_0=0`$). Using the index set $`J`$ above, we can write $`𝒜=\{\mu (e_1e_2)\}`$ for $`A_1`$ and $`𝒜=\{e_\mu |\mu J)\}`$ for $`B_n,C_n`$ and $`D_n`$.
Denoting by $`E_{\mu \nu }\mathrm{End}V_\mathrm{\Lambda }`$ the matrix unit, having all elements zero, except at row $`\mu `$ and column $`\nu `$, the projector can be written as
$$𝒫_0=\frac{1}{Q+Q^1}\underset{\mu ,\nu J}{}ϵ(\mu )ϵ(\nu )q^{<e_\mu +e_\nu ,\rho >}E_{\mu \nu }E_{\mu \nu }.$$
(2.6)
In the following we will refer to all models generically as higher spin models for simplicity, even when not talking about $`A_1`$. The dimensions of $`B_n`$ are related with the dimensions of $`A_1`$ for $`s`$ integer and the dimensions of $`C_n`$ and $`D_n`$ are related with the dimensions of $`A_1`$ for $`s`$ semi-integer.
If we consider then a one-dimensional chain of length $`N`$ with a ”spin” at each site, the ”spin variables ” range over the set of weight vectors $`𝒜=\left\{v_\mu \right|\mu J\}`$ and our Hilbert space is an $`N`$-fold tensor product $`V_\mathrm{\Lambda }\mathrm{}V_\mathrm{\Lambda }`$. For $`A_1`$, these are the $`q`$-analogs of the usual spin states. The periodic Hamiltonians associated with the TL representations are given by the following sum over $`N`$ sites
$$=\underset{k=1}{\overset{N}{}}U_k.$$
(2.7)
Here $`\mathrm{𝐛𝐭}=U_{N,N+1}=U_{N1}`$ and the Hamiltonian densities acting on two neighboring sites are then given by:
$`\nu ,\lambda |U|\mu ,\kappa `$ $`=`$ $`ϵ(\mu )ϵ(\nu )q^{<e_\mu +e_\nu ,\rho >}\delta _{\mu +\kappa ,0}\delta _{\nu +\lambda ,0}`$
$`\mu ,\nu ,\kappa ,\lambda `$ $``$ $`J`$ (2.8)
Using representations of the TL algebra, one can also build solvable vertex models whose Hamiltonian limit leads to the previously mentioned quantum spin chains. To do so, we introduce an operator $`R(u)\mathrm{End}(V_\mathrm{\Lambda }V_\mathrm{\Lambda })`$ by
$`R_k(u)`$ $`=`$ $`I\mathrm{}I\stackrel{\underset{}{R\left(u\right)}}{k,k+1}I\mathrm{}I`$
$`R(u)`$ $`=`$ $`{\displaystyle \underset{\mu ,\nu ,\kappa ,\lambda J}{}}R_{\mu \kappa }^{\nu \lambda }(u,\eta )E_{\mu \nu }^{(k)}E_{\kappa \lambda }^{(k+1)}`$
$$R_{\mu \kappa }^{\nu \lambda }(u,\eta )=\frac{\mathrm{sinh}(\eta u)}{\mathrm{sinh}\eta }\delta _{\mu \nu }\delta _{\kappa \lambda }+\frac{\mathrm{sinh}u}{\mathrm{sinh}\eta }ϵ(\mu )ϵ(\nu )q^{<e_\mu +e_\nu ,\rho >}\delta _{\mu +\kappa ,0}\delta _{\nu +\lambda ,0}$$
(2.9)
where $`u`$ is the spectral parameter and the anisotropy parameter $`\eta `$ is chosen so that
$$2\mathrm{cosh}\eta =Q+Q^1$$
(2.10)
The $`R`$ matrix commutes with the quantum group action and the Yang-Baxter equation
$$R_k(u)R_{k+1}(u+v)R_k(v)=R_{k+1}(v)R_k(u+v)R_{k+1}(u)$$
(2.11)
is valid owing to the TL relations (1.1).
Due to the TL algebra these models are equivalent to the $`6`$-vertex and the ” $`p`$-state ” self-dual Potts models ( $`\sqrt{p}=lim_{q1}(Q+Q^1)`$) through the argument in . In fact, the cases $`𝒰_q(A_1)`$ with $`s=1/2`$ and $`𝒰_q(C_1)`$ yield the $`6`$-vertex model itself. When $`q=1`$, the vertex models here reduce to those discussed in , where the number of states is equal to $`\sqrt{p}`$.
The case $`A_1`$ have been studied by several authors. When $`q=1`$ and $`s=1/2`$ the operator $`U_k`$ is essentially the Heisenberg interaction term $`\sigma _k^x\sigma _{k+1}^x+\sigma _k^y\sigma _{k+1}^y+\sigma _k^z\sigma _{k+1}^z`$. Besides the explicit matrix elements , available from (2.8), the local Hamiltonian $`U_k`$ is in principle also expressible in terms of the usual representation matrices for $`SU(2)`$ generators. The resulting $`q`$-deformed Hamiltonian has been written down for $`s=1`$ in . However for $`s>1`$, writing the Hamiltonian in terms of the usual $`SU(2)`$ operators becomes very cumbersome. These Hamiltonians have alternative expressions in terms of Casimir operators .
The limit $`q1`$ has discussed for general $`s`$ . As for related spin-$`1`$ biquadratic model , they are massive for $`s1`$ and of relevance to the dimerization transition on $`SU(n)`$ antiferromagnetic chains . The case $`s=1/2`$ has also been investigated in some detail in .
Having now built common ground for all models, whose salient feature is they being spin zero projectors, we may now follow the steps of reference to find their spectra.
## 3 Bethe Ansatz: Periodic boundary conditions
All the above Hamiltonians are $`U(1)`$ invariant and we can classify their spectra according to sectors. For $`A_1`$ the commuting operator is the total spin $`S^z=_{k=1}^NS_k^z`$ where $`S_k^z=`$diag$`(s,s1,\mathrm{},s+1,s)_k`$ and we set the conserved quantum number $`r=sNS^z`$. We extend these quantum numbers for the other algebras as $`r=N\omega 𝐒^z`$ where $`𝐒^z=_k𝐒_k^z`$ with $`𝐒_k^z=`$diag$`(\omega ,\omega 1,\mathrm{},\omega +1,\omega )_k`$ and $`\omega =\mathrm{max}\{J\}`$. Eigenvalues of the operator $`r`$ can be used to collect the eigenstates of $``$ in sectors, $`\mathrm{\Psi }_r`$. Therefore, there exists a reference state $`\mathrm{\Psi }_0`$, satisfying $`\mathrm{\Psi }_0=E_0\mathrm{\Psi }_0`$, with $`E_0=0`$. We take $`\mathrm{\Psi }_0`$ to be $`\mathrm{\Psi }_0=|\omega \omega \omega \mathrm{}\omega `$. It is the only eigenstate in the sector $`r=0`$ and all other energies will be measured relative to this state.
In every sector $`r`$ there are eigenstates degenerate with $`\mathrm{\Psi }_0`$. They contain a set of impurities. We call impurity any state obtained by lowering some of the $`|\omega ,k`$’s, such that the sum of any two neighboring spins is non-zero. Since $``$ is a projector on spin zero, all these states are annihilated by $``$. In particular, they do not move under the action of $``$, which is the reason for their name.
Nothing interesting happens in sectors with $`r<2\omega `$. Sector $`r=2\omega `$, we encounter the situation where the states $`|\alpha ,k`$ and $`|\alpha ,k\pm 1`$, $`\alpha J`$, occur in neighboring pairs. They move under the action of $``$, i.e., the sector $`r=2\omega `$ contains one free pseudoparticle. In general, for a sector $`r`$ we may have $`p`$ pseudoparticles and $`N_{\omega 1},N_{\omega 2},\mathrm{},N_{\omega +1}`$ impurities of the type $`\omega 1,\omega 2,\mathrm{},\omega +1,`$ respectively, such that
$$r=2\omega p+\underset{\alpha =1}{\overset{2\omega 1}{}}\alpha N_{\omega \alpha }.$$
(3.1)
The main result of this section is to show that $``$ can be diagonalized in a convenient basis, constructed from products of single pseudoparticle wavefunctions. The energy eigenvalues will be parametrized as a sum of single pseudoparticle contributions.
The first nontrivial sector is $`r=2\omega `$ and the correspondent eigenspace is spanned by the states $`|k(\alpha ,\alpha )>=|\omega \omega \mathrm{}\omega \stackrel{\alpha }{k}\alpha \omega \mathrm{}\omega >`$ , where $`k=1,2,\mathrm{},N1`$and $`\alpha J.`$ We seek eigenstates of $``$ which are linear combinations of these vectors. It is very convenient to consider the linear combination
$$|\mathrm{\Omega }(k)=\underset{\alpha =\omega }{\overset{\omega }{}}ϵ(\omega )ϵ(\alpha )q^{<e_\omega +e_\alpha ,\rho >}|k(\alpha ,\alpha ),$$
(3.2)
which is an eigenstate of $`U_k`$:
$$U_k|\mathrm{\Omega }(k)=(Q+Q^1)|\mathrm{\Omega }(k).$$
(3.3)
Moreover, the action of $`U_{k\pm 1}`$ on $`|\mathrm{\Omega }(k)`$ is very simple
$$\begin{array}{ccc}U_{k1}|\mathrm{\Omega }(k)=ϵ|\mathrm{\Omega }(k1)\hfill & & U_{k+1}|\mathrm{\Omega }(k)=ϵ|\mathrm{\Omega }(k+1)\hfill \\ & & \\ U_k|\mathrm{\Omega }(m)=0\hfill & & k\{m\pm 1,m\}\hfill \end{array}$$
(3.4)
where $`ϵ=1`$ for $`B_n`$, $`D_n`$ and $`A_1`$ ( $`s`$ integer) and $`ϵ=1`$ for $`C_n`$ and $`A_1`$ ( $`s`$ semi-integer).
It should be emphasized that although the linear combination (3.2) is different for each model, the action of $`U_k`$ is always given by (3.3) and (3.4). Therefore, all Hamiltonians can be treated in a similar way and it affords a considerable simplification in their diagonalizations when we compare with the calculus used in the usual spin basis .
### 3.1 One-pseudoparticle
We will now start to diagonalize $``$ in every sector. Let us consider one free pseudoparticle as a highest weight state which lies in the sector $`r=2\omega `$
$$\mathrm{\Psi }_{2\omega }=\underset{k}{}A(k)|\mathrm{\Omega }(k).$$
(3.5)
Using the eigenvalue equation $``$ $`\mathrm{\Psi }_{2\omega }=E_{2\omega }\mathrm{\Psi }_{2\omega }`$, one can derive a complete set of equations for the wavefunctions $`A(k)`$.
When the bulk of $``$ acts on $`|\mathrm{\Omega }(k)`$ it sees the reference configuration, except in the vicinity of $`k`$ where we use (3.3) and (3.4) to get
$`|\mathrm{\Omega }(k)=(Q+Q^1)|\mathrm{\Omega }(k)+ϵ|\mathrm{\Omega }(k1)+ϵ|\mathrm{\Omega }(k+1)`$
$`2kN2`$ (3.6)
Substituting (3.6) in the eigenvalue equation, we get
$`(E_{2\omega }QQ^1)A(k)=ϵA(k1)+ϵA(k+1)`$
$`2kN2`$ (3.7)
Here we will treat periodic boundary conditions . They demand $`U_{N,N+1}=U_{N,1}`$, implying $`A(k+N)=A(k)`$. This permits us to complete the set of equations (3.7) for $`A(k)`$ by including the equations for $`k=1`$ and $`k=N1`$. Now we parametrize $`A(k)`$ by plane wave $`A(k)=A\xi ^k`$ to get the energy of one free pseudoparticle as:
$`E_{2\omega }=Q+Q^1+ϵ\left(\xi +\xi ^1\right)`$
$`\xi ^N=1`$ (3.8)
Here $`\xi =\mathrm{e}^{i\theta }`$, $`\theta `$ being the momenta determined from the periodic boundary to be $`\theta =2\pi l/N`$, with $`l`$ an integer.
### 3.2 One-pseudoparticle and impurities
Let us consider the state with one pseudoparticle and one impurity of type $`(\omega 1)`$, which lies in the sector $`r=2\omega +1`$. We seek eigenstates in the form
$$\mathrm{\Psi }_{2\omega +1}(\xi _1,\xi _2)=\underset{k_1<k_2}{}\left\{A_1(k_1,k_2)|\mathrm{\Omega }_1(k_1,k_2)+A_2(k_1,k_2)|\mathrm{\Omega }_2(k_1,k_2)\right\}$$
(3.9)
We try to build these eigenstates out of translational invariant products of one pseudoparticle excitation with parameter $`\xi _2`$ and one impurity with parameter $`\xi _1`$:
$$\mathrm{\Psi }_{2\omega +1}(\xi _1,\xi _2)=|(\omega 1)(\xi _1)\times \mathrm{\Psi }_{2\omega }(\xi _2)+\mathrm{\Psi }_{2\omega }(\xi _2)\times |(\omega 1)(\xi _1)$$
(3.10)
Using one-pseudopaticle eigenstate solution (3.5) and comparing this with (3.9) we get
$`|\mathrm{\Omega }_1(k_1,k_2)`$ $`=`$ $`{\displaystyle \underset{\alpha =\omega }{\overset{\omega }{}}}ϵ(\omega )ϵ(\alpha )q^{<e_\omega +e_\alpha ,\rho >}|k_1(\omega 1),k_2(\alpha ,\alpha )`$
$`|\mathrm{\Omega }_2(k_1,k_2)`$ $`=`$ $`{\displaystyle \underset{\alpha =\omega }{\overset{\omega }{}}}ϵ(\omega )ϵ(\alpha )q^{<e_\omega +e_\alpha ,\rho >}|k_1(\alpha ,\alpha ),k_2(\omega 1)`$
and
$$A_1(k_1,k_2)=A_1\xi _1^{k_1}\xi _2^{k_2},A_2(k_1,k_2)=A_2\xi _2^{k_1}\xi _1^{k_2}.$$
(3.12)
Periodic boundary conditions $`A_1(k_2,N+k_1)=A_2(k_1,k_2)`$ and $`A_i(N+k_1,N+k_2)=A_i(k_1,k_2)`$, $`i=1,2`$ imply that
$$A_1\xi _2^N=A_2,\xi ^N=(\xi _1\xi _2)^N=1$$
(3.13)
When $``$ now acts on $`\mathrm{\Psi }_{2\omega +1}`$, we will get a set of coupled equations for $`A_i(k_1,k_2),`$ $`i=1,2`$. We split the equations into far equations, when the pseudoparticle do not meet the impurity and near equations, containing terms when they are neighbors.
Since the impurity is annihilated by $``$, the action of $``$ on (3.9) in the case far (i.e., $`(k_2k_1)3`$), can be write down directly from (3.7) :
$`\left(E_{2\omega +1}QQ^1\right)A_1(k_1,k_2)=ϵA_1(k_1,k_21)+ϵA_1(k_1,k_2+1)`$ (3.14)
$`\left(E_{2\omega +1}QQ^1\right)A_2(k_1,k_2)=ϵA_2(k_11,k_2)+ϵA_2(k_1+1,k_2)`$ (3.15)
Using the parametrization (3.12), these equations will give us the energy eigenvalues
$$E_{2\omega +1}=Q+Q^1+ϵ(\xi _2+\xi _2^1)$$
(3.16)
To find $`\xi _2`$ we must consider the near equations. First, we compute the action of $``$ on the coupled near states $`|\mathrm{\Omega }_1(k,k+1)`$ and $`|\mathrm{\Omega }_2(k,k+2)`$:
$`|\mathrm{\Omega }_1(k,k+1)=(Q+Q^1)|\mathrm{\Omega }_1(k,k+1)+ϵ|\mathrm{\Omega }_1(k,k+2)+ϵ|\mathrm{\Omega }_2(k,k+2)`$
(3.17)
$`|\mathrm{\Omega }_2(k,k+2)=(Q+Q^1)|\mathrm{\Omega }_2(k,k+2)+ϵ|\mathrm{\Omega }_2(k1,k+2)+ϵ|\mathrm{\Omega }_1(k,k+1)`$
(3.18)
The last terms in these equations tell us that a pseudoparticle can propagate past the isolated impurity, but in so doing causes a shift in its position by two lattice site. Substituting (3.17) and (3.18) into the eigenvalue equation, we get
$`\left(E_{2\omega +1}QQ^1\right)A_1(k,k+1)=ϵA_1(k,k+2)+ϵA_2(k,k+2)`$ (3.19)
$`\left(E_{2\omega +1}QQ^1\right)A_2(k,k+2)=ϵA_2(k1,k+2)+ϵA_1(k,k+1)`$ (3.20)
These equations, which are not automatically satisfied by the ansatz (3.12), are equivalent to the conditions
$$A_1(k,k)A_2(k,k+2),A_1(k,k+1)A_2(k+1,k+2).$$
(3.21)
obtained by subtracting Eq. (3.19) from Eq.(3.14) for $`k_1=k`$ , $`k_2=k+1`$ and by subtracting Eq. (3.20) from Eq.(3.15) for $`k_1=k`$ , $`k_2=k+2`$, respectively. The conditions (3.21) require a modification of the amplitude relation (3.13):
$$\frac{A_2}{A_1}=\xi _1^2=\xi _2^N\xi _2^N\xi _1^2=1\mathrm{or}\xi _2^{N2}\xi ^2=1$$
(3.22)
Putting $`\xi _i=\mathrm{e}^{i\theta _i}`$, $`i=1,2`$, it means $`\mathrm{cos}(N2)\theta _2=\mathrm{cos}2\theta _1`$. Hence
$$\theta _2=\frac{2\pi m\pm 4\pi l/N}{N2},l\mathrm{and}m\mathrm{integers}.$$
(3.23)
In other words, $`\mathrm{\Psi }_{2\omega +1}(\xi _1,\xi _2)`$ are eigenstates of $``$ with energy eigenvalues given by $`E_{2\omega +1}=Q+Q^1+2ϵ\mathrm{cos}\theta _2`$. Note that when $`Nm\pm 2l`$ is a multiple of $`(N2)`$ we get states which are degenerate with the one-pseudoparticle states $`\mathrm{\Psi }_{2\omega }`$, which lie in the sector $`r=2\omega `$.
In the sectors $`r=2\omega +l`$ we also will find states, which consist of one pseudoparticle with parameter $`\xi _{l+1}`$ interacting with $`l`$ impurities, distributing according to (3.1), with parameters $`\xi _i,i=1,2\mathrm{},l`$.
The energy of these states is parametrized as in (3.16) and $`\xi _{l+1}`$ satisfies the condition (3.22) with $`\xi =\xi _1\mathrm{}\xi _l\xi _{l+1}`$. It involves only $`\xi _{l+1}`$ and $`\xi _{\mathrm{imp}}=\xi _1\xi _2\mathrm{}\xi _l`$, being therefore highly degenerate, i.e.
$$\xi _{l+1}^N\xi _1^2\xi _2^2\mathrm{}\xi _l^2=1$$
(3.24)
This is to be expected due to the irrelevance of the relative distances, up to jumps of two positions via exchange with a pseudoparticle. Moreover, these results do not depend on impurity type.
### 3.3 Two-pseudoparticles
The sector $`r=4\omega `$ contains, in addition to the cases discussed above, states which consist of two interacting pseudoparticles. We seek eigenstates in the form
$$\mathrm{\Psi }_{4\omega }(\xi _1,\xi _2)=\underset{k_1+1<k_2}{}A(k_1,k_2)|\mathrm{\Omega }(k_1,k_2)$$
(3.25)
Again, we try to build two-pseudoparticle eigenstates out of translational invariant products of one-pseudoparticle excitations at $`k_1`$ and $`k_2`$ ($`k_2`$ $`k_1`$ $`+2`$) :
$$\mathrm{\Psi }_{4\omega }(\xi _1,\xi _2)=\mathrm{\Psi }_{2\omega }(\xi _1)\times \mathrm{\Psi }_{2\omega }(\xi _2)+\mathrm{\Psi }_{2\omega }(\xi _2)\times \mathrm{\Psi }_{2\omega }(\xi _1)$$
(3.26)
Using again (3.5) and comparing (3.26) with (3.25) we get
$$|\mathrm{\Omega }(k_1,k_2)=\underset{\alpha ,\beta =\omega }{\overset{\omega }{}}ϵ(\alpha )ϵ(\beta )q^{<2e_\omega +e_\alpha +e_\beta ,\rho >}|k_1(\alpha ,\alpha ),k_2(\beta ,\beta )$$
(3.27)
for $`k_2k_1+3`$ and
$$A(k_1,k_2)=A_{12}\xi _1^{k_1}\xi _2^{k_2}+A_{21}\xi _2^{k_1}\xi _1^{k_2},$$
(3.28)
for $`k_2`$ $`k_1`$ $`+2.`$
Periodic boundary conditions $`A(k_2,N+k_1)=A(k_1,k_2)`$ and $`A(N+k_1,N+k_2)=A(k_1,k_2)`$ imply
$$A_{12}\xi _2^N=A_{21}\mathrm{and}\xi ^N=1$$
(3.29)
where $`\xi =\xi _1\xi _2`$ ($`\xi _i=\mathrm{e}^{i\theta _i},i=1,2`$) and the total momentum is $`\theta _1+\theta _2=2\pi l/N`$, with $`l`$ integer.
Applying $``$ to the state of (3.25), we obtain a set of equations for the wavefunctions $`A(k_1,k_2)`$. When the two pseudoparticles are separated, ($`k_2k_1+3`$) these are the following far equations:
$$\begin{array}{ccc}\left(E_{4\omega }2Q2Q^1\right)A(k_1,k_2)\hfill & =\hfill & ϵA(k_11,k_2)+ϵA(k_1+1,k_2)\hfill \\ & & \\ & +\hfill & ϵA(k_1,k_21)+ϵA(k_1,k_2+1)\hfill \end{array}$$
(3.30)
We already know them to be satisfied, if we parametrize $`A(k_1,k_2)`$ by plane waves (3.28). The corresponding energy eigenvalue is
$$E_{4\omega }=2Q+2Q^1+ϵ\left(\xi _1+\xi _1^1+\xi _2+\xi _2^1\right)$$
(3.31)
The real problem arises of course, when pseudoparticles are neighbors, so that they interact and we have no guarantee that the total energy is sum of single pseudoparticle energies.
Acting of $``$ on the state (3.27) gives the following set of equations for the near states
$$\begin{array}{ccc}|\mathrm{\Omega }(k,k+2)\hfill & =\hfill & 2\left(Q+Q^1\right)|\mathrm{\Omega }(k,k+2)+ϵ|\mathrm{\Omega }(k1,k+2)\hfill \\ & & \\ & +\hfill & ϵ|\mathrm{\Omega }(k,k+3)+U_{k+1}|\mathrm{\Omega }(k,k+2)\hfill \end{array}$$
(3.32)
Before we substitute this result into the eigenvalue equation, we observe that some new states are appearing. In order to incorporate these new states in the eigenvalue problem, we define
$$U_{k+1}|\mathrm{\Omega }(k,k+2)ϵ|\mathrm{\Omega }(k,k+1)+ϵ|\mathrm{\Omega }(k+1,k+2)$$
(3.33)
Here we underline that we are using the same notation for these new states. Applying $``$ to them we obtain
$$\begin{array}{ccc}|\mathrm{\Omega }(k,k+1)\hfill & =\hfill & \left(Q+Q^1\right)|\mathrm{\Omega }(k,k+1)+ϵ|\mathrm{\Omega }(k1,k+1)\hfill \\ & & \\ & +\hfill & ϵ|\mathrm{\Omega }(k,k+2)\hfill \end{array}$$
(3.34)
Now, we extend (3.25), the definition of $`\mathrm{\Psi }_{4\omega }`$ , to
$$\mathrm{\Psi }_{4\omega }(\xi _1,\xi _2)=\underset{k_1<k_2}{}A(k_1,k_2)|\mathrm{\Omega }(k_1,k_2)$$
(3.35)
Substituting (3.32) and (3.34) into the eigenvalue equation, we obtain the following set of near equations
$$\left(E_{4\omega }QQ^1\right)A(k,k+1)=ϵA(k1,k+1)+ϵA(k,k+2)$$
(3.36)
Using the same parametrization (3.28) for these new wavefunctions, the equation (3.36) gives us the phase shift produced by the interchange of the two interacting pseudoparticles
$$\frac{A_{21}}{A_{12}}=\frac{ϵ(1+\xi )+(Q+Q^1)\xi _2}{ϵ(1+\xi )+(Q+Q^1)\xi _1}$$
(3.37)
We thus arrive to the Bethe Ansatz equations which fix the values of $`\xi _1`$ and $`\xi _2`$ in the energy equation (3.31)
$`\xi _2^N={\displaystyle \frac{1+\xi +ϵ(Q+Q^1)\xi _2}{1+\xi +ϵ(Q+Q^1)\xi _1}}`$
$`\xi ^N=(\xi _1\xi _2)^N=1`$ (3.38)
### 3.4 Two-pseudoparticles and impurities
In the sectors $`r>6\omega `$, in addition the cases already discussed, we find states with two interacting particles and impurities. Let us now consider two pseudoparticles with one impurity of type $`\omega 1`$. Theses eigenstates lie in the sector $`r=4\omega +1`$ and we seek them in the form
$`\mathrm{\Psi }_{4\omega +1}(\xi _1,\xi _2,\xi _3)`$ $`=`$ $`{\displaystyle \underset{k_1+1<k_2<k_32}{}}A_\mathrm{𝟏}(k_1,k_2,k_3)|\mathrm{\Omega }_1(k_1,k_2,k_3)`$ (3.39)
$`+{\displaystyle \underset{k_1+1<k_2<k_3}{}}A_\mathrm{𝟐}(k_1,k_2,k_3)|\mathrm{\Omega }_2(k_1,k_2,k_3)`$
$`+{\displaystyle \underset{k_1+1<k_2<k_31}{}}A_\mathrm{𝟑}(k_1,k_2,k_3)|\mathrm{\Omega }_3(k_1,k_2,k_3)`$
In $`A_𝐢(k_1,k_2,k_3)`$ the index $`𝐢=1,2,3`$ characterizes the impurity position. Comparing (3.39) with the state build from the translational invariant products of two-pseudoparticles with parameters $`\xi _2`$ and $`\xi _3`$ and one-impurity with parameter $`\xi _1`$:
$`\mathrm{\Psi }_{4\omega +1}(\xi _1,\xi _2,\xi _3)`$ $`=`$ $`|(\omega 1)(\xi _1)\times \mathrm{\Psi }_{2\omega }(\xi _2)\times \mathrm{\Psi }_{2\omega }(\xi _3)`$ (3.40)
$`+|(\omega 1)(\xi _1)\times \mathrm{\Psi }_{2\omega }(\xi _3)\times \mathrm{\Psi }_{2\omega }(\xi _2)`$
$`+\mathrm{\Psi }_{2\omega }(\xi _2)\times |(\omega 1)(\xi _1)\times \mathrm{\Psi }_{2\omega }(\xi _3)`$
$`+\mathrm{\Psi }_{2\omega }(\xi _3)\times |(\omega 1)(\xi _1)\times \mathrm{\Psi }_{2\omega }(\xi _2)`$
$`+\mathrm{\Psi }_{2\omega }(\xi _2)\times \mathrm{\Psi }_{2\omega }(\xi _3)\times |(\omega 1)(\xi _1)`$
$`+\mathrm{\Psi }_{2\omega }(\xi _3)\times \mathrm{\Psi }_{2\omega }(\xi _2)\times |(\omega 1)(\xi _1)`$
we get
$`|\mathrm{\Omega }_1(k_1,k_2,k_3)={\displaystyle \underset{\alpha ,\beta =\omega }{\overset{\omega }{}}}W(\alpha ,\beta ,\omega )|k_1(\omega 1),k_2(\alpha ,\alpha ),k_3(\beta ,\beta )`$
$`|\mathrm{\Omega }_2(k_1,k_2,k_3)={\displaystyle \underset{\alpha ,\beta =\omega }{\overset{\omega }{}}}W(\alpha ,\beta ,\omega )|k_1(\alpha ,\alpha ),k_2(\omega 1),k_3(\beta ,\beta )`$
$`|\mathrm{\Omega }_3(k_1,k_2,k_3)={\displaystyle \underset{\alpha ,\beta =\omega }{\overset{\omega }{}}}W(\alpha ,\beta ,\omega )|k_1(\alpha ,\alpha ),k_2(\beta ,\beta ),k_3(\omega 1)`$ (3.41)
where
$$W(\alpha ,\beta ,\omega )=ϵ(\alpha )ϵ(\beta )q^{<2e_\omega +e_\alpha +e_\beta ,\rho >}$$
(3.42)
and the wavefunctions $`A_𝐢(k_1,k_2,k_3)`$ which are parametrized by plane waves as
$`A_\mathrm{𝟏}(k_1,k_2,k_3)=A_{\mathrm{𝟏}23}\xi _1^{k_1}\xi _2^{k_2}\xi _3^{k_3}+A_{\mathrm{𝟏}32}\xi _1^{k_1}\xi _2^{k_3}\xi _3^{k_2}`$
$`A_\mathrm{𝟐}(k_1,k_2,k_3)=A_{\mathrm{𝟐}13}\xi _1^{k_2}\xi _2^{k_1}\xi _3^{k_3}+A_{\mathrm{𝟐}31}\xi _1^{k_2}\xi _2^{k_3}\xi _3^{k_1}`$
$`A_\mathrm{𝟑}(k_1,k_2,k_3)=A_{\mathrm{𝟑}12}\xi _1^{k_3}\xi _2^{k_1}\xi _3^{k_2}+A_{\mathrm{𝟑}21}\xi _1^{k_3}\xi _2^{k_2}\xi _3^{k_1}.`$ (3.43)
Periodic boundary conditions read now
$`A_𝐢(k_1,k_2,k_3)=A_𝐢(N+k_1,N+k_2,N+k_3),`$
$`A_𝐢(k_2,k_3,N+k_1)=A_{𝐢+\mathrm{𝟏}}(k_1,k_2,k_3),𝐢=1,2,3mod3`$ (3.44)
which imply that
$`\xi _1^N={\displaystyle \frac{A_{\mathrm{𝟏}23}}{A_{\mathrm{𝟑}12}}}={\displaystyle \frac{A_{\mathrm{𝟏}32}}{A_{\mathrm{𝟑}21}}},\xi _2^N={\displaystyle \frac{A_{\mathrm{𝟑}12}}{A_{\mathrm{𝟐}31}}}={\displaystyle \frac{A_{\mathrm{𝟐}13}}{A_{\mathrm{𝟏}32}}},`$
$`\xi _3^N={\displaystyle \frac{A_{\mathrm{𝟑}21}}{A_{\mathrm{𝟐}13}}}={\displaystyle \frac{A_{\mathrm{𝟐}31}}{A_{\mathrm{𝟏}23}}},\xi ^N=(\xi _1\xi _2\xi _3)^N=1`$ (3.45)
Action of $``$ on the state $`\mathrm{\Psi }_{4\omega +1}`$ gives the following set of far equations:
$`\left(E_{4\omega +1}2Q2Q^1\right)A_1(k_1,k_2,k_3)=ϵA_1(k_1,k_21,k_3)+ϵA_1(k_1,k_2+1,k_3)`$
$`+ϵA_1(k_1,k_2,k_31)+ϵA_1(k_1,k_2,k_3+1)`$
(3.46)
and a similar set of eigenvalue equations for $`A_\mathrm{𝟐}(k_1,k_2,k_3)`$ and $`A_\mathrm{𝟑}(k_1,k_2,k_3)`$. The parametrization (3.43) solves these far equations provided that
$$E_{4\omega +1}=2Q+2Q^1+ϵ(\xi _2+\xi _2^1+\xi _3+\xi _3^1)$$
(3.47)
Taking into account the near equations we must split them in three different neighborhood: (i) impurity neighbors of separated pseudoparticles and, (ii) impurity far from neighbors pseudoparticles and (iii) when impurity and pseudoparticles share the same neighborhood.
In the case (i) we consider the second pseudoparticle far and follow the steps for the case of one-pseudoparticle with impurity eigenstates. Thus, the near equations can be read off from (3.36)
$`(E_{4\omega +1}2Q2Q^1)A_\mathrm{𝟏}(k,k+1,k_3)=ϵA_\mathrm{𝟏}(k,k+1,k_31)+ϵA_\mathrm{𝟏}(k,k+1,k_3+1)`$
$`+ϵA_\mathrm{𝟏}(k,k+2,k_3)+ϵA_\mathrm{𝟐}(k,k+2,k_3)`$ (3.48)
and a similar set of equations coupling $`A_\mathrm{𝟐}`$ and $`A_\mathrm{𝟑}`$. It follows from the consistency between (3.46) and (3.48) that
$$A_\mathrm{𝟏}(k,k,k_3)A_\mathrm{𝟐}(k,k+2,k_3)$$
(3.49)
and similar identification between $`A_\mathrm{𝟐}`$ and $`A_\mathrm{𝟑}`$. The plane waves (3.43) solve these identifications provided
$$\xi _1^2=\frac{A_{\mathrm{𝟏}23}}{A_{\mathrm{𝟐}13}}=\frac{A_{\mathrm{𝟏}32}}{A_{\mathrm{𝟐}31}}.=\frac{A_{\mathrm{𝟐}31}}{A_{\mathrm{𝟑}21}}=\frac{A_{\mathrm{𝟐}13}}{A_{\mathrm{𝟑}12}}$$
(3.50)
For the case (ii) we can derive the near equations from those of two-pseudoparticles case. Keeping the impurity far and following the steps (3.30)–(3.37) we get
$$\left(E_{4\omega +1}QQ^1\right)A_1(k_1,k,k+1)=ϵA_1(k_1,k1,k+1)+ϵA_1(k_1,k,k+2)$$
(3.51)
and a similar set of equations for $`A_\mathrm{𝟐}`$ and $`A_\mathrm{𝟑}`$. The case (iii) is obtained from (3.51) for $`k_1=k1`$.
The parametrization (3.43) solves this provided that
$$\frac{A_{\mathrm{𝟏}23}}{A_{\mathrm{𝟏}32}}=\frac{ϵ(1+\xi _3\xi _2)+(Q+Q^1)\xi _2}{ϵ(1+\xi _2\xi _3)+(Q+Q^1)\xi _3}$$
(3.52)
Matching the constraint equations (3.52), (3.50) and (3.45) we arrive to the Bethe equations
$`\xi _a^N\xi _1^2={\displaystyle \frac{1+\xi _b\xi _a+ϵ(Q+Q^1)\xi _a}{1+\xi _a\xi _b+ϵ(Q+Q^1)\xi _b}},ab=2,3`$ (3.53)
$`\xi ^N=(\xi _1\xi _2\xi _3)^N=1,\xi _1^{N4}=1.`$ (3.54)
The origin of the exponent ($`N4`$) in the impurity parameter can be understood by saying that after the two pseudoparticles propagate past the impurity, the position of impurity is shifted by four lattice sites.
Next, we can also find eigenstates with two pseudoparticles and more than one impurities. They can be described in the following way: Let us consider an eigenstate with $`l>1`$ impurities with parameters $`\xi _1,\xi _2,\mathrm{},\xi _l`$ and two pseudoparticles with parameters $`\xi _{l+1}`$ and $`\xi _{l+2}`$. The energy eigenvalue is
$$E_r=2Q+2Q^1+ϵ\left(\xi _{l+1}+\xi _{l+1}^1+\xi _{l+2}+\xi _{l+2}^1\right)$$
(3.55)
and the Bethe equations
$`\xi _{l+1}^N\xi _1^2\xi _2^2\mathrm{}\xi _l^2={\displaystyle \frac{1+\xi _{l+1}\xi _{l+2}+ϵ(Q+Q^1)\xi _{l+1}}{1+\xi _{l+1}\xi _{l+2}+ϵ(Q+Q^1)\xi _{l+2}}}`$
$`\xi _a^{N4}=1,a=1,2,\mathrm{},l`$ (3.56)
Moreover, $`\xi ^N=1`$ with $`\xi =\xi _1\xi _2\mathrm{}\xi _{l+2}`$.
### 3.5 Three-pseudoparticle eigenstates
In the sector $`r=6\omega `$, in addition to the previously discussed eigenstates of one and two pseudoparticles with impurities, one can find eigenstates with three interacting pseudoparticles with parameters $`\xi _1,\xi _2`$ and $`\xi _3`$. We start seek them in the form
$$\mathrm{\Psi }_{6\omega }(\xi _1,\xi _2,\xi _3)=\underset{k_1+2k_2k_32}{}A(k_1,k_2,k_3)|\mathrm{\Omega }(k_1,k_2,k_3)$$
(3.57)
where $`|\mathrm{\Omega }(k_1,k_2,k_3)=_{i=1}^3|\mathrm{\Omega }(k_i)`$ . The corresponding wavefunctions
$$\begin{array}{ccc}A(k_1,k_2,k_3)\hfill & =\hfill & A_{123}\xi _1^{k_1}\xi _2^{k_2}\xi _3^{k_3}+A_{132}\xi _1^{k_1}\xi _2^{k_3}\xi _3^{k_2}+A_{213}\xi _1^{k_2}\xi _2^{k_1}\xi _3^{k_3}\hfill \\ & & \\ & & +A_{231}\xi _1^{k_2}\xi _2^{k_3}\xi _3^{k_1}+A_{312}\xi _1^{k_3}\xi _2^{k_1}\xi _3^{k_2}+A_{321}\xi _1^{k_3}\xi _2^{k_2}\xi _3^{k_1}\hfill \end{array}$$
(3.58)
satisfy the periodic boundary conditions
$$A(k_2,k_3,N+k_1)=A(k_1,k_2,k_3),A(N+k_1,N+k_2,N+k_3)=A(k_1,k_2,k_3)$$
(3.59)
which imply that
$`\xi _1^N`$ $`=`$ $`{\displaystyle \frac{A_{123}}{A_{312}}}={\displaystyle \frac{A_{132}}{A_{321}}},\xi _2^N={\displaystyle \frac{A_{312}}{A_{231}}}={\displaystyle \frac{A_{213}}{A_{132}}},`$
$`\xi _3^N`$ $`=`$ $`{\displaystyle \frac{A_{321}}{A_{213}}}={\displaystyle \frac{A_{231}}{A_{123}}},\xi ^N=(\xi _1\xi _2\xi _3)^N=1`$ (3.60)
These relations show us that the interchange of two-pseudoparticles is independent of the position of the third particle.
Applying $``$ to (3.57), we obtain a set of equations for $`A(k_1,k_2,k_3)`$. When the three pseudoparticles are separated, $`(k_1+2<k_2<k_32)`$, we get the following far equations:
$`(E_{6s}3Q3Q^1)A(k_1,k_2,k_3)`$ $`=`$ $`ϵA(k_11,k_2,k_3)+ϵA(k_1+1,k_2,k_3)`$
$`+ϵA(k_1,k_21,k_3)+ϵA(k_1,k_2+1,k_3)`$
$`+ϵA(k_1,k_2,k_31)+ϵA(k_1,k_2,k_3+1)`$
It is simple verify that the wavefunctions (3.58) satisfy these far equations provided
$$E_{6\omega }=\underset{n=1}{\overset{3}{}}\left\{Q+Q^1+\xi _n+\xi _n^1\right\}$$
(3.62)
Applying $``$ on the near states we get the following set equations:
$`|\mathrm{\Omega }(k_1,k_1+2,k_3)`$ $`=`$ $`(2Q+2Q^1)|\mathrm{\Omega }(k_1,k_1+2,k_3)+ϵ|\mathrm{\Omega }(k_11,k_1+2,k_3)`$
$`+ϵ|\mathrm{\Omega }(k_1,k_1+3,k_3)+ϵ|\mathrm{\Omega }(k_1,k_1+2,k_31)`$
$`+ϵ|\mathrm{\Omega }(k_1,k_1+2,k_3+1)+U_{k_1+1}|\mathrm{\Omega }(k_1,k_1+2,k_3)`$
for $`k_3>k_1+4`$, which correspond to the meeting of two pseudoparticles at the left of the third pseudoparticle, which is far from of the meeting position.
$`|\mathrm{\Omega }(k_1,k_2,k_2+2)`$ $`=`$ $`(2Q+2Q^1)|\mathrm{\Omega }(k_1,k_2,k_2+2)+ϵ|\mathrm{\Omega }(k_11,k_2,k_2+2)`$
$`+ϵ|\mathrm{\Omega }(k_1+1,k_2,k_2+2)+ϵ|\mathrm{\Omega }(k_1,k_21,k_2+2)`$
$`+ϵ|\mathrm{\Omega }(k_1,k_2,k_2+3)+U_{k_2+1}|\mathrm{\Omega }(k_1,k_2,k_2+2)`$
for $`k_2>k_1+2`$, which correspond to the meeting of two pseudoparticles at the right of the far pseudoparticle. Moreover, there is one set of equations which correspond to the meeting of three pseudoparticles
$`|\mathrm{\Omega }(k,k+2,k+4)`$ $`=`$ $`(Q+Q^1)|\mathrm{\Omega }(k,k+2,k+4)+ϵ|\mathrm{\Omega }(k1,k+2,k+4)`$ (3.65)
$`+ϵ|\mathrm{\Omega }(k,k+2,k+5)+U_{k+1}|\mathrm{\Omega }(k,k+2,k+4)`$
$`+U_{k+3}|\mathrm{\Omega }(k,k+2,k+4)`$
In deriving these equations new states made their debut. In order to incorporate these new states in the eigenvalue problem we define:
$`U_{k_1+1}|\mathrm{\Omega }(k_1,k_1+2,k_3)`$ $`=`$ $`ϵ|\mathrm{\Omega }(k_1,k_1+1,k_3)+ϵ|\mathrm{\Omega }(k_1+1,k_1+2,k_3)`$
$`U_{k_2+1}|\mathrm{\Omega }(k_1,k_2,k_2+2)`$ $`=`$ $`ϵ|\mathrm{\Omega }(k_1,k_2+1,k_2+2)+ϵ|\mathrm{\Omega }(k_1,k_2,k_2+1)`$
$`U_{k+1}|\mathrm{\Omega }(k,k+2,k+4)`$ $`=`$ $`ϵ|\mathrm{\Omega }(k,k+1,k+4)+ϵ|\mathrm{\Omega }(k+1,k+2,k+4)`$
$`U_{k+3}|\mathrm{\Omega }(k,k+2,k+4)`$ $`=`$ $`ϵ|\mathrm{\Omega }(k,k+3,k+4)+ϵ|\mathrm{\Omega }(k,k+2,k+3)`$ (3.66)
Applying $``$ to these new states the result can be incorporated to the eigenvalue problem provided the definition of $`\mathrm{\Psi }_{6\omega }`$ (3.57) is extended to
$$\mathrm{\Psi }_{6\omega }(\xi _1,\xi _2,\xi _3)=\underset{k_1<k_2<k_3}{}A(k_1,k_2,k_3)|\mathrm{\Omega }(k_1,k_2,k_3)$$
(3.67)
After this we are left with three meeting equations
$`(E_{6\omega }2Q2Q^1)A(k_1,k_1+1,k_3)=ϵA(k_11,k_1+1,k_3)+ϵA(k_1,k_1+2,k_3)`$
$`+ϵA(k_1,k_1+1,k_31)+ϵA(k_1,k_1+1,k_3+1)`$
for $`k_3>k_1+2`$,
$`(E_{6\omega }2Q2Q^1)A(k_1,k_2,k_2+1)=ϵA(k_1,k_21,k_2+1)+ϵA(k_1,k_2,k_2+2)`$
$`+ϵA(k_11,k_2,k_2+1)+ϵA(k_1+1,k_2,k_2+1)`$
for $`k_1+2<k_2`$ and
$$(E_{6\omega }QQ^1)A(k,k+1,k+2)=ϵA(k1,k+1,k+2)+ϵA(k,k+1,k+3)$$
It is easy to verify that the parametrization (3.58) and (3.62) solve these equations provided
$`{\displaystyle \frac{A_{123}}{A_{213}}}`$ $`=`$ $`{\displaystyle \frac{A_{231}}{A_{321}}}={\displaystyle \frac{1+\xi _1\xi _2+ϵ(Q+Q^1)\xi _1}{1+\xi _1\xi _2+ϵ(Q+Q^1)\xi _2}}`$
$`{\displaystyle \frac{A_{132}}{A_{231}}}`$ $`=`$ $`{\displaystyle \frac{A_{213}}{A_{312}}}={\displaystyle \frac{1+\xi _1\xi _3+ϵ(Q+Q^1)\xi _1}{1+\xi _1\xi _3+ϵ(Q+Q^1)\xi _3}}`$
$`{\displaystyle \frac{A_{312}}{A_{321}}}`$ $`=`$ $`{\displaystyle \frac{A_{123}}{A_{132}}}={\displaystyle \frac{1+\xi _2\xi _3+ϵ(Q+Q^1)\xi _2}{1+\xi _2\xi _3+ϵ(Q+Q^1)\xi _3}}`$ (3.70)
Matching these constraints and the periodic boundary conditions (3.60) we get the Bethe Ansatz equations
$`\xi _a^N`$ $`=`$ $`{\displaystyle \underset{ba=1}{\overset{3}{}}}\left\{{\displaystyle \frac{1+\xi _a\xi _b+ϵ(Q+Q^1)\xi _a}{1+\xi _a\xi _b+ϵ(Q+Q^1)\xi _b}}\right\},a=1,2,3`$
$`(\xi _1\xi _2\xi _3)^N`$ $`=`$ $`1`$ (3.71)
### 3.6 General eigenstates
The generalization to any $`r`$ is now immediate. Since the Yang-Baxter equations are satisfied, there is only two-pseudoparticle scattering (using the $`S`$-matrix language). Therefore, neighbor equations, where more then two pseudoparticles become neighbors, are nor expected to give any new restrictions. For instance, in the sector $`r=6\omega `$, we saw that the interchange of two-pseudoparticles is independent of the position of the third particle. Thus, in a sector with $`p`$ pseudoparticles we expect that the $`p`$-pseudoparticle phase shift will be a sum of $`p(p1)/2`$ two-pseudoparticle phase shift. The energy is given by the sum of single pseudoparticle energies. The corresponding Bethe Ansatz equations depend on the phase shift of two pseudoparticles and on the number of impurity. For a generic sector one can verify that no different neighborhood those discussed above can appear. So, no additional meeting conditions will be encountered. Thus, we can extend the previous results to the $`p`$ -pseudoparticle states in the following way: In a generic sector $`r`$ with $`l`$ impurities parametrized by $`\xi _1\xi _2\mathrm{}\xi _l`$ and $`p`$ pseudoparticles with parameters $`\xi _{l+1}\xi _{l+2}\mathrm{}\xi _{l+p}`$, the energy is
$$E_r=\underset{n=l+1}{\overset{p}{}}\left\{Q+Q^1+ϵ\left(\xi _n+\xi _n^1\right)\right\}$$
(3.72)
with $`\xi _n`$ determined by the Bethe ansatz equations
$`\xi _a^N\xi _1^2\xi _2^2\mathrm{}\xi _l^2`$ $`=`$ $`{\displaystyle \underset{ba=l+1}{\overset{l+p}{}}}\left\{{\displaystyle \frac{1+\xi _b\xi _a+ϵ(Q+Q^1)\xi _a}{1+\xi _a\xi _b+ϵ(Q+Q^1)\xi _b}}\right\}`$
$`\xi _c^{N2p}`$ $`=`$ $`1,c=1,2,\mathrm{},l`$
$`\xi ^N`$ $`=`$ $`1,\xi =\xi _1\xi _2\mathrm{}\xi _l\xi _{l+1}\xi _{l+2}\mathrm{}\xi _{l+p}.`$ (3.73)
The energy eigenvalues and the Bethe equations depend on the deformation parameter $`q`$, through the relation (2.1):
$$Q+Q^1=\{\begin{array}{ccc}[2s+1]\hfill & \mathrm{for}\hfill & A_1\hfill \\ [2n1][n+1/2]/[n1/2]\hfill & \mathrm{for}\hfill & B_n(n2)\hfill \\ [n][2n+2]/[n+1]\hfill & \mathrm{for}\hfill & C_n(n1)\hfill \\ [n][2n2]/[n1]\hfill & \mathrm{for}\hfill & D_n(n3)\hfill \end{array}$$
(3.74)
We obtained thus the spectra with periodic boundary conditions of quantum spin-chain models, arising as representations of the Temperley-Lieb algebra. As expected, all these models have equivalent spectra up to degeneracies of their eigenvalues. From a suitable sorting of the parameters $`\xi _i`$, one can insure that the spectra of lower-$`r`$ sectors are contained entirely in the higher-$`r`$ sectors.
## 4 Bethe Ansatz: Non-local boundary conditions
It is the purpose of this section to present and solve, via coordinate Bethe ansatz, the quantum group invariant closed TL Hamiltonians which can be written as :
$$=\underset{k=1}{\overset{N1}{}}U_k+𝒰_0$$
(4.1)
where $`U_k`$ is a Temperley-Lieb operator and $`\mathrm{𝐛𝐭}=𝒰_0`$ is non-local term defined through of a operator $`G`$ which plays the role of the translation operator
$$𝒰_0=GU_{N1}G^1,G=(QU_1)(QU_2)\mathrm{}(QU_{N1})$$
(4.2)
satisfying $`[,G]=0`$ and additionally invariance with respect to the quantum algebra. The operator $`G`$ shifts the $`U_k`$ by one unit $`GU_kG^1=U_{k+1}`$ and maps $`𝒰_0`$ into $`U_1`$, which manifest the translational invariance of $``$. In this sense the Hamiltonian (4.1) is periodic.
For the $`q`$-deformed $`A`$-$`D`$ Temperley-Lieb algebra, the matrix elements of $`U_k`$ is again given by (2.8), i.e. :
$`\nu ,\lambda |U|\mu ,\kappa =ϵ(\mu )ϵ(\nu )q^{<e_\mu +e_\nu ,\rho >}\delta _{\mu +\kappa ,0}\delta _{\nu +\lambda ,0}`$
$`\mu ,\nu ,\kappa ,\lambda J`$ (4.3)
From (4.3) we choose an particular Bethe state
$$|\mathrm{\Omega }(k)=\underset{\alpha =\omega }{\overset{\omega }{}}ϵ(\omega )ϵ(\alpha )q^{<e_\omega +e_\alpha ,\rho >}|k(\alpha ,\alpha )$$
(4.4)
which is an eigenstate of $`U_k`$ and it is shifted by one unit under the action of $`U_{k\pm 1}`$
$`U_k|\mathrm{\Omega }(k)=(Q+Q^1)|\mathrm{\Omega }(k)`$
$`U_{k\pm 1}|\mathrm{\Omega }(k)=ϵ|\mathrm{\Omega }(k\pm 1),U_k|\mathrm{\Omega }(k\pm )=ϵ|\mathrm{\Omega }(k)`$
$`U_k|\mathrm{\Omega }(l)=0fork\{l1,l,l+1\}`$ (4.5)
where $`ϵ=1`$ for $`B_n`$, $`D_n`$ and $`A_1`$ ( $`s`$ integer) and $`ϵ=1`$ for $`C_n`$ and $`A_1`$ ( $`s`$ semi-integer).
The action of the operator $`G`$ on the states $`|\mathrm{\Omega }(k)`$ can be easily computed using (4.5): It is simple on the bulk and at the left boundary
$$G|\mathrm{\Omega }(k)=ϵQ^{N2}|\mathrm{\Omega }(k+1),1kN2$$
(4.6)
but manifests its nonlocality at the right boundary
$$G|\mathrm{\Omega }(N1)=ϵQ^{N2}\underset{k=1}{\overset{N1}{}}(ϵQ)^k|\mathrm{\Omega }(Nk)$$
(4.7)
Similarly, the action of the operator $`G^1=(Q^1U_{N1})\mathrm{}(Q^1U_1)`$ is simple on the bulk and at the right boundary
$$G^1|\mathrm{\Omega }(k)=ϵQ^{N+2}|\mathrm{\Omega }(k1),2kN1$$
(4.8)
and non-local at the left boundary
$$G^1|\mathrm{\Omega }(1)=ϵQ^{N+2}\underset{k=1}{\overset{N1}{}}(ϵQ)^k|\mathrm{\Omega }(k).$$
(4.9)
Now we proceed the diagonalization of $``$ as was made for the periodic case. As (4.1) and (1.2) have the same bulk, i.e., differences came from the boundary terms, we will keep all results relating to the bulk of the periodic case presented in the previous section.
### 4.1 One-pseudoparticle eigenstates
Let us consider one free pseudoparticle which lies in the sector $`r=2\omega `$
$$\mathrm{\Psi }_{2\omega }=\underset{k=1}{\overset{N1}{}}A(k)|\mathrm{\Omega }(k).$$
(4.10)
The action of the operator $`𝒰=_{k=1}^{N1}U_k`$ on the states $`|\mathrm{\Omega }(k)`$ can be computed using (4.5):
$`𝒰|\mathrm{\Omega }(1)=(Q+Q^1)|\mathrm{\Omega }(1)+ϵ|\mathrm{\Omega }(2)`$
$`𝒰|\mathrm{\Omega }(k)=(Q+Q^1)|\mathrm{\Omega }(k)+ϵ|\mathrm{\Omega }(k1)+ϵ|\mathrm{\Omega }(k+1)`$
$`\mathrm{for}2kN2`$
$`𝒰|\mathrm{\Omega }(N1)=(Q+Q^1)|\mathrm{\Omega }(N1)+ϵ|\mathrm{\Omega }(N2).`$ (4.11)
and using (4.6)–(4.9) one can see that the action of $`𝒰_0=GU_{N1}G^1`$ vanishes on the bulk
$$𝒰_0|\mathrm{\Omega }(k)=0,2kN2$$
(4.12)
and is nonlocal at the boundaries
$$𝒰_0|\mathrm{\Omega }(1)=ϵ\underset{k=1}{\overset{N1}{}}(ϵQ)^k|\mathrm{\Omega }(k),𝒰_0|\mathrm{\Omega }(N1)=ϵ\underset{k=1}{\overset{N1}{}}(ϵQ)^{N+k}|\mathrm{\Omega }(k).$$
(4.13)
which are connected by
$$𝒰_0|\mathrm{\Omega }(N1)=(ϵQ)^N𝒰_0|\mathrm{\Omega }(1).$$
(4.14)
From these equations we can understood the role of $`𝒰_0`$: Although the Hamiltonian (4.1) is a global operator, it manifests the property of essential locality. From the physical point of view, this type of models exhibit behavior similar to closed chains with twisted boundary conditions.
Before we substitute these results into the eigenvalue equation, we will define two new states
$$ϵ|\mathrm{\Omega }(0)=𝒰_0|\mathrm{\Omega }(1),ϵ|\mathrm{\Omega }(N)=𝒰_0|\mathrm{\Omega }(N1)$$
(4.15)
to include the cases $`k=0`$ and $`k=N`$ into the definition of $`\mathrm{\Psi }_{2\omega }`$, equation (4.10). Finally, the action of $`=𝒰+𝒰_0`$ on the states $`|\mathrm{\Omega }(k)`$ is
$`|\mathrm{\Omega }(0)=(Q+Q^1)|\mathrm{\Omega }(0)+(ϵQ)^Nϵ|\mathrm{\Omega }(N1)+ϵ|\mathrm{\Omega }(1)`$
$`|\mathrm{\Omega }(k)=(Q+Q^1)|\mathrm{\Omega }(k)+ϵ|\mathrm{\Omega }(k1)+ϵ|\mathrm{\Omega }(k+1)`$
$`\mathrm{for}1kN2`$
$`|\mathrm{\Omega }(N1)=(Q+Q^1)|\mathrm{\Omega }(N1)+ϵ|\mathrm{\Omega }(N2)+(ϵQ)^Nϵ|\mathrm{\Omega }(0)`$
$`|\mathrm{\Omega }(N)=(Q+Q^1)|\mathrm{\Omega }(N)+ϵ|\mathrm{\Omega }(N1)+(ϵQ)^Nϵ|\mathrm{\Omega }(1)`$ (4.16)
Substituting these results into the eigenvalue equation $`\mathrm{\Psi }_{2\omega }=E_{2\omega }\mathrm{\Psi }_{2\omega }`$ we get a complete set of eigenvalue equations for the wavefunctions
$`E_{2s}A(k)=(Q+Q^1)A(k)+ϵA(k1)+ϵA(k+1)`$
$`\mathrm{for}1kN1`$ (4.17)
provided the following boundary conditions
$$(ϵQ)^NA(k)=A(N+k)$$
(4.18)
are satisfied.
The plane wave parametrization $`A(k)=A\xi ^k`$ solves these eigenvalue equations and the boundary conditions provided that:
$`E_{2\omega }=Q+Q^1+ϵ(\xi +\xi ^1)`$
$`\xi ^N=(ϵQ)^N`$ (4.19)
where $`\xi =\mathrm{e}^{i\theta }`$ and $`\theta `$ being the momentum.
### 4.2 Two-pseudoparticle eigenstates
Let us now consider the sector $`r=4\omega `$, where we can find an eigenstate with two interacting pseudoparticles. We seek the corresponding eigenfunction as products of single pseudoparticles eigenfunctions, i.e.
$$\mathrm{\Psi }_{4\omega }=\underset{k_1+1<k_2}{}A(k_1,k_2)|\mathrm{\Omega }(k_1,k_2)$$
(4.20)
where
$$|\mathrm{\Omega }(k_1,k_2)=\underset{\alpha ,\beta =\omega }{\overset{\omega }{}}ϵ(\alpha )ϵ(\beta )q^{<2e_\omega +e_\alpha +e_\beta ,\rho >}|k_1(\alpha ,\alpha ),k_2(\beta ,\beta )$$
(4.21)
To solve the eigenvalue equation $`\mathrm{\Psi }_{4\omega }=E_{4\omega }\mathrm{\Psi }_{4\omega }`$, we recall (4.5) to get the action of $`𝒰`$ and $`𝒰_0`$ on the states $`|\mathrm{\Omega }(k_1,k_2)`$. Here we have to consider four cases: (i) when the two pseudoparticles are separated in the bulk, the action of $`𝒰`$ is
$`𝒰|\mathrm{\Omega }(k_1,k_2)=2(Q+Q^1)|\mathrm{\Omega }(k_1,k_2)+ϵ|\mathrm{\Omega }(k_11,k_2)+ϵ|\mathrm{\Omega }(k_1+1,k_2)`$
$`+ϵ|\mathrm{\Omega }(k_1,k_21)+ϵ|\mathrm{\Omega }(k_1,k_2+1)`$ (4.22)
i.e., for $`k_1`$ $`2`$ and $`k_1+3k_2N2`$; (ii) when the two pseudoparticles are separated but one of them or both are at the boundaries
$`𝒰|\mathrm{\Omega }(1,k_2)=2(Q+Q^1)|\mathrm{\Omega }(1,k_2)+ϵ|\mathrm{\Omega }(2,k_2)+ϵ|\mathrm{\Omega }(1,k_21)`$
$`+ϵ|\mathrm{\Omega }(1,k_2+1)`$ (4.23)
$`𝒰|\mathrm{\Omega }(k_1,N1)=2(Q+Q^1)|\mathrm{\Omega }(k_1,N1)+ϵ|\mathrm{\Omega }(k_11,N1)`$
$`+ϵ|\mathrm{\Omega }(k_1+1,N1)+ϵ|\mathrm{\Omega }(k_1,N2)`$ (4.24)
$$𝒰|\mathrm{\Omega }(1,N1)=2(Q+Q^1)|\mathrm{\Omega }(1,N1)+ϵ|\mathrm{\Omega }(2,N1)+ϵ|\mathrm{\Omega }(1,N2)$$
(4.25)
where $`2k_1N4`$ and $`4k_2N2`$; (iii) when the two pseudoparticles are neighbors in the bulk
$`𝒰|\mathrm{\Omega }(k,k+2)=2(Q+Q^1)|\mathrm{\Omega }(k,k+2)+ϵ|\mathrm{\Omega }(k1,k+2)+ϵ|\mathrm{\Omega }(k,k+3)`$
$`+U_{k+1}|\mathrm{\Omega }(k,k+2)`$ (4.26)
for $`2kN4`$ and (iv) when the two pseudoparticles are neighbors and at the boundaries
$$𝒰|\mathrm{\Omega }(1,3)=2(Q+Q^1)|\mathrm{\Omega }(1,3)+ϵ|\mathrm{\Omega }(1,4)+U_2|\mathrm{\Omega }(1,3)$$
(4.27)
$`𝒰|\mathrm{\Omega }(N3,N1)=2(Q+Q^1)|\mathrm{\Omega }(N3,N1)+ϵ|\mathrm{\Omega }(N4,N1)`$
$`+U_{N2}|\mathrm{\Omega }(N3,N1)`$ (4.28)
Moreover, the action of $`𝒰_0`$ does not depend on the pseudoparticles are neither separated nor neighbors. It is vanishes in the bulk
$$𝒰_0|\mathrm{\Omega }(k_1,k_2)=0\mathrm{for}k_11\mathrm{and}k_2N1,$$
(4.29)
and different of zero at the boundaries:
$`𝒰_0|\mathrm{\Omega }(1,k_2)`$ $`=`$ $`ϵ{\displaystyle \underset{k=1}{\overset{k_22}{}}}(ϵQ)^k|\mathrm{\Omega }(k,k_2)(ϵQ)^{k_21}U_{k_2}|\mathrm{\Omega }(k_21,k_2+1)`$ (4.30)
$`ϵ{\displaystyle \underset{k=k_2+2}{\overset{N1}{}}}(ϵQ)^{k2}|\mathrm{\Omega }(k_2,k)`$
$$𝒰_0|\mathrm{\Omega }(k_1,N1)=(ϵQ)^{N+2}𝒰_0|\mathrm{\Omega }(1,k_2)$$
(4.31)
where $`2k_1N3`$ and $`3k_2N2`$.
Following the same procedure of one-pseudoparticle case we again define new states in order to have consistency between bulk and boundaries terms
$`𝒰_0|\mathrm{\Omega }(1,k_2)`$ $`=`$ $`ϵ|\mathrm{\Omega }(0,k_2),𝒰_0|\mathrm{\Omega }(k_1,N1)=ϵ|\mathrm{\Omega }(k_1,N)`$
$`𝒰_0|\mathrm{\Omega }(1,N1)`$ $`=`$ $`ϵ|\mathrm{\Omega }(0,N1)+ϵ|\mathrm{\Omega }(1,N)`$
$`U_{k+1}|\mathrm{\Omega }(k,k+2)`$ $`=`$ $`ϵ|\mathrm{\Omega }(k,k+1)+ϵ|\mathrm{\Omega }(k+1,k+2)`$ (4.32)
Acting with $``$ on these new states, we get
$`|\mathrm{\Omega }(0,k_2)`$ $`=`$ $`2(Q+Q^1)|\mathrm{\Omega }(0,k_2)+ϵ|\mathrm{\Omega }(0,k_21)+ϵ|\mathrm{\Omega }(0,k_2+1)`$ (4.33)
$`+ϵ|\mathrm{\Omega }(1,k_2)+(ϵQ)^{N2}ϵ|\mathrm{\Omega }(k_2,N1)`$
$`|\mathrm{\Omega }(k_1,N)`$ $`=`$ $`2(Q+Q^1)|\mathrm{\Omega }(k_1,N)+ϵ|\mathrm{\Omega }(k_11,N)+ϵ|\mathrm{\Omega }(k_1+1,N)`$ (4.34)
$`+ϵ|\mathrm{\Omega }(k_1,N1)+(ϵQ)^{N+2}ϵ|\mathrm{\Omega }(1,k_1)`$
$$|\mathrm{\Omega }(k,k+1=(Q+Q^1)\left|\mathrm{\Omega }(k,k+1+ϵ\right|\mathrm{\Omega }(k1,k+1+ϵ|\mathrm{\Omega }(k,k+2$$
(4.35)
Substituting these results into the eigenvalue equation, we get the following equations for wavefunctions corresponding to the separated pseudoparticles.
$`(E_{4\omega }2Q2Q^1)A(k_1,k_2)`$ $`=`$ $`ϵA(k_11,k_2)+ϵA(k_1+1,k_2)`$ (4.36)
$`+ϵA(k_1,k_21)+ϵA(k_1,k_2+1)`$
i.e., for $`k_11`$ and $`k_1+3k_2N1`$. The boundary conditions read now
$$A(k_2,N+k_1)=(ϵQ)^{N2}A(k_1,k_2).$$
(4.37)
The parametrization for the wavefunctions
$$A(k_1,k_2)=A_{12}\xi _1^{k_1}\xi _2^{k_2}+A_{21}\xi _1^{k_2}\xi _2^{k_1}$$
(4.38)
solves the equation (4.36) provided that
$$E_{4s}=2(Q+Q^1)+ϵ(\xi _1+\xi _1^1+\xi _2+\xi _2^1)$$
(4.39)
and the boundary conditions (4.37) provided that
$$\xi _2^N=(ϵQ)^{N2}\frac{A_{21}}{A_{12}},\xi _1^N=(ϵQ)^{N2}\frac{A_{12}}{A_{21}}\xi ^N=(ϵQ)^{2(N2)}$$
(4.40)
where $`\xi =\xi _1\xi _2=e^{i(\theta _1+\theta _2)}`$, $`\theta _1+\theta _2`$ being the total momenta.
Now we include the new states (4.32) into the definition of $`\mathrm{\Psi }_{4\omega }`$ in order to extend (4.20) to
$$\mathrm{\Psi }_{4\omega }=\underset{k_1<k_2}{}A(k_1,k_2)|\mathrm{\Omega }(k_1,k_2.$$
(4.41)
Here we have used the same notation for separated and neighboring states.
Substituting (4.26) and (4.35) into the eigenvalue equation, we get
$$(E_{4\omega }QQ^1)A(k,k+1)=ϵA(k1,k+1)+ϵA(k,k+2)$$
(4.42)
which gives us the phase shift produced by the interchange of the two pseudoparticles
$$\frac{A_{21}}{A_{12}}=\frac{1+\xi +ϵ(Q+Q^1)\xi _2}{1+\xi +ϵ(Q+Q^1)\xi _1}.$$
(4.43)
We thus arrive to the Bethe ansatz equations which fix the values of $`\xi _1`$ and $`\xi _2`$:
$`\xi _2^N`$ $`=`$ $`(ϵQ)^{N2}\left\{{\displaystyle \frac{1+\xi +ϵ(Q+Q^1)\xi _2}{1+\xi +ϵ(Q+Q^1)\xi _1}}\right\},`$
$`\xi _1^N\xi _2^N`$ $`=`$ $`(ϵQ)^{2(N2)}`$ (4.44)
### 4.3 General eigenstates
Thus in the sector $`r=2\omega p`$, we expect that the $`p`$-pseudoparticle phase shift will be a sum of two-pseudoparticle phase shifts and the energy is given by
$$E_{p(2s)}=\underset{n=1}{\overset{p}{}}\left\{Q+Q^1+ϵ(\xi _n+\xi _n^1)\right\}$$
(4.45)
where
$$\xi _a^N=(ϵ_sQ)^{N2p+2}\underset{ba}{\overset{p}{}}\left\{\frac{1+\xi _a\xi _b+ϵ(Q+Q^1)\xi _a}{1+\xi _a\xi _b+ϵ(Q+Q^1)\xi _b}\right\},a=1,\mathrm{},p$$
$$\left(\xi _1\xi _2\mathrm{}\xi _p\right)^N=(ϵQ)^{p(N2p+2)}$$
(4.46)
The corresponding eigenstates are
$$\mathrm{\Psi }_r(\xi _1,\xi _2,\mathrm{}\xi _p)=\underset{1k_1<\mathrm{}<k_pN1}{}A(k_1,k_{2,}\mathrm{},k_p)|\mathrm{\Omega }(k_1,k_2,\mathrm{},k_p)$$
(4.47)
where $`|\mathrm{\Omega }(k_1,k_2,\mathrm{},k_p)=_{i=1}^p|\mathrm{\Omega }(k_i)`$ and the wavefunctions satisfy the following boundary conditions
$$A(k_1,k_{2,}\mathrm{},k_p,N+k_1)=(ϵQ)^{N2p+2}A(k_1,k_{2,}\mathrm{},k_p)$$
(4.48)
It is not all, in a sector $`r`$ we may have $`p`$ pseudoparticle and $`N_{\omega 1},N_{\omega 2},\mathrm{},N_{\omega +1}`$ impurities of the type $`(\omega 1),(\omega 2),\mathrm{},(\omega +1)`$, respectively, such that
$$N_{\omega 1}+2N_{\omega 2}+\mathrm{}+(2\omega 1)N_{\omega +1}=r2\omega p$$
(4.49)
We called impurity a state $`|\alpha ,k`$ flanked by at least two states $`|\beta ,k\pm 1`$ such that $`\alpha +\beta 0`$. Since $``$ is a sum of projectors on spin zero, these states are also annihilated by $`𝒰_0`$ . Therefore the impurities play here the same role as in the periodic case. It means that for a sector $`r`$ with $`l`$ impurities with parameters $`\xi _1,\mathrm{},\xi _l`$ and $`p`$ pseudoparticles with parameters $`\xi _{l+1},\mathrm{},\xi _{l+p}`$ the energy is given by (4.46), and the Bethe equations do not depend on impurity type and are given by
$$\xi _a^N\xi _1^2\xi _2^2\mathrm{}\xi _l^2=(ϵQ)^{N2p+2}\underset{ba=l+1}{\overset{l+p}{}}\left\{\frac{1+\xi _a\xi _b+ϵ(Q+Q^1)\xi _a}{1+\xi _a\xi _b+ϵ(Q+Q^1)\xi _b}\right\}$$
(4.50)
with $`a=l+1,l+2,\mathrm{},l+p,p1`$, and
$$\xi ^{2p}(\xi _{l+1}\mathrm{}\xi _{l+p})^{N2p}=(ϵQ)^{p(N2p+2)}$$
(4.51)
where $`\xi =\xi _1\xi _2\mathrm{}\xi _l\xi _{l+1}\mathrm{}\xi _{l+p}`$.
We have shown that these closed Temperley-Lieb quantum invariant spin chains can be solved by the coordinate Bethe ansatz. A consequence of the nonlocal terms $`𝒰_0`$ is the arising of boundary conditions depending on the quantum group parameter $`q`$ via the relation $`Q+Q^1=`$Tr$`_{V_\mathrm{\Lambda }}`$ $`(q^{2\rho }).`$ It is also $`p`$-pseudoparticle dependent (which is equal to spin sector $`r`$ for $`A_1`$, when $`s=1/2)`$.
For the algebra $`A_1`$, and $`s=1/2`$ , Q=q and $`U_k`$ are 4x4 matrices giving a nearest-neighbour interaction $`U_k=q\sigma _k^+\sigma _{k+1}^{}\sigma _k^{}\sigma _{k+1}^+(q+q^1)/4(\sigma _k^z\sigma _{k+1}^z+1)+(q+q^1)/4(\sigma _k^z\sigma _{k+1}^z2).`$ This Hamiltonian was investigated in Ref. and we summarize some basic results regarding this case. Assuming q=$`\mathrm{exp}(i\phi )`$ some interesting properties were found. For instance, the spin $`L`$ of the ground state becomes $`\phi `$ \- dependent. For any N (even), $`L`$ depends on the value of $`\phi `$ according to :
$`L`$ $`=`$ $`0\text{ }\text{for}\text{ }{\displaystyle \frac{\pi }{2}}<\phi <\pi `$
$`L`$ $`=`$ $`l\text{ }\text{for}\text{ }{\displaystyle \frac{\pi }{2(l+1)}}<\phi <{\displaystyle \frac{\pi }{2l}}`$
$`L`$ $`=`$ $`{\displaystyle \frac{N}{2}}\text{ }\text{for}\text{ }0<\phi <{\displaystyle \frac{\pi }{N}}`$
The ground state is non-degenerate (up to the trivial SU<sub>q</sub> degeneracy). At the edges of the intervals, $`\phi `$ = $`\pi /2l`$, additional degeneracies occur. These transitions to higher spins resemble the incommensurate transition obtained in various other models.
From the statistical mechanics point of view the Hamiltonian presents critical behavior and it is conformal invariant. The central charge ( or conformal anomaly ) is :
$$c=1\frac{6(\pi \phi )^2}{\pi \phi },\phi [\pi /2,\pi ].$$
In particular, if we choose the rational form :
$$\phi =\frac{\pi m}{(m+1)},\text{ }m=3,4\mathrm{},$$
then
$$c=1\frac{6}{m(m+1)},$$
which give us the conformal anomalies of the minimal unitary models.
The connection between the Hamiltonian of the closed SU<sub>q</sub>(2) invariant chain and the unitary minimal series was explored in . For a generic irrational $`\phi `$ one can decompose the space of states into the direct sum of irreducible representations of the quantum group which are in one-to-one correspondence with the usual SU(2) representations.
We have constructed and diagonalized numerically the Hamiltonian for small values of N and s=1/2, 1. We checked that, for a given N , s=1/2 and 1 have the same spectra up to degeneracies.
## 5 Bethe Ansatz: Free boundary conditions
It is for free boundary conditions that the Hamiltonian (2.7) naturally commutes with the quantum group $`𝒰_q(X_n)`$. Since the our linear combination (3.2) left all models with the same status, which concern to the coordinate Bethe ansatz, we expect that all procedure developed for the coordinate Bethe ansatz with free boundary conditions in for the case $`A_1`$ ($`s=1/2`$). can be used here. To show this we recall the previous section, taking into account $`𝒰_0=0`$, where almost all equations can be seized for the free boundary conditions eigenvalue problem.
### 5.1 One-pseudoparticle
In this sector, the eigenstate is given by (4.10):
$$\mathrm{\Psi }_{2\omega }(\xi )=\underset{k=1}{\overset{N1}{}}A(k)|\mathrm{\Omega }(k)$$
(5.1)
where $`|\mathrm{\Omega }(k)`$ is again given by (4.4).
The action of $``$ on the states $`|\mathrm{\Omega }(k)`$ is given by (4.11), which gives us the following eigenvalue equations
$$(E_{2\omega }QQ^1)A(k)=ϵA(k1)+ϵA(k+1),2kN2$$
(5.2)
At the boundaries, we get more two slightly different equations
$`(E_{2\omega }QQ^1)A(1)=ϵA(2)`$
$`(E_{2\omega }QQ^1)A(N1)=ϵA(N2)`$ (5.3)
where $`ϵ=1`$ for $`B_n`$, $`D_n`$ and $`A_1`$ ( $`s`$ integer) and $`ϵ=1`$ for $`C_n`$ and $`A_1`$ ( $`s`$ semi-integer). We now try as a solution
$$A(k)=\mathrm{A}(\theta )\xi ^k\mathrm{A}(\theta )\xi ^k$$
(5.4)
where $`\xi =e^{i\theta }`$, $`\theta `$ being the momenta. Substituting this in equation (5.2) we obtain the energy eigenvalue associated with a free pseudoparticle with free boundary conditions
$$E_{2\omega }=Q+Q^1+ϵ\left(\xi +\xi ^1\right)$$
(5.5)
We want equations (5.2) to be valid for $`k=1`$ and $`k=N1`$ also, where $`A(0)`$ and $`A(N)`$ are defined by (5.4). Matching (5.2) and (5.3) we get the end conditions
$$A(0)=0\mathrm{and}A(N)=0$$
(5.6)
implying that $`A(\theta )=A(\theta )`$ and $`\xi ^{2N}=1`$, respectively. $`A(\theta )`$ it now determined ( up to a factor that is invariant under $`\theta \theta `$), to be equal to $`\xi ^N`$.
### 5.2 One pseudoparticle and impurities
Differently from the previous cases, due to the lack of periodicity, the impurity positions are fixed. So, they have a different role in the eigenvalue problem with free boundary conditions. For instance, let us consider the case of one impurity of the type $`\omega 1`$, with parameter $`\xi _1`$ and one pseudoparticle with parameter $`\xi _2`$. This eigenstate lies in the sector $`r=2\omega +1`$ and we can write
$$\mathrm{\Psi }_{2\omega +1}(\xi _1,\xi _2)=\underset{k_1<k_2}{}\left\{A_1(k_1,k_2)|\mathrm{\Omega }_1(k_1,k_2)+A_2(k_1,k_2)|\mathrm{\Omega }_2(k_1,k_2)\right\}$$
(5.7)
where $`|\mathrm{\Omega }_i(k_1,k_2),i=1,2`$ are given by (LABEL:pba10).
For this case we obtain the following eigenvalue equations
$`(E_{2\omega +1}QQ^1)A_1(k_1,k_2)`$ $`=`$ $`ϵA_1(k_11,k_2)+ϵA_1(k_1+1,k_2)`$
$`(E_{2\omega +1}QQ^1)A_2(k_1,k_2)`$ $`=`$ $`ϵA_2(k_1,k_21)+ϵA_2(k_1,k_2+1)`$ (5.8)
We also have two meeting conditions that arise because pseudoparticle and impurity may be neighbors (see (3.21))
$$A_1(k,k)=A_2(k,k+2),A_2(k+1,k+2)=A_1(k,k+1)$$
(5.9)
in addition to the two conditions to be satisfied at the free ends
$$A_1(k_1,N)=0,A_2(0,k_2)=0$$
(5.10)
Now we try the following ansatz for the wavefunctions
$`A_1(k_1,k_2)`$ $`=`$ $`\mathrm{A}_1(\theta _1,\theta _2)\xi _1^{k_1}\xi _2^{k_2}\mathrm{A}_1(\theta _1,\theta _2)\xi _1^{k_1}\xi _2^{k_2}`$
$`A_2(k_1,k_2)`$ $`=`$ $`\mathrm{A}_2(\theta _1,\theta _2)\xi _2^{k_1}\xi _1^{k_2}\mathrm{A}_2(\theta _1,\theta _2)\xi _2^{k_1}\xi _1^{k_2}`$ (5.11)
From (5.8) we get the energy eigenvalue
$$E_{2\omega +1}=Q+Q^1+ϵ\left(\xi _2+\xi _2^1\right)$$
(5.12)
and from (5.9) and (5.10) the following relations between the coefficients $`A_i`$
$`\mathrm{A}_1(\theta _1,\theta _2)\xi _2^N`$ $`=`$ $`\mathrm{A}_1(\theta _1,\theta _2)\xi _2^N,\mathrm{A}_2(\theta _1,\theta _2)=\mathrm{A}_2(\theta _1,\theta _2)`$
$`\mathrm{A}_1(\theta _1,\theta _2)`$ $`=`$ $`\mathrm{A}_2(\theta _1,\theta _2)\xi _1^2,\mathrm{A}_1(\theta _1,\theta _2)=\mathrm{A}_2(\theta _1,\theta _2)\xi _1^2`$ (5.13)
from this we get
$$\xi _2^{2N}=1$$
(5.14)
as the Bethe equation of (5.12). The coefficients $`A_i`$ are determined up to a factor that is invariant under $`\theta _2\theta _2`$ as:
$$\mathrm{A}_1(\theta _1,\theta _2)=\xi _1^2\xi _2^N\mathrm{and}\mathrm{A}_2(\theta _1,\theta _2)=\xi _2^N.$$
(5.15)
In general, for the eigenstate with $`l`$ impurities with parameters $`\xi _1,\mathrm{},\xi _l`$ and one pseudoparticle with parameter $`\xi _{l+1}`$, which lies in a sector $`r`$, we can write
$$\mathrm{\Psi }_r(\xi _1,\mathrm{},\xi _{l+1})=\underset{j=1}{\overset{l+1}{}}\left\{\underset{1k_1<\mathrm{}<k_{l+1}N1}{}A_j(k_1,\mathrm{},k_{l+1})|\mathrm{\Omega }_j(k_1,\mathrm{},k_{l+1})\right\}$$
(5.16)
The corresponding eigenvalue is given by (5.5) , with $`\xi =\xi _{l+1}`$, and the ansatz for the coefficients of the wavefunctions becomes
$$\mathrm{A}_j(\theta _1,\mathrm{},\theta _{l+1})=\left(\underset{i=1}{\overset{l+1j}{}}\xi _i^2\right)\xi _{l+1}^N$$
(5.17)
Here we notice that the index $`j`$ in the wavefunctions $`A_j(k_1,\mathrm{},k_{l+1})`$ means that the pseudoparticle is at the position $`k_{l+2j}`$.
### 5.3 Two-pseudoparticles
For the sector $`r=4\omega `$, beside eigenstates with impurities, we have an eigenstate with two pseudoparticles. We obtain the following eigenvalue equations
$`(E_{4\omega }2Q2Q^1)A(k_1,k_2)=ϵA(k_11,k_2)+ϵA(k_1+1,k_2)`$
$`+ϵA(k_1,k_21)+ϵA(k_1,k_2+1)`$ (5.18)
We have again two conditions to be satisfied at the ends of the chain
$$A(0,k_2)=0\mathrm{and}A(k_1,N)=0$$
(5.19)
In addition to this we have a meeting condition
$$ϵA(k,k)+ϵA(k+1,k+1)+(Q+Q^1)A(k,k+1)=0$$
(5.20)
Now we try the ansatz
$$\begin{array}{ccc}A(k_1,k_2)\hfill & =\hfill & \mathrm{A}(\theta _1,\theta _2)\xi _1^{k_1}\xi _2^{k_2}\mathrm{A}(\theta _2,\theta _1)\xi _1^{k_2}\xi _2^{k_1}\hfill \\ & \hfill & \mathrm{A}(\theta _1,\theta _2)\xi _1^{k_1}\xi _2^{k_2}+\mathrm{A}(\theta _2,\theta _1)\xi _1^{k_2}\xi _2^{k_1}\hfill \\ & \hfill & \mathrm{A}(\theta _1,\theta _2)\xi _1^{k_1}\xi _2^{k_2}+\mathrm{A}(\theta _2,\theta _1)\xi _1^{k_2}\xi _2^{k_1}\hfill \\ & +\hfill & \mathrm{A}(\theta _1,\theta _2)\xi _1^{k_1}\xi _2^{k_2}\mathrm{A}(\theta _2,\theta _1)\xi _1^{k_2}\xi _2^{k_1}\hfill \end{array}$$
(5.21)
Here we observe the permutations and negations of $`\theta _1`$ and $`\theta _2`$. Substituting this ansatz in (5.18) we obtain the energy eigenvalue for the sector with two pseudoparticles
$$E_{4\omega }=2Q+2Q^1+ϵ\left(\xi _1+\xi _1^1+\xi _2+\xi _2^1\right)$$
(5.22)
The ansatz (5.21) satisfy equations (5.19) provided that
$`\mathrm{A}(\theta _1,\theta _2)=\mathrm{A}(\theta _1,\theta _2),\mathrm{A}(\theta _2,\theta _1)=\mathrm{A}(\theta _2,\theta _1)`$
$`\mathrm{A}(\theta _1,\theta _2)=\mathrm{A}(\theta _1,\theta _2),\mathrm{A}(\theta _2,\theta _1)=\mathrm{A}(\theta _2,\theta _1)`$ (5.23)
and
$$\xi _2^{2N}=\frac{\mathrm{A}(\theta _1,\theta _2)}{\mathrm{A}(\theta _1,\theta _2)}=\frac{\mathrm{A}(\theta _1,\theta _2)}{\mathrm{A}(\theta _1,\theta _2)},\xi _1^{2N}=\frac{\mathrm{A}(\theta _2,\theta _1)}{\mathrm{A}(\theta _2,\theta _1)}=\frac{\mathrm{A}(\theta _2,\theta _1)}{\mathrm{A}(\theta _2,\theta _1)}$$
(5.24)
Moreover, the meeting conditions are satisfied provided that
$`{\displaystyle \frac{\mathrm{A}(\theta _1,\theta _2)}{\mathrm{A}(\theta _2,\theta _1)}}={\displaystyle \frac{\mathrm{A}(\theta _2,\theta _1)}{\mathrm{A}(\theta _1,\theta _2)}}={\displaystyle \frac{1+\xi _1\xi _2+ϵ\left(Q+Q^1\right)\xi _2}{1+\xi _1\xi _2+ϵ\left(Q+Q^1\right)\xi _1}}`$
$`{\displaystyle \frac{\mathrm{A}(\theta _1,\theta _2)}{\mathrm{A}(\theta _2,\theta _1)}}={\displaystyle \frac{\mathrm{A}(\theta _2,\theta _1)}{\mathrm{A}(\theta _1,\theta _2)}}={\displaystyle \frac{1+\xi _1^1\xi _2+ϵ\left(Q+Q^1\right)\xi _2}{1+\xi _1^1\xi _2+ϵ\left(Q+Q^1\right)\xi _1^1}}`$ (5.25)
Matching these conditions we get
$$\xi _1^{2N}=\frac{B(\theta _1,\theta _2)}{B(\theta _1,\theta _2)},\xi _2^{2N}=\frac{B(\theta _2,\theta _1)}{B(\theta _2,\theta _1)}$$
(5.26)
and
$$\mathrm{A}(\theta _1,\theta _2)=\xi _1^N\xi _2^NB(\theta _1,\theta _2)\xi _2^1.$$
(5.27)
Here we have used the usual free boundary notations
$$B(\theta _a,\theta _b)=s(\theta _a,\theta _b)s(\theta _b,\theta _a)$$
(5.28)
where
$$s(\theta _a,\theta _b)=1+\xi _a\xi _b+ϵ\left(Q+Q^1\right)\xi _b.$$
(5.29)
Now let us consider the eigenstates with two pseudoparticle and impurities. The energy eigenvalue is the same of the two pseudoparticles pure state. The parameters associated with impurities are embraced in the definition of the coefficients of the wavefunctions. For instance, when we have an eigenstate of two pseudoparticles with parameters $`\xi _2`$ and $`\xi _3`$ and one impurity of parameter $`\xi _1`$, the energy is given by (5.22) and the Bethe equations by (5.26), with $`\xi _1\xi _3`$ and $`\theta _1\theta _3`$. But now the wavefunctions are different
$`\mathrm{A}_1(\theta _1,\theta _2,\theta _3)=\left\{\xi _1^4\right\}\xi _2^N\xi _3^NB(\theta _2,\theta _3)\xi _3^1`$
$`\mathrm{A}_2(\theta _1,\theta _2,\theta _3)=\left\{\xi _1^2\right\}\xi _2^N\xi _3^NB(\theta _2,\theta _3)\xi _3^1`$
$`\mathrm{A}_3(\theta _1,\theta _2,\theta _3)=\xi _2^N\xi _3^NB(\theta _2,\theta _3)\xi _3^1`$ (5.30)
where $`B(\theta _2,\theta _3)`$ is given by (5.28)
### 5.4 General eigenstates
The generalization follows as in the previous cases. In a sector $`r`$ with
$`p`$ pseudoparticles, we get
$$E_r=\underset{n=1}{\overset{p}{}}\left[Q+Q^1+ϵ\left(\xi _n+\xi _n^1\right)\right]$$
(5.31)
and the Bethe equations
$$\xi _a^{2N}=\underset{ba=l+1}{\overset{l+p}{}}\frac{B(\theta _a,\theta _b)}{B(\theta _a,\theta _b)},a=1,2,\mathrm{},p$$
(5.32)
The corresponding eigenfunction can be written as
$$\mathrm{\Psi }_r(\xi _1,\mathrm{},\xi _p)=\underset{k_1<\mathrm{}<k_{l+p}}{}A(k_1,k_2,\mathrm{},k_{l+p})|\mathrm{\Omega }(k_1,k_2,\mathrm{},k_p)$$
(5.33)
with
$$A(k_1,k_2,\mathrm{},k_p)=\underset{P}{}\epsilon _P\mathrm{A}(\theta _1,\theta _2,\mathrm{},\theta _p)\xi _1^{k_1}\xi _2^{k_2}\mathrm{}\xi _p^{k_p}$$
(5.34)
where the sum extends over all permutations and negations of $`\theta _1,\mathrm{},\theta _p`$ and $`\epsilon _P`$ changes sign at each such interchange. The coefficients in the wavefunctions are given by
$$\mathrm{A}(\theta _1,\theta _2,\mathrm{},\theta _p)=\underset{j=1}{\overset{p}{}}\xi _j^N\underset{l+1j<il+p}{}B(\theta _j,\theta _i)\xi _j^1$$
(5.35)
where $`B(\theta _j,\theta _i)`$ are defined in (5.28).
For a sector $`r`$ with $`l`$ impurities with parameters $`\xi _1,\mathrm{},\xi _l`$ and $`p`$ pseudoparticles with parameters $`\xi _{l+1},\mathrm{},\xi _{l+p}`$ the energy is given by (5.31) and the Bethe equations by (5.32). Only the coefficients of the wave functions are modified
$$\mathrm{A}_j(\theta _1,\theta _2,\mathrm{},\theta _{l+p})=A_j(\xi _1\xi _2\mathrm{}\xi _l)\mathrm{A}(\theta _{l+1},\theta _2,\mathrm{},\theta _{l+p}).$$
(5.36)
The functions $`A_j(\xi _1\xi _2\mathrm{}\xi _l)=\xi _1^{a_1}\xi _2^{a_2}\mathrm{}\xi _l^{a_l}`$ where the index $`j`$ characterizes the possible configurations of $`l`$ impurities relative to the $`p`$ pseudoparticles. Here $`a_i`$ are numbers which depend on the position of corresponding impurity relative to the pseudoparticles.
Here we observe again the valid of these results for all Temperley-Lieb spin chain Hamiltonians defined as projector of spin zero on the representations of the quantum groups $`𝒰_q(X_n)`$, characterized by the values of $`Q+Q^1=\mathrm{Tr}_{V_\mathrm{\Lambda }}(q^{2\rho }).`$
## 6 Conclusion
We have applied in a systematic way the coordinate Bethe ansatz to find the spectra of a series of ”spin ” Hamiltonians arising as representations of the Temperley–Lieb algebra. We consider several boundary conditions in order to include all previously known cases.
Due to $`𝒰_k`$ be a projector of spin zero, there is a linear combination of eigenstates of $`𝐒_T^Z=_k𝐒_k^Z`$ where $`𝐒_k^Z=\mathrm{diag}(\omega ,\mathrm{},\omega ),\omega =\mathrm{max}(J)`$, which beside simplify the calculus permits us a unified treatment for all models. We find that for a given set of boundary conditions, all models have equivalent spectra, i.e. they differ at most in their degeneracies. Moreover, for the closed cases, the spectra of the lower-dimensional representations are entirely contained in the higher-dimensional ones (see Eq.(3.23)).
Here we notice that this spectrum equivalence is, of course, a consequence of the TL algebra. Nevertheless there is in the literature a large class of Hamiltonians which are not derived from representations of the TL algebra which share the same property. The authors of reference developed a technique for construction of spin chain Hamiltonians which affine quantum group symmetry whose spectra coincides with the spectra of spin chain Hamiltonians which have non-affine quantum group symmetry.
The energy eigenvalues are given by
$`E`$ $`=`$ $`{\displaystyle \underset{n=1}{\overset{p}{}}}\left(Q+Q^1+2ϵ\mathrm{cos}\theta _n\right)`$
$`Q+Q^1`$ $`=`$ $`\mathrm{Tr}_{V_\mathrm{\Lambda }}(q^{2\rho })`$ (6.1)
where $`\rho `$ is half the sum of the positive roots of $`𝒰_q(X_n),X_n=A_1,B_n,C_n`$and$`D_n`$. $`\theta _n`$ are solutions of the Bethe ansatz equations (3.73), (4.50) and (5.32).
The Hamiltonians for the cases $`X_n=B_n,C_n`$and$`D_n`$ appear to be new although due to the Temperley–Lieb equivalence , they are expected to possess the same thermodynamic properties as the $`A_1(s=1/2)`$ case, i.e., the spin-$`1/2`$ XXZ chain with appropriate coupling.
There are several issues left for future work. In particular, one would like clarify from an algebraic point of view, the equality of the spectra, for instance, of the biquadratic model ($`A_1,s=1`$) and XXZ model ($`A_1,s=1/2`$ ) for free boundary conditions and the inclusion of the XXZ spectrum in the one of the biquadratic Hamiltonian for periodic boundary conditions. Furthermore, the completeness and complete characterization as highest weight states of the Bethe ansatz eigenstates here presented are not considered.
Using our solutions, one can derive partition functions in the finite–size scaling limit and find the operator content of the systems constructed from these quantum chains.
Finally, we remark here that although the Hamiltonian (4.1) is a global operator, it manifests the property of essential locality . From the physical point of view, this type of models exhibit behavior similar to closed chains with twisted boundary conditions, however now the boundary conditions become sector dependent.
Acknowledgments: We would like to thank A. Lima–Santos for helpful comments.This work was supported in part by CAPES-PICD, Conselho Nacional de Desenvolvimento Científico e Tecnológico (CNPq-Brazil) (RCTG), and by Fundação de Amparo à Pesquisa do Estado de São Paulo (FAPESP-Brazil) (ALM).
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# Concentration Dependent Sedimentation of Colloidal Rods
## I Introduction
There is extensive literature concerned with sedimentation behaviour of spherically shaped colloidal particles (for a review see ref. ). Essentially exact predictions can be made for the sedimentation velocity of spherical colloids to first order in concentration . For non-spherical colloids a similar exact prediction is non-existent. The only attempt to calculate the first order concentration dependence of the sedimentation velocity for rod like colloids we are aware of is due to Peterson . This theory is based on approximate, orientationally pre-averaged hydrodynamic interactions between the colloidal rods and a rather crude estimate of certain multiple integrals that represent the ensemble averaged velocity. As yet there are no accurate expressions for hydrodynamic interaction tensors for rods. In the first part of the present paper, in section II, we calculate these interaction tensors in a mean-field approximation. In section III we use this approximate expression for the hydrodynamic interaction functions to derive an explicit expression for the first order in concentration coefficient of the sedimentation velocity as a function of the aspect ratio of the rods. This expression is found to agree remarkably well with Peterson’s result for aspect ratios less than about $`30`$. For larger aspect ratios our result for the first order in concentration coefficient is much larger then Peterson’s prediction. In the second part of this paper, section IV, sedimentation experiments on fd-virus are discussed. Experiments are done at low concentration to find the first order concentration dependence, which is compared to the theory mentioned above. In addition, sedimentation experiments at larger concentrations, including the nematic phase are performed.
## II Hydrodynamic Interaction between Long and Thin Rods
In order to calculate sedimentation velocities, the connection between translational and angular velocities, and hydrodynamic forces and torques must be found. In the present section such a relation will be established for two rods on the Rodne-Prager level, that is, with the neglect of reflection contributions between the rods. Considering only two rods limits the discussion on the sedimentation velocity to first order in concentration. Reflection contributions to the two-rod hydrodynamic interaction functions and multi-body rod interactions are both probably small in comparison to the Rodne-Prager terms, due to the fact that the distance between segments of different rods is of the order of the length of the rods, at least in the isotropic state. A Rodne-Prager approximation could therefore work quite well for long and thin rods, although explicit results for reflection contributions should be obtained to confirm this intuition.
For the low Reynolds numbers under consideration, the translational velocities $`𝐯_j`$, $`j=1,2`$, and the angular velocities $`𝛀_j`$ are linearly related to the hydrodynamic forces $`𝐅_j^h`$ and torques $`𝐓_j^h`$ that the fluid exerts on the rods,
$`\left(\begin{array}{c}𝐯_1\\ 𝐯_2\\ 𝛀_1\\ 𝛀_2\end{array}\right)=\left(\begin{array}{cccc}𝐌_{11}^{TT}& 𝐌_{12}^{TT}& 𝐌_{11}^{TR}& 𝐌_{12}^{TR}\\ 𝐌_{21}^{TT}& 𝐌_{22}^{TT}& 𝐌_{21}^{TR}& 𝐌_{22}^{TR}\\ 𝐌_{11}^{RT}& 𝐌_{12}^{RT}& 𝐌_{11}^{RR}& 𝐌_{12}^{RR}\\ 𝐌_{21}^{RT}& 𝐌_{22}^{RT}& 𝐌_{21}^{RR}& 𝐌_{22}^{RR}\end{array}\right)\left(\begin{array}{c}𝐅_1^h\\ 𝐅_2^h\\ 𝐓_1^h\\ 𝐓_2^h\end{array}\right).`$ (13)
The superscripts “T” and “R” refer to translation and rotation, respectively, while the superscript “h” on the forces and torques refer to their hydrodynamic origin. On the Brownian time scale, the $`3\times 3`$-dimensional mobility matrices $`𝐌`$ are functions of the positions of the centers of the two rods and their orientation.
As will turn out, in order to find the sedimentation velocity, we need expressions for $`𝐌_{1j}^{TT}`$, for which approximations are obtained in subsection II B. As a first step, the fluid flow field generated by a translating rod must be calculated. This is the subject of subsection II A. Rotation of rods also plays a role in sedimentation, but as will turn out, to first order in concentration and with the neglect of hydrodynamic reflection contributions, these do not contribute to the sedimentation velocity. Explicit expressions pertaining to the hydrodynamics of rotating rods are derived in the same spirit as for translating rods in appendix A. Subsection II C contains some concluding remarks.
For the hydrodynamic calculations the rods will be thought of as a rigid string of spherical beads with diameter $`D`$. The length of the rods is $`L`$, and there are $`n+1=L/D`$ beads per rod, with $`n`$ an even integer.
### A Flow field generated by a translating rod
The flow field generated by a rod that consists of $`n+1`$ beads is given by,
$`𝐮(𝐫)={\displaystyle \underset{j=n/2}{\overset{n/2}{}}}{\displaystyle _{V_j}}𝑑S^{}𝐓(𝐫𝐫^{})𝐟_j(𝐫^{}),`$ (14)
where $`𝐓`$ is the Oseen tensor,
$`𝐓(𝐫)={\displaystyle \frac{1}{8\pi \eta _0r}}\left[\widehat{𝐈}+\widehat{𝐫}\widehat{𝐫}\right],`$ (15)
with $`\eta _0`$ the shear viscosity of the solvent and $`\widehat{𝐫}=𝐫/r`$ the unit position vector. Furthermore, $`𝐟_j(𝐫^{})`$ is the force per unit area that a surface element at $`𝐫^{}`$ of bead $`j`$ exerts on the fluid, and $`V_j`$ is the spherical surface of bead $`j`$. For long and thin rods, the distances $`𝐫`$ of interest, relative to the positions of the beads, are those for which $`rD`$, with $`D`$ the diameter of the beads. Now write $`𝐫^{}=𝐫_j+𝐑^{}`$ with $`𝐫_j`$ the position coordinate of the $`j^{th}`$ bead, so that $`R^{}=D/2`$, and Taylor expand the Oseen tensor in eq.(2) with respect to $`𝐑^{}`$. Keeping only the first term in this Taylor expansion leads to relative errors of the order $`R^{}/rD/L`$. Up to that order we then find,
$`𝐮(𝐫)={\displaystyle \underset{j=n/2}{\overset{n/2}{}}}𝐓(𝐫𝐫_j)𝐅_j^h,`$ (16)
with,
$`𝐅_j^h={\displaystyle 𝑑S^{}𝐟_j(𝐫^{})},`$ (17)
the total force that the fluid exerts on bead $`j`$. With the neglect of end-effects this force is equal for each bead, $`𝐅_j𝐅^h/(n+1)=\frac{D}{L}𝐅^h`$, with $`𝐅^h`$ the total force on the rod. Eq.(4) thus reduces to,
$`𝐮(𝐫)={\displaystyle \frac{D}{L}}{\displaystyle \underset{j=n/2}{\overset{n/2}{}}}𝐓(𝐫𝐫_j)𝐅^h.`$ (18)
The force $`𝐅^h`$ is calculated in terms of the translational velocity of the rod self-consistently from eq.(6) using Faxén’s theorem for translational motion for each spherical bead, where the velocity $`𝐯_j`$ of bead $`j`$ is expressed in terms of the force $`𝐅_j^h`$ on bead $`j`$ and the velocity $`𝐮_0(𝐫_j)`$ at the center of the bead that would have existed without that bead being present,
$`𝐯_j={\displaystyle \frac{1}{3\pi \eta _0D}}𝐅_j^h+𝐮_0(𝐫_j)+{\displaystyle \frac{1}{24}}D^2_j^2𝐮_0(𝐫_j),`$ (19)
where $`_j`$ is the gradient operator with respect to $`𝐫_j`$. The first term on the right hand-side is just Stokes friction of a single bead in an unbounded fluid, while the second term accounts for hydrodynamic interaction between the beads. The fluid flow field $`𝐮_0`$ in turn is equal to,
$`𝐮_0(𝐫)={\displaystyle \underset{i=n/2,ij}{\overset{n/2}{}}}{\displaystyle _{V_i}}𝑑S^{}𝐓(𝐫𝐫^{})𝐟_i^{}(𝐫^{}),`$ (20)
where $`𝐟_i^{}`$ is the force per unit area that a surface element of bead $`i`$ exerts on the fluid in the absence of bead $`j`$. For very long rods, consisting of many beads, the difference between $`𝐟_i`$ (the corresponding force for the intact rod) and $`𝐟_i^{}`$ may be neglected: there are only a few neighbouring beads for which the difference is significant, but there are many more beads further away from bead $`j`$ for which the difference is insignificant. To within the same approximations involved to arrive at eq.(4), eq.(8) can then be written as,
$`𝐮_0(𝐫_j)={\displaystyle \underset{i=n/2,ij}{\overset{n/2}{}}}𝐓(𝐫_j𝐫_i)𝐅_i^h.`$ (21)
Substitution of this expression into Faxén’s theorem (7), and using that $`𝐫_j𝐫_i=(ji)D\widehat{𝐮}`$, with $`\widehat{𝐮}`$ the orientation of the rod, leads to,
$`𝐯_j`$ $`=`$ $`{\displaystyle \frac{1}{3\pi \eta _0D}}𝐅_j^h{\displaystyle \frac{1}{8\pi \eta _0D}}\widehat{𝐮}\widehat{𝐮}{\displaystyle \underset{i=n/2,ij}{\overset{n/2}{}}}\left[{\displaystyle \frac{2}{ij}}{\displaystyle \frac{1}{6ij^3}}\right]𝐅_i^h`$ (23)
$`{\displaystyle \frac{1}{8\pi \eta _0D}}\left[\widehat{𝐈}\widehat{𝐮}\widehat{𝐮}\right]{\displaystyle \underset{i=n/2,ij}{\overset{n/2}{}}}\left[{\displaystyle \frac{1}{ij}}+{\displaystyle \frac{1}{12ij^3}}\right]𝐅_i^h,`$
where eq.(3) has been used, together with,
$`^2𝐓(𝐫)={\displaystyle \frac{1}{4\pi \eta _0r^3}}\left[\widehat{𝐈}3\widehat{𝐫}\widehat{𝐫}\right].`$ (24)
For pure translational motion, the velocity $`𝐯_j`$ of each bead is equal to the velocity $`𝐯`$ of the rod, so that both sides of eq.(10) can be summed over $`j`$, yielding for the left hand-side $`𝐯L/D`$. Neglecting end-effects and replacing sums by integrals (which is allowed for long and thin rods), it is found that,
$`𝐯={\displaystyle \frac{\mathrm{ln}\{L/D\}}{4\pi \eta _0L}}\left[\widehat{𝐈}+\widehat{𝐮}\widehat{𝐮}\right]𝐅^h.`$ (25)
Notice that the Stokes friction contribution (the first term on the right hand-side in eq.(10)) is logarithmically small in comparison to the friction contribution due to hydrodynamic interaction between the beads. In fact, the Stokes contribution is neglected in eq.(12). A matrix inversion, in order to express $`𝐅^h`$ in terms of $`𝐯`$, and subsequent substitution into eq.(6), after rewriting the sum over beads as an integral over the center line of the rod, yields,
$`𝐮(𝐫)={\displaystyle \frac{4\pi \eta _0}{\mathrm{ln}\{L/D\}}}{\displaystyle _{L/2}^{L/2}}𝑑l𝐓(𝐫𝐫_pl\widehat{𝐮})\left[\widehat{𝐈}{\displaystyle \frac{1}{2}}\widehat{𝐮}\widehat{𝐮}\right]𝐯,`$ (26)
with $`𝐫_p`$ the position coordinate of the rod. This is the approximate expression for the fluid flow generated by a translating long and thin rod that will be used in the following subsection to obtain an expression for the mobility matrices $`𝐌_{1j}^{TT}`$, $`j=1,2`$.
### B Calculation of $`𝐌^{TT}`$
In order to calculate the velocity $`𝐯_2`$ that rod $`2`$ acquires in the flow field (13) generated by a translating rod $`1`$, one should in principle perform a reflection calculation up to very high order : the field generated by rod $`1`$ is scattered by each bead of rod $`2`$ and subsequently reflected hence and forth between the different beads within rod 2. Such a calculation is hardly feasible analytically. Here, the field generated by rod $`1`$ that is incident on rod $`2`$ is approximated by a constant fluid flow field $`\overline{𝐮}`$ equal to the average of the incident field over the center line of rod $`2`$. This “hydrodynamic mean-field approximation” is accurate for distances of the order $`L`$ or larger, for which separations the incident field indeed becomes equal to a constant. For smaller distances between the rods this procedure provides a semi-quantitative approximation. Within this approximation, the velocity of rod $`2`$ immediately follows from eq.(12), with $`𝐯=𝐯_2\overline{𝐮}`$, $`𝐅^h=𝐅_2^h`$, the total force of the fluid on rod $`2`$, and $`\widehat{𝐮}=\widehat{𝐮}_2`$, the orientation of rod $`2`$,
$`𝐯_2=\overline{𝐮}{\displaystyle \frac{\mathrm{ln}\{L/D\}}{4\pi \eta _0L}}\left[\widehat{𝐈}+\widehat{𝐮}_2\widehat{𝐮}_2\right]𝐅_2^h.`$ (27)
The average incident flow field follows from eqs.(13) and (12), with $`𝐯=𝐯_1`$, the velocity of rod $`1`$ and $`𝐅^h=𝐅_1^h`$, the force on rod $`1`$,
$`\overline{𝐮}`$ $`=`$ $`{\displaystyle \frac{4\pi \eta _0L}{\mathrm{ln}\{L/D\}}}𝐀\left[\widehat{𝐈}{\displaystyle \frac{1}{2}}\widehat{𝐮}_1\widehat{𝐮}_1\right]𝐯_1`$ (28)
$`=`$ $`𝐀\left[\widehat{𝐈}{\displaystyle \frac{1}{2}}\widehat{𝐮}_1\widehat{𝐮}_1\right]\left[\widehat{𝐈}+\widehat{𝐮}_1\widehat{𝐮}_1\right]𝐅_1^h`$ (29)
$`=`$ $`𝐀𝐅_1^h,`$ (30)
where,
$$𝐀=\frac{1}{L^2}_{L/2}^{L/2}𝑑l_1_{L/2}^{L/2}𝑑l_2𝐓(𝐫_{21}+l_2\widehat{𝐮}_2l_1\widehat{𝐮}_1),$$
(31)
with $`𝐫_{21}=𝐫_2𝐫_1`$ the distance between the centers of the two rods. Notice that for distances $`𝐫_{21}`$ between the centers of the rods larger than $`L`$, the matrix $`𝐀`$ asymptotes to $`𝐓(𝐫_{21})`$. By definition the following “mean-field” expressions for the translational mobility matrices are thus obtained (after an interchange of the indices $`1`$ and $`2`$),
$`𝐌_{11}^{TT}`$ $`=`$ $`{\displaystyle \frac{\mathrm{ln}\{L/D\}}{4\pi \eta _0L}}\left[\widehat{𝐈}+\widehat{𝐮}_1\widehat{𝐮}_1\right],`$ (32)
$`𝐌_{12}^{TT}`$ $`=`$ $`{\displaystyle \frac{1}{L^2}}{\displaystyle _{L/2}^{L/2}}𝑑l_1{\displaystyle _{L/2}^{L/2}}𝑑l_2𝐓(𝐫_{12}+l_1\widehat{𝐮}_1l_2\widehat{𝐮}_2).`$ (33)
One might try to device approximate expressions for the matrix $`𝐀`$. However, sedimentation velocities are obtained as ensemble averages, also with respect to orientations, giving rise to integrals with respect to $`𝐫_{12}`$ and $`\widehat{𝐮}_{1,2}`$, which can be evaluated by numerical integration.
### C Concluding Remarks
The approximations involved in the above discussion are justified for very long rods, since $`𝒪`$(1)-constants are neglected against terms of order $`\mathrm{ln}\{L/D\}`$, both by neglecting end-effects and replacing sums over beads by integrals (for the evaluation of the sums in Faxén’s theorem in eq.(10)). Such approximations are most important for the diagonal mobility matrix $`𝐌_{11}^{TT}`$ (notice that factors $`\mathrm{ln}\{L/D\}`$ do not appear in the off-diagonal matrice $`𝐌_{12}^{TT}`$, due to the resubstitution of velocities in terms of forces). Both end-effects and the mathematical approximations involved in the calculation of the diagonal mobility matrix $`𝐌_{11}^{TT}`$ in eq.(17) may be accurately accounted for by the replacement,
$`\mathrm{ln}\{L/D\}\mathrm{ln}\{L/D\}\nu ,`$ (34)
with $`\nu =\nu _{}`$ and $`\nu =\nu _{}`$ a constant, pertaining to translational motion perpendicular and parallel to the rods orientation, respectively. This correction is experimentally significant for somewhat shorter rods ($`L/D<20`$, say), but vanishes relatively to the logarithmic term for very long rods. The actual values of $`\nu _{}`$ and $`\nu _{}`$ for cylindrical rods are equal to ,
$`\nu _{}`$ $`=`$ $`0.84,`$ (35)
$`\nu _{}`$ $`=`$ $`0.21.`$ (36)
A more accurate expression for $`𝐌_{11}^{TT}`$ than in eq.(17) is,
$`𝐌_{11}^{TT}={\displaystyle \frac{\mathrm{ln}\{L/D\}}{4\pi \eta _0L}}\left[\widehat{𝐈}+\widehat{𝐮}_1\widehat{𝐮}_1\right]{\displaystyle \frac{1}{4\pi \eta _0L}}\left[\nu _{}\widehat{𝐈}+(2\nu _{}\nu _{})\widehat{𝐮}_1\widehat{𝐮}_1\right].`$ (37)
In the sequel we will use this expression for the mobility matrix $`𝐌_{11}^{TT}`$ instead of eq.(17). The approximations involved in the off-diagonal mobility matrices in eqs.(34,38) are primarily due to the mean-field treatment of the incident flow field. It is probably a formidable task to improve on these expressions.
The use of the more accurate expression (22) for the diagonal translational mobility matrix also circumvents the practical problem of calculating volume fractions of colloidal rod material from given values for $`L`$, $`D`$ and number density. For the bead model it is not so clear how the volume of a rod must be expressed in terms of $`L`$ and $`D`$. The volume of the cylindrical rod is simply equal to $`\frac{\pi }{4}D^2L`$.
## III An Expression for the Sedimentation Velocity of Rods
The sedimentation of colloidal material induces, through the presence of the walls of the container, backflow of solvent. The backflow velocity is inhomogeneous, and varies on the length scale of the container. On a local scale, however, the backflow may be considered homogeneous, and the sedimentation velocity can be calculated relative to the local backflow velocity. This relative sedimentation velocity is a constant throughout the container (except possibly in a small region of extent $`L`$ near the walls of the container, where gradients of the backflow velocity are large), and depends only on the properties of the suspension. A formal evaluation of the sedimentation velocity directly from eq.(1), by ensemble averaging, leads to spurious divergences, which are the result of the neglect of the hydrodynamic effects of the walls of the container which lead to solvent backflow. Batchelor was the first to deal with these divergences correctly, and we will use his arguments here .
Ensemble averaging of $`𝐯_1`$ in eq.(1) gives the sedimentation velocity $`𝐯_s`$, which is thus found to be equal to,
$`𝐯_s=<𝐌_{11}^{TT}𝐅_1^h+\overline{\rho }V𝐌_{12}^{TT}𝐅_2^h+𝐌_{11}^{TR}𝐓_1^h+\overline{\rho }V𝐌_{12}^{TR}𝐓_2^h>,`$ (38)
where the brackets $`<\mathrm{}>`$ denote ensemble averaging with respect to positions and orientations of the rods. The factors $`\overline{\rho }V=NN1`$ account for the presence of $`N1`$ rods which all interact with rod $`1`$ under consideration. The divergence problems mentioned above arising in the explicit evaluation of the ensemble averages will be dealt with later.
In order to be able to calculate these ensemble averages, the forces and torques must be expressed in terms of the positions and orientations of the rods. On the Brownian time scale there is a balance between all the forces and torques on each of the rods, that is, the total force and torque are equal to zero. The total force in turn is equal to the sum of the force $`𝐅_j^h`$ that the fluid exerts on the rod, the interaction force $`𝐅_j^I=_j\mathrm{\Phi }`$ (with $`\mathrm{\Phi }`$ the total interaction energy of the rods), the Brownian force $`𝐅_j^{Br}=k_BT_j\mathrm{ln}\{P\}`$ (with $`k_B`$ Boltzmann’s constant, $`T`$ the temperature and $`P`$ the probability density function for positions and orientations), and the external force $`𝐅^{ext}`$ due to the gravitational field. Hence,
$`𝐅_j^h=_j\mathrm{\Phi }+k_BT_j\mathrm{ln}\{P\}𝐅^{ext}.`$ (39)
Similarly, the total torque is the sum of the hydrodynamic torque $`𝐓_j^h`$, the interaction torque $`\widehat{}_j\mathrm{\Phi }`$, the Brownian torque $`k_BT\widehat{}_j\mathrm{ln}\{P\}`$, while the torque on each rod due to the homogeneous external force vanishes. Hence,
$`𝐓_j^h=\widehat{}_j\mathrm{\Phi }+k_BT\widehat{}_j\mathrm{ln}\{P\},`$ (40)
where the rotation operator is defined as,
$`\widehat{}_j(\mathrm{})\widehat{𝐮}_j\times _{u_j}(\mathrm{}),`$ (41)
with $`_{u_j}`$ the gradient operator with respect to $`\widehat{𝐮}_j`$. Substitution of eqs.(24,25) into eq.(23) for the sedimentation velocity yields,
$`𝐯_s`$ $`=`$ $`<𝐌_{11}^{TT}[_1\mathrm{\Phi }+k_BT_1\mathrm{ln}\{P\}𝐅^{ext}]`$ (44)
$`\overline{\rho }V𝐌_{12}^{TT}\left[_2\mathrm{\Phi }+k_BT_2\mathrm{ln}\{P\}𝐅^{ext}\right]`$
$`+𝐌_{11}^{TR}\left[\widehat{}_1\mathrm{\Phi }+k_BT\widehat{}_1\mathrm{ln}\{P\}\right]+\overline{\rho }V𝐌_{12}^{TR}\left[\widehat{}_2\mathrm{\Phi }+k_BT\widehat{}_2\mathrm{ln}\{P\}\right]>.`$
The next step in the explicit evaluation of these ensemble averages is to determine the stationary probability density function $`PP(𝐫_1,𝐫_2,\widehat{𝐮}_1,\widehat{𝐮}_2)`$ for the positions and orientations of two rods. At this point it is convenient to introduce the pair-correlation function $`g`$, defined as,
$`P(𝐫_1,𝐫_2,\widehat{𝐮}_1,\widehat{𝐮}_2)P(𝐫_1,\widehat{𝐮}_1)P(𝐫_2,\widehat{𝐮}_2)g(𝐫_{12},\widehat{𝐮}_1,\widehat{𝐮}_2),`$ (45)
where $`P(𝐫_j,\widehat{𝐮}_j)`$ is the probability density function for the position and orientation of a single rod. For spherical particles in a homogeneous external gravitational field, the probability density function for the position coordinates is simply the equilibrium function, without an external field, provided that the particles are identical. The probability density function differs from the equilibrium function only in case the relative sedimentation velocity of two spheres is different, for example due to differing masses and/or sizes. For rods, things are somewhat more complicated. Even if two rods are identical their relative sedimentation velocity generally differs as a result of the fact that the translational friction constant of rods is orientation dependent (see eq.(17)). The probability density function is generally dependent on the external force due to the fact that rods with different orientations overtake each other during sedimentation. Suppose, however, that the sedimentation velocity is so slow, that during a relative displacement of two rods in the gravitational field of the order of the length $`L`$ of the rods, each rod rotated many times due to their Brownian motion. The relative sedimentation velocity of two rods then averages out to zero. For such a case, the pair-correlation function is only weakly perturbed by the external field, so that we may use its equilibrium Boltzmann form,
$`g(𝐫_1,𝐫_2,𝐮_1,𝐮_2)=\mathrm{exp}\{\beta \mathrm{\Phi }(𝐫_1,𝐫_2,\widehat{𝐮}_1,\widehat{𝐮}_2)\},`$ (46)
Let us derive the inequality that should be satisfied for eq.(29) to be valid. According to eq.(17), the largest relative sedimentation velocity $`\mathrm{\Delta }𝐯_s`$ of two rods is approximately equal to,
$`\mathrm{\Delta }𝐯_s{\displaystyle \frac{\mathrm{ln}\{L/D\}}{4\pi \eta _0L}}𝐅^{ext}.`$ (47)
On the other hand, the time $`\tau _{rot}`$ required for a rotational revolution is equal to,
$`\tau _{rot}={\displaystyle \frac{\pi \eta _0L^3}{3k_BT\mathrm{ln}\{L/D\}}}.`$ (48)
The condition under which eq.(29) for the pair-correlation function is a good approximation is therefore,
$`{\displaystyle \frac{L/\mathrm{\Delta }𝐯_s}{\tau _{rot}}}{\displaystyle \frac{12k_BT}{𝐅^{ext}L}}\mathrm{\hspace{0.25em}1}.`$ (49)
Hence, the work provided by the external force to displace a colloidal rod over a distance equal to its length should be much smaller than a few times the thermal energy of the rods. Substitution of typical numbers shows that this inequality is satisfied under normal practical circumstances.<sup>*</sup><sup>*</sup>* For example, for rods with a length of 100 nm, the sedimentation velocity should be much less than 1 mm/min, in order that the inequality (32) is satisfied. Furthermore, when alignment of the rods during sedimentation in a homogeneous suspension is of no importance, the one-particle probability density functions in eq.(28) are both constants equal to,
$`P(𝐫_j,\widehat{𝐮}_j)={\displaystyle \frac{1}{4\pi V}}.`$ (50)
We will restrict ourselves here to the most common situation where the inequality (32) is satisfied, and assume negligible alignment, so that the probability density function is well approximated by eqs.(28,29,33). In that case many of the terms in eq.(27) for the sedimentation velocity cancel : the interaction contributions cancel against the Brownian terms. The sedimentation velocity reduces simply to,
$`𝐯_s`$ $`=`$ $`\left[<𝐌_{11}^{TT}>+\overline{\rho }V<𝐌_{12}^{TT}>\right]𝐅^{ext}.`$ (51)
Since $`𝐌_{12}^{TT}(𝐫_{12},\widehat{𝐮}_1,\widehat{𝐮}_2)1/r_{12}`$ for large distances, its ensemble average with respect to position coordinates diverges. Such a spurious divergence is also found for spherical particles, and is the result of the neglect of the hydrodynamic effect of the walls of the container. Batchelor showed that a formally divergent quantity, which is unambiguously finite valued on physical grounds, can be subtracted from the ensemble average, rendering a perfectly well defined sedimentation velocity. This finite valued quantity is formally divergent for the same reason that the ensemble average of $`𝐌_{12}^{TT}`$ is divergent, and subtraction accounts for the local hydrodynamic effects of the wall. To within the approximations made here, valid for long and thin rods, we do not encounter a conditionally divergent contribution due to terms $`1/r_{12}^3`$, as for spheres. Such a divergence can be dealt with by noting that the ensemble average of the deviatoric part of the stress tensor vanishes. We will not go into the extension of Batchelor’s arguments to rods to deal with this conditionally divergence problem. Batchelor’s argument is as follows. First define the velocity $`𝐮(𝐫𝐫_1,\mathrm{},𝐫_N,\widehat{𝐮}_1,\mathrm{},\widehat{𝐮}_N)`$ as the velocity at a point $`𝐫`$ (either in the fluid or inside a colloidal rod), given the positions $`𝐫_1,\mathrm{}𝐫_N`$ and orientations $`\widehat{𝐮}_1,\mathrm{},\widehat{𝐮}_N`$ of $`N`$ rods. In the laboratory reference frame the net flux of material through a cross sectional area must be zero. This means that the ensemble average of $`𝐮`$ must be zero. Hence,
$`\mathrm{𝟎}`$ $`=`$ $`<𝐮(𝐫𝐫_1,\mathrm{},𝐫_N,\widehat{𝐮}_1,\mathrm{},\widehat{𝐮}_N)>`$ (53)
$`={\displaystyle }d𝐫_1\mathrm{}{\displaystyle }d𝐫_N{\displaystyle }d\widehat{𝐮}_1\mathrm{}{\displaystyle }d\widehat{𝐮}_N𝐮(𝐫𝐫_1,\mathrm{},𝐫_N,\widehat{𝐮}_1,\mathrm{},\widehat{𝐮}_N)P(𝐫_1,,\mathrm{},𝐫_N,\widehat{𝐮}_1,\mathrm{},\widehat{𝐮}_N),`$
where $`P(𝐫_1,\mathrm{},\widehat{𝐮}_N)`$ is the probability density function for $`\{𝐫_1,\mathrm{},\widehat{𝐮}_N\}`$. Formally, this ensemble average diverges for the same reason that the sedimentation velocity diverges : the flow field $`𝐮`$ varies like $`1/r`$ for large distances due to its Oseen contribution. The formally divergent expression (35), that must be zero for physical reasons, is subtracted from eq.(34) for the sedimentation velocity to render this expression convergent. The field $`𝐮`$ can generally be written as the sum of a two terms : a term to which the field would be equal to in the absence of reflection contributions plus a term that accounts for the reflection contributions. To within our approach reflection contributions are neglected so that only the former term survives here. Hence,
$`𝐮(𝐫𝐫_1,\mathrm{},𝐫_N,\widehat{𝐮}_1,\mathrm{},\widehat{𝐮}_N)`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{N}{}}}𝐮(𝐫𝐫_j),for𝐫inthefluid,`$ (54)
$`=`$ $`𝐯_s,for𝐫inacore,`$ (55)
where the field $`𝐮(𝐫𝐫_j)`$, for $`𝐫`$ in the fluid, is the sum of the flow fields in eq.(13), (with $`𝐯`$ replaced by $`𝐯_s𝐮_s`$), and eq.(66) in appendix A, (both considered as a function of the relative distance to the $`j^{th}`$ rod under consideration), that is,
$`𝐮(𝐫𝐫_j)=𝐮_T(𝐫𝐫_j)+𝐮_R(𝐫𝐫_j),`$ (56)
where $`𝐮_T(𝐫𝐫_j)`$ is the fluid flow generated by a translating rod as given in eq.(13),
$`𝐮_T(𝐫𝐫_j)={\displaystyle \frac{1}{L}}{\displaystyle _{L/2}^{L/2}}𝑑l_j𝐓(𝐫𝐫_jl_j\widehat{𝐮}_j)𝐅^{ext},`$ (57)
where the inverse of eq.(12) is used, together with $`𝐅^h=𝐅^{ext}`$, and $`𝐮_R(𝐫𝐫_j)`$ is the field generated by a rotating rod, which is similarly given by eq.(66) in appendix A. Operating on both sides of eq.(35) with $`𝑑𝐫𝑑\widehat{𝐮}P(𝐫,\widehat{𝐮})`$, where $`P(𝐫,\widehat{𝐮})`$ is the constant specified in eq.(33), and subtraction of the resulting equation from eq.(34) for the sedimentation velocity yields, for identical rods and to first order in concentration (rename $`𝐫_1=𝐫_2,𝐫=𝐫_1,\widehat{𝐮}_1=\widehat{𝐮}_2`$ and $`\widehat{𝐮}=\widehat{𝐮}_1`$),
$`𝐯_s`$ $`=`$ $`\phi 𝐯_s+\left[<𝐌_{11}^{TT}>+\overline{\rho }V<𝐌_{12}^{TT}>\right]𝐅^{ext}`$ (59)
$`{\displaystyle \frac{\overline{\rho }}{(4\pi )^2}}{\displaystyle 𝑑𝐫_{12}𝑑\widehat{𝐮}_1𝑑\widehat{𝐮}_2\left[𝐮_T(𝐫_{12})+𝐮_R(𝐫_{12})\right]\chi _f(𝐫_1𝐫_2,\widehat{𝐮}_2)},`$
where $`\chi _f`$ is the characteristic function that restricts the integrations to points $`𝐫`$ which are in the fluid, not inside the core of rod $`2`$,
$`\chi _f(𝐫𝐫_2,\widehat{𝐮}_2)`$ $`=`$ $`1,for𝐫inthefluid,`$ (60)
$`=`$ $`0,for𝐫inthecoreofrod\mathrm{\hspace{0.33em}2}.`$ (61)
Without interactions the angular velocity of each rod is simply proportional to the hydrodynamic torque on the same rod (see eqs.(59,64) in appendix A), which hydrodynamic torques are zero. Since the integral in the above expression for the sedimentation velocity is multiplied by the concentration $`\overline{\rho }`$, it follows that the rotational field $`𝐮_R`$ does not contribute to first order in the density. For the same reason, in each term that is multiplied by the concentration, $`𝐅^{ext}`$ may be expressed with eq.(22) in terms of the sedimentation velocity $`𝐯_s^0`$ without interactions, at infinite dilution, as,
$`𝐯_s^0=\left[{\displaystyle \frac{\mathrm{ln}\{L/D\}}{3\pi \eta _0L}}{\displaystyle \frac{1}{6\pi \eta _0L}}(\nu _{}+\nu _{})\right]𝐅^{ext}.`$ (62)
We thus find the following expression for the sedimentation velocity, valid to first order in volume fraction,
$`𝐯_s=𝐯_s^0\left[1\left({\displaystyle \frac{f_1+f_2}{2\mathrm{ln}\{L/D\}(\nu _{}+\nu _{})}}+𝒪(D/L)\right){\displaystyle \frac{L}{D}}\phi \right],`$ (63)
where the functions $`f_1`$ and $`f_2`$ are equal to,
$`f_1`$ $`=`$ $`{\displaystyle \frac{1}{4\pi ^3DL^3}}{\displaystyle 𝑑𝐫_{12}𝑑\widehat{𝐮}_1𝑑\widehat{𝐮}_2\left[g(𝐫_{12},\widehat{𝐮}_1,\widehat{𝐮}_2)\chi _f(𝐫_1𝐫_2,\widehat{𝐮}_2)\right]}`$ (65)
$`\times {\displaystyle _{L/2}^{L/2}}dl_1{\displaystyle _{L/2}^{L/2}}dl_2{\displaystyle \frac{1}{𝐫_{12}+l_1\widehat{𝐮}_1l_2\widehat{𝐮}_2}},`$
$`f_2`$ $`=`$ $`{\displaystyle \frac{1}{4\pi ^3DL^3}}{\displaystyle 𝑑𝐫_{12}𝑑\widehat{𝐮}_1𝑑\widehat{𝐮}_2\chi _f(𝐫_1𝐫_2,\widehat{𝐮}_2)}`$ (67)
$`\times {\displaystyle _{L/2}^{L/2}}dl_1{\displaystyle _{L/2}^{L/2}}dl_2[{\displaystyle \frac{1}{𝐫_{12}+l_1\widehat{𝐮}_1l_2\widehat{𝐮}_2}}{\displaystyle \frac{1}{𝐫_{12}l_2\widehat{𝐮}_2}}].`$
where the expressions (13) and (38) for the mobility matrix $`𝐌_{12}^{TT}`$ and the field $`𝐮_T`$ are used, respectively. We also used that integrals over the Oseen tensor must be proportional to the identity tensor, so that in these integrals the Oseen tensor may be replaced by the trace $`\frac{1}{3}Tr\{𝐓\}`$ of that tensor, which, according to its defining equation (3), is equal to,
$`Tr\{𝐓(𝐫)\}={\displaystyle \frac{1}{2\pi \eta _0r}}.`$ (68)
For rods interacting only via a hard-core repulsion, it is shown in appendix B how to reduce the number of integrations, leading to the following results,
$`f_1`$ $`=`$ $`{\displaystyle \frac{8}{\pi ^3}}{\displaystyle _0^{\mathrm{}}}𝑑x{\displaystyle _1^1}𝑑z_1{\displaystyle _1^1}𝑑z_2{\displaystyle _0^\pi }𝑑\mathrm{\Psi }`$ (71)
$`\times j_0^2(xz_1)j_0^2(xz_2)\left[1\left(z_1z_2+\sqrt{(1z_1^2)(1z_2^2)}\mathrm{cos}\{\mathrm{\Psi }\}\right)^2\right]^{1/2}+𝒪(D/L)`$
$`=\mathrm{\hspace{0.25em}6.4}\mathrm{}+𝒪(D/L),`$
$`f_2`$ $`=`$ $`{\displaystyle \frac{2}{9}}{\displaystyle \frac{L}{D}}+𝒪(D/L).`$ (72)
The numerical value for $`f_1`$ has been obtained by numerical integration and applies for hard-core interactions, where $`g`$ is equal to $`0`$ when two cores overlap, and equal to $`1`$ otherwise. The result for $`f_2`$ is independent of the kind of direct interaction between the rods.
Substitution of the numerical values for $`f_1`$ and $`f_2`$ from eqs.(46,47) and into eq.(42) gives our final result for the sedimentation velocity up to $`𝒪(D/L)`$ contributions,
$`𝐯_s=𝐯_s^0\left[1{\displaystyle \frac{6.4+\frac{2}{9}\frac{L}{D}}{2\mathrm{ln}\{L/D\}(\nu _{}+\nu _{})}}{\displaystyle \frac{L}{D}}\phi \right],`$ (73)
The volume fraction prefactor,
$`\alpha ={\displaystyle \frac{6.4+\frac{2}{9}\frac{L}{D}}{2\mathrm{ln}\{L/D\}(\nu _{}+\nu _{})}}{\displaystyle \frac{L}{D}},`$ (74)
is plotted as a function of $`L/D`$ in Fig.1 (where both $`\nu _{}`$ and $`\nu _{}`$ are taken equal to $`0`$). Also plotted is an older result due to Peterson who predicted,
$`\alpha ={\displaystyle \frac{8(3/8)^{2/3}(L/D)^{1/3}}{2\mathrm{ln}\{L/D\}}}{\displaystyle \frac{L}{D}}.`$ (75)
In this latter theory back flow is not correctly accounted for, hydrodynamic interactions are orientationally preaveraged and certain integrals are not precisely calculated but only estimated. The data points shown in Fig. 1a are experimental results for silica rods, coated with stearyl alcohol and disolved in cyclohexane. The data point $`\times `$ is taken from Ref. . The point $``$ is an unpublished result from the same author’s. The data point if Fig. 1b is a data point for fd-virus at high salt concentration, as obtained in the experimental section of the present paper. As can be seen from Fig. 1, the present prediction is virtually equal to that of Peterson for $`L/D<30`$, but large differences are found for large aspect ratios. For large aspect ratios, $`\alpha `$ is predicted to vary like $`\left(L/D\right)^2/\mathrm{ln}\{L/D\}`$, in contrast to Peterson’s result $`\left(L/D\right)^{4/3}/\mathrm{ln}\{L/D\}`$. For smaller aspect ratios $`\alpha `$ approaches approximately Batchelor’s value for spheres $`\alpha =6.55`$, which is probably fortuitous in view of the approximations made here which limit the results to be meaningful only for long and thin rods.
## IV Experimental Results
The concentration dependence of the sedimentation velocity predicted by Eq. (74) differs significantly from Peterson’s result Eq. (75) only for rods with large L/D. In our experiments we have used filamentous bacteriophage fd which is a rod-like virus with L/D $``$ 130. Other relevant physical characteristics of fd, are it’s length $`L=880`$nm, it’s diameter $`D=6.6`$ nm , and it’s density of 1.285 mg/ml . Because of it’s large L/D ratio the virus is a semi-flexible rather then a rigid rod characterized with persistence length of 2.2 $`\mu `$m . It’s linear charge density is 10 e<sup>-</sup>/nm at pH 8.2 .
We have grown fd virus according to standard procedures of molecular biology described in Maniatis . The virus suspension was first purified in a cesium chloride density gradient and then extensively dialyzed against tris buffer at pH 8.15 and at the desired ionic strength for the sedimentation experiments. After that the virus was concentrated by ultracentrifugation and from this stock solution a series of samples with different concentrations were prepared. The sedimentation velocity was measured on a Beckman XL-A analytical ultracentrifuge equipped with UV absorbance optics. Most of the experiments were done at 25 C and at a centrifugal force equal to 45,500 g’s (25,000 rpm). Before each sedimentation experiment the sample and rotor were allowed to equilibrate at the desired temperature for a few hours. Sedimentation data showed some unexpected features, interfering with straightforward calculation of the sedimentation coefficient. For this reason we have added appendix C, where the detailed analysis of our data is given.
The measured sedimentation velocity for a range of volume fractions of fd from dilute solution up to a stable nematic phase are shown in Fig. 2. All the samples in these measurements were kept at 8mM ionic strength. The sedimentation velocity of rods in the isotropic phase uniformly decreases with increasing concentration. After Ref. we have tried to fit our experimental data to a functional form $`S_\varphi =S_0(1p\varphi )^\nu `$ where $`S_0`$ is the sedimentation velocity at infinite dilution and $`\varphi `$ is the volume fraction of rods. The experimental values of sedimentation velocity are reported in Svedbergs where 1 $`S=10^{13}s^1`$. As seen from Fig. 2 we obtain a reasonable fit to the experimental data in the isotropic phase and find that the sedimentation velocity at infinite dilution is $`S_0=46.0`$ for the value of constants $`\nu =1/3`$ and p=3600. After linearizing our fitted formula we find that the volume prefactor $`\alpha 1200`$ is much larger then predicted by Eq. (74). The reason for such a high value of slope $`\alpha `$ is the low ionic strength at which the experiment was performed. The same increase in $`\alpha `$ with decreasing ionic strength is observed in sedimentation of spherical particles . Also we note that in this case the region where the sedimentation velocity varies linearly with rod concentration is limited to very low volume fraction of rods.
It is a well known fact that elongated rods at high volume fractions undergo a first order phase transition to a liquid crystalline nematic phase . The nematic phase is characterized by a short range liquid-like positional order and long range solid-like orientational order of rods. fd virus forms a cholesteric phase instead of the nematic phase . Locally the cholesteric phase is equivalent to nematic, however on a macroscopic scale the average direction of molecules in a cholesteric phase forms a helix. The free energy difference between a cholesteric phase and a nematic phase is very small and although our experiments are performed on the cholestric phase only, we expect that our results are generic and would hold for a nematic phase of hard rods as well.
Bacteriophage fd in the cholesteric phase exhibits qualitatively new behavior when placed in centrifugal field. Instead of a single sedimenting boundary and single plateau we observe two boundaries with two plateau’s sedimenting at different velocities as shown in Fig. 3. To confirm that this change in sedimentation behavior is indeed due to the formation of the nematic phase we have made a sample which is co-existing between the isotropic and nematic phase. After the sample had phase separated into macroscopically distinct co-existing phases, each phase was sedimented separately. In the isotropic phase there was no sign of a second boundary, while in the nematic (cholesteric) phase we observed a fast sedimenting second boundary that was slightly more concentrated then the first component. On one hand, the slow component had a plateau concentration and a sedimentation velocity that was almost independent of the average concentration. On the other hand, the sedimentation velocity of the faster moving component rapidly decreases with increasing $`\mathrm{𝑓𝑑}`$ concentration and at the same time the difference between the plateau concentrations of the two components increases with increasing average concentration.
The unstable sedimentation of colloidal rods in the nematic phase has a similar origin as the self-sharpening effect described in Appendix C. The reason for the instability is the discontinous jump in sedimentation velocity that occurs at the isotropic-nematic phase transition as is shown in Fig. 2 . The denser nematic phase sediments at a significantly higher velocity then a more dilute isotropic phase. Initially a stable nematic phase occupies the whole sample length. When the centrifugal field is turned on a sharp sedimenting boundary starts moving towards the bottom of the container. Below this boundary (to the right) the rods are still in the nematic phase while above it the concentration of rods is very low and therefore they are in the isotropic state. This occurs because some particles will inevitably diffuse against the centrifugal field from a highly concentrated plateau into the dilute region. As this happens they simultaneously undergo a transition from the nematic to the isotropic phase. Since the sedimentation velocity of rods in the isotropic phase is much lower then in nematic phase, the probability of the rods diffusing from the isotropic phase back into the nematic phase is virtually zero. It is this asymmetry that results in a continuous flux of particles from the nematic into the isotropic phase and contributes to the formation of the second plateau in the isotropic phase that is moving at a slower speed. The concentration of the rods in the isotropic plateau will be very close to the concentration of isotropic rods co-existing with the nematic phase due to the self-sharpening effect because dilute isotropic rods will catch up with more concentrated isotropic rods. However, the highest concentration the isotropic rods can attain is the co-existence concentration between the isotropic and nematic phases, as long as the nematic phase sediments faster then the isotropic phase. Indeed, this is very close to what we observe in Fig. 3. Another experimental observation corroborating our explanation is that the sedimentation velocity and concentration of the slower isotropic plateau does not change significantly with average concentration of rods as seen in Fig 2.
Since the theory presented in this paper is valid only to first order in concentration of rods, to obtain an accurate value of the prefactor $`\alpha `$ in Eq. 74 we have made additional measurements in the dilute to semi-dilute range. Our results are presented in Fig. 4. We note that the overlap to semi-dilute concentration for fd with $`L=0.88\mu `$m is at volume fraction of $`5.910^5`$. Unlike the previous measurements, we have done these measurements at high ionic strength where the behavior of charged rods is expected to approach the behavior of hard rods. Additionally at high ionic strength we expect the sedimentation velocity to have a linear dependence on volume fraction of rods up to higher values of volume fraction. The results for ionic strength of 50mM and 100mM ionic strength are shown in figure 4. The volume prefactor in Eq. (74) at 50 mM ionic strength is $`\alpha =450\pm 40`$ and at 100mM ionic strength $`\alpha =440\pm 60`$ . We have repeated the experiment at 100mM ionic strength on a different analytical Beckman Xl-A ultracentrifuge and obtained the following result $`\alpha =490\pm 50`$. We conclude that $`\alpha =470\pm 50`$ which is the result plotted in Fig. 1b. Since the values of the coefficient $`\alpha `$ do not change much with changing ionic strength from 50 mM to 100mM we conclude that the charged rods have approached the hard rod limit. Note that because of it’s large L/D ratio $`fd`$ is slightly flexible with a persistence length which is 2.5 times it’s contour length . Still for the experimentally determined parameters of $`fd`$, which are $`L=880`$nm and $`D=6.6`$nm, our experimental results compare favorably to the Eq. 74, which predicts the value of $`\alpha =488`$ (see Fig. 1). In contrast, the previous result due to Peterson in Eq. 75 predicts a lower value of $`\alpha =288`$
## V Acknowledgment
We acknowledge valuable discussions with R. B. Meyer. This research was supported by National Science Foundations grants No. DMR-9705336, INT-9113312 and by the Netherlands Foundations of Fundamental Research(FOM). Additional information is avaliable online: www.elsie.brandeis.edu.
## Appendix A
Flow field generated by a rotating rod
Consider a rod with its center at the origin, which rotates with an angular velocity $`𝛀`$. The angular velocity is decomposed in its component perpendicular and parallel to the rods center line,
$`𝛀_{}`$ $`=`$ $`\left[\widehat{𝐈}\widehat{𝐮}\widehat{𝐮}\right]𝛀,`$ (76)
$`𝛀_{}`$ $`=`$ $`\widehat{𝐮}\widehat{𝐮}𝛀.`$ (77)
Due to the linearity of the governing hydrodynamic equations, the flow fields generated by a rod rotating along $`𝛀_{}`$ and $`𝛀_{}`$ may be calculated separately and added to obtain the flow field of the rod rotating along $`𝛀`$.
Let us first consider a rod rotating with an angular velocity $`𝛀_{}`$. The flow field that is generated by this rotating rod is given by the general equation (2). The relative change of the velocity of the beads is $`1/j`$. For beads further away from the origin one may therefore consider the velocity over larger groups of beads as being virtually constant. The force on bead $`j`$ is then proportional to its own velocity,
$`𝐅_j^h=C𝛀_{}\times 𝐫_j=CDj𝛀_{}\times \widehat{𝐮}.`$ (78)
where $`C`$ is an as yet unknown proportionality constant. This expression is not valid for beads close to the center of the rod : for these beads the forces may have a different direction than their velocity. The fluid flow field generated by a long and thin rod, however, is primarily determined by the relatively large velocities of the beads further away from its center. We may therefore use eq.(53), except for relatively few beads close to the center and near the tips of the rod. Since $`𝐫_j=jD\widehat{𝐮}`$, the torque is thus found, to leading order in $`D/L`$, to be equal to,
$`𝐓_{}^h={\displaystyle \underset{j=n/2}{\overset{n/2}{}}}𝐫_j\times 𝐅_j^h=CD^2{\displaystyle \frac{1}{12}}\left({\displaystyle \frac{L}{D}}\right)^3\widehat{𝐮}\times (𝛀_{}\times \widehat{𝐮})=CD^2{\displaystyle \frac{1}{12}}\left({\displaystyle \frac{L}{D}}\right)^3𝛀_{},`$ (79)
since $`𝛀_{}`$ is perpendicular to $`\widehat{𝐮}`$. It is used here that $`_{j=1}^kj^2=\frac{1}{6}k(k+1)(2k+1)`$. First of all, the constant $`C`$ is calculated self-consistently from Faxén’s theorem in the form of eq.(10). Multiplying both sides of eq.(10) by $`𝐫_j\times `$, using that $`𝐫_j\times 𝐯_j=j^2D^2𝛀_{}`$, and summation over beads, leads to,
$`{\displaystyle \frac{1}{12}}\left({\displaystyle \frac{L}{D}}\right)^3D^2𝛀_{}={\displaystyle \frac{1}{3\pi \eta _0D}}𝐓_{}^h+{\displaystyle \frac{CD}{8\pi \eta _0}}\left({\displaystyle \frac{L}{D}}\right)^3g(L/D)𝛀_{},`$ (80)
where the function $`g`$ is defined as,
$`g(L/D)={\displaystyle \frac{1}{(n+1)^3}}{\displaystyle \underset{j=n/2}{\overset{n/2}{}}}{\displaystyle \underset{i=n/2,ij}{\overset{n/2}{}}}ij\left[{\displaystyle \frac{1}{ij}}+{\displaystyle \frac{1}{12}}{\displaystyle \frac{1}{ij^3}}\right].`$ (81)
For long and thin rods the summations may be replaced by integrals, leading to,
$`g(L/D)={\displaystyle \frac{1}{6}}\mathrm{ln}\{L/D\},`$ (82)
up to leading order in $`D/L`$. Substitution of eq.(17) for the torque yields a single equation for $`C`$, yielding, again up to leading order,
$`C={\displaystyle \frac{4\pi \eta _0D}{\mathrm{ln}\{L/D\}}}.`$ (83)
Hence, from eq.(17),
$`𝛀_{}={\displaystyle \frac{3\mathrm{ln}\{L/D\}}{\pi \eta _0L^3}}𝐓_{}^h.`$ (84)
The flow field $`𝐮_{}`$ that is generated by a rotating rod may now be obtained from eq.(4), to within the same approximations that were discussed in the section II A, as,
$`𝐮_{}(𝐫)={\displaystyle \underset{j=n/2}{\overset{n/2}{}}}𝐓(𝐫𝐫_j)𝐅_j^h={\displaystyle \frac{4\pi \eta _0D^2}{\mathrm{ln}\{L/D\}}}{\displaystyle \underset{j=n/2}{\overset{n/2}{}}}𝐓(𝐫𝐫_j)(𝛀_{}\times j\widehat{𝐮}).`$ (85)
Replacing the sum over beads by a line integral, we thus find,
$`𝐮_{}(𝐫)={\displaystyle \frac{4\pi \eta _0}{\mathrm{ln}\{L/D\}}}{\displaystyle _{L/2}^{L/2}}𝑑l𝐓(𝐫𝐫_pl\widehat{𝐮})(𝛀_{}\times l\widehat{𝐮}),`$ (86)
where $`𝐫_p`$ is the position coordinate of the rod.
Next consider a rod rotating with an angular velocity $`𝛀_{}`$. For this case we have to resort to Faxén’s theorem for rotational motion of a bead, which reads,
$`𝛀_{}={\displaystyle \frac{1}{\pi \eta _0D^3}}𝐓_j^h+{\displaystyle \frac{1}{2}}_j\times 𝐮_0(𝐫_j),`$ (87)
where, as in the translational Faxén’s theorem (7), $`𝐮_0`$ is the fluid flow velocity that would have existed in the absence of bead $`j`$. The first term on the right hand-side is just Stokes rotational friction of a single bead in an unbounded fluid, while the second term accounts for hydrodynamic interaction between the beads. The important thing to note here is that the fluid flow generated by a single rotating bead is now equal to,
$`𝐮_j(𝐫)=\left({\displaystyle \frac{D/2}{𝐫𝐫_j}}\right)^3𝛀_{}\times (𝐫𝐫_j),`$ (88)
so that this fluid flow is $`\mathrm{𝟎}`$ along the entire center line of the rod. This implies that hydrodynamic interaction between the beads is unimportant for this case. For a long and thin rod rotating along its center line, each bead experiences a rotational friction that is practically equal to the Stokes friction, as if each bead where alone in an unbounded fluid. As a result, the total torque on the rod is simply the sum of the Stokesian torques on the beads, so that it follows immediately from Faxén’s theorem (62) that,
$`𝛀_{}={\displaystyle \frac{1}{\pi \eta _0D^2L}}𝐓_{}^h.`$ (89)
Furthermore, the total fluid flow $`𝐮_{}`$ is simply the sum of the fluid flows (63) generated by the rotating beads as if they were alone in an unbounded fluid, since hydrodynamic interaction between the beads is unimportant in the present case. Replacing the sum by a line integral thus yields,
$`𝐮_{}(𝐫)={\displaystyle \frac{D^2}{8}}{\displaystyle _{L/2}^{L/2}}𝑑l{\displaystyle \frac{1}{𝐫𝐫_pl\widehat{𝐮}^3}}\left(𝛀_{}\times (𝐫𝐫_p)\right).`$ (90)
The fluid flow $`𝐮=𝐮_{}+𝐮_{}`$ generated by a rotating rod with an arbitrary angular velocity $`𝛀=𝛀_{}+𝛀_{}`$ follows by combining eqs.(51,52) and (60,65),
$$𝐮(𝐫)=\frac{4\pi \eta _0}{\mathrm{ln}\{L/D\}}_{L/2}^{L/2}𝑑l𝐓(𝐫l\widehat{𝐮})(𝛀\times l\widehat{𝐮})+\frac{D^2}{8}_{L/2}^{L/2}𝑑l\frac{1}{𝐫𝐫_pl\widehat{𝐮}^3}\left((\widehat{𝐮}\widehat{𝐮}𝛀)\times (𝐫𝐫_p)\right).$$
(91)
This approximate expression will be used in the following paragraph to obtain an expression for the mobility matrices $`𝐌_{1j}^{TR}`$, $`j=1,2`$.
Calculation of $`𝐌^{TR}`$
In order to calculate the velocity $`𝐯_2`$ that rod $`2`$ acquires in the flow field (66) generated by a rotating rod $`1`$, we apply, without further discussion, the same “mean-field” approach as in the previous section. The velocity $`𝐯_2`$ is approximated by taking the fluid flow field generated by the rotating rod as a constant, equal to the average of the actual field over the center line of the rod. Hence,
$`𝐯_2=\overline{𝐮}{\displaystyle \frac{\mathrm{ln}\{L/D\}}{4\pi \eta _0L}}\left[\widehat{𝐈}+\widehat{𝐮}_2\widehat{𝐮}_2\right]𝐅_2^h,`$ (92)
where the average flow field in terms of the torque on rod $`1`$ follows from eqs.(66) and (59,64), with $`𝛀=𝛀_1`$, the angular velocity of rod $`1`$ and $`𝐓_1^h`$ the torque on rod 1,
$`\overline{𝐮}`$ $`=`$ $`{\displaystyle \frac{12}{L^4}}{\displaystyle _{L/2}^{L/2}}𝑑l_1{\displaystyle _{L/2}^{L/2}}𝑑l_2𝐓(𝐫_{21}+l_2\widehat{𝐮}_2l_1\widehat{𝐮}_1)\left(l_1\widehat{𝐮}_1\times 𝐓_1^h\right)`$ (93)
$`+`$ $`{\displaystyle \frac{1}{8\pi \eta _0L^2}}{\displaystyle _{L/2}^{L/2}}𝑑l_1{\displaystyle _{L/2}^{L/2}}𝑑l_2{\displaystyle \frac{1}{𝐫_{21}+l_2\widehat{𝐮}_2l_1\widehat{𝐮}_1^3}}(𝐫_{21}+l_2\widehat{𝐮}_2)\times \left(\widehat{𝐮}_1\widehat{𝐮}_1𝐓_1^h\right).`$ (94)
By definition the following “mean-field” expression for the translational-rotational mobility matrices are thus obtained (after an interchange of the indices $`1`$ and $`2`$), The outer product $`𝐚\times 𝐀`$ of a vector $`𝐚`$ and a matrix $`𝐀`$ is defined as the matrix with column vectors equal to the outer product of $`𝐚`$ and the column vectors of $`𝐀`$. The outer product is thus taken with respect to the first index on $`𝐀`$.
$`𝐌_{11}^{TR}`$ $`=`$ $`\mathrm{𝟎},`$ (95)
$`𝐌_{12}^{TR}`$ $`=`$ $`{\displaystyle \frac{12}{L^4}}{\displaystyle _{L/2}^{L/2}}𝑑l_1{\displaystyle _{L/2}^{L/2}}𝑑l_2l_2\widehat{𝐮}_2\times 𝐓(𝐫_{12}+l_1\widehat{𝐮}_1l_2\widehat{𝐮}_2)`$ (96)
$`+`$ $`{\displaystyle \frac{1}{8\pi \eta _0L^2}}{\displaystyle _{L/2}^{L/2}}𝑑l_1{\displaystyle _{L/2}^{L/2}}𝑑l_2{\displaystyle \frac{1}{𝐫_{12}+l_1\widehat{𝐮}_1l_2\widehat{𝐮}_2^3}}\left[\widehat{𝐮}_2\times (𝐫_{12}+l_1\widehat{𝐮}_1)\right]\widehat{𝐮}_2.`$ (97)
A non-zero contribution to $`𝐌_{11}^{TR}`$ stems entirely from reflection contributions, since a pure rotation of a single rod in an unbounded fluid does not induce a translational velocity of the same rod. As mentioned before, reflection contributions are small in the isotropic state, since the typical distance between the beads of different rods is of the order $`L`$.
## Appendix B
As a first step in the evaluation of the integrals in eq.(43) for $`f_1`$, the Fourier transform of the Oseen tensor ($`𝐓(𝐤)=\frac{1}{\eta _0k^2}\left[\widehat{𝐈}\widehat{𝐤}\widehat{𝐤}\right]`$, with $`\widehat{𝐤}=𝐤/k`$) is substituted, and the integrations with respect to $`l_1`$ and $`l_2`$ are performed, with the result,
$`f_1`$ $`=`$ $`{\displaystyle \frac{1}{8\pi ^5DL}}{\displaystyle 𝑑𝐤k^2𝑑\widehat{𝐮}_1𝑑\widehat{𝐮}_2𝑑𝐫_{12}\left[g(𝐫_{12},\widehat{𝐮}_1,\widehat{𝐮}_2)\chi _f(𝐫_1𝐫_2,\widehat{𝐮}_2)\right]}`$ (99)
$`\times \mathrm{exp}\{i𝐤𝐫_{12}\}j_0({\displaystyle \frac{1}{2}}L𝐤\widehat{𝐮}_1)j_0({\displaystyle \frac{1}{2}}L𝐤\widehat{𝐮}_2),`$
where,
$`j_0(x){\displaystyle \frac{\mathrm{sin}\{x\}}{x}}.`$ (100)
Consider the integral with respect to $`𝐫_{12}`$,
$`I{\displaystyle 𝑑𝐫_{12}\left[g(𝐫_{12},\widehat{𝐮}_1,\widehat{𝐮}_2)\chi _f(𝐫_1𝐫_2,\widehat{𝐮}_2)\right]\mathrm{exp}\{i𝐤𝐫_{12}\}}.`$ (101)
Replace the expression in the square brackets by $`(g1)+(1\chi _f)`$. The integral over $`1\chi _f`$ is easily found to be equal to,
$`{\displaystyle 𝑑𝐫_{12}\left[1\chi _f(𝐫_1𝐫_2,\widehat{𝐮}_2)\right]\mathrm{exp}\{i𝐤𝐫_{12}\}}={\displaystyle \frac{\pi }{4}}D^2Lj_0({\displaystyle \frac{1}{2}}L𝐤\widehat{𝐮}_2),`$ (102)
while the integral over $`g1`$ is equal to,
$`{\displaystyle 𝑑𝐫_{12}\left[g(𝐫_{12},\widehat{𝐮}_1,\widehat{𝐮}_2)1\right]\mathrm{exp}\{i𝐤𝐫_{12}\}}=2DL^2\widehat{𝐮}_1\times \widehat{𝐮}_2j_0({\displaystyle \frac{1}{2}}L𝐤\widehat{𝐮}_1)j_0({\displaystyle \frac{1}{2}}L𝐤\widehat{𝐮}_2).`$ (103)
These results are valid for $`kD1`$ (say $`kD<0.2`$), while, in addition, eq.(75) is valid for orientations where $`\frac{D}{L}\widehat{𝐮}_1\times \widehat{𝐮}_2`$. As will turn out, the $`kD`$-dependence is of no importance for long and thin rods, since convergence of the wavevector integral is assured by the $`kL`$-dependent functions, which tend to zero for wavevectors for which, indeed, $`kD1`$. Moreover, the angular integration range, pertaining to orientations where $`\frac{D}{L}/\widehat{𝐮}_1\times \widehat{𝐮}_2`$ is not small, vanishes for long and thin rods. Substitution of the results (74,75) into eq.(71) for $`f_1`$, and noting that after integration over orientations the dependence on the direction $`\widehat{𝐤}`$ of the wavevector is lost, so that its direction may be chosen along the $`z`$-direction, yields (with $`x=\frac{1}{2}kL`$),
$`f_1={\displaystyle \frac{2}{\pi ^4}}{\displaystyle _0^{\mathrm{}}}𝑑x{\displaystyle 𝑑\widehat{𝐮}_1𝑑\widehat{𝐮}_2j_0^2(xz_2)\left[\widehat{𝐮}_1\times \widehat{𝐮}_2j_0^2(xz_1)\frac{\pi }{8}\frac{D}{L}j_0(xz_1)\right]}.`$ (104)
with $`z_j`$, $`j=1,2`$, is the $`z`$-component of $`\widehat{𝐮}_j`$. The second term between the square brackets is an $`𝒪(D/L)`$ contribution as compared to the first term and may be neglected. Transforming the orientational integrals to spherical coordinates, for which $`z_j=\mathrm{cos}\{\mathrm{\Theta }_j\}`$, and using that (with $`\mathrm{\Psi }=\phi _1\phi _2`$),
$`\widehat{𝐮}_1\times \widehat{𝐮}_2=\left[1\left(\mathrm{cos}\{\mathrm{\Theta }_1\}\mathrm{cos}\{\mathrm{\Theta }_2\}+\mathrm{sin}\{\mathrm{\Theta }_1\}\mathrm{sin}\{\mathrm{\Theta }_2\}\mathrm{cos}\{\mathrm{\Psi }\}\right)^2\right]^{1/2},`$ (105)
finally yields eq.(46) for $`f_1`$.
Next consider the evaluation of the integrals in eq.(44) for $`f_2`$. That the integrals are convergent follows from the Taylor expansion,
$`{\displaystyle \frac{1}{𝐫_{12}𝐚}}={\displaystyle \frac{1}{r_{12}}}+𝐚{\displaystyle \frac{1}{r_{12}}}+{\displaystyle \frac{1}{2}}\mathrm{𝐚𝐚}:{\displaystyle \frac{1}{r_{12}}}+\mathrm{}.`$ (106)
Using this in eq.(44) and integration with respect to $`\widehat{𝐮}_1`$ shows that the integrand varies like $`r_{12}^4`$ for large $`r_{12}`$, since $`^2r_{12}^1=0`$ for $`r_{12}0`$. Following the same procedure as above one finds,
$`f_2`$ $`=`$ $`{\displaystyle \frac{1}{2\pi ^3DL}}{\displaystyle 𝑑𝐤_1^1𝑑z_1_1^1𝑑z_2\left[(2\pi )^3\delta (𝐤)\frac{\pi }{4}D^2Lj_0(\frac{1}{2}Lkz_2)\right]}`$ (108)
$`\times j_0({\displaystyle \frac{1}{2}}Lkz_2){\displaystyle \frac{j_0(\frac{1}{2}Lkz_1)1}{k^2}},`$
where $`\delta `$ is the delta distribution. The second term in the square brackets is easily seen to be $`𝒪(D/L)`$, using the same integration tricks as above for the evaluation of $`f_1`$. For the evaluation of the delta distribution contribution, the integrand can be expanded in a power series expansion in $`k`$. Using that $`j_0(x)=1x^2/6+\mathrm{}`$, results in eq.(47) for $`f_2`$.
## Appendix C
An analytical centrifuge measures the concentration of sedimenting colloid along the centrifugal field. From a single run in an analytical ultracentrifuge we obtain a time sequence of plots usually taken every few minutes. A representative sequence of these plots is shown in Fig. 5. Each plot in the series indicates the fd concentration as a function of radial position in the cell at that particular time. The concentrations of the dilute virus solutions were determined with the extinction coefficients of 3.84 mg<sup>-1</sup>cm<sup>2</sup> at 270 nm . For samples with higher concentration the solution is optically opaque at 270 nm and therefore we measure absorbance at progressively higher wavelengths, which corresponds to a lower extinction coefficient of fd. The sedimenting particles in Fig. 5 move from left to right. The water-air interface is indicated by a sharp peak located at radial position of 5.95 cm that is due to refraction by the meniscus. Note that this peak does not move as a function of time indicating that the container does not leak. As the rods start sedimenting towards the cell bottom, the region at the top of the solution (to the right of the air-water interface and to the left of the sedimentation front in Fig. 5) is depleted of virus as indicated by absence of absorbtion. Also the value of the concentration of rods in the plateau region, always to the right of the depleted region, is decreasing as the bulk of the sample moves towards the bottom of the container. The reason for this is that the walls of the cell are not parallel to each other, but instead follow the lines of centrifugal field in order to minimize convective disturbances, an effect refered to as “radial dilution” . Between the flat plateau region and the depleted region there is a sharp boundary.
At higher concentrations of fd we observed the appearance of a sharp peak at the sedimenting boundary as is shown in Fig. 6. The peak height increases with increasing concentration while the magnitude of the peak is independent of the wavelength and thus this peak cannot be due to absorption of the fd, which is wavelength dependent. The probable cause of the peak is the refraction of incident light due to the steep gradient in the virus concentration and hence the refractive index at the sedimenting boundary. As the incident light is refracted away from the detector, less light is collected by it and this results in apparent increased absorption of the sample. The peak at the water/air meniscus has the same origin.
Two factors that determine the shape of the sedimenting boundary are the diffusion constant and the self-sharpening effect . The diffusion of the particles leads to gradual spreading of the initially very sharp boundary. This diffusion of particles is countered by the self-sharpening effect, which is due to the concentration dependence of the sedimentation velocity. On one hand, any molecule lagging behind the boundary is in a more dilute environment and will therefore sediment at an enhanced velocity. On the other hand, the particles in the plateau region are in a more concentrated environment and their sedimentation will be retarded. As a consequence the boundaries will self-sharpen. In a suspension of elongated particles the self-sharpening effect will be much stronger then in a suspension of globular particles because the volume prefactor $`\alpha `$ in Eq. 74 is much larger for elongated particles then for globular particles. The pronounced self-sharpening effect leads to hyper-sharp boundaries, which result in a steep gradient of refractive index which in turn causes the artifacts shown in Fig. 6. In globular colloids these effects are usually not observed.
In sedimentation analysis it is assumed that the rate of movement of the sedimentation boundary is approximately equivalent to the sedimentation velocity of the particles in the plateau (bulk) region. To compare results from different runs it is usual to express the sedimentation velocity in units independent of centrifugal force as follows:
$$S=\frac{1}{\omega ^2r}\frac{dr}{dt}=\frac{1}{\omega ^2}\frac{d\mathrm{ln}r}{dt}$$
(109)
The sedimentation velocity unit is called a Svedberg (S), with 1 S = $`[10^{13}\text{sec}^{}1]`$. We define $`r`$ as the radial position at the sedimentation boundary where the virus concentration is equal to half the concentration of the plateau region. This quantity is easily obtained from experimental data for samples at low concentration. For samples at higher concentration, where we observe a peak at the sedimenting boundary due to refraction of light, we define $`r`$ as the radial position of the highest point of the peak. A typical plot of the logarithm of $`r`$ against $`\omega ^2t`$ used in the determination of the sedimentation constant is shown in Fig. 7. Surprisingly, we found out that a linear function provided an inadequate fit to our data. When the sedimentation data are collected between radial positions of 6.1 cm and 6.8 cm a polynomial of second order fits the data much better:
$$\mathrm{ln}r=A+B\omega ^2t+C\omega ^4t^2$$
(110)
We introduce the experimentally observed sedimentation constant $`S^r`$ by combining Eq. 110 and Eq. 109:
$$S^r=\frac{\text{d}\mathrm{ln}r}{\text{d}\omega ^2t}=B+2C\omega ^2t$$
(111)
From this equation we see that the experimentally measured sedimentation velocity is not a constant but depends on the position of the measurement $`r`$ of equivalently time $`\omega ^2t`$ at which the sedimentation front is found at position $`r`$. The reason for this unexpected behavior is not clear, but we assume it is an instrumental artifact. There is no physical reason to believe that the sedimentation velocity is a function of time or of radial position in the cell. In table 1 we see that the coefficient C, obtained when the quadratic polynomial in Eq. 110 is fitted to data in Fig. 4a is independent of concentration. This is another indication that this artifact is due to the instrument.
Theoretically, the constant $`B`$ in eq. (111) should be equal to the concentration dependent Svedberg constant:
$$S_\varphi =S_0(1\alpha \varphi )$$
(112)
and the constant $`C`$ should be zero. Instead, we have found that the experimental Svedberg (eq. 111) is described by:
$$S^r=S_\varphi +\text{offset}=S_0+\text{offset}S_0\alpha \varphi $$
(113)
where “offset” depends on position in the centrifuge, but is independent of colloid concentration. $`S_0`$ is the Svedberg constant of the rods in the limit of zero concentration. However, the value of slope
$$\frac{\text{d}S^r}{\text{d}\varphi }=\alpha S_0$$
(114)
is independent of radial position (or equivalently $`\omega ^2t`$) where we evaluate Eq. 111 as is shown in table I. The measurement artifact only introduces a position dependent offset in the sedimentation velocity which affects the measured value of $`S_0`$. From a few measurements where we did not observe measurement artifacts ($`C=0`$) we obtained the value of $`S_0=47`$. Since this is in good agreement with previous measurements we use this value throughout our analysis .
It is important to note that the dependence of $`S^r`$ on position $`r`$ shown in Fig. 7 is not due to the decreasing concentration of rods in plateau, which in turn is due to radial dilution. To show this we have made two measurements. In a first measurement we evaluated the sedimentation velocity at the point where the sedimenting boundary is close to the bottom of the container. At this time, due to radial dilution, the plateau concentration is about 70% of the initial concentration. In the second run our initial concentration was 70% of the concentration of rods in the first run. In this run we evaluated the sedimentation velocity right at the beginning of the run. We find that sedimentation velocities obtained in these two ways are vastly different, which indicates that the systematic errors described are not due to radial dilution.
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# CMB 𝐵-polarization to map the Large-scale Structures of the Universe
## I Introduction
In the new era of precision cosmology we are entering in, the forthcoming experiments will provides us with accurate data on Cosmic Microwave Background anisotropies. This should lead to accurate determinations of the cosmological parameters, provided the large-scale structures of the Universe indeed formed from gravitational instabilities of initial adiabatic scalar perturbations. It has been soon realized however that even with the most precise experiments, the cosmological parameter space is degenerate when the primary CMB anisotropies alone are considered. Complementary data, that may be subject to more uncontrollable systematics are thus required, such as supernovae surveys (but see ) or constraints derived from the large-scale structure properties. Among the latter, weak lensing surveys are probably the safer, but still have not yet proved to be accurate enough with the present day observations.
Secondary CMB anisotropies (i.e. induced by a subsequent interaction of the photons with the mass or matter fluctuations) offer opportunities for raising this degeneracy. Lens effects are particularly attractive since they are expected to be one of the most important.They also are entirely driven by the properties of the dark matter fluctuations, the physics of which involve only gravitational dynamics, and are therefore totally controlled by the cosmological parameters and not by details on galaxy or star formation rates. More importantly an unambiguous detection of the lens effects on CMB maps would be a precious confirmation of the gravitational instability picture. Methods to detect the lens effects on CMB maps have been proposed recently. High order correlation functions, peak ellipticities or large scale lens induced correlators have been proposed for detecting such effects. All of them are however very sensitive to cosmic variance since lens effect is only a sub-dominant alteration of the CMB temperature patterns. The situation is different when one considers the polarization properties. The reason is that in standard cosmological models temperature fluctuations at small scale are dominated by scalar perturbations. Therefore the pseudo-scalar part, the so called $`B`$ component, of the polarization is negligible compared to its scalar part (the $`E`$ component) and can only be significant when CMB lens couplings are present. This mechanism has been recognized in earlier papers. The aim of this paper is to study systematically the properties of the lens induced B field and uncover its properties.
In section II, we perturbatively compute the lens effect on the CMB polarization $`E`$ and $`B`$ field. This first order equation is illustrated by numerical experiments. Possibility of direct reconstruction of the projected mass distribution is also examined. As it has already been noted a significant fraction of the potential wells that deflect the CMB photons can actually be mapped in local weak lensing surveys. This feature has been considered so far in relation to the CMB temperature fluctuations. We extend in Section III these studies to the CMB polarization exploiting the specificities of the field found in previous section. In particular we propose two quantities that can be built from weak lensing and Cosmic Microwave Background polarization surveys, the average value of which does not vanish in presence of CMB lens effects. Compared to direct analysis of the CMB polarization, such tools have the joint advantage of being less sensitive to systematics –systematic errors coming from CMB mapping on one side and weak lensing measurement on the other side have no reason to correlate!– and so emerge even in presence of noisy data, and of being an efficient probe of the cosmological constant. Indeed the expected amplitude of correlation is directly sensitive to the relative length of the optical bench, from the galaxy source plane to the CMB plane, which is mainly sensitive to the cosmological constant. Filtering effects and cosmic variance estimation of such quantities are considered in this section as well.
## II Lens effects on CMB polarization
### II.1 First order effect
Photons emerging from the last scattering surface are deflected by the large scale structures of the Universe that are present on the line-of-sights. Therefore photons observed from apparent direction $`\stackrel{}{\alpha }`$ must have left the last scattering surface from a slightly different direction, $`\stackrel{}{\alpha }+\stackrel{}{\xi }(\stackrel{}{\alpha })`$, where $`\stackrel{}{\xi }`$ is the lens induced apparent displacement at that distance. The displacement field is related to the angular gradient of the projected gravitational potential. In the following, the lens effect will be described by the deformation effects it induces, encoded in the amplification matrix,
$`𝒜^1`$ $`=`$ $`\left(\begin{array}{cc}1\kappa \gamma _1& \gamma _2\\ \gamma _2& 1\kappa +\gamma _1\end{array}\right)`$ (3)
$`=`$ $`\delta _i^j+\xi _{,i}^j`$ (4)
so that
$`\kappa `$ $`=`$ $`{\displaystyle \frac{1}{2}}(\xi _{,x}^x+\xi _{,y}^y)`$
$`\gamma _1`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\xi _{,x}^x\xi _{,y}^y)`$
$`\gamma _2`$ $`=`$ $`\xi _{,x}^y=\xi _{,y}^x.`$ (5)
The lens effect affects the local polarization just by moving the apparent direction of the line of sight. Thus, if we use the Stokes parameters $`Q`$ and $`U`$ to describe the local polarization vector,
$$\stackrel{}{P}=\left(\begin{array}{c}Q\\ U\end{array}\right)$$
we can relate the observed polarization $`\widehat{\stackrel{}{P}}`$ to the primordial one by the relation
$$\widehat{Q}(\stackrel{}{\alpha })=Q(\stackrel{}{\alpha }+\stackrel{}{\xi }),\widehat{U}(\stackrel{}{\alpha })=U(\stackrel{}{\alpha }+\stackrel{}{\xi }).$$
(6)
From now on we will denote $`\widehat{x}`$ an observed quantity and $`x`$ the primordial one. $`\stackrel{}{\alpha }^{}=\stackrel{}{\alpha }+\stackrel{}{\xi }`$ is the sky coordinate system for the observer, therefore the amplification matrix $`𝒜`$ is also the Jacobian of the transformation between the source plane and the image plane. We will restrain here our computation to the weak lensing effect so observed quantity will not take into account any other secondary effect. It is very important at this point to note that the lensing effect does not produce any polarization nor rotate the Stokes parameter. In this regime its effect reduces to a simple deformation of the polarization patterns, similar to the temperature maps. This is the mechanism by which the geometrical properties of the polarization field are changed.
To see that we have to consider the *electric* ($`E`$) and *magnetic* ($`B`$) components instead of the Stokes parameters. At small angular scales (we assume that a small fraction of the sky can be described by a plane), these two quantities are defined as,
$`E`$ $``$ $`\mathrm{\Delta }^1\left[\left(_x^2_y^2\right)Q+2_x_yU\right]`$ (7)
$`B`$ $``$ $`\mathrm{\Delta }^1\left[\left(_x^2_y^2\right)U2_x_yQ\right].`$
This fields reflect non-local geometrical properties of the polarization field. The electric component accounts for the scalar part of the polarization and the magnetic one, the pseudo-scalar part: by parity change $`E`$ is conserved, whereas $`B`$ sign is changed. As it has been pointed out in previous papers, lens effects partly redistribute polarization power in these two fields.
We explicit this latter effect in the weak lensing regime where distortions, $`\kappa `$ and $`\gamma _i`$ components are small. This is indeed expected to be the case when lens effects by the large-scale structures are considered, for which the typical value of the convergence field $`\kappa `$ is expected to be $`1\%`$ at 1 degree scale. The leading order effect is obtained by simply pluging (6) in (7) and by expanding the result at leading order in $`\xi `$ , $`\kappa `$, and $`\gamma `$ . Noting that (these calculations are very similar to those done in ),
$`_i\widehat{X}`$ $`=`$ $`\widehat{_kX}\left(\delta _i^k+\xi _{,i}^k\right)`$ (8)
$`_i_j\widehat{X}`$ $`=`$ $`\widehat{_k_lX}(\delta _i^k+\xi _{,i}^k)(\delta _j^l+\xi _{,j}^l)`$
$`+\widehat{_kX}\xi _{,ij}^k`$
we can write a perturbation description of the lensing effect on electric and magnetic components of the polarization. At leading order one obtains:
$`\mathrm{\Delta }\widehat{E}`$ $`=`$ $`\mathrm{\Delta }E+\xi ^i_i\mathrm{\Delta }E2\kappa \mathrm{\Delta }E`$
$`2\delta _{ij}\left(\gamma ^i\mathrm{\Delta }P^j+\gamma _{,k}^iP^{j,k}\right)+O(\gamma ^2)`$
$`\mathrm{\Delta }\widehat{B}`$ $`=`$ $`\mathrm{\Delta }B+\xi ^i_i\mathrm{\Delta }B2\kappa \mathrm{\Delta }B`$
$`2ϵ_{ij}\left(\gamma ^i\mathrm{\Delta }P^j+\gamma _{,k}^iP^{j,k}\right)+O(\gamma ^2),`$
Where we used the fact that $`\widehat{\mathrm{\Delta }X}=\mathrm{\Delta }X+\xi ^i_i\mathrm{\Delta }X`$ at the leading order. The formulas for $`E`$ and $`B`$ are alike. The only difference stands in the $`\delta _{ij}`$ and $`ϵ_{ij}`$ (the latter is the totally antisymmetric tensor, $`ϵ_{11}=ϵ_{22}=0,`$ $`ϵ_{12}=ϵ_{21}=1`$) that reflects the geometrical properties of the two fields. The first three terms of each of these equations represent the naive effect: the lens induced deformation of the $`E`$ or $`B`$ fields. This effect is complemented by an enhancement effect (respectively $`\kappa \mathrm{\Delta }E`$ and $`\kappa \mathrm{\Delta }B`$) and by shear-polarization mixing terms.The latter effects consist in two parts. One which we will call the $`\mathrm{\Delta }`$-term that couples the shear with second derivative of the polarization field. The other one, hereafter the $``$-term, mixes gradient of the shear and polarization. Although terms like $``$ have been neglected in similar computations we cannot do that here a priori. We will indeed show later that these two terms have similar amplitudes.
One consequence of standard inflationary models on CMB anisotropies is the unbalanced distribution of power between the electric ($`E`$) and magnetic ($`B`$) component of its polarization. Adiabatic scalar fluctuations do not induce $`B`$-type polarization and they dominate at small scales over the tensor perturbations (namely the gravity waves). So, even though gravity waves induce $`E`$ and $`B`$ type polarization in a similar amount, *primary* CMB sky is expected to be completely dominated by $`E`$ type polarization at small scales. Then for this class of models the actual magnetic component of the polarization field is generated by the corrective part of eq. (II.1),
$$\mathrm{\Delta }\widehat{B}=2ϵ_{ij}\left(\gamma ^i\mathrm{\Delta }\widehat{P}^j+\gamma _{,k}^i\widehat{P}^{j,k}\right)$$
(10)
This result extends the direct lens effects described in Benabed & Bernardeau who focused their analysis on the lens effect due to the discontinuity of the polarization field in case of cosmic strings. Previous studies of the weak lensing effect on CMB showed that with lensing, the $`B`$ component becomes important at small scales. We obtain here the same result but with a different method; eq. (10) means that the polarization signal $`P`$ is redistributed by the lensing effect in a way that breaks the geometrical properties of the primordial field. Note here that it is mathematically possible to build a shear field that preserves these geometrical properties and that does not create any $`B`$ signal at small scales. We will discuss this problem in Sec. II.3. It also means that $`B`$ directly reflects the properties of the shear map. We will take advantage of this feature to probe the correlation properties of $`B`$ with the projected mass distribution in next sections.
### II.2 Lens-induced $`B`$ maps
We show examples of lens induced $`B`$ maps. These maps have been calculated using “CMBSlow” code developed by A. Riazuelo (see ) to compute primordial polarization maps (we use realizations of standard CDM model to illustrate lens effects). Then various shear maps are applied. We present both true distortions, (obtained by Delaunay triangulation used to shear the $`Q`$ and $`U`$ fields), and the first order calculations given by eq. (10).
Fig. 1 presents the shear effect induced by an isothermal sphere with finite core radius (and the lens edges have been suppressed by an exponential cutoff to minimize numerical noise). The agreement between true distortion (central panel ) and first order formula (right panel) is good. However, a close examination of the maps reveals that some structures in the true map are slightly wider than their counterparts in the first order map. This error is more severe in the center, where the distortion is bigger, which is to be expected since the limits of the validity region of first order calculations are reached.
Fig. 2 shows the $`B`$ field induced by a *realistic* distortion. We use second order Lagrangian dynamics to create a $`2.5\times 2.5`$ degree map that mimics a realistic projected mass density up to $`z=1000`$ and used its gravitational distortion to compute a typical weak lensing-induced $`B`$ map. Again we compare the *exact* effect (i.e. left panel where Delaunay triangulation is used) and the first order formula (middle panel). Right panel shows the difference between the two maps. It reveals the locations where the two significantly disagree. In fact most disagreements are due to slight mismatch of the $`B`$ patch positions, which lead to dipole like effects in this map.
We also show here a comparison of the two parts of the first order formula eq. (10) in order to see which of the $`\mathrm{\Delta }`$ or $``$ terms dominates. It would be more comfortable if one of the two terms was dominant, however, Fig. 3 shows that it is not the case. Even if the $`\mathrm{\Delta }`$-term dominates at low ($`<1000)`$ $`\mathrm{}`$, it is only twice bigger than $``$-one at this scale. The inverse is true for higher ($`30005000)`$ $`\mathrm{}`$s. This can be seen by looking at Fig. 4 where we show the relative amplitudes of the $`\mathrm{\Delta }`$ and $``$ contributions. The $`\mathrm{\Delta }`$ part gives birth to large patches (around $`10^{}`$) while $``$ panel shows a lot more of small features.
### II.3 Direct reconstruction – Kernel problem
The fact that the observable $`B`$ is at leading order proportional to the weak lensing signal invites us to try a direct reconstruction, similar to the lensing mass reconstruction. In fact, we can write
$$\mathrm{\Delta }\widehat{B}=2ϵ_{ij}\left(\gamma ^i\mathrm{\Delta }\widehat{P}^j+\gamma _{,k}^i\widehat{P}^{j,k}\right)\mathrm{F}[\gamma ]$$
(11)
and our reconstruction problem becomes an inversion problem for the operator $`\mathrm{F}`$. Unfortunately, one can prove that this problem has no unique solution. It is due to the fact that $`\mathrm{F}`$ admits a huge kernel, in the sense that, given a polarization map, there is a wide class of shear fields that will conserve a null $`B`$ polarization. The demonstration of this property is sketched in the following.
Since the unlensed polarization is only electric in our approximation, we can describe it by the Laplacian of a scalar field ;
$$E\mathrm{\Delta }\phi \text{ so }\{\begin{array}{c}Q=\left(_x^2_y^2\right)\phi \hfill \\ U=2_x_y\phi \hfill \end{array}.$$
(12)
The same holds for the shear and convergence fields
$$\kappa \frac{\mathrm{\Delta }\psi }{2},\gamma _1=\frac{1}{2}\left(_x^2_y^2\right)\psi \text{ },\gamma _2=_x_y\psi .$$
(13)
Thus we need to know, for a given $`\phi `$ field, whether there is any $`\psi `$ that fulfills the equation
$$\gamma _2\mathrm{\Delta }Q\gamma _1\mathrm{\Delta }U+_i\gamma _2^iQ_i\gamma _1^iU=0.$$
(14)
$`\phi `$ and $`\psi `$ can be written as polynomial decompositions
$`\phi (x,y)`$ $`=`$ $`{\displaystyle \underset{n,l}{}}a_{nl}x^ny^l`$
$`\psi (x,y)`$ $`=`$ $`{\displaystyle \underset{m,k}{}}b_{mk}x^my^k.`$ (15)
Using (15) in (14) we are left with a new polynomial whose coefficients $`c_{ij}`$ are sums of $`a_{nl}\times b_{mk}`$ and have to be all put to zero. With the coefficient equations in hand, it is easy to prove that assuming all the $`b_{mk}`$ coefficient up to $`m+k=N`$ are known and writing the equations $`i+j=(N+1)3,c_{ij}=0`$, we can compute out of all the $`a_{nl}`$ all but three $`b_{mk}`$ with $`m+k=N+1`$. This is somewhat similar to mass reconstruction problems from galaxy surveys where one cannot avoid the mass sheet degeneracy. The situation is however worse in our case since not only constant convergence but also translations and a whole class of $`a_{nk}`$ realization dependent complex deformations are indiscernible. Thus, with the only knowledge of the $`B`$ component of the polarization one cannot, with the first order eq. (10), recover the projected mass distribution.
## III Cross-correlating CMB maps and weak lensing surveys
### III.1 Motivations
Even with the most precise experiments it is clear that clean detections of $`B`$ component will be difficult to obtain. The magnetic polarization amplitude induced with such a mechanism is expected to be one order of magnitude below the electric one. Besides even if we know that there is a window in angular scale where the other secondary effects will not interfere too much with the detection of the lens-induced $`B`$, few is known about removing the foregrounds to obtain clean maps reconstruction algorithms would require.
These considerations lead us to look for complementary data sets to compare $`B`$ with. Although the source plane for weak lensing surveys is much closer than for the lensed CMB fluctuations, we expect to have a significant overlapping region in the two redshift lens distributions, so that weak lensing surveys can map a fair fraction of the line-of-sight CMB lenses. Consequently, weak lensing surveys can potentially provide us with shear maps correlated with $`B`$, but which have different geometrical degeneracy, noise sources and systematics than the polarization field.
The correlation strength between the lensing effects at two different redshifts can be evaluated. We define $`r`$ as the cross-correlation coefficient between two lens planes:
$$r(z_{\text{gal}})=\frac{\kappa \kappa _{\text{gal}}}{\sqrt{\kappa ^2\kappa _{\text{gal}}^2}}.$$
(16)
In a broad range of realistic cases (see tab. 1), $`r40\%`$. To take advantage of this large overlapping we will consider quantity that cross correlates the CMB $`B`$ field and galaxy surveys. Moreover, cross-correlation observations are expected to be insensitive to noises in weak lensing surveys and in CMB polarization maps. This idea has already been explored for temperature maps. We extend this study here taking advantage of the specific geometrical dependences uncovered in the previous section.
### III.2 Definition of $`b_\mathrm{\Delta }`$ and $`b_{}`$.
The magnetic component of the polarization in eq. (10) appears to be built from a pure CMB part, which comes from the primordial polarization, and a gravitational lensing part. It is natural to define $`b`$, in such a way that mimics the $`\mathrm{\Delta }\widehat{B}`$ fonction dependance, by replacing the CMB shear field by the galaxy one.
$`b`$ $`=`$ $`ϵ_{ij}\left(\gamma _{\text{gal}}^i\mathrm{\Delta }\widehat{P}^j+\gamma _{\text{gal},k}^i\widehat{P}^{j,k}\right)`$
$`=`$ $`ϵ_{ij}\left(\gamma _{\text{gal}}^i\mathrm{\Delta }P^j+\gamma _{\text{gal},k}^iP^{j,k}\right)+O(\kappa ^2).`$
In the following, we will label local lensing quantities, such as what one can obtain from lensing reconstruction on galaxy surveys, with a gal index. This new quantity can be viewed as a guess for the CMB polarization $`B`$ component if lensing was turned on only in a redshift range matching the depth of galaxy surveys. The correlation coefficient of this guess with the true $`\mathrm{\Delta }B`$ field, that is $`\mathrm{\Delta }\widehat{B}b`$, is expected to be quadratic both in $`P`$ and in $`\gamma `$ and to be proportional to the cross-coefficient $`r`$.
For convenience, and in order to keep the objects we manipulate as simple as possible, we will not exactly implement this scheme, as it will lead to uneven angular derivative degrees in the two terms of resulting equations. We can, instead, decompose the effect in the $`\mathrm{\Delta }`$ and $``$-part. These two are not correlated, since their components do not share the same degrees of angular derivation<sup>1</sup><sup>1</sup>1generically, a random field and its derivative at the same point are not correlated. . Hence, we can play the proposed game, considering the two terms of eq. (10) as if they were two different fields, creating two guess-quantities that should correlate independently with the observed $`B`$ field. Following this idea we build $`b_\mathrm{\Delta }`$ as,
$`b_\mathrm{\Delta }`$ $``$ $`ϵ_{ij}\gamma _{\text{gal}}^i\mathrm{\Delta }\widehat{P}^j`$
$`=`$ $`ϵ_{ij}\gamma _{\text{gal}}^i\mathrm{\Delta }P^j+O(\kappa ^2)`$
which corresponds to the $`\mathrm{\Delta }`$-term in eq. (10). The amplitude of the cross-correlation between $`\mathrm{\Delta }B`$ and $`b_\mathrm{\Delta }`$ can easily be estimated. At leading order, we have
$$\mathrm{\Delta }\widehat{B}b_\mathrm{\Delta }=2ϵ_{ij}ϵ_{kl}\gamma ^k\gamma _{\text{gal}}^i\mathrm{\Delta }P^l\mathrm{\Delta }P^j.$$
(19)
The corresponding $``$ correlation is
$$\mathrm{\Delta }\widehat{B}b_{}=2ϵ_{ij}ϵ_{kl}_m\gamma ^k_n\gamma _{\text{gal}}^i_mP^l_nP^j$$
(20)
where we have defined
$$b_{}ϵ_{ij}_k\gamma _{\text{gal}}^i_k\widehat{P}^j.$$
(21)
Fig. 4 shows numerical simulations presenting maps of first order $`\mathrm{\Delta }\widehat{B}`$, its $`\mathrm{\Delta }`$ and $``$ contributions and the corresponding guess maps one can build with a low $`z`$ shear map. The similarities between the top maps and the bottom maps are not striking. Yet, under close examination one can recognize individual patterns shared between the maps. This is confirmed by the computation of the correlation coefficient between the maps, that shows significant overlapping, between 50% and 15%, depending correlation and filtering strategy. The calculations hereafter will evaluate the theoretical correlation structure between maps given in figs. 4-b and 4-g & h.
For galaxy surveys, the amplification matrix is,
$`𝒜_{\text{gal}}^1(\stackrel{}{\alpha })\text{Id}={\displaystyle _0^{z_{\text{gal}}}}d\chi w_{\text{gal}}(\chi )`$ (22)
$`\times {\displaystyle }{\displaystyle \frac{\mathrm{d}^3k}{(2\pi )^{\frac{3}{2}}}}\delta (\stackrel{}{k})\mathrm{e}^{\mathrm{i}\left(k_r\chi +\stackrel{}{k}_{}D(\chi )\stackrel{}{\alpha }\right)}`$
$`\times \left(\begin{array}{cc}1+\mathrm{cos}(2\varphi _k_{})& \mathrm{sin}(2\varphi _k_{})\\ \mathrm{sin}(2\varphi _k_{})& 1\mathrm{cos}(2\varphi _k_{})\end{array}\right)`$ (25)
where $`\delta (k)`$ is the Fourier transform of the density contrast at redshift $`z(\chi )`$, $`w`$ is the lens efficiency function, $`\mathrm{D}`$ is the angular distance, and $`\varphi _k_{}`$is the position angle of the transverse wave-vector $`k_{}`$in the $`k_{}=(k_x,k_y)`$ plane. Assuming a Dirac source distribution the efficiency function is given by
$$w_{\text{gal}}(z)=\frac{3}{2}\mathrm{\Omega }_\mathrm{o}\frac{\mathrm{D}_z\mathrm{D}_{zz_{\text{gal}}}}{a\mathrm{D}_{z_{\text{gal}}}}.$$
(26)
Note that the Fourier components $`\delta (k)`$ include the density time evolution. They are thus proportional to the growth factor in the linear theory. The time evolution of these components is much more complicated in the nonlinear regime (see ).
Then, $`b_{\mathrm{}}`$ is
$$b_{\mathrm{}}(\stackrel{}{\alpha })=^{\chi _{\text{gal}}}𝒟(\chi ,\stackrel{}{l},\stackrel{}{k})\stackrel{~}{E}(l)\delta (k)𝒢_{\mathrm{}}^{\mathrm{Ker}}(\stackrel{}{l},\stackrel{}{k}_{})$$
(27)
with the integration element defined as,
$$𝒟(\chi ,\stackrel{}{l},\stackrel{}{k})=\mathrm{d}\chi w_{\text{gal}}(\chi )\frac{\mathrm{d}^3k}{(2\pi )^{3/2}}\frac{\mathrm{d}^2l}{2\pi }\mathrm{e}^{\mathrm{i}\left[k_r\chi +\left(\stackrel{}{k}_{}\mathrm{D}(\chi )+\stackrel{}{l}\right)\stackrel{}{\alpha }\right]},$$
(it actually depends on the position of the source plane through the efficiency function $`w(z)`$) and where $`\mathrm{}`$ stands for either $`\mathrm{\Delta }`$ or $``$. The geometrical kernel $`𝒢^{\mathrm{Ker}}`$ is given by (using eq. (12))
$`𝒢_\mathrm{\Delta }^{\mathrm{Ker}}(\stackrel{}{l},\stackrel{}{k})`$ $``$ $`l^2\mathrm{sin}2\left(\varphi _k\varphi _l\right)`$ (28)
$`𝒢_{}^{\mathrm{Ker}}(\stackrel{}{l},\stackrel{}{k})`$ $``$ $`lk\mathrm{cos}\left(\varphi _k\varphi _l\right)\mathrm{sin}2\left(\varphi _k\varphi _l\right).`$ (29)
This function contains all the geometrical structures of the $`\mathrm{\Delta }`$ and $``$ terms. We can write the same kind of equation for $`\mathrm{\Delta }\widehat{B}`$. Then, the cross-correlation is
$`\mathrm{\Delta }\widehat{B}b_{\mathrm{}}(\stackrel{}{\alpha })`$ $`=`$ $`2{\displaystyle ^{\chi _{\text{gal}}}}𝒟(\chi _{\text{gal}},\stackrel{}{l}_{\text{gal}},\stackrel{}{k}_{\text{gal}})`$
$`\times {\displaystyle ^{\chi _{\text{cmb}}}}𝒟(\chi _{\text{cmb}},\stackrel{}{l}_{\text{cmb}},\stackrel{}{k}_{\text{cmb}})𝒢_{\mathrm{}}^{\mathrm{Ker}}(\stackrel{}{l}_{\text{gal}},\stackrel{}{k}_{\text{gal}})`$
$`\times 𝒢_{\mathrm{}}^{\mathrm{Ker}}(\stackrel{}{l}_{\text{cmb}},\stackrel{}{k}_{\text{cmb}})\delta (\stackrel{}{k}_{\text{gal}})\delta (\stackrel{}{k}_{\text{cmb}})\stackrel{~}{E}(\stackrel{}{l}_{\text{gal}})\stackrel{~}{E}(\stackrel{}{l}_{\text{cmb}}).`$
The completion of this calculation requires the use of the small angle approximation,
$`\delta (\stackrel{}{k}_{\text{gal}})\delta (\stackrel{}{k}_{\text{cmb}})`$ $`=`$ $`P(k)\delta ^3(k_{\text{gal}}+k_{\text{cmb}})`$
$`P(k_{})\delta ^2(k_{\text{gal}_{}}+k_{\text{cmb}_{}})\delta (k_{\text{gal}_r}+k_{\text{cmb}_r})`$
which implies
$$\stackrel{}{k}_{\text{gal}}=\stackrel{}{k}_{\text{cmb}}=\stackrel{}{k}$$
(32)
and after the radial components have been integrated out,
$$\chi _{\text{gal}}=\chi _{\text{cmb}}=\chi .$$
(33)
We also define the $`C_E(l)`$ as the angular power spectrum of the $`E`$ field,
$$\stackrel{~}{E}(\stackrel{}{l}_{\text{gal}})\stackrel{~}{E}(\stackrel{}{l}_{\text{cmb}})=C_E(l)\delta ^2(l_{\text{gal}}l_{\text{cmb}})$$
(34)
Eventually one gets,
$`\mathrm{\Delta }\widehat{B}b_{\mathrm{}}(\stackrel{}{\alpha })`$ $`=`$ $`2{\displaystyle ^{z_{\text{gal}}}}d\chi w_{\text{gal}}w_{\text{cmb}}{\displaystyle \frac{\mathrm{d}^2k\mathrm{d}^2l}{(2\pi )^4}}`$
$`\times C_E(l)P(k)𝒢_{\mathrm{}}^{\mathrm{Ker}}(\stackrel{}{l},\stackrel{}{k})^2`$
Then, integrating on the geometrical dependencies in $`𝒢_{\mathrm{}}^{\mathrm{Ker}}`$, we have
$`\mathrm{\Delta }\widehat{B}b_\mathrm{\Delta }(\stackrel{}{\alpha })`$ $`=`$ $`2{\displaystyle ^{z_{\text{gal}}}}d\chi w_{\text{gal}}w_{\text{cmb}}`$ (36)
$`\times {\displaystyle }{\displaystyle \frac{\mathrm{d}k\mathrm{d}l}{2(2\pi )^2}}kl^5C_E(l)P(k)`$
$`=`$ $`\mathrm{\Delta }E^2\kappa \kappa _{\text{gal}},`$
and
$`\mathrm{\Delta }\widehat{B}b_{}(\stackrel{}{\alpha })`$ $`=`$ $`{\displaystyle ^{z_{\text{gal}}}}d\chi w_{\text{gal}}w_{\text{cmb}}`$ (37)
$`\times {\displaystyle }{\displaystyle \frac{\mathrm{d}k\mathrm{d}l}{2(2\pi )^2}}k^3l^3C_E(l)P(k)`$
$`=`$ $`{\displaystyle \frac{1}{2}}(\stackrel{}{}E)^2\stackrel{}{}\kappa \stackrel{}{}\kappa _{\text{gal}},`$
implying that, ignoring filtering effects, we are able to measure directly the correlation between lensing effect at $`z_{\text{cmb}}`$ and any $`z_{\text{gal}}`$ a weak lensing survey can access. Since $`\mathrm{\Delta }\widehat{E}=\mathrm{\Delta }E(1+O(\kappa ))`$ we get, for the $`\mathrm{\Delta }`$ type quantity,
$`\mathrm{\Delta }\widehat{E}^2`$ $`=`$ $`\mathrm{\Delta }E^2\left(1+O(\kappa )\right)^2`$
$`=`$ $`\mathrm{\Delta }E^2\left(1+O\left(\kappa ^2\right)\right).`$
The same holds for $`.`$ We are then able to construct two quantities insensitive to the normalization of CMB and $`\sigma _8`$
$`𝒳_\mathrm{\Delta }`$ $``$ $`{\displaystyle \frac{\mathrm{\Delta }\widehat{B}b_\mathrm{\Delta }(\stackrel{}{\alpha })}{\mathrm{\Delta }\widehat{E}^2\kappa _{\text{gal}}^2}}={\displaystyle \frac{\kappa \kappa _{\text{gal}}}{\kappa _{\text{gal}}^2}}`$
$``$ $`r\sqrt{{\displaystyle \frac{\kappa ^2}{\kappa _{\text{gal}}^2}}}.`$
and
$`𝒳_{}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Delta }\widehat{B}b_{}(\stackrel{}{\alpha })}{(\stackrel{}{}\widehat{E})^2(\stackrel{}{}\kappa _{\text{gal}})^2}}={\displaystyle \frac{1}{2}}{\displaystyle \frac{\stackrel{}{}\kappa \stackrel{}{}\kappa _{\text{gal}}}{\stackrel{}{}\kappa _{\text{gal}}^2}}`$
$``$ $`{\displaystyle \frac{1}{2}}r_{}\sqrt{{\displaystyle \frac{\kappa ^2}{\kappa _{\text{gal}}^2}}}.`$
We implicitly defined $`r_{}`$ like $`r`$ but with $`\kappa `$ instead of $`\kappa `$
$$r_{}(z_{\text{gal}})=\frac{\stackrel{}{}\kappa \stackrel{}{}\kappa _{\text{gal}}}{\sqrt{(\kappa )^2(\kappa _{\text{gal}})^2}}.$$
(41)
We will see in Sect. III.4 that they behave very much alike. This result is to be compared with the formula for $`\mathrm{cos}(\theta _g)`$ established in where the obtained quantity was going like $`r\sqrt{\kappa ^2}`$. These calculations however have neglected the filtering effects that may significantly affect our conclusions. These effects are investigated in next section.
### III.3 Filtering effects
In above section we conduct our calculations assuming no filtering. Obviously we have to take it into account! We will show here that the results obtained before hold, in certain limits, when one adds filtering effects.
In the following, we consider, for simplicity, top-hat filters only. It is expected that other window functions will show very similar behaviors and this simplification does not restrain the generality of our results. Let us call $`W(x)`$ the top-hat filter function in Fourier space
$$W(x)2\frac{\mathrm{J}_1(x)}{x}.$$
(42)
$`\mathrm{J}_1`$ is the first $`\mathrm{J}`$-Bessel function. We will also define $`W_i(x)`$ a general function
$$W_i(x)2\frac{\mathrm{J}_i(x)}{x}$$
(43)
where $`\mathrm{J}_i`$ is the $`i^{\mathrm{th}}`$ $`\mathrm{J}`$-Bessel function, so that $`W=W_1`$. Then, if $`X(\stackrel{}{\alpha })`$ is the value of any quantity $`X`$ at position $`\stackrel{}{\alpha }`$ on the sky, its top-hat filtered value can be computed as
$$X_{(\theta )}(\stackrel{}{\alpha })=\frac{\mathrm{d}^2k}{2\pi }\stackrel{~}{X}_kW(k\theta )\mathrm{e}^{\mathrm{i}\stackrel{}{k}\stackrel{}{\alpha }},$$
(44)
where $`\stackrel{~}{X}`$ is $`X`$ Fourier transform. In the following we will note $`X_{(\theta )}`$ the filtered quantity at scale $`\theta `$.
The tricky thing for $`\mathrm{\Delta }\widehat{B}b_{\mathrm{}}`$ is that the CMB part and the low-redshift weak lensing part are *a priori* filtered at different scale. For $`\mathrm{\Delta }\widehat{B}`$, which is a measured value, its pure CMB part and its weak lensing part are filtered at the same scale $`\theta `$. Hence, $`\widehat{B}`$ reads,
$`\mathrm{\Delta }\widehat{B}(\stackrel{}{\alpha })_{(\theta )}`$ $`=`$ $`2{\displaystyle ^{\chi _{\text{cmb}}}}𝒟(\chi ,\stackrel{}{l},\stackrel{}{k})\stackrel{~}{E}(l)\delta (k)`$ (45)
$`\times \left[𝒢_\mathrm{\Delta }^{\mathrm{Ker}}(\stackrel{}{l},\stackrel{}{k}_{})+𝒢_\mathrm{\Delta }^{\mathrm{Ker}}(\stackrel{}{l},\stackrel{}{k}_{})\right]W\left(|\stackrel{}{k}_{}\mathrm{D}+\stackrel{}{l}|\theta \right)`$
*A contrario* $`b_{\mathrm{}}`$ is a composite value. The CMB part is still filtered at $`\theta `$ whereas the weak lensing part (which comes from a weak lensing survey of galaxies) is filtered independently at another scale which we denote $`\theta _{\text{gal}}`$. It implies that,
$`b_{\mathrm{}}(\stackrel{}{\alpha })_{(\theta )}`$ $`=`$ $`2{\displaystyle ^{\chi _{\text{gal}}}}𝒟(\chi ,\stackrel{}{l},\stackrel{}{k})\stackrel{~}{E}(l)\delta (k)`$ (46)
$`\times 𝒢_{\mathrm{}}^{\mathrm{Ker}}(\stackrel{}{l},\stackrel{}{k}_{})W(k\mathrm{D}\theta _{\text{gal}})W(l\theta ).`$
Taking filtering into account, the cross-correlation coefficient becomes,
$`\mathrm{\Delta }\widehat{B}_{(\theta )}b_{\mathrm{}(\theta ,\theta _{\text{gal}})}`$ $`=`$ $`2{\displaystyle ^{z_{\text{gal}}}}d\chi w_{\text{gal}}w_{\text{cmb}}`$
$`\times {\displaystyle }{\displaystyle \frac{\mathrm{d}^2k\mathrm{d}^2l}{(2\pi )^4}}C_E(l)P(k)𝒢_{\mathrm{}}^{\mathrm{Ker}}(\stackrel{}{l},\stackrel{}{k})`$
$`\times W(k\mathrm{D}\theta _{\text{gal}})W(l\theta )W(|\stackrel{}{k}\mathrm{D}+\stackrel{}{l}|\theta ).`$
It can be shown (from the summation theorems of the Bessel functions) that,
$`W_1(|\stackrel{}{k}\mathrm{D}+\stackrel{}{l}|\theta )=`$ (48)
$`{\displaystyle \underset{i=1}{}}iW_i(k\mathrm{D}\theta )W_i(l\theta )(1)^i{\displaystyle \frac{\mathrm{sin}i(\varphi _k\varphi _l)}{\mathrm{sin}(\varphi _k\varphi _l)}}`$
It is then possible to break the $`W\left(|\stackrel{}{k}\mathrm{D}+\stackrel{}{l}|\theta \right)`$ into a sum of $`W_i(k\mathrm{D}\theta )W_i(l\theta )`$ with coefficients that depend on the geometrical properties of our problem. Integrating over the geometrical dependencies of $`𝒢_{\mathrm{}}^{\mathrm{Ker}}`$, leaves us with only a few non vanishing terms in our sum,
$$d\varphi \mathrm{sin}^2(2\varphi )\frac{\mathrm{sin}(i\varphi )}{\mathrm{sin}\varphi }=\{\begin{array}{cc}\pi \hfill & i\text{ }=\text{ }1\text{ or }i=3\hfill \\ 0\hfill & \text{elsewhere}\hfill \end{array},$$
(49)
for the $`\mathrm{\Delta }`$-term and
$$d\varphi \mathrm{cos}\varphi \mathrm{sin}^2(2\varphi )\mathrm{sin}(i\varphi )\mathrm{sin}\varphi =\{\begin{array}{cc}\pi /2& i=1\hfill \\ 3\pi /4& i=3\hfill \\ \pi /4& i=5\hfill \\ 0& \text{elsewhere}\hfill \end{array},$$
(50)
for the $``$-term. Each term can be computed exactly, and it turns out that the terms built from $`W_i,i>1`$ are always negligible compared to the ones coming from $`W_1`$. It implies that we can safely ignore the $`W_3`$ and $`W_5`$ in both $`\mathrm{\Delta }`$ and $``$ expressions, therefore it is reasonable to assume that $`W\left(|\stackrel{}{k}\mathrm{D}+\stackrel{}{l}|\theta \right)=W(k\mathrm{D}\theta )W(l\theta ).`$ It is expected that other windows, in particular the Gaussian window function, share similar properties. Then, taking into accounts the filtering effects, the equations for the cross-correlations reduce to
$$\mathrm{\Delta }\widehat{B}_{(\theta )}b_{\mathrm{\Delta }(\theta ,\theta _{\text{gal}})}=\mathrm{\Delta }E_{(\theta )}^2\kappa _{(\theta )}\kappa _{\text{gal}(\theta _{\text{gal}})}$$
(51)
and
$$\mathrm{\Delta }\widehat{B}_{(\theta )}b_{(\theta ,\theta _{\text{gal}})}=\frac{1}{2}E_{(\theta )}^2\kappa _{(\theta )}\kappa _{\text{gal}(\theta _{\text{gal}})},$$
(52)
so that our correlation coefficients can be written,
$$𝒳_{\mathrm{\Delta }(\theta ,\theta _{\text{gal}})}=r_{(\theta ,\theta _{\text{gal}})}\sqrt{\frac{\kappa _{(\theta )}^2}{\kappa _{\text{gal}(\theta _{\text{gal}})}^2}}$$
(53)
and
$$𝒳_{(\theta ,\theta _{\text{gal}})}=\frac{1}{2}r_{(\theta ,\theta _{\text{gal}})}\sqrt{\frac{\kappa _{(\theta )}^2}{\kappa _{\text{gal}(\theta _{\text{gal}})}^2}}.$$
(54)
The results obtained in eqs. (III.2-III.2) are thus still formally valid. Actually, eqs. (53-54) simply tell that filtering effects can simply be assumed to act independently on the lensing effects and on the primary Cosmic Microwave Background maps. We are left with two quantities that only reflects the line-of-sight overlapping effects of lensing distortions.
### III.4 Sensitivity to the cosmic parameters
We quickly explore here the behavior of $`𝒳_{\mathrm{}}`$ in different sets of cosmological parameters. These quantities only depend on weak lensing quantities. Ignoring the $`\mathrm{\Omega }_0`$ dependence in the angular distances and growing factor, one would expect $`\kappa ^2`$ to scale like $`\mathrm{\Omega }_0^2`$. Yet, because of the growth factor, the convergence field exhibits a weaker sensitivity to $`\mathrm{\Omega }_0`$. Assuming $`\mathrm{\Lambda }=0`$ and a power law spectrum, we know from that $`\kappa _{\text{gal}}^2\mathrm{\Omega }_0^{1.66}`$ for $`z_{\text{gal}}=1`$. The same calculation leads to $`\kappa _{\text{cmb}}\kappa _{\text{gal}}\mathrm{\Omega }_0^{1.68}`$, $`(\kappa _{\text{gal}})^2\mathrm{\Omega }_0^{1.91}`$ and $`\stackrel{}{}\kappa _{\text{cmb}}\stackrel{}{}\kappa _{\text{gal}}\mathrm{\Omega }_0^{1.915}`$. Then, in this limit, the quantities $`𝒳_{\mathrm{}}`$ have a very low dependence on $`\mathrm{\Omega }_0`$ :
$$𝒳_\mathrm{\Delta }\mathrm{\Omega }_0^{0.02}\text{ and }𝒳_{}\mathrm{\Omega }_0^{0.005}.$$
Eventually, the $`𝒳_{\mathrm{}}`$ quantities should exhibit a seizable sensitivity to $`\mathrm{\Lambda }`$; changing $`\mathrm{\Lambda }`$ increases or reduces the size of the optic bench and accordingly the overlapping between $`\kappa _{\text{cmb}}`$ and $`\kappa _{\text{gal}}`$.
Figs. 5 and 6 present contour plots of the amplitude of $`𝒳_\mathrm{\Delta }`$ and $`𝒳_{}`$ in the $`(\mathrm{\Omega }_0,\mathrm{\Lambda })`$ plane for CDM models. They show the predicted low $`\mathrm{\Omega }_0`$ sensitivity and the expected $`\mathrm{\Lambda }`$ dependency. Both figures are very alike. This is due to the fact that the dominant features are contained in the efficiency function dependences on the angular distances.
### III.5 Cosmic variance
In previous sections we looked at the sensitivity of observable quantities which mixed galaxy weak lensing surveys and CMB polarization detection. It is very unlikely that both surveys will be able to cover, with a good resolution and low foreground contamination, a large fraction of the sky. It seems however reasonable to expect to have at our disposal patches of at least a few hundreds square degrees. The issue we address in this section is to estimate the cosmic variance of such a detection in joint surveys in about 100 square degrees.
The computation of cosmic variance is a classical problem in cosmological observation . A natural estimate for an ensemble average $`X`$ is its geometrical average. If the survey has size $`\mathrm{\Sigma }`$ then,
$$\overline{X}=\frac{1}{\mathrm{\Sigma }}_\mathrm{\Sigma }\mathrm{d}^2\alpha X(\stackrel{}{\alpha })$$
(55)
For a compact survey with circular shape of radius $`\mathrm{\Xi }`$ we formally have,
$$\overline{X}=\frac{\mathrm{d}^2k}{2\pi }\stackrel{~}{X}(\stackrel{}{k})W(k\mathrm{\Xi }).$$
(56)
For sake of simplicity this is what we use in the following but we will see that the shape of the survey has no significant consequences.
Taking $`\overline{X}`$ as an estimate of $`X`$ (the ensemble average of $`X`$), leads to an error of the order $`\sqrt{\overline{X}^2\overline{X}^2}`$which usually scales like $`1/\sqrt{\mathrm{\Sigma }}`$ if the survey is large enough.
When we are measuring $`𝒳_{\mathrm{}}`$ on a small patch of the sky, we are apart from the statistical value by the same kind of error. We can neglect the errors on $`\mathrm{\Delta }\widehat{E}^2`$, $`(\widehat{E})^2`$, $`(\kappa _{\text{gal}})^2`$and $`\kappa _{\text{gal}}^2`$; those may not be the dominant source of discrepancy and can even be measured on wider and independent samples. The biggest source of error is the measure of $`\mathrm{\Delta }\widehat{B}b_{\mathrm{}}`$. It is given by,
$$C_{\mathrm{}}=\sqrt{\left(\overline{\mathrm{\Delta }\widehat{B}b_{\mathrm{}}}\overline{\mathrm{\Delta }\widehat{B}}\overline{b_{\mathrm{}}}\right)^2\overline{\mathrm{\Delta }\widehat{B}b_{\mathrm{}}}\overline{\mathrm{\Delta }\widehat{B}}\overline{b_{\mathrm{}}}^2}.$$
(57)
The computation of (57) is made easier if we write explicitly the geometrical average as a summation over $`N`$ measurement points ($`N`$ can be as large as we want),
$$\overline{X}=\frac{1}{N}\underset{i=1}{\overset{N}{}}X(\theta _i),$$
(58)
we then developed (57), and replace the ensemble average of the summation sign by the geometrical average over the survey size. We are left with a sum of correlators containing $`8`$ fields taken at 2, 3 and 4 different points. The calculations can be carried out analytically if we assume that all our fields follow Gaussian statistics, which is reasonable at the scale we are working on. In that case indeed, we can take advantage of the Wick theorem to contract each of the $`8`$ fields correlators in products of $`2`$ points correlation functions. By definition, (57) contains only connected correlators, moreover the ensemble averages $`\mathrm{\Delta }\widehat{B}`$ and $`b_{\mathrm{}}`$ vanish, therefore only a small fraction of correlators among all the possible combination of the $`8`$ fields survive. We can use a simple diagrammatic representations to describe their geometrical shape. All the non vanishing terms in $`C_{\mathrm{}}`$ are given in Fig. 7. Each line between two vertex represents a $`2`$ points correlation function such as $`X(\stackrel{}{\alpha }_1)X(\stackrel{}{\alpha }_2)`$, and the different symbols at the vertex correspond to different $`X`$ fields (the cross stands for $`\mathrm{\Delta }P`$, the dot for $`\gamma _{\text{cmb}}`$, and the open dot stands for $`\gamma _{\text{gal}}`$). The $`𝒜`$-terms represent terms where the two top (and the two bottom) $`\mathrm{\Delta }B`$ and $`b_{\mathrm{}}`$ are taken at the same point, but top and bottom fields are not at the same place. The $``$-terms are three points diagrams: the top $`\mathrm{\Delta }B`$ and $`b_{\mathrm{}}`$ are at the same point whereas the right and left bottom vertexes are at two different locations. The $`𝒞`$ terms are four-points diagrams, where each vertex is at a different point. To illustrate our notations, let us write $`_{2c}^{\mathrm{}}`$ as an example,
$$\begin{array}{ccc}_{2c}^{\mathrm{}}& =& \gamma _{\text{cmb}}(\stackrel{}{\alpha }_1)\gamma _{\text{gal}}(\stackrel{}{\alpha }_2)\gamma _{\text{gal}}(\stackrel{}{\alpha }_3)\gamma _{\text{cmb}}(\stackrel{}{\alpha }_1)\times \\ & & \mathrm{\Delta }P(\stackrel{}{\alpha }_1)\mathrm{\Delta }P(\stackrel{}{\alpha }_1)\mathrm{\Delta }P(\stackrel{}{\alpha }_2)\mathrm{\Delta }P(\stackrel{}{\alpha }_3)\end{array}$$
We only focus on the calculation of the $`𝒜`$ terms because we can use the approximation that
$$𝒜𝒞.$$
(59)
Indeed, in perturbative theory, if the survey is large enough, the $`n`$-points correlation functions naturally dominates over the $`n+1`$-points correlation function. This is true as long as the local variance is much bigger than the autocorrelation at survey scale and we assume the surveys are still large enough to be in this case.
The general expression for any $`𝒜`$ diagram is
$`𝒜_i^{\mathrm{}}`$ $`=`$ $`4{\displaystyle ^{\text{cmb}}}𝒟(\chi _{\text{cmb}1},\stackrel{}{l}_{\text{cmb}1},\stackrel{}{k}_{\text{cmb}1})𝒟(\chi _{\text{cmb}2},\stackrel{}{l}_{\text{cmb}2},\stackrel{}{k}_{\text{cmb}2})`$
$`\times {\displaystyle ^{\text{gal}}}𝒟(\chi _{\text{gal}1},\stackrel{}{l}_{\text{gal}1},\stackrel{}{k}_{\text{gal}1})𝒟(\chi _{\text{gal}2},\stackrel{}{l}_{\text{gal}2},\stackrel{}{k}_{\text{gal}2})`$
$`\times 𝒢_{\mathrm{}}^{\mathrm{Ker}}(\stackrel{}{l}_{\text{cmb}1},\stackrel{}{k}_{\text{cmb}1_{}})𝒢_{\mathrm{}}^{\mathrm{Ker}}(\stackrel{}{l}_{\text{cmb}2},\stackrel{}{k}_{\text{cmb}2_{}})`$
$`\times 𝒢_{\mathrm{}}^{\mathrm{Ker}}(\stackrel{}{l}_{\text{gal}1},\stackrel{}{k}_{\text{gal}1_{}})𝒢_{\mathrm{}}^{\mathrm{Ker}}(\stackrel{}{l}_{\text{gal}2},\stackrel{}{k}_{\text{gal}2_{}})_i\stackrel{}{k}_i\stackrel{}{l}_j`$
$`\times W(|\stackrel{}{k}_{\text{cmb}1_{}}\mathrm{D}+\stackrel{}{l}_{\text{cmb}1}|\theta )W(|\stackrel{}{k}_{\text{cmb}2_{}}\mathrm{D}+\stackrel{}{l}_{\text{cmb}2}|\theta )`$
$`\times W(k_{\text{gal}1_{}}\mathrm{D}_1\theta _{\text{gal}})W(l_{\text{gal}1}\theta )W(k_{\text{gal}2_{}}\mathrm{D}_2\theta _{\text{gal}})W(l_{\text{gal}2}\theta )`$
$`\times W\left(|\stackrel{}{k}_{\text{gal}1_{}}\mathrm{D}_1+\stackrel{}{l}_{\text{gal}1}+\stackrel{}{k}_{\text{cmb}1_{}}\mathrm{D}_1+\stackrel{}{l}_{\text{cmb}1}|\mathrm{\Xi }\right)`$
$`\times W\left(|\stackrel{}{k}_{\text{gal}2_{}}\mathrm{D}_2+\stackrel{}{l}_{\text{gal}2}+\stackrel{}{k}_{\text{cmb}2_{}}\mathrm{D}_2+\stackrel{}{l}_{\text{cmb}2}|\mathrm{\Xi }\right)`$
where $`_i`$ gives the 2-point correlations associated with the lines of the diagram. For example :
$`_1`$ $`=`$ $`\delta (\stackrel{}{k}_{\text{gal}1})\delta (\stackrel{}{k}_{\text{cmb}1})\delta (\stackrel{}{k}_{\text{gal}2})\delta (\stackrel{}{k}_{\text{cmb}2})`$ (61)
$`\times \stackrel{~}{E}(l_{\text{gal}1})\stackrel{~}{E}(l_{\text{gal}2})\stackrel{~}{E}(l_{\text{cmb}1})\stackrel{~}{E}(l_{\text{cmb}2})`$
We explicit in the following the computation of $`𝒜_1^{\mathrm{}}`$. The other terms follow the same treatment or can be neglected. The lines in the $`𝒜_1^{\mathrm{}}`$ diagram give us the relations
$`\stackrel{}{k}_{\text{cmb}1}`$ $`=`$ $`\stackrel{}{k}_{\text{gal}1}=\stackrel{}{k}_1`$
$`\stackrel{}{k}_{\text{cmb}2}`$ $`=`$ $`\stackrel{}{k}_{\text{gal}2}=\stackrel{}{k}_2`$ (62)
$`\stackrel{}{l}_{\text{cmb}1}`$ $`=`$ $`\stackrel{}{l}_{\text{cmb}2}=\stackrel{}{l}_{\text{cmb}}`$
$`\stackrel{}{l}_{\text{gal}1}`$ $`=`$ $`\stackrel{}{l}_{\text{gal}2}=\stackrel{}{l}_{\text{gal}}`$
Then, using these relations and the small angular approximation, we have :
$`𝒜_1^{\mathrm{}}`$ $`=`$ $`4{\displaystyle ^{\text{gal}}}d\chi _1d\chi _2w_{\text{cmb}1}w_{\text{gal}1}w_{\text{cmb}2}w_{\text{gal}2}`$
$`\times {\displaystyle }{\displaystyle \frac{\mathrm{d}^2k_1\mathrm{d}^2k_2}{(2\pi )^4}}{\displaystyle \frac{\mathrm{d}^2l_{\text{gal}}\mathrm{d}^2l_{\text{cmb}}}{(2\pi )^4}}`$
$`\times C_E(l_{\text{gal}})C_E(l_{\text{cmb}})P(k_1)P(k_2)`$
$`\times 𝒢_{\mathrm{}}^{\mathrm{Ker}}(\stackrel{}{l}_{\text{cmb}},\stackrel{}{k}_1)𝒢_{\mathrm{}}^{\mathrm{Ker}}(\stackrel{}{l}_{\text{cmb}},\stackrel{}{k}_2)`$
$`\times 𝒢_{\mathrm{}}^{\mathrm{Ker}}(\stackrel{}{l}_{\text{gal}},\stackrel{}{k}_1)𝒢_{\mathrm{}}^{\mathrm{Ker}}(\stackrel{}{l}_{\text{gal}},\stackrel{}{k}_2)`$
$`\times W(|\stackrel{}{k}_1\mathrm{D}+\stackrel{}{l}_{\text{cmb}}|\theta )W(|\stackrel{}{k}_2\mathrm{D}+\stackrel{}{l}_{\text{cmb}}|\theta )`$
$`\times W(k_1\mathrm{D}\theta _{\text{gal}})W(k_2\mathrm{D}\theta _{\text{gal}})W^2(l_{\text{gal}}\theta )`$
$`\times W^2(|\stackrel{}{l}_{\text{gal}}+\stackrel{}{l}_{\text{cmb}}|\mathrm{\Xi }).`$
We apply the decomposition of $`W_1\left(|\stackrel{}{k}\mathrm{D}(\chi )+\stackrel{}{l}|\theta \right)`$ we used in eq. (48). The geometry of our problem is the same and the result (49) still holds for the terms in $`W_1\left(|\stackrel{}{k}_1\mathrm{D}(\chi _1)+\stackrel{}{l}_{\text{cmb}}|\theta \right)`$ and $`W_1\left(|\stackrel{}{k}_2\mathrm{D}(\chi _2)+\stackrel{}{l}_{\text{cmb}}|\theta \right)`$. This however is not true for $`W_1^2\left(|\stackrel{}{l}_{\text{gal}}+\stackrel{}{l}_{\text{cmb}}|\mathrm{\Xi }\right)`$ for which the application of the re-summation theorem does not bring any simplification. Then, neglecting all the $`W_3`$ parts and after integration on the $`\varphi _{k_i}`$, for the $`\mathrm{\Delta }`$-term, we have,
$`𝒜_1^\mathrm{\Delta }`$ $`=`$ $`{\displaystyle ^{\text{gal}}}d\chi _1d\chi _2w_{\text{cmb}1}w_{\text{gal}1}w_{\text{cmb}2}w_{\text{gal}2}`$
$`\times {\displaystyle }{\displaystyle \frac{\mathrm{d}k_1\mathrm{d}k_2}{(2\pi )^2}}{\displaystyle \frac{\mathrm{d}^2l_{\text{gal}}\mathrm{d}^2l_{\text{cmb}}}{(2\pi )^4}}l_{\text{gal}}^4l_{\text{cmb}}^4k_1k_2`$
$`\times C_E(l_{\text{gal}})C_E(l_{\text{cmb}})P(k_1)P(k_2)`$
$`\times W^2(|\stackrel{}{l}_{\text{gal}}+\stackrel{}{l}_{\text{cmb}}|\mathrm{\Xi })\mathrm{cos}^22\left(\varphi _{l_{\text{cmb}}}\varphi _{l_{\text{gal}}}\right)`$
$`\times W(k_1\mathrm{D}\theta _{\text{gal}})W(k_2\mathrm{D}\theta _{\text{gal}})W^2(l_{\text{gal}}\theta )`$
$`\times W(k_1\mathrm{D}\theta )W(k_2\mathrm{D}\theta )W^2(l_{\text{cmb}}\theta ).`$
Note that for the evaluation of the $``$ part, using the same kind of method, we obtain the same equation as eq. (III.5) where $`l_{\text{gal}}^4l_{\text{cmb}}^4`$ is replaced by $`l_{\text{gal}}^2l_{\text{cmb}}^2k_1^2k_2^2/2`$.
We can get rid of the remaining $`W^2\left(|\stackrel{}{l}_{\text{gal}}+\stackrel{}{l}_{\text{cmb}}|\mathrm{\Xi }\right)`$ with another approximation. The power spectrum $`C_E(l)`$ favors large values of $`l`$ whereas $`W^2(|\stackrel{}{l}_{\text{gal}}+\stackrel{}{l}_{\text{cmb}}|\mathrm{\Xi })`$ will be non-zero for $`|\stackrel{}{l}_{\text{gal}}+\stackrel{}{l}_{\text{cmb}}|1/\mathrm{\Xi }`$. Then for typical survey size of about one hundred square-degrees, $`|\stackrel{}{l}_{\text{gal}}+\stackrel{}{l}_{\text{cmb}}|l_i`$ and we can assume $`\stackrel{}{l}_{\text{gal}}\stackrel{}{l}_{\text{cmb}}`$ and $`\stackrel{}{l}_{\text{gal}}+\stackrel{}{l}_{\text{cmb}}=\stackrel{}{ϵ}`$. In this limit, $`\mathrm{cos}^22\left(\varphi _{l_{\text{cmb}}}\varphi _{l_{\text{gal}}}\right)=1`$ and $`𝒜_1^{\mathrm{}}`$ can be written
$`𝒜_1^\mathrm{\Delta }`$ $`=`$ $`{\displaystyle \frac{l\mathrm{d}l}{(2\pi )^2}l^8C_E^2(l)W^4(l\theta )\frac{\mathrm{d}^2ϵ}{2\pi }l^8W_1^2(ϵ\mathrm{\Xi })}`$
$`\times \left[{\displaystyle ^{\text{gal}}}d\chi w_{\text{cmb}}w_{\text{gal}}{\displaystyle \frac{k\mathrm{d}k}{2\pi }P(k)W(k\mathrm{D}\theta )W(k\mathrm{D}\theta _{\text{gal}})}\right]^2`$
which is essentially the cosmic variance of $`\mathrm{\Delta }E^2`$, for the $`\mathrm{\Delta }`$ part and of $`(E)^2`$ for the $``$ one (where $`l^8`$ in eq. (III.5) is replaced by $`l^4k_1^2k_2^2/2`$). Finally we have,
$`{\displaystyle \frac{𝒜_1^\mathrm{\Delta }}{B_{(\theta )}b_{\mathrm{\Delta }(\theta ,\theta _{\text{gal}})}^2}}`$ $`=`$ $`{\displaystyle \frac{2\pi }{\mathrm{\Sigma }}}{\displaystyle \frac{dll^9C_E^2(l)W_1^4(l\theta )}{\left(dll^5C_E(l)W_1^2(l\theta )\right)^2}}`$
$``$ $`\text{Cosmic variance of }\mathrm{\Delta }E^2`$
where $`\mathrm{\Sigma }=\pi \mathrm{\Xi }^2`$ in case of a disc shape survey. We show in Fig. 8 numerical results for a $`100\mathrm{deg}^2`$ survey although the numerical calculations were done with a Gaussian window function instead of a top-hat.
Numerically, for $`\theta =10^{}`$, we get
$$\frac{𝒜_1^{\mathrm{}}}{B_{(\theta )}b_{\mathrm{}(\theta ,\theta _{\text{gal}})}^2}\frac{(3.7\%)^2}{\mathrm{\Sigma }/100\text{ deg}^2}.$$
(67)
We expect that for the same reasons, the $`𝒜_2^{\mathrm{}}`$ terms will be dominated by the weak lensing variance. Yet a correct evaluation here is harder to reach. We have made this estimation within the framework of a power law $`P(k)`$. With this simplification in hand, we can write for $`𝒜_{2n}^{\mathrm{}}`$ (we focus only the $`\mathrm{\Delta }`$ part, but the same discussion holds for the $``$ observable.)
$`{\displaystyle \frac{𝒜_{2n}^\mathrm{\Delta }}{B_{(\theta )}b_{\mathrm{\Delta }(\theta ,\theta _{\text{gal}})}^2}}`$ $`=`$
$`{\displaystyle \frac{1}{r^2}}{\displaystyle \mathrm{d}^2k_1\mathrm{d}^2k_2P(k_1)P(k_2)\mathrm{cos}^2\left(\varphi _{k_1}\varphi _{k_2}\right)}`$
$`\times {\displaystyle \frac{W_1^2(k_1\theta )W_1^2(k_2\theta _{\text{gal}})W_1^2(|\stackrel{}{k}_1+\stackrel{}{k}_2|\mathrm{\Xi })}{\left[\mathrm{d}^2kP(k)W_1(k\theta )W_1(k\theta _{\text{gal}})\right]^2}}.`$
The last integral behaves essentially like the cosmic variance of $`\kappa ^2`$. More exactly, it goes like $`1/\sqrt{2}`$ this variance. It should even be smaller, because of the extra $`\mathrm{cos}^2`$ factor. We evaluated this cosmic variance using the ray-tracing simulations described in . These simulations provide us with realistic convergence maps (for the cosmological models we are interested in) with a resolution of 0.1’, and a survey size of 9 square degrees. The sources have been put at a redshift unity, and the ray-lights are propagated through a simulated Universe whose the density field has been evolved from an initial CDM power spectrum. The measured cosmic variance of $`\kappa _{(\theta )}\kappa _{(\theta _{\text{gal}})}`$ is about $`3\%`$ (see Table 2) when filtered at scales $`\theta _{\text{gal}}=5^{}`$ and $`\theta =10^{}`$ for a $`\mathrm{\Omega }_0=0.3`$ cosmology.
An estimation of $`𝒜_{2n}^\mathrm{\Delta }`$ is then given by,
$$\frac{𝒜_{2n}^\mathrm{\Delta }}{B_{(\theta )}b_{\mathrm{\Delta }(\theta ,\theta _{\text{gal}})}^2}\left(\frac{2.12\%}{r}\right)^2\frac{1}{\mathrm{\Sigma }/100\text{ deg}^2}.$$
(70)
Since $`r_{}`$ is very comparable to $`r`$, we very roughly estimate $`𝒜_{2n}^{}`$
$$\frac{𝒜_{2n}^{}}{B_{(\theta )}b_{(\theta ,\theta _{\text{gal}})}^2}\left(\frac{2.12\%}{r}\right)^2\frac{1}{\mathrm{\Sigma }/100\text{ deg}^2}.$$
(71)
The same considerations gives
$`{\displaystyle \frac{𝒜_{2n}^{\mathrm{}}}{B_{(\theta )}b_{\mathrm{}(\theta ,\theta _{\text{gal}})}^2}}`$ $`=`$ $`{\displaystyle \frac{(2.12\%)^2}{\mathrm{\Sigma }/100\text{ deg}^2}}.`$ (72)
There is no $`r`$ dependency here; the diagram cross-correlates $`\kappa _{\text{cmb}}`$ and $`\kappa _{\text{gal}}`$.
We can approximate the remaining $`𝒜`$-terms. They should be smaller than the former. We have
$`𝒜_{3n}^{\mathrm{}}`$ $``$ $`{\displaystyle \frac{1}{r_{\mathrm{}}^2}}{\displaystyle \frac{(2.12\%\times 3.7\%)^2}{\mathrm{\Sigma }/100\text{ deg}^2}}B_{(\theta )}b_{\mathrm{}(\theta ,\theta _{\text{gal}})}^2`$
$``$ $`𝒜_{2n}^{\mathrm{}}`$
and
$`𝒜_{3c}^{\mathrm{}}`$ $``$ $`{\displaystyle \frac{(2.12\%\times 3.7\%)^2}{\mathrm{\Sigma }/100\text{ deg}^2}}B_{(\theta )}b_{\mathrm{}(\theta ,\theta _{\text{gal}})}^2`$
$``$ $`𝒜_{2c}^{\mathrm{}}.`$
Then, only the $`𝒜_1^{\mathrm{}}`$ and $`𝒜_2^{\mathrm{}}`$ terms (boxed on Fig. 7) contribute substantially to the cosmic variance of $`𝒳_{\mathrm{}}`$. Since $`𝒜_1^{\mathrm{}}`$ and $`𝒜_2^{\mathrm{}}`$ are respectively the cosmic variance of $`\mathrm{\Delta }E^2`$ (resp. $`(\stackrel{}{}E)^2)`$ and of $`\kappa ^2`$ (resp. $`(\stackrel{}{}\kappa )^2`$), we can write the variance of $`𝒳_{\mathrm{}}`$ as
$`\mathrm{C}osVar(𝒳_\mathrm{\Delta })`$ $`=`$
$`\mathrm{C}osVar\left(\mathrm{\Delta }E^2\right)+\left({\displaystyle \frac{1+r^2}{2r^2}}\right)\mathrm{C}osVar\left(\kappa ^2\right).`$
and
$`\mathrm{C}osVar(𝒳_{})`$ $`=`$
$`\mathrm{C}osVar\left((\stackrel{}{}E)^2\right)+\left({\displaystyle \frac{1+r_{}^2}{2r_{}^2}}\right)\mathrm{C}osVar\left((\stackrel{}{}\kappa )^2\right).`$
Table 3 presents numerical results for various filtering scenarii and models.
The two quantities, $`b_\mathrm{\Delta }`$ and $`b_{}`$, lead to similar cosmic variance that are rather small. Obviously it would be even better to use $`b=b_\mathrm{\Delta }+b_{}`$. For such a quantity the resulting cosmic variance for the cross-correlation coefficient should even be smaller, by a factor $`\sqrt{2}`$, although a detailed analysis is made complicated because of the complex correlation patterns it contains.
## IV Conclusion
We have computed a first order mapping that describes, in real space, the weak lensing effects on the CMB polarization. In particular we derived the explicit mathematical relation between the primary CMB polarization and the shear field at leading order in lens effect. It demonstrates that a $`B`$-component of the polarization field can be induced by lens couplings. We have shown however that the $`B`$-map alone cannot lead to a non-ambiguous reconstruction of the projected mass map.
Nonetheless, the $`B`$-component can potentially exhibit a significant correlation signal with local weak lensing surveys. This opens a new window for detecting lens effects on CMB maps. In particular, and contrary to previous studies involving the temperature maps alone, we found that such a correlation can be measured with a rather high signal to noise ratio even in surveys of rather modest size and resolution. Anticipating data sets that should be available in the near future, ($`100\text{deg}^2`$ survey, with $`5^{}`$ resolution for galaxy survey and $`10^{}`$ Gaussian beam size for CMB polarization detection), we have obtained a cosmic variance around $`8\%`$. Needless is to say that this estimation does not take into account systematics and possible foreground contaminations. It shows anyway that Cosmic Microwave Background polarization contains a precious window for studying the large scale mass distribution and consequently putting new constraints on the cosmological parameters.
In this paper we have investigated specific quantities that would accessible to observations. They both would permit to put constraint on the cosmological constant. The simulated maps we presented here are only of illustrative interest. We plan to complement this study with extensive numerical experiments to validate our results (in particular on the cosmic variance), and explore the effect of realistic ingredients we did not include in our simple analytical framework, a shear non-gaussianity, lens-lens coupling and so forth.
###### Acknowledgements.
We thank B. Jain, U. Seljak and S. White for the use of their ray-tracing simulations. KB and FB thank CITA for hospitality and LvW is thankful to SPhT Saclay for hospitality. We are all grateful to the TERAPIX data center located at IAP for providing us computing facilities.
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# Instabilities of vortices in binary mixture of trapped Bose-Einstein condensates: Role of collective excitations with positive and negative energies
## I Introduction
Recent observations of the quantized vortices in a dilute ultracold gas of $`{}_{}{}^{87}Rb`$ atoms are spectacular evidences of the superfluid properties of atomic gases below transition temperature for formation of Bose-Einstein condensate (BEC). First experimental results were obtained, using method suggested by Williams and Holland , in a binary mixture of the hyperfine states of $`{}_{}{}^{87}Rb`$ and more recent experiments have demonstrated vortices and vortex lattices in an optically stirred single component condensate.
Gross-Pitaevskii (GP) equation is a widely used approximation to describe dynamics of the quantized vortices in the superfluid component of Bose gases at temperatures below critical. The GP equation falls into a general class of Hamiltonian nonlinear systems. The theory of stability of stationary solutions (equilibria) of such systems is well established at the moment, see e.g. . However, a certain gap between formal mathematical knowledge and applications to physical examples still exists and BEC field is not an exception. One of the controversial examples in the BEC context is complex or imaginary eigenfrequencies in the spectrum of elementary excitations. Existence of these frequencies was first pointed out by Bogoliubov in his original work , for more recent references see, e.g. . Unjustified negligence by the nonselfadjoitness of the Bogoliubov equations often lead to the association between frequencies and energies in a way standard for quantum mechanics based on selfadjoint operators, thus admitting possibility of complex energies in a conservative system.
However, rigorously proved theorem allowing clear physical interpretation of complex frequencies is known from the general theory of Hamiltonian systems . It states that excitations with complex frequencies can appear only as a result of the resonance between two excitations with positive and negative energies . The roots of this theorem go back to 19th century and it has been used for interpretation of instabilities with complex eigenvalues in plasma physics and fluid dynamics . Motivated by recent experiments on observation of vortices in binary mixture of condensates we will demonstrate that these vortices can have complex frequencies in their spectrum, thereby giving good practical ground for selfconsistent theoretical interpretation of this phenomenon in the BEC context.
Properties of vortices in the single component magnetically trapped ultracold gases have been subject to the intensive theoretical investigations, see e.g. \[9,17-32\], which significantly extended classical works dealing with spatially unbounded case. The properties of unit vortices in the two component condensates have also been studied, in parallel and independently from this work, by Garcia-Ripoll and Pérez-Garcia . Richness of the dynamics of the two-condensate system and different approaches to the problem have led only to few overlaps which are outlined where appropriate.
In the next section we introduce coupled GP equations and briefly describe their general properties. Then, in section III, we derive Bogoliubov equations for excitations, clearly specifying differences between frequency and energy spectra of the excitations. We also explain scenario of appearance of the excitations with complex frequencies and show that they have zero energies. In section IV we verify validity of general results presented in Sec. III using perturbation theory in the limit of weak interaction and direct numerical study of Bogoliubov equations. In Sections IV, V we describe how long term dynamics of unit and higher order vortices can be interpreted using linear Bogoliubov theory.
## II Gross-Pitaevskii equations
Studies of superfluid mixtures using coupled GP equations have long history and attracted significant recent activities, see e.g. and references therein. Following these works we assume that wave functions $`\psi _{1,2}`$ of two-species condensate inside an axial harmonic trap obey equations
$`i\mathrm{}_t\psi _1={\displaystyle \frac{\mathrm{}^2}{2m}}\stackrel{}{}^2\psi _1+{\displaystyle \frac{1}{2}}m\mathrm{\Omega }^2(r^2+\sigma ^2z^2)\psi _1`$ (1)
$`+(u_{11}|\psi _1|^2+u_{12}|\psi _2|^2)\psi _1,`$ (2)
$`i\mathrm{}_t\psi _2={\displaystyle \frac{\mathrm{}^2}{2m}}\stackrel{}{}^2\psi _2+{\displaystyle \frac{1}{2}}m\mathrm{\Omega }^2(r^2+\sigma ^2z^2)\psi _2`$ (3)
$`+(u_{22}|\psi _2|^2+u_{21}|\psi _1|^2)\psi _2,`$ (4)
where for simplicity we have neglected by possible differences of atomic masses $`m_{1,2}=m`$ and trap frequencies, $`\mathrm{\Omega }_{1,2}=\mathrm{\Omega }`$, $`\sigma _{1,2}=\sigma `$, $`\stackrel{}{}=\stackrel{}{i}_x_x+\stackrel{}{i}_y_y+\stackrel{}{i}_z_z`$, $`r^2=x^2+y^2`$. Coefficients $`u_{ij}=4\pi \mathrm{}^2a_{ij}/m`$ characterise intra- and inter-species interaction with corresponding two-body scattering lengths $`a_{11}a_{22}`$ and $`a_{12}=a_{21}`$.
At this point we introduce dimensionless time and space variables $`\stackrel{~}{t}=\mathrm{\Omega }t`$ and $`(\stackrel{~}{x},\stackrel{~}{y},\stackrel{~}{z})=(x,y,z)/a_{ho}`$ and normalisation for the wave functions $`\stackrel{~}{\psi }_{1,2}=a_{ho}^{3/2}\psi _{1,2}/\sqrt{N_1}`$, where $`a_{ho}=\sqrt{\mathrm{}/(m\mathrm{\Omega })}`$ is the harmonic oscillator strength and $`N_{1,2}=𝑑V|\psi _{1,2}|^2`$ are the numbers of particles. We will consider the quasi-2D model to simplify our numerical study. This approximation was previously used in several works, see e.g. , and it is applicable not only for pancake traps, $`\sigma 1`$, but also captures main qualitative features of spherical traps. To make further reduction of Eqs. (2) we redefine the wave functions once more:
$$\stackrel{~}{\psi }_{1,2}=\left[\frac{\sigma }{2\pi }\right]^{1/4}\mathrm{\Psi }_{1,2}(\stackrel{~}{x},\stackrel{~}{y},\stackrel{~}{t})e^{\sigma \stackrel{~}{z}^2/4}e^{i\mu _{1,2}\stackrel{~}{t}i\sigma \stackrel{~}{t}/2},$$
(5)
where $`\mu _{1,2}`$ are the chemical potentials. Dropping tilde we find that equations for $`\mathrm{\Psi }_{1,2}`$ and $`\mathrm{\Psi }_{1,2}^{}`$ can be put into Hamiltonian form
$$i_t\stackrel{}{\psi }+\widehat{\eta }\frac{\delta H}{\delta \stackrel{}{\psi }^{}}=0,$$
(6)
$$\stackrel{}{\psi }=\left[\begin{array}{c}\mathrm{\Psi }_1\\ \mathrm{\Psi }_1^{}\\ \mathrm{\Psi }_2\\ \mathrm{\Psi }_2^{}\end{array}\right],\widehat{\eta }=\left[\begin{array}{cccc}1& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& 1\end{array}\right],\frac{\delta H}{\delta \stackrel{}{\psi }^{}}=\left[\begin{array}{c}\delta H/\delta \mathrm{\Psi }_1^{}\\ \delta H/\delta \mathrm{\Psi }_1\\ \delta H/\delta \mathrm{\Psi }_2^{}\\ \delta H/\delta \mathrm{\Psi }_2\end{array}\right],$$
$`H`$ $`={\displaystyle }dxdy(|\stackrel{}{}\mathrm{\Psi }_1|^2+|\stackrel{}{}\mathrm{\Psi }_2|^2`$ (9)
$`+(\widehat{V}\mu _1)|\mathrm{\Psi }_1|^2+(\widehat{V}\mu _2)|\mathrm{\Psi }_2|^2`$
$`+{\displaystyle \frac{g}{2}}[\beta _{11}|\mathrm{\Psi }_1|^4+\beta _{22}|\mathrm{\Psi }_2|^4+2\beta _{12}|\mathrm{\Psi }_1|^2|\mathrm{\Psi }_2|^2]).`$
where $`H`$ is the energy functional (or Hamiltonian), $`\stackrel{}{}=\stackrel{}{i}_x_x+\stackrel{}{i}_y_y`$, $`\widehat{V}=r^2/4`$ is the harmonic potential, $`g`$ is the interaction parameter $`g=8\sqrt{\pi \sigma }N_1a_{11}/a_{ho}`$, $`\beta _{12}=a_{12}/a_{11}`$, and $`\beta _{22}=a_{22}/a_{11}`$. $`\beta _{11}=1`$ and it is left in the equations for the sake of the symmetry.
Invariancies of $`H`$ with respect to the infinitesimal rotations and two parameter phase transformation
$$(\mathrm{\Psi }_1,\mathrm{\Psi }_2)(\mathrm{\Psi }_1e^{i\varphi _1},\mathrm{\Psi }_2e^{i\varphi _2}),$$
(10)
result in conservation of the total angular momentum and of the total number of particles in each component.
Radially symmetric stationary states of the condensate (equilibria) can be presented in the form
$$\mathrm{\Psi }_j=A_j(r)e^{iL_j\theta },j=1,2$$
(11)
where $`\theta `$ is the polar angle and $`A_j`$ are real functions. Using method suggested in only states with vortex in one of the components can be created and therefore we will consider below only cases with $`L_2>0`$ and $`L_1=0`$. Functions $`A_j(r)`$ were found numerically using Newton method. Chemical potentials $`\mu _{1,2}`$ were found from the normalization conditions
$$2\pi r𝑑rA_1^2=1,2\pi r𝑑rA_2^2=\frac{N_2}{N_1}N.$$
(12)
## III Frequency and energy spectra of collective excitations
### A Bogoliubov equations and frequency spectrum
To study spectrum of BEC at equilibrium we linearise Eqs. (6) using substitutions
$$\mathrm{\Psi }_j=(A_j(r)+f_j(r,\theta ,t))e^{iL_j\theta },$$
(13)
where $`f_j`$ are small and complex. Assuming that excitation are periodic in $`\theta `$ with period $`2\pi `$ we expand $`f_j`$ into Fourier series:
$$f_j=\underset{l}{}\left(U_{lj}(r,t)e^{il\theta }+V_{lj}^{}(r,t)e^{il\theta }\right),$$
(14)
where $`l=0,\pm 1,\pm 2,\mathrm{}`$ Then
$$i_t\stackrel{}{W}_l+\widehat{\eta }\widehat{}_l\stackrel{}{W}_l=0$$
(15)
is the set of linear partial differential equations resulting from the substitution of (13), (14) into Eq. (6). Here $`\stackrel{}{W}_l=(U_{l1},V_{l1},U_{l2},V_{l2})^T`$,
$$\widehat{}_l=\left[\begin{array}{cccc}\widehat{}_{l,1}& g\beta _{11}A_1^2& g\beta _{12}A_1A_2& g\beta _{12}A_1A_2\\ g\beta _{11}A_1^2& \widehat{}_{l,1}& g\beta _{12}A_1A_2& g\beta _{12}A_1A_2\\ g\beta _{12}A_1A_2& g\beta _{12}A_1A_2& \widehat{}_{l,2}& g\beta _{22}A_2^2\\ g\beta _{12}A_1A_2& g\beta _{12}A_1A_2& g\beta _{22}A_2^2& \widehat{}_{l,2}\end{array}\right],$$
is a self-adjoint operator and
$`\widehat{}_{l,j}={\displaystyle \frac{1}{r}}{\displaystyle \frac{}{r}}r{\displaystyle \frac{}{r}}+{\displaystyle \frac{1}{r^2}}(L_j+l)^2+\widehat{V}\mu _j`$ (16)
$`+2g\beta _{jj}A_j^2+g\beta _{12}A_j^{}^2;j=1,2;j^{}=2,1.`$ (17)
Frequency spectra of $`\widehat{\eta }\widehat{}_l`$ are discrete, providing $`\widehat{V}0`$, therefore phonons, strictly speaking, are absent in the trap geometry. In accord with standard terminology , all spatially bounded elementary excitations can be called collective excitations (or collective modes) of an equilibrium under consideration. Linearised equations for excitations in a Bose gas, similar to Eq. (15), were first derived by Bogoliubov and expansion (14) was first applied in the context of the vortex excitations by Pitaevskii . To find frequency spectrum and collective modes we need to solve set of the eigenvalue problems
$$\widehat{\eta }\widehat{}_l\stackrel{}{w}_{ln}=\omega _{ln}\stackrel{}{w}_{ln}.$$
(18)
$`\widehat{\eta }\widehat{}_l`$ are non-self-adjoint operators and therefore complex frequencies are not forbidden. If $`\stackrel{}{w}_{ln}`$ is an eigenvector of $`\widehat{\eta }\widehat{}_l`$ with eigenvalue $`\omega _{ln}`$ it is selfevident that $`\stackrel{}{w}_{ln}^{}`$ is also an eigenvector with eigenvalue $`\omega _{ln}^{}`$ and it can be shown that $`\widehat{\eta }\widehat{}_l`$ has eigenvectors $`\stackrel{}{w}_{ln}=\widehat{\tau }\stackrel{}{w}_{ln}`$ and $`\stackrel{}{w}_{ln}^{}`$ with eigenvalues $`\omega _{ln}`$ and $`\omega _{ln}^{}`$, respectively. Here
$$\widehat{\tau }=\left[\begin{array}{cccc}0& 1& 0& 0\\ 1& 0& 0& 0\\ 0& 0& 0& 1\\ 0& 0& 1& 0\end{array}\right].$$
Thus spectrum of $`\widehat{\eta }\widehat{}_l`$ can be obtained by reflection of the spectrum of $`\widehat{\eta }\widehat{}_l`$ with respect to the line $`Re\omega =0`$ in the plane $`(Re\omega ,Im\omega )`$. In other words, it means that purely real or purely imaginary frequencies of the elementary excitation exist in pairs and complex frequencies exist in quartets and that
$$Tr\{\widehat{\eta }(\widehat{}_l+\widehat{}_l)\}=\underset{n}{}(\omega _{ln}+\omega _{ln})=0.$$
(19)
Any equilibrium state of the condensate is spectrally stable if its spectrum is real. If there is at least one frequency with negative imaginary part then corresponding collective mode will grow in time destabilising the equilibrium, which is called spectrally unstable.
### B Biorthogonality
Eigenmodes of $`\widehat{\eta }\widehat{}_l`$ are biorthogonal to the modes of the adjoint eigenvalue problem , i.e.
$$\stackrel{}{w}_{ln},\stackrel{}{a}_{ln^{}}=0,$$
(20)
where $`nn^{}`$ and $`\stackrel{}{a}_{ln}`$ obey
$$(\widehat{\eta }\widehat{}_l)^{}\stackrel{}{a}_{ln}=\widehat{}_l\widehat{\eta }\stackrel{}{a}_{ln}=\omega _{ln}^{}\stackrel{}{a}_{ln},$$
(21)
and $`\stackrel{}{b},\stackrel{}{c}=2\pi _k_0^{\mathrm{}}r𝑑rb_k^{}c_k`$ for any $`\stackrel{}{b}`$, $`\stackrel{}{c}`$. Factor $`2\pi `$ is introduced to mimic integration over $`\theta `$. The key feature of our model, originating in its Hamiltonian structure, is that transformation linking $`\stackrel{}{w}_{ln}`$ and its adjoint $`\stackrel{}{a}_{ln}`$ can be found in explicit and simple form. If $`\stackrel{}{w}_{ln}`$ is an eigenmode of $`\widehat{\eta }\widehat{}_l`$ with frequency $`\omega _{ln}`$ then it can be checked that $`\widehat{\eta }\stackrel{}{w}_{ln}^{}`$ and $`\widehat{\eta }\stackrel{}{w}_{ln}`$ are eigenmodes of $`\widehat{}_l\widehat{\eta }`$ with eigenvalues, respectively, $`\omega _{ln}^{}`$ and $`\omega _{ln}`$. The mode adjoint to $`\stackrel{}{w}_{ln}`$ is $`\widehat{\eta }\stackrel{}{w}_{ln}^{}`$, therefore if $`Im\omega _{ln}0`$ then biorthogonality condition (20) implies
$$\stackrel{}{w}_{ln},\widehat{\eta }\stackrel{}{w}_{ln}=0.$$
(22)
If $`\omega _{ln}`$ is real, then $`\stackrel{}{w}_{ln}`$ is also real and $`\stackrel{}{w}_{ln},\widehat{\eta }\stackrel{}{w}_{ln}0`$. Normalisation constant can always be chosen in such a way that
$$|\stackrel{}{w}_{ln},\widehat{\eta }\stackrel{}{w}_{ln}|=2.$$
(23)
Convenience of making the left hand side of Eq. (23) equal to $`2`$ will become clear below, when Eqs. (19), (20) for energies of elementary excitations are derived. Eq. (23) makes it explicit that inner product $`\stackrel{}{w}_{ln},\widehat{\eta }\stackrel{}{w}_{ln}`$ can be either positive or negative. This point often remains silent if one derives conditions similar to (23) as part of the diagonalisation procedure of the second-quantised Hamiltonian disregarding eigenmodes with negative and zero values of $`\stackrel{}{w}_{ln},\widehat{\eta }\stackrel{}{w}_{ln}`$ . The fact that inner product of a real eigenmode with its adjoint can be negative, is different from standard quantum mechanics based on the selfadjoint operators. These difference can have series of consequences and one of them is that energy levels are not necessarily linked to the eigenfrequencies according to the standard rule $`\mathrm{}\omega `$, see below. Note, that origin of the nonselfadjoitness in our case is the nonlinearity of GP equations. If particles in the condensate do not interact, $`g=0`$, then $`\widehat{\eta }\widehat{}_l`$ become diagonal and self-adjoint.
### C Energy spectrum
As a prelude to calculation of energies of the elementary excitations it is instructive to introduce notion of the nonlinear stability, i.e. stability under the full nonlinear dynamics. According to Dirichlet’s theorem , nonlinear stability is ensured if an equilibrium state under consideration is either minimum or maximum of the functional $`H`$. Note here, that excitations change number of particles in the equilibrium state, therefore it was convenient to introduce energy functional $`H`$, which is actually the so-called modified energy , i.e. it is the energy functional for Eq. (1) modified by addition of the number of particles integrals, see terms proportional to $`\mu _j`$ in Eq. (9). Let us assume that $`\stackrel{}{w}_l(r)e^{il\theta }`$ is a small initial perturbation of an equilibrium state of the condensate, then
$$H=H_0+\frac{1}{2}\stackrel{}{w}_l,\widehat{}_l\stackrel{}{w}_l+\mathrm{},$$
(24)
where $`H`$ is the energy of the perturbed equilibrium and $`H_0`$ is calculated at the exact equilibrium. The equilibrium is nonlinearly stable if eigenvalue problems
$$\widehat{}_l\stackrel{}{\beta }_{lm}=\alpha _{lm}\stackrel{}{\beta }_{lm}$$
(25)
have all their eigenvalues either negative or positive, except zero eigenvalues generated by continuous symmetries. Spectral instability implies nonlinear one, and nonlinear stability implies spectral one, but not vise versa . In the nonrotating traps all the higher order states of the condensate including vortices are not the local extrema of the energy . Let us stress, however, that their nonlinear instability can not be guaranteed by this fact alone and requires separate consideration.
If $`\stackrel{}{w}_l`$ in Eq. (24) is an eigenmode $`\stackrel{}{w}_{ln}`$ of $`\widehat{\eta }\widehat{}_l`$ then
$$ϵ_{ln}=\frac{1}{2}\stackrel{}{w}_{ln},\widehat{}_l\stackrel{}{w}_{ln}=\frac{1}{2}\omega _{ln}\stackrel{}{w}_{ln},\widehat{\eta }\stackrel{}{w}_{ln}$$
(26)
measures energy of this mode. Assuming that $`\omega _{ln}`$ is real and using biorthogonality conditions (23) one gets
$$ϵ_{ln}=\omega _{ln}sign\stackrel{}{w}_{ln},\widehat{\eta }\stackrel{}{w}_{ln}.$$
(27)
If $`Im\omega _{ln}0`$ then Eq. (22) implies
$$ϵ_{ln}=0.$$
(28)
Energy in physical units is given by $`ϵ_{ln}\mathrm{}\mathrm{\Omega }`$.
It is readily demonstrated that
$$\stackrel{}{w}_{ln},\widehat{\eta }\stackrel{}{w}_{ln}=\widehat{\tau }\stackrel{}{w}_{ln},\widehat{\eta }\widehat{\tau }\stackrel{}{w}_{ln}=\stackrel{}{w}_{ln},\widehat{\eta }\stackrel{}{w}_{ln},$$
(29)
therefore modes with frequencies $`\omega _{ln}`$ and $`\omega _{ln}=\omega _{ln}`$ have their energies equal in sign and value.
Eigenfunctions $`\stackrel{}{\beta }_{ml}`$ are orthogonal $`\stackrel{}{\beta }_{lm},\stackrel{}{\beta }_{lm^{}}=\delta _{mm^{}}`$ and form complete basis. Therefore $`\stackrel{}{w}_{ln}`$ in Eq. (26) can be expanded in terms of $`\stackrel{}{\beta }_{lm}`$, which gives $`ϵ_{ln}=_m\alpha _{lm}\stackrel{}{w}_{ln},\stackrel{}{\beta }_{lm}^2`$. Thus an equilibrium can have collective modes carrying energy with opposite signs only if it is not a local extremum of the energy functional.
### D Resonances and instabilities, crossings and avoided crossings
Instabilities with imaginary eigenvalues have been found to dominate dynamics of homogeneous BEC with attractive interaction and mixture of two condensates with repulsive interaction . Higher order vortices and dark solitons in single component trapped condensates can be unstable with respect to modes with complex frequencies . In our problem instabilities with complex eigenvalues are the most important ones also. Therefore it is desirable to have criterion or theorem allowing interpretation and prediction of these instabilities. Such theorem is actually known from the general theory of the Hamiltonian dynamical systems and, in our context, it can be reformulated as: (i) Sign of the energy of a collective excitation preserves as parameters vary as long as there is no frequency resonance with another excitation. (ii) Condensate at equilibrium can lose spectral stability as parameters vary only in two ways: either via frequency resonance of two elementary excitations with positive and negative energies or by resonance at zero frequency. Proof of these results is based on the fact, that transition from spectral stability to instability can not violate energy conservation law: $`_tH=0`$.
From this theorem and preceding considerations one can conclude that complex eigenvalues in the spectrum of the vortices in Bose-Einstein condensate can appear only due to mutual annihilation of the collective excitations with positive and negative energies. The latter can coexist only for nonlinearly unstable equilibria. The theorem does not forbid for frequencies of two excitations with either the same or opposite signs of the energies simply cross each other without change of the spectral stability. In our example it typically happens when corresponding eigenmodes remain orthogonal at the exact resonance. If, however, eigenmodes of the colliding excitations start to compete for the same direction in the functional phase space and become degenerate at the resonance then crossing is not a generic scenario. If the signs of the energies of the excitations are opposite, then quartet of complex frequencies appears upon passing the resonance. Alternatively, if the signs are the same, the exact resonance can not be achieved and it becomes replaced by the so-called avoided crossing of the energy levels .
Figs. 1(a),(b) show numerically calculated frequency and energy spectra of the collective excitations with $`l=\pm 1`$ for spectrally stable (left panel) and spectrally unstable (right panel) unit vortices $`(L_1=0,L_1=1)`$. The negative energy excitation is clearly seen in the spectrally stable situation. At the transition threshold to spectral instability this excitation and another one, having energy with the same absolute value but opposite sign, annihilate each other. This transition is accompanied by appearance of the zero energy excitations. Examples of frequency and energy evolution under the parameter variation resulting in instabilities, crossings and avoided crossing are given in the section IV. Figs. 1(c) show spectra of $`\widehat{}_{\pm 1}`$ illustrating that there are no qualitative changes in these spectra after appearance of complex frequencies.
Possibility of observation of the collective modes with complex frequencies has caused some concerns and discussions . However, initial perturbations of the equilibrium state having nonzero projections on the adjoint mode $`\widehat{\eta }\stackrel{}{w}_{ln}^{}`$ with $`Im\omega _{ln}^{}>0`$ will lead to the ultimate growth of the corresponding mode $`\stackrel{}{w}_{ln}`$ because it has $`Im\omega _{ln}<0`$. The consequences of this growth is in any way diminished by the fact that the energy of this mode is zero. Wealth of references and examples of instabilities with complex eigenvalues existing in other physical contexts can be found in .
### E Goldstone and dipole modes
Infinitesimal variations of $`\varphi _j`$, see Eq. (10), generate two zero-energy eigenmodes (Goldstone modes) $`(A_1,A_1,0,0)^T`$ and $`(0,0,A_2,A_2)^T`$ belonging to the null-eigenspace of $`\widehat{\eta }\widehat{}_0`$. Harmonic trapping modifies spectrum in such a way that $`\widehat{\eta }\widehat{}_{\pm 1}`$ acquire couple of parameter independent eigenvalues $`\omega =\pm 1`$ with eigenfunctions
$$\stackrel{}{w}_{\pm 1d}=\left[\begin{array}{c}\\ \omega \left(\frac{dA_1}{dr}\frac{1}{r}L_1A_1\right)+\frac{1}{2}rA_1\\ \\ \omega \left(\frac{dA_1}{dr}\pm \frac{1}{r}L_1A_1\right)\frac{1}{2}rA_1\\ \\ \omega \left(\frac{dA_2}{dr}\frac{1}{r}L_2A_2\right)+\frac{1}{2}rA_2\\ \\ \omega \left(\frac{dA_2}{dr}\pm \frac{1}{r}L_2A_2\right)\frac{1}{2}rA_2\end{array}\right],$$
(30)
where $`\omega `$ can take values $`\pm 1`$ for both eigenmodes. $`\stackrel{}{w}_{\pm 1d}`$ are often called dipole modes and $`ϵ_{\pm 1d}=1`$. Eqs. (30) generalise expressions previously derived for single-species condensates . Existence of the dipole modes can also be associated with the Kohn’s theorem .
## IV Collective excitations of unit vortices: $`L_1=0,L_1=1`$
We begin our analysis considering weakly interacting condensate: $`g1`$, $`N1`$. In this limit potential energy due to harmonic trapping $`\widehat{V}`$ strongly dominates over the interaction energy, which allows to make explicit calculations of the frequency and energy spectra of the elementary excitations. Calculations of the excitation spectra of the vortices and dark solitons in the singly component weakly interacting condensates have been recently done by several groups of authors . However, these calculations lack analysis of the energy sign in a sense explained in the preceding section.
We substitute asymptotic expansions $`\mu _{1,2}=\mu _{1,2}^{(0)}+g\mu _{1,2}^{(1)}+O(g^2)`$, $`A_{1,2}=A_{1,2}^{(0)}+gA_{1,2}^{(1)}+O(g^2)`$ into the stationary ($`_t=0`$) version of Eqs. (6) and derive recurrent system of linear equations. In the zero approximation we have two uncoupled harmonic oscillator problems with eigenmodes
$$A_1^{(0)}=\frac{1}{\sqrt{2\pi }}e^{r^2/4},A_2^{(0)}=\sqrt{\frac{N}{\pi }}\frac{r}{2}e^{r^2/4}.$$
(31)
Using solvability condition of the first order problem we find asymptotic expressions for the chemical potentials: $`\mu _1=1+g(2\beta _{11}+n\beta _{12})/(8\pi )+O(g^2)`$, $`\mu _2=2+g(n\beta _{22}+\beta _{12})/(8\pi )+O(g^2)`$.
Then we expand operators $`\widehat{}_l`$, eigenmodes $`\stackrel{}{w}_l`$ and frequencies $`\omega _l`$ into the series $`\widehat{}_l=\widehat{}_l^{(0)}+g\widehat{}_l^{(1)}+O(g^2)`$, $`\stackrel{}{w}_l=\stackrel{}{w}_l^{(0)}+g\stackrel{}{w}_l^{(1)}+O(g^2)`$, $`\omega _l=\omega _l^{(0)}+g\omega _l^{(1)}+O(g^2)`$. After substitution into Eq. (18) in the first order we find standard equation
$$(\widehat{\eta }\widehat{}_l^{(0)}\omega _l^{(0)})\stackrel{}{w}_l^{(1)}=(\omega _l^{(1)}\widehat{\eta }\widehat{}_l^{(1)})\stackrel{}{w}_l^{(0)}.$$
(32)
$`\widehat{\eta }\widehat{}_l^{(0)}`$ is diagonal and self-adjoint and all its eigenmodes and eigenvalues can be found explicitly. Then using solvability condition for Eq. (32) one can find corrections for all frequencies and coefficients in the linear superposition of the zero approximation eigenmodes. Computer algebra makes technical realization of this plan by a straightforward exercise. We will present and analyse explicit analytical results only for the operator $`\widehat{\eta }\widehat{}_1`$ considering vicinity of the spectral point $`\omega =1`$, because it contains information about origin of spectral instabilities of unit vortices. Equivalent analysis of $`\widehat{\eta }\widehat{}_1`$ near $`\omega =1`$ gives the same results.
$`\widehat{\eta }\widehat{}_1^{(0)}`$ has three unit frequencies, $`\omega _1^{(0)}=1`$. Corresponding eigenmodes are
$`\stackrel{}{b}_1=(0,1,0,0)^T{\displaystyle \frac{r}{2\sqrt{\pi }}}e^{r^2/4},`$ (33)
$`\stackrel{}{b}_2=(0,0,1,0)^T{\displaystyle \frac{1}{\sqrt{2\pi }}}e^{r^2/4},`$ (34)
$`\stackrel{}{b}_3=(0,0,0,1)^T{\displaystyle \frac{r^2}{4\sqrt{\pi }}}e^{r^2/4}.`$ (35)
Solvability conditions for Eq. (32) lead to the characteristic determinant of the three by three matrix. We find that one of the three frequencies is associated with dipole mode $`\stackrel{}{w}_{1d}`$ and that the other two are
$`\omega _1^\pm =1+{\displaystyle \frac{g}{32\pi }}\left(3\beta _{12}N\beta _{22}\pm \sqrt{R}\right)+O(g^2),`$ (36)
$`R(3\beta _{12}+N\beta _{22})^28\beta _{12}^2(N+1).`$ (37)
Corresponding unnormalised eigenmodes are
$`\stackrel{}{w}_1^\pm =`$ $`\sqrt{2N}\left(N\beta _{22}\beta _{12}\pm \sqrt{R}\right)\stackrel{}{b}_1`$ (40)
$`+\sqrt{2}\left(3N\beta _{22}+[14N]\beta _{12}\sqrt{R}\right)\stackrel{}{b}_2`$
$`+4N\left(\beta _{12}\beta _{22}\right)\stackrel{}{b}_3+O(g).`$
Pair of frequencies (36) becomes complex if $`R<0`$, signaling spectral instability of the unit vortex. If $`R>0`$, then
$`\stackrel{}{w}_1^\pm ,\widehat{\eta }\stackrel{}{w}_1^\pm =`$ (41)
$`4\left[R(N1)\pm \sqrt{R}\left\{(15N)\beta _{12}+N(3+N)\beta _{22}\right\}\right]`$ (42)
and one can check that
$$ϵ_1^\pm =\pm \omega _1^\pm ,$$
(43)
i.e. modes with frequencies $`\omega _1^+`$ and $`\omega _1^{}`$ have, respectively, positive and negative energies, see Fig. 2(c). One can also verify biorthogonality $`\stackrel{}{w}_1^\pm ,\widehat{\eta }\stackrel{}{w}_1^\pm =0=ϵ_1^\pm `$ for $`R<0`$, see Eq. (22), and $`\stackrel{}{w}_1^\pm ,\widehat{\eta }\stackrel{}{w}_1^{}=0`$ for $`R>0`$, see Eq. (23).
Assuming $`\beta _{ij}>0`$ and rewriting instability condition $`R<0`$ in the form
$$\frac{\beta _{22}}{\beta _{12}}<\frac{1}{N}(2\sqrt{2(N+1)}3),$$
(44)
one can see that unit vortices are more stable if intraspecies interaction of atoms in the vortex containing part of the condensate is somewhat larger than interspecies interaction. The choices of scattering lengths corresponding to the experiment are: $`\beta _{22}=1.03/0.97`$, $`\beta _{12}=1/0.97`$ ($`\beta _{22}>\beta _{11}`$) and $`\beta _{22}=0.97/1.03`$, $`\beta _{12}=1/1.03`$ ($`\beta _{22}<\beta _{11}`$). The former case corresponds to the vortex in the spin state $`\{F=1,m_f=1\}`$ of $`{}_{}{}^{87}Rb`$ and the latter to the vortex in the state $`\{2,2\}`$. These two states will be called – state $`1`$ and state $`2`$. It is clear that for $`N=1`$ Eq. (44) predicts instability for the vortex in the state $`2`$ and stability for the vortex in the state $`1`$, which supports results of the experimental observations .
Fig. 2 shows frequency resonance accompanied by the simultaneous mutual annihilation of excitations with positive and negative energies happening at some critical value of $`\beta _{12}`$, see Eq. (44). In fact it models transition from the situation with vortex in the state 1 to the case with vortex in the state 2. Performing numerical studies for wide range of parameters, outside the validity region of analytical considerations, we have not been able to find regions of spectral instability of the unit vortex in state $`1`$. Contrary, existence of instabilities of the vortex in state $`2`$ due to exactly the same scenario, which is predicted in the weak interaction limit, can be readily demonstrated, see Fig. 3.
It is important to stress, that if two condensates are decoupled, $`\beta _{12}=0`$, then the negative energy mode $`\stackrel{}{w}_1^{}`$ belongs to the vortex containing component and the positive energy mode $`\stackrel{}{w}_1^+`$ belongs to the vortex free component. Because of this separation instability is not possible for any values of $`g`$, which agrees with spectral stability of unit vortices in the singly component case reported in . Thus we can conclude that the vortex instability in our example has essentially two-component nature and its analog may also exist in the case when second component is a non-condensate one.
Criterion similar to (44), but without corresponding energy analysis, has also been independently obtained in using two mode approach, i.e. condensate wave functions $`\mathrm{\Psi }_1`$ and $`\mathrm{\Psi }_2`$ have been presented as linear superposition of modes (31) with time dependent coefficients and GP equations have been reduced to the set of ordinary differential equations for these coefficients . However, this method fails taking into account an eigenmode proportional to $`r^2`$, see $`\stackrel{}{b}_3`$ in (34), which makes an important contribution to the expressions for frequencies. Therefore Eqs. (36) and (44) are different from the corresponding results presented in .
To illustrate crossings and avoided crossings in the spectrum of unit vortices we show in Fig. 4 positive part of the frequency spectra of $`\widehat{\eta }\widehat{}_{\pm 1}`$ and $`\widehat{\eta }\widehat{}_{\pm 2}`$. One can see numerous points, which on the first glance can be interpreted as crossings in the frequency spectrum. However, under the close investigation those of them which are marked by the open circles, turns to be the avoided crossing of the excitations with equal energy signs, see Fig. 5.
## V Drift of unit vortices
The unstable modes of the vortex in the state $`2`$ are of the dipole type, i.e. $`|l|=1`$, therefore their growth leads to the initial displacement of the vortex from the trap center. Displaced vortex carries less of the total angular momentum, compare to the momentum of the vortex positioned at the trap center. Lack of the angular momentum is compensated by two factors. First, vortex acquires a nonzero tangential velocity and therefore its trajectory is actually a spiral, which is similar to the dynamics of an optical vortex displaced from the center of a gaussian beam . Second, angular momentum, and vortex itself, become gradually transferred into the second condensate component, which was first demonstrated in .
If drift instability is absent, then dissipative effects still can result in the vortex drift, which was predicted for the single component condensates by several authors . Presence of the second condensate opens a channel for the energy and angular momentum transfer from the vortex into the vortex free component. It can be seen that vector $`\stackrel{}{b}_1`$, being excited by the instability, provides a channel for this transfer. Thus drift instability can also be interpreted as due to dissipation of the energy and momentum by the vortex free condensate component.
In the limits $`N1`$ and $`N1`$ our model can be approximately considered as a single component condensate with ($`N1`$) or without ($`N1`$) vortex. Unit vortex and ground state of the single component condensate are known to be spectrally stable. Therefore drift instability disappears in both limits, see Fig. 3. Increase of $`g`$ for fixed $`N`$ also results in the suppression of the instability, see Fig. 3, which means that not only relative, but also absolute increase of the number of atoms in the vortex free component stabilizes the condensate.
## VI Drift and splitting of higher order vortices
Considering higher order vortices one can expect to find instability scenario resulting in their splitting into unit vortices. This scenarion is expected to be due to growth of the collective modes with $`|l|>1`$. However, as we will see below the drift instability linked to the dipole like modes, $`|l|=1`$, also can be presented. It leads to displacement of the whole vortex from the trap center without splitting, at least at the onset of the instability.
Both drift and splitting instabilities of the higher-order vortices appear as a result of the frequency resonances of the elementary excitations with negative and positive energies, similar to the case of unit vortices. $`N1`$ corresponds to the vortex free condensate and therefore both instabilities disappear in this limit. In the limit $`N1`$ only drift instability is suppressed and one can recover periodic in $`N`$ bands of the instabilities with complex frequencies and $`|l|>1`$ similar to the results reported for higher order vortices in a singly component condensate . Note, that splitting itself was not explicitly demonstrated in . It is also interesting to note that higher-order vortices in the free, $`\widehat{V}=0`$, single component condensate are spectrally stable . Thus splitting can be considered as induced by the trapping. The vortex free condensate component plays crucial role in the drift instability, but will the latter one be presented without trapping or not remains an open problem.
### Double vortices: $`L_1=0,L_2=2`$
Considering double vortex we have found that it can be unstable with respect to the $`|l|=1,2`$ excitations. The vortex in the state $`1`$ has been found surprisingly stable. One has to take relatively small values of $`g`$ and large $`N`$ to find splitting instability. Vortex in the state $`2`$ is more unstable in a sense that splitting exists already for $`N1`$, see Fig. 6. As it is evident from Fig. 6 either drift or splitting instability can dominate vortex dynamics. If drift instability is dominant, then vortex first gets displaced from the trap center and only then splits into the unit ones, see Fig. 7. The latter happens due to the curved background which breaks cylindrical symmetry with respect to the vortex axis. After the splitting vortices remain close to each other and move towards the condensate periphery. Results of the numerical simulation of GP equations (1) presented in this Section were obtained starting from equilibrium states perturbed by random noise with amplitude $`0.05A_{1,2}`$.
The dynamics is quite different when splitting instability is dominant, see Fig. 8. In this case unit vortices appear straight at the onset of the instability development and spiral out of the condensate center being always positioned symmetrically with respect to it. After a certain period of time vortices move back to the trap center and condensate state close to the initial one is restored, see Fig. 8, then the cycle is repeated with gradually worsening degree of periodicity.
During development of the instability angular momentum and vorticity become partially transferred into the second condensate. Analysis of the transverse profiles of the phases corresponding to the density profiles shown in Fig. 8, has revealed that black spots appearing in the second condensate are indeed unit vortices, not the density holes without topological structure.
### Triple vortices: $`L_1=0,L_2=3`$
Reach variety of beautiful vortex lattices can be found considering instabilities of vortices of the order 3 and higher. This reachness can be understood in terms of spatial profiles of the unstable collective modes. E.g. triple vortex has been found unstable with respect to the perturbations with $`|l|=1,2,3`$. All components of the $`|l|=2`$ excitations are equal to zero at the trap center. Therefore one can expect that growth of this mode will develop into a spatial structure preserving vortex at the trap center. Contrary some of the components of the $`|l|=3`$ modes have humps at the center. Therefore their growth should repel all unit vortices out of the center, which leads to the breaking of the triple vortex into a triangular structure of the unit vortices moving away from the trap center. Both instabilities can be found for the same parameters values and can have close growth rates. Therefore the winning mode is selected through the process of complex competition, see Figs. 9,10.
## VII Summary
We have described general approach to stability of equilibrium in BEC using Bogoliubov theory and GP equation. Biorthogonality conditions (22), (23) and correspondence between frequency and energy spectra of the elementary excitations (27), (28) have been derived selfconsistently from first principals revealing several novel and conceptually important aspects originating in nonseldjoitness of the Bogoliubov operator.
It has been demonstrated that frequency resonances of the excitations with positive and negative energies can lead to their mutual annihilation and appearance of the collective modes with complex frequencies and zero energies. Conditions for the avoided crossing of energy levels have also been discussed. General theory has been verified both numerically and analytically in the weak interaction limit considering an example of vortices in a binary mixture of condensates.
Growth of excitations with complex frequencies leads to the two main scenarios of the instability development. First one is the spiraling of unit and double vortices out of the condensate center to its periphery. Second scenario is the splitting of the double and higher order vortices into unit ones. Absolute and/or relative increase of the number of particles in the vortex free condensate component have been found to have stabilizing effect.
###### Acknowledgements.
Author is particularly grateful to referees for their critical and very helpful comments. He also acknowledges discussions with S.M. Barnett and W.J. Firth. Numerical part of the work was significantly speeded up due to access to the computer equipment obtained via U.K. EPSRC grant GR/M31880 and assistance of G. Harkness and R. Martin.
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# FIELDS ON THE POINCARÉ GROUP: Arbitrary Spin Description and Relativistic Wave Equations
## I Introduction
Traditionally in field theory particles with different spins are described by multicomponent spin-tensor fields on Minkowski space. However, it is possible to use for this purpose scalar functions as well, which depend on both Minkowski space coordinates and on some continuous bosonic variables corresponding to spin degrees of freedom. For the first time, such fields were introduced in in connection with the problem of constructing relativistic wave equations (RWE). Fields of this type may be treated as ones on homogeneous spaces of the Poincaré group. A systematic development of such point of view was given by Finkelstein . He also gave a classification and explicit constructions of homogeneous spaces of the Poincaré group, which contain Minkowski space. The next logical step was done by Lurçat who suggested to construct quantum field theory on the Poincaré group. One of the motivations was to give a dynamical role to the spin. Some development of these ideas was given in . For example, different homogeneous spaces were described, as well as possibilities to introduce interactions in spin phase space, and to construct Lagrangian formulations were studied. The authors of arrived at the conclusion that eight is the lowest dimension of a homogeneous space suitable for a description of both half-integer and integer spins. However, no convinced physical motivation for the choice of homogeneous spaces was presented, and the interpretation of additional degrees of freedom and of corresponding quantum numbers remained an open problem.
In this paper, starting from pure group-theoretical point of view, we develop a regular approach to describing particles with different spins in the framework of a theory of scalar fields on the Poincaré group. Such fields can be considered as generating functions for conventional spin-tensor fields. In this language the problem of constructing RWE of different types is formulated from a unique position.
In our consideration, we use scalar fields on the proper Poincaré group, i.e., fields on the ten-dimensional manifold; this manifold is a direct product of Minkowski space and of the manifold of the Lorentz subgroup. These fields arise in our constructions in course of the study of a generalized regular representation (GRR). That provides a possibility to analyze then all the representations of the Poincare group. The study of GRR implies a wide use of harmonic analysis method . In a sense, this method is an alternative to one of induced representations suggested by Wigner (see ). It turns out that the fields on the Poincaré group can be considered as generating functions for usual spin-tensor fields on Minkowski space, thus we naturally obtain all results for the latter fields. However, sometimes it is more convenient to formulate properties and equations for spin-tensor fields in terms of the generating functions. Moreover, the problem of constructing RWE looks very natural in the language of the scalar fields on the group. We show that this problem can be formulated as a problem of a classification of different scalar fields. For this purpose, in accordance with the general theory of harmonic analysis, we consider various sets of commuting operators and identify constructing RWE with eigenvalue problems for this operators. We succeeded to define discrete transformations for the scalar fields using some automorphisms of the proper Poincare group. The space of scalar fields on the group turns out to be closed with respect to the discrete transformations. One ought to say that the latter transformations are of fundamental importance for constructing RWE and for their analysis. Consideration of the discrete transformations helps us to give right physical interpretation for quantum numbers which appear in course of the classification of the scalar fields.
The paper is organized as follows.
In Sect. 2 we introduce the basic objects of our study, namely, scalar fields $`f(x,𝐳)`$. The scalar fields depend on $`x`$, which are coordinates on Minkowski space, and on $`𝐳`$, which are coordinates on the Lorentz subgroup. The complex coordinates $`𝐳`$ describe spin degrees of freedom. It is shown that these fields are generating functions for usual spins-tensor fields. Classifying the scalar fields with the help of various sets of commuting operators on the group, we get description of irreps of the group. We formulate a general scheme of constructing RWE in this language in any dimensions. We introduce discrete transformations in the space of the scalar functions and we relate these transformations to automorphisms of the proper Poincaré group.
In Sect. 3 we apply the above general scheme to detailed study of scalar fields on two-dimensional Poincaré and Euclidean groups. In particular, we construct RWE and analyze their solutions.
Three-dimensional Poincaré and Euclidean group case is considered in Sect. 4. Besides finite-component equations, we also construct positive energy RWE assotiated with unitary infinite-dimensional irreps of 2+1 Lorentz group. These equations, in particular, describe particles with fractional spins.
In Sect. 5 we study scalar fields on the $`3+1`$ proper Poincaré group. A connection of the present consideration with other approaches to RWE theory is considered in detail. In particular, we pay significant attention to equations with subsidiary conditions. General first-order Gel’fand–Yaglom equations (including Bhabha equations), Dirac–Fierz–Pauli equations, and Rarita–Schwinger equations arise in the present consideration as well. This give a regular base for comparison of properties of various RWE.
Doing the classification of scalar functions in 2, 3, and 4 dimensions, we obtain equations describing fields with fixed mass and spin. In Sect. 6 we consider the general features of these equations.
One ought to say that the construction of RWE is elaborated in detail only for the massive case. We plane to discuss the massless case in a later article.
## II Fields on the proper Poincaré group and spin description
### A Parametrization of the Poincaré group
Consider Poincaré group transformations
$$x^\nu =\mathrm{\Lambda }_\mu ^\nu x^\mu +a^\nu $$
(1)
of coordinates $`x=(x^\mu ,\mu =0,\mathrm{},D)`$ in $`d=D+1`$-dimensional Minkowski space, $`ds^2=\eta _{\mu \nu }dx^\mu dx^\nu `$, $`\eta _{\mu \nu }=diag(1,1,\mathrm{},1)`$. The matrices $`\mathrm{\Lambda }`$ define rotations in Minkovski space and belong to the vector representation of $`O(D,1)`$ group. We are also going to consider $`D`$-dimensional Euclidean case in which $`ds^2=\eta _{ik}dx^idx^k`$, and $`\eta _{ik}=diag(1,1,\mathrm{},1)`$, $`i,k=1,\mathrm{},D`$. Here the matrices $`\mathrm{\Lambda }`$ belong to the vector representation of $`O(D)`$ group.
The transformations (1) which can be obtained continuously from the identity form the proper Poincaré group $`M_0(D,1)`$ with the elements $`g=(a,\mathrm{\Lambda })`$. Corresponding homogeneous transformations ($`a=0`$) form the proper Lorentz group $`SO_0(D,1)`$. In the Euclidean case we deal with $`M_0(D)`$ and $`SO(D)`$ respectively. The composition law and the inverse element of these groups have the form
$$(a_2,\mathrm{\Lambda }_2)(a_1,\mathrm{\Lambda }_1)=(a_2+\mathrm{\Lambda }_2a_1,\mathrm{\Lambda }_2\mathrm{\Lambda }_1),g^1=(\mathrm{\Lambda }^1a,\mathrm{\Lambda }^1).$$
(2)
Thus, the groups $`M_0(D,1)`$ and $`M_0(D)`$ are semidirect products
$`M_0(D,1)=T(d)\times )SO_0(D,1),M_0(D)=T(D)\times )SO_(D),`$where T(d) is $`d`$-dimensional translation group.
There exists one-to-one correspondence between the vectors $`x`$ and $`2\times 2`$ Hermitian matrices $`X`$ in pseudo-Euclidean spaces of 2,3 and 4 dimensions,<sup>*</sup><sup>*</sup>* We use two sets of $`2\times 2`$ matrices $`\sigma _\mu =(\sigma _0,\sigma _k)`$ and $`\overline{\sigma }_\mu =(\sigma _0,\sigma _k)`$, $`\sigma _0=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),\sigma _1=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),\sigma _2=\left(\begin{array}{cc}0& i\\ i& 0\end{array}\right),\sigma _3=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right).`$
$$xX,X=x^\mu \sigma _\mu .$$
(3)
Namely:
$`d=3+1:X=\left(\begin{array}{cc}x^0+x^3& x^1ix^2\\ x^1+ix^2& x^0x^3\end{array}\right),`$ (6)
$`d=2+1:X=\left(\begin{array}{cc}x^0& x^1ix^2\\ x^1+ix^2& x^0\end{array}\right),`$ (9)
$`d=1+1:X=\left(\begin{array}{cc}x^0& x^1\\ x^1& x^0\end{array}\right).`$ (12)
In all the above cases
$$detX=\eta _{\mu \nu }x^\mu x^\nu ,x^\mu =\frac{1}{2}Tr(X\overline{\sigma }^\mu ).$$
(13)
In Euclidean spaces of 2 and 3 dimensions a similar correspondence has the form
$`D=3:X=\left(\begin{array}{cc}x^3& x^1ix^2\\ x^1+ix^2& x^3\end{array}\right),`$ (16)
$`D=2:X=\left(\begin{array}{cc}\hfill x^2& \hfill x^1\\ \hfill x^1& \hfill x^2\end{array}\right).`$ (19)
If $`x`$ is subjected to a transformation (1), then $`X`$ transforms as follows (see, for example, :
$$X^{}=UXU^{}+A,$$
(20)
where $`A=a^\mu \sigma _\mu `$ and $`U`$ are some $`2\times 2`$ complex matrices obeying the conditions
$$\sigma _\nu \mathrm{\Lambda }_\mu ^\nu =U\sigma _\mu U^{}.$$
(21)
Eq. (21) relates the matrices $`\mathrm{\Lambda }`$ and $`U`$. There are many $`U`$ which correspond to the same $`\mathrm{\Lambda }`$. We may fix this arbitrariness imposing the condition
$$detU=1,$$
(22)
which does not contradict to the relation $`detU=e^{i\varphi }`$, which follows from (21). However, even after that, there is no one-to-one correspondence between $`\mathrm{\Lambda }`$ and $`U`$, namely two matrices ($`U`$,$`U`$) correspond to one $`\mathrm{\Lambda }`$. Considering both $`U`$ and $`U`$ as representatives for $`\mathrm{\Lambda }`$, we in fact go over from $`SO_0(D,1)`$ to its double covering group $`\mathrm{Spin}(D,1)`$, or, in the Euclidean case, from $`SO(D)`$ to its double covering group $`\mathrm{Spin}(D)`$. In the dimensions under consideration the groups $`\mathrm{Spin}(D,1)`$ and $`\mathrm{Spin}(D)`$ are isomorphic to the following ones:We denote the complex conjugation by $``$ above the quantities.
$`d=3+1:USL(2,C),U=\left(\begin{array}{cc}u_1^1& u_2^1\\ u_1^2& u_2^2\end{array}\right),u_1^1u_2^2u_1^2u_2^1=1,`$ (25)
$`d=2+1:USU(1,1),U=\left(\begin{array}{cc}u_1& u_2\\ \stackrel{}{}u_2& \stackrel{}{}u_1\end{array}\right),|u_1|^2|u_2|^2=1,`$ (28)
$`D=3:USU(2),U=\left(\begin{array}{cc}\hfill u_1& \hfill u_2\\ \hfill \stackrel{}{}u_2& \hfill \stackrel{}{}u_1\end{array}\right),|u_1|^2+|u_2|^2=1,`$ (31)
$`d=1+1:USO(1,1),U=\left(\begin{array}{cc}\mathrm{cosh}\frac{\varphi }{2}& \mathrm{sinh}\frac{\varphi }{2}\\ \mathrm{sinh}\frac{\varphi }{2}& \mathrm{cosh}\frac{\varphi }{2}\end{array}\right),`$ (34)
$`D=2:USO(2),U=\left(\begin{array}{cc}\hfill \mathrm{cos}\frac{\varphi }{2}& \hfill \mathrm{sin}\frac{\varphi }{2}\\ \hfill \mathrm{sin}\frac{\varphi }{2}& \hfill \mathrm{cos}\frac{\varphi }{2}\end{array}\right).`$ (37)
Considering nonhomogeneous transformations and retaining both elements $`U`$ and $`U`$ in the consideration, we go over from the groups $`M_0(D,1)`$ and $`M_0(D)`$ to the groups
$`M(D,1)=T(d)\times )\mathrm{Spin}(D,1),M(D)=T(D)\times )\mathrm{Spin}(D)`$respectively. As it is known, that allows one to avoid double-valued representations for half integer spins. Thus, there exists one-to-one correspondence between the elements $`g`$ of the groups $`M(D,1)`$, $`M(D)`$ and two $`2\times 2`$ matrices, $`g(A,U)`$. The first one $`A`$ corresponds to translations and the second one $`U`$ corresponds to rotations. Eq. (20) describes the action of $`M(D,1)`$ on Minkowski space (the latter is coset space $`M(D,1)/\mathrm{Spin}(D,1)`$). As a consequence of (20), one can obtain the composition law and the inverse element of the groups $`M(D,1)`$, $`M(D)`$:
$$(A_2,U_2)(A_1,U_1)=(U_2A_1U_2^{}+A_2,U_2U_1),g^1=(U^1A(U^1)^{},U^1).$$
(38)
The matrices $`U`$ in the dimensions under consideration satisfy the following identities:
$`USL(2,C):\sigma _2U\sigma _2=(U^T)^1;`$ (39)
$`USU(1,1):\sigma _1U\sigma _1=\stackrel{}{}U,\sigma _2U\sigma _2=(U^T)^1,\sigma _3U\sigma _3=(U^{})^1,`$ (40)
$`USU(2):\sigma _2U\sigma _2=(U^T)^1=\stackrel{}{}U.`$ (41)
An equivalent picture arise in terms of the matrices $`\overline{X}=x^\mu \overline{\sigma }_\mu `$. Using the relation $`\overline{X}=\sigma _2X^T\sigma _2`$, the transformation law for $`X`$ (20), and the identity (39), one can get
$$\overline{X}^{}=(U^{})^1\overline{X}U^1+\overline{A}.$$
(42)
Thus, $`\overline{X}`$ are transformed by means of the elements $`(\overline{A},(U^{})^1)`$. The relation $`(A,U)(\overline{A},(U^{})^1)`$ defines an automorphism of the Poincaré group $`M(D,1)`$. In Euclidean case the matrices $`U`$ are unitary, and the latter relation is reduced to $`(A,U)(A,U)`$.
The representation of the Poincaré transformations in the form (20) is closely related to a representation of finite rotations in $`^d`$ in terms of the Clifford algebra. In higher dimensions the transformation law has the same form, where $`A`$ is a vector element and $`U`$ corresponds to an invertible element (spinor element) of the Clifford algebra . Besides, the representation of the finite transformations in the form (20) can be useful for spin description by means of Grassmannian variables $`\xi `$, since $`\xi `$ and $`\xi `$ give a realization of the Clifford algebra .
### B Regular representation and scalar functions on the group
It is well known that any irrep of a group $`G`$ is contained (up to the equivalence) in a decomposition of a GRR. Thus, the study of GRR is an effective method for the analysis of irreps of the group. Consider, first, the left GRR $`T_L(g)`$, which is defined in the space of functions $`f(g_0)`$, $`g_0G`$, on the group, as
$$T_L(g)f(g_0)=f^{}(g_0)=f(g^1g_0),gG.$$
(43)
As a consequence of the relation (43), we can write
$$f^{}(g_0^{})=f(g_0),g_0^{}=gg_0.$$
(44)
Let $`G`$ be the group $`M(3,1)`$, and we use the parametrization of its elements by two $`2\times 2`$ matrices (one hermitian and another one from $`SL(2,C)`$), which was described in the previous Sect. At the same time, using such a parametrization, we choose the following notations:
$$g(A,U),g_0(X,Z),$$
(45)
where $`A,X`$ are $`2\times 2`$ hermitian matrices and $`U,ZSL(2,C).`$ The map $`g_0(X,Z)`$ creates the correspondence
$`g_0(x,z,\underset{¯}{z}),\text{where}x=(x^\mu ),z=(z_\alpha ),\underset{¯}{z}=(\underset{¯}{z}_\alpha ),`$ (46)
$`\mu =0,1,2,3,\alpha =1,2,z_1\underset{¯}{z}_2z_2\underset{¯}{z}_1=1,`$ (47)
by virtue of the relations
$$X=x^\mu \sigma _\mu ,Z=\left(\begin{array}{cc}z_1& \underset{¯}{z}_1\\ z_2& \underset{¯}{z}_2\end{array}\right)SL(2,C).$$
(48)
On the other hand, we have the correspondence $`g_0^{}(x^{},z^{},\underset{¯}{z}^{})`$,
$`g_0^{}=gg_0(X^{},Z^{})=(A,U)(X,Z)=(UXU^++A,UZ)(x^{},z^{},\underset{¯}{z}^{}),`$ (49)
$`x^\mu \sigma _\mu =X^{}=UXU^++Ax^\mu =(\mathrm{\Lambda }_0)_\nu ^\mu x^\nu +a^\mu ,\mathrm{\Lambda }USL(2,C),`$ (50)
$`\left(\begin{array}{cc}z_1^{}& \underset{¯}{z}_1^{}\\ z_2^{}& \underset{¯}{z}_2^{}\end{array}\right)=Z^{}=UZz_\alpha ^{}=U_\alpha ^\beta z_\beta ,\underset{¯}{z}_\alpha ^{}=U_\alpha ^\beta \underset{¯}{z}_\beta ,U=(U_\alpha ^\beta ),z_1^{}\underset{¯}{z}_2^{}z_2^{}\underset{¯}{z}_1^{}=1.`$ (53)
Then the relation (44) takes the form
$`f^{}(x^{},z^{},\underset{¯}{z}^{})=f(x,z,\underset{¯}{z}),`$ (54)
$`x^\mu =(\mathrm{\Lambda }_0)_\nu ^\mu x^\nu +a^\mu ,\mathrm{\Lambda }USL(2,C),`$ (55)
$`z_\alpha ^{}=U_\alpha ^\beta z_\beta ,\underset{¯}{z}_\alpha ^{}=U_\alpha ^\beta \underset{¯}{z}_\beta ,z_1\underset{¯}{z}_2z_2\underset{¯}{z}_1=z_1^{}\underset{¯}{z}_2^{}z_2^{}\underset{¯}{z}_1^{}=1.`$ (56)
The relations (54)-(56) admit a remarkable interpretation. We may treat $`x`$ and $`x^{}`$ in these relations as position coordinates in Minkowski space (in different Lorentz refrence frames) related by proper Poincare transformations, and the sets $`(z,\underset{¯}{z})`$ and $`(z,\underset{¯}{z}^{})`$ may be treated as spin coordinates in these Lorentz frames. They are transformed according to the formulas (56). Carrying two-dimensional spinor representation of the Lorentz group, the variables $`z`$ and $`\underset{¯}{z}`$ are invariant under translations as one can expect for spin degrees of freedom. Thus, we may treat sets $`(x,z,\underset{¯}{z})`$ as points in a position-spin space with the transformation law (55), (56) under the change from one Lorentz reference frame to another. In this case equations (54)-(56) present the transformation law for scalar functions on the position-spin space.
On the other hand, as we have seen, the sets $`(x,z,\underset{¯}{z})`$ are in one-to-one correspondence to the group $`M(3,1)`$ elements. Thus, the functions $`f(x,z,\underset{¯}{z})`$ are still functions on this group. That is why we often call them scalar functions on the group as well, remembering that the term ”scalar” came from the above interpretation.
Remember now that different functions of such type correspond to different representations of the group $`M(3,1)`$. Thus, the problem of classification of all irreps of this group is reduced to the problem of a classification of all scalar functions on position-spin space. However, for the purposes of such classification, it is natural to restrict ourselves by the scalar functions which are analytic both in $`z,\underset{¯}{z}`$ and in $`\stackrel{}{}z,\stackrel{}{}\underset{¯}{z}`$ (or, simply speaking, which are differentiable with respect to these arguments). Further such functions are denoted by $`f(x,z,\underset{¯}{z},\stackrel{}{}z,\stackrel{}{}\underset{¯}{z})=f(x,𝐳)`$, $`𝐳=(z,\underset{¯}{z},\stackrel{}{}z,\stackrel{}{}\underset{¯}{z})`$.
Consider now the right GRR $`T_R(g)`$. This representation is defined in the space of functions $`f(g_0)`$, $`g_0G`$ as
$$T_R(g)f(g_0)=f^{}(g_0)=f(g_0g),gG,$$
(57)
As a consequence of the relation (57), we can write
$$f^{}(g_0^{})=f(g_0),g_0^{}=g_0g^1.$$
(58)
In the case of the proper Poincare group, the right transformations act on $`g_0(X,Z)`$ according to the formula
$$g_0^{}=g_0g^1(X^{},Z^{})=(X+Z^1A(Z^1)^{},ZU^1),$$
(59)
hence $`x^\mu =x^\mu +L_\nu ^\mu a^\nu `$, where the matrix $`L`$ depends on $`𝐳`$, $`\sigma _\nu L_\mu ^\nu =Z^1\sigma _\mu (Z^1)^{}`$. The transformations for $`x,𝐳`$ do not admit similar to the left GRR case interpretation. In particular, the transformation low for $`x`$ does not look as a Lorentz transformation. On the other hand, the study of the right GRR is useful for the purposes of the classification of the Poincare group irreps, since the generators of the right GRR are used to construct complete sets of commuting operators on the group.
### C Generators of generalized regular representations
Generators of the left GRR correspond to translations and rotations. They can be written as
$$\widehat{p}_\mu =i/x^\mu ,\widehat{J}_{\mu \nu }=\widehat{L}_{\mu \nu }+\widehat{S}_{\mu \nu },$$
(60)
where $`\widehat{L}_{\mu \nu }=i(x_\mu _\nu x_\nu _\mu )`$ are angular momentum operators, and $`\widehat{S}_{\mu \nu }`$ are spin operators depending on $`𝐳`$ and $`/𝐳`$. An explicit form of spin operators is given in the Appendix.
The algebra of the generators (60) has the form
$`[\widehat{p}_\mu ,\widehat{p}_\nu ]=0,[\widehat{J}_{\mu \nu },\widehat{p}_\rho ]=i(\eta _{\nu \rho }\widehat{p}_\mu \eta _{\mu \rho }\widehat{p}_\nu ),`$ (61)
$`[\widehat{J}_{\mu \nu },\widehat{J}_{\rho \sigma }]=i\eta _{\nu \rho }\widehat{J}_{\mu \sigma }i\eta _{\mu \rho }\widehat{J}_{\nu \sigma }i\eta _{\nu \sigma }\widehat{J}_{\mu \rho }+i\eta _{\mu \sigma }\widehat{J}_{\nu \rho }.`$ (62)
In the space of Fourier transforms
$$\phi (p,𝐳)=(2\pi )^{d/2}f(x,𝐳)e^{ipx}𝑑x$$
(63)
the left GRR acts as (one has to use (43)):
$$T_L(g)\phi (p,𝐳)=e^{iap^{}}\phi (p^{},g^1𝐳),p^{}=g^1pP^{}=U^1P(U^1)^{},P=p_\mu \sigma ^\mu .$$
(64)
One can see that $`detZ`$ and $`detP=p^2`$ are invariant under the transformations $`p^2=\eta ^{\mu \nu }p_\mu p_\nu `$. Since we do not use $`p`$ with upper indices, this does not lead to a misunderstanding. (64) and that $`p^2`$ is an eigenvalue of the Casimir operator $`\widehat{p}^2`$.
For the groups $`M(D)`$ there are two types of representations depending on $`p^2`$: 1) $`p^20`$; 2) $`p^2=0`$; then all $`p_i=0`$, and irreps are labelled by eigenvalues of Casimir operators of the rotation subgroup.
For the groups $`M(D,1)`$ there are four types of representations depending on the eigenvalues $`m^2`$ of the Casimir operator $`\widehat{p}^2`$: 1) $`m^2>0`$; 2) $`m^2<0`$ (tachyon); 3) $`m^2=0`$, $`p_00`$ (massless particle); 4) $`m^2=p_0=0`$, irreps are labelled by eigenvalues of the Casimir operators of the Lorentz subgroup, and the corresponding functions do not depend on $`x`$.
For decomposing the left GRR we are going to construct a complete set of commuting operators in the space of functions on the group. Together with the Casimir operators some functions of right generators<sup>§</sup><sup>§</sup>§The physical meaning of the right generators is not so transparent. However, one can remember that the right generators of $`SO(3)`$ in the nonrelativistic rotator theory are interpreted as operators of angular momentum in a rotating body-fixed reference frame . may be included in such a set. Therefore it is necessary to know the explicit form of right generators. As a consequence of the formulas
$`T_R(g)f(x,𝐳)=f(xg,𝐳g),xgX+ZAZ^{},𝐳gZU,`$ (65)
$`T_R(g)\phi (p,𝐳)=e^{ia^{}p}\phi (p,𝐳g),a^{}A^{}=ZAZ^{}`$ (66)
one can obtain
$$\widehat{p}_\mu ^R=(L^1(𝐳))_\mu ^\nu p_\nu ,\widehat{J}_{\mu \nu }^R=\widehat{S}_{\mu \nu }^R,$$
(67)
where $`LSO(D,1)`$ (or $`LSO(D,1)`$ in the Euclidean case). The operators of right translations can also be written in the form $`\widehat{P}^R=Z^1\widehat{P}(Z^1)^{}`$; operators $`\widehat{S}_{\mu \nu }`$ and $`\widehat{S}_{\mu \nu }^R`$ are left and right generators of $`\mathrm{Spin}(D,1)`$ (or $`\mathrm{Spin}(D)`$) and depend on $`𝐳`$ only. All the right generators (67) commute with all the left generators (60) and obey the same commutation relations (62).
In accordance with theory of harmonic analysis on Lie groups there exists a complete set of commuting operators, which includes Casimir operators, a set of the left generators and a set of right generators (both sets contain the same number of the generators). The total number of commuting operators is equal to the number of parameters of the group. In a decomposition of the left GRR the nonequivalent representations are distinguished by eigenvalues of the Casimir operators, equivalent representations are distinguished by eigenvalues of the right generators, and the states inside the irrep are distinguished by eigenvalues of the left generators.
In particular, Casimir operators of spin Lorentz subgroup are functions of $`\widehat{S}_{\mu \nu }^R`$ (or $`\widehat{S}_{\mu \nu }`$) and commute with all the left generators (with left translations and rotations), but do not commute with generators of the right translations. These operators distinguish equivalent representations in the decomposition of the left GRR. Notice that some aspects of the theory of harmonic analysis on the 3+1 and 2+1 Poincaré groups were considered in and respectively.
If GRR acts in the space of all functions on the group $`G`$, then a regular representation acts in the space of functions $`L^2(G,\mu )`$, such that the norm
$$\stackrel{}{}f(g)f(g)d\mu (g)$$
(68)
is finite , where $`d\mu (g)`$ is an invariant measure on the group. The regular representation is unitary, as it follows from (68) and from the invariance of the measure. However we will also use nonunitary representations (in particular, finite-dimensional representations of the Lorentz group). Therefore we consider the GRR as a more useful concept.
### D Fields on the Poincaré group
As we have shown, the relations associated with the left GRR (43) define the transformation law for coordinates ($`x,𝐳`$) on the position-spin space under the change from one Lorentz reference frame to another. The equations
$`f^{}(x^{},𝐳^{})=f(x,𝐳),`$ (69)
$`x^{}=gx=\mathrm{\Lambda }x+aUXU^{}+A,𝐳^{}=g𝐳UZ.`$ (70)
define a scalar field on this space (i.e. a scalar field on the Poincaré group). In contrast to scalar field on Minkowski space, this field is reducible with respect to both mass and spin.
Consider the transformation laws of $`x`$ and $`𝐳`$ in various dimensions more detail.
In two-dimensional case matrices $`Z`$ depend on only one parameter (angle or hyperbolic angle, see (34),(37)). The functions on the group depend on $`x=(x^\mu )`$ and $`z=e^\alpha `$ (or $`x=(x^k)`$ and $`z=e^{i\alpha }`$ in Euclidean case); it is appropriate to consider these functions as functions of real parameter $`\alpha `$ directly.
In three-dimensional case according to (28),(31)
$$D=3:Z=\left(\begin{array}{cc}\hfill z_1& \hfill \stackrel{}{}z_{\dot{2}}\\ \hfill z_2& \hfill \stackrel{}{}z_{\dot{1}}\end{array}\right),d=2+1:Z=\left(\begin{array}{cc}z_1& \stackrel{}{}z_{\dot{2}}\\ z_2& \stackrel{}{}z_{\dot{1}}\end{array}\right),detZ=1.$$
(71)
Functions $`f(x,𝐳)`$ depend on $`x=(x^\mu )`$ (in Euclidean case $`x=(x^k)`$) and $`𝐳=(z,\stackrel{}{}z)`$, where $`z`$ are the elements of the first column of matrix (71). Let us write the relation (70) for $`d=2+1`$ in component-wise form
$`x^\nu \sigma _{\nu \alpha \dot{\alpha }}=U_\alpha ^\beta x^\mu \sigma _{\mu \beta \dot{\beta }}\stackrel{}{}U_{\dot{\alpha }}^{\dot{\beta }}+a^\mu \sigma _{\mu \alpha \dot{\alpha }},`$ (72)
$`z_\alpha ^{}=U_\alpha ^\beta z_\beta ,\stackrel{}{}z_{\dot{\alpha }}^{}=\stackrel{}{}U_{\dot{\alpha }}^{\dot{\beta }}\stackrel{}{}z_{\dot{\beta }},z^\alpha =(U^1)_\beta ^\alpha z^\beta ,\stackrel{}{}z^{\dot{\alpha }}=(\stackrel{}{}U^1)_{\dot{\beta }}^{\dot{\alpha }}\stackrel{}{}z^{\dot{\beta }}.`$ (73)
Undotted and dotted indices correspond to spinors transforming by means of matrix $`U`$ and complex conjugate matrix $`\stackrel{}{}U`$. Invariant tensor $`\sigma _{\nu \alpha \dot{\alpha }}`$ has one vector index and two spinor indices of distinct types.
For the group $`M(3,1)`$ matrices $`Z`$, $`detZ=1`$, has the form (48); the elements $`z^\alpha `$ and $`\underset{¯}{z}^\alpha `$ of first and second columns of matrix (48) are subjected to the same transformation law. The functions $`f(x,𝐳)`$ depend on $`x=(x^\mu )`$ and $`𝐳=(z,\stackrel{}{}z,\underset{¯}{z},\stackrel{}{}\underset{¯}{z})`$. The main reason to consider not real parameters (for example, real and imaginary parts of $`z,\underset{¯}{z}`$), but of $`z,\underset{¯}{z}`$ and $`\stackrel{}{}z,\stackrel{}{}\underset{¯}{z}`$, is the fact that the complex variables are subjected to simple transformation rule. Besides, the use of spaces of analytic and antianalytic functions is suitable for the problem of decomposition of GRR.
According to (70) and (42) one may write the transformation law of $`x^\mu `$, $`z_\alpha `$, $`\stackrel{}{}z_{\dot{\alpha }}`$ in component-wise form
$`x^\nu \sigma _{\nu \alpha \dot{\alpha }}=U_\alpha ^\beta x^\mu \sigma _{\mu \beta \dot{\beta }}\stackrel{}{}U_{\dot{\alpha }}^{\dot{\beta }}+a^\mu \sigma _{\mu \alpha \dot{\alpha }},x^\nu \overline{\sigma }_\nu ^{\dot{\alpha }\alpha }=(\stackrel{}{}U^1)_{\dot{\beta }}^{\dot{\alpha }}x^\mu \overline{\sigma }_\mu ^{\dot{\beta }\beta }(U^1)_\beta ^\alpha +a^\mu \overline{\sigma }_\mu ^{\dot{\alpha }\alpha },`$ (74)
$`z_\alpha ^{}=U_\alpha ^\beta z_\beta ,\stackrel{}{}z_{\dot{\alpha }}^{}=\stackrel{}{}U_{\dot{\alpha }}^{\dot{\beta }}\stackrel{}{}z_{\dot{\beta }},z^\alpha =(U^1)_\beta ^\alpha z^\beta ,\stackrel{}{}z^{\dot{\alpha }}=(\stackrel{}{}U^1)_{\dot{\beta }}^{\dot{\alpha }}\stackrel{}{}z^{\dot{\beta }}.`$ (75)
It is easy to see from (74) that the tensors
$$\sigma _{\mu \alpha \dot{\alpha }}=(\sigma _\mu )_{\alpha \dot{\alpha }},\overline{\sigma }_\mu ^{\dot{\alpha }\alpha }=(\overline{\sigma }_\mu )^{\dot{\alpha }\alpha }$$
(76)
are invariant. These tensors are usually used to convert vector indices into spinor ones and vice versa or to construct vector from two spinors of different types:
$$x^\mu =\frac{1}{2}\overline{\sigma }^{\mu \dot{\alpha }\alpha }x_{\dot{\alpha }\alpha },x_{\alpha \dot{\alpha }}=\sigma _{\mu \alpha \dot{\alpha }}x^\mu ,q^\mu =\frac{1}{2}\overline{\sigma }^{\mu \dot{\alpha }\alpha }z_\alpha \stackrel{}{}z_{\dot{\alpha }}.$$
(77)
In consequence of the unimodularity of $`2\times 2`$ matrices $`U`$ there exist invariant antisymmetric tensors $`\epsilon ^{\alpha \beta }=\epsilon ^{\beta \alpha }`$, $`\epsilon ^{\dot{\alpha }\dot{\beta }}=\epsilon ^{\dot{\beta }\dot{\alpha }}`$, $`\epsilon ^{12}=\epsilon ^{\dot{1}\dot{2}}=1`$, $`\epsilon _{12}=\epsilon _{\dot{1}\dot{2}}=1`$. Now spinor indices are lowered and raised according to the rules
$$z_\alpha =\epsilon _{\alpha \beta }z^\beta ,z^\alpha =\epsilon ^{\alpha \beta }z_\beta ,$$
(78)
and in particular one can get $`\sigma _{\mu \alpha \dot{\alpha }}=\overline{\sigma }_{\mu \dot{\alpha }\alpha }`$. Below we will also use the notations $`_\alpha =/z^\alpha `$, $`^{\dot{\alpha }}=/\stackrel{}{}z_{\dot{\alpha }}`$, and correspondingly $`^\alpha =/z_\alpha `$, $`_{\dot{\alpha }}=/\stackrel{}{}z^{\dot{\alpha }}`$.
In the framework of theory of the scalar functions on the Poincaré group a standard spin description in terms of multicomponent functions arises under the separation of space and spin variables.
Since $`𝐳`$ is invariant under translations, any function $`\varphi (𝐳)`$ carry a representation of the Lorentz group. Let a function $`f(h)=f(x,𝐳)`$ allows the representation
$$f(x,𝐳)=\varphi ^n(𝐳)\psi _n(x),$$
(79)
where $`\varphi ^n(𝐳)`$ form a basis in the representation space of the Lorentz group. The latter means that one may decompose the functions $`\varphi ^n(𝐳^{})`$ of transformed argument $`𝐳^{}=g𝐳`$ in terms of the functions $`\varphi ^n(z)`$:
$$\varphi ^n(𝐳^{})=\varphi ^l(𝐳)L_l^n(U).$$
(80)
An action of the Poincaré group on a line $`\varphi ^n(𝐳)\varphi ^n(𝐳)`$ is reduced to a multiplication by matrix $`L(U)`$, where $`U\mathrm{Spin}(D,1)`$, $`\varphi (𝐳^{})=\varphi (𝐳)L(U)`$.
Comparing the decompositions of the function $`f^{}(x^{},𝐳^{})=f(x,𝐳)`$ over the transformed basis $`\varphi (𝐳^{})`$ and over the initial basis $`\varphi (𝐳)`$,
$$f^{}(x^{},𝐳^{})=\varphi (𝐳^{})\psi ^{}(x^{})=\varphi (𝐳)L(U)\psi ^{}(x^{})=\varphi (𝐳)\psi (x),$$
where $`\psi (x)`$ is a column with components $`\psi _n(x)`$, one may obtain
$$\psi ^{}(x^{})=L(U^1)\psi (x),$$
(81)
i.e. the transformation law of a tensor field on Minkowski space. This law correspond to the representation of the Poincaré group acting in a linear space of tensor fields as follows $`T(g)\psi (x)=L(U^1)\psi (\mathrm{\Lambda }^1(xa))`$. According to (80) and (81), the functions $`\varphi (z)`$ and $`\psi (x)`$ are transformed under contragradient representations of the Lorentz group.
For example, let us consider scalar functions on the Poincaré group $`f_1(x,𝐳)=\psi _\alpha (x)z^\alpha `$ and $`f_2(x,𝐳)=\overline{\psi }_\alpha (x)\stackrel{}{}z^\alpha `$, which correspond to spinor representations of Lorentz group. According to (79) and (81)
$$\psi _\alpha ^{}(x^{})=U_\alpha ^\beta \psi _\beta (x),\overline{\psi }_{\dot{\alpha }}^{}(x^{})=\stackrel{}{}U_{\dot{\alpha }}^{\dot{\beta }}\overline{\psi }_{\dot{\beta }}(x).$$
(82)
The product $`\psi _\alpha (x)\overline{\psi }^\alpha (x)`$ is Poincaré invariant.
Thus tensor fields of all spins are contained in the decomposition of the field (69) on the Poincaré group, and the problems of their classification and construction of explicit realizations are reduced to problem of the decomposition of the left GRR.
Notice that above we reject the phase transformations, which correspond to $`U=e^{i\varphi }`$. This transformations of $`U(1)`$ group do not change space-time coordinates $`x`$, but change the phase of $`𝐳`$. According to (80) and (81) that leads to the transformation of phase of tensor field components $`\psi _n(x)`$. Taking account of this transformations means the consideration of the functions on the group $`T(d)\times )\mathrm{Spin}(D,1)\times U(1)`$.
### E Automorphisms of the Poincaré group and discrete transformations: P,C,T
Let us consider elements $`g(A,U)`$, $`g_0(X,Z)`$ of the Poincaré group $`M(D,1)`$. It is easy to see that transformations
$`(A,U)(\overline{A},(U^{})^1),(X,Z)(\overline{X},(Z^{})^1),`$ (83)
$`(A,U)(\stackrel{}{}A,\stackrel{}{}U),(X,Z)(\stackrel{}{}X,\stackrel{}{}Z),`$ (84)
$`(A,U)(A,U),(X,Z)(X,Z)`$ (85)
are outer involutory automorphisms of the group and generate finite group consisting of eight elements.
The automorphisms (83)-(85) define discrete transformations of space-time and spin coordinates $`x,𝐳`$. The substitution of transformed coordinates into the functions $`f(x,𝐳)`$ (or into the generators (60)) leads to change signs of some physical variables. (Notice that the substitution both into the functions and into the generators leaves signs unaltered.)
The space reflection (or parity transformation $`P`$) is defined by the relations $`x^0x^0`$, $`x^kx^k`$, or $`X\overline{X}`$. If $`X`$ is transformed by means of the group element $`(A,U)`$, then $`\overline{X}`$ is transformed by means of the group element $`(\overline{A},(U^{})^1)`$, see (42). Therefore the space reflection represents a realization of the automorphism (83) of the Poincaré group
$$(X,Z)\stackrel{P}{}(\overline{X},(Z^{})^1).$$
(86)
Thus, under the space reflection $`x`$ and $`𝐳`$ have to be changed in all the constructions according to (86). In particular, for the momentum $`P=p_\mu \sigma ^\mu `$ we obtain $`P\overline{P}`$, where $`\overline{P}=p_\mu \overline{\sigma }^\mu `$. The generators of the rotations are not changed and the generators of the boosts change their signs only.
The time reflection transformation $`T^{}`$ is defined by the relation $`x^\mu (1)^{\delta _{0\mu }}x^\mu `$, or $`X\overline{X}`$, and corresponds to the composition of automorphisms (83) and (85):
$$(X,Z)\stackrel{T^{}}{}(\overline{X},(Z^{})^1).$$
(87)
Inversion $`PT^{}`$, $`(X,Z)\stackrel{PT^{}}{}(X,Z)`$, corresponds to the automorphism (85).
Automorphism of complex conjugation (84) means substitution $`ii`$,
$$f(x,𝐳)\stackrel{C}{}\stackrel{}{}f(x,𝐳).$$
(88)
One can show that in the framework of the characteristics related to the Poincaré group this transformation corresponds to the charge conjugation. Both the transformation (88) and charge conjugation change signs of all the generators, $`\widehat{p}_\mu \widehat{p}_\mu `$, $`\widehat{L}_{\mu \nu }\widehat{L}_{\mu \nu }`$, $`\widehat{S}_{\mu \nu }\widehat{S}_{\mu \nu }`$. Below, considering RWE, we will see that transformation (88) change also the sign of current vector $`j^\mu `$.
The time reversal $`T`$ is defined by the relation $`X\overline{X}`$ (the time reflection transformation $`T^{}`$), with the supplementary condition of energy sign conservation that means $`P\overline{P}`$. Therefore, the conditions $`\widehat{p}_\mu (1)^{\delta _{0\mu }}\widehat{p}_\mu `$, $`\widehat{L}_{\mu \nu }(1)^{\delta _{0\mu }+\delta _{0\nu }}\widehat{L}_{\mu \nu }`$, $`\widehat{S}_{\mu \nu }(1)^{\delta _{0\mu }+\delta _{0\nu }}\widehat{S}_{\mu \nu }`$ take place. The transformation $`CT^{}`$ obeys these conditions.
However, it is known that it is possible to give two distinct definitions of time reversal transformation obeying conditions mentioned above. Wigner time reversal $`T_w`$ leaves the total charge (and correspondingly $`j^0`$) unaltered, and reverses the direction of current $`j^k`$. Schwinger time reversal $`T_{sch}`$ leaves the current $`j^k`$ invariant and reverses the charge.
The transformation $`CT^{}`$ changes the sign of $`j^0`$ and therefore can be identify with Schwinger time reversal, $`T_{sch}=CT^{}`$. $`CPT_{sch}`$-transformation corresponds to the inversion $`(X,Z)(X,Z)`$. Wigner time reversal $`T_w`$ and $`CPT_w`$-transformation can be defined considering both outer and inner automorphisms of the proper Poincaré group . Namely, $`CPT_w=I_xI_z`$, where $`I_z`$ is defined as
$$(X,Z)\stackrel{I_z}{}(X,Z(i\sigma _2))$$
(89)
and is a composition of the inner automorphism $`(X,Z)(\overline{X}^T,(Z^T)^1)`$ and of the rotation by the angle $`\pi `$. Wigner time reversal is the composition of above considered transformations, $`T_w=I_zCT=I_zT_{sch}`$.
The improper Poincaré group is defined as a group, which includes continuous transformations of the proper Poincaré group $`gM(D,1)`$ and the space reflection $`P`$.
In the Euclidean case the space reflection is reduced to the substitution $`(X,Z)\stackrel{P}{}(X,Z)`$. The charge conjugation inverts the momentum and spin orientation.
### F Equivalent representations
In the decomposition of scalar field (69) on the Poincaré group (or, that is the same, of the left GRR) there are equivalent representations distinguished by the right generators.
Remember that representations $`T_1(g)`$ and $`T_2(g)`$ acting in linear spaces $`L_1`$ and $`L_2`$ respectively are equivalent if there exists an invertible linear operator $`A:L_1L_2`$ such that
$$AT_1(g)=T_2(g)A.$$
(90)
In particular, the left and the right GRR of a Lee group G are equivalent. The operator $`(Af)(g)=f(g^1)`$ realizes the equivalence .
Let us consider functions $`f(x,𝐳)`$ belonging to two equivalent representations in the decomposition of the left GRR of the group $`M(D,1)`$ (or $`M(D)`$). If the representations $`T_1(g)`$ and $`T_2(g)`$ acting in the different subspaces $`L_1`$ and $`L_2`$ of the space of functions on the group are equivalent, then
$$AT_1(g)f_1(x,𝐳)=T_2(g)Af_1(x,𝐳),f_2(x,𝐳)=Af_1(x,𝐳),$$
where $`f_1(x,𝐳)L_1`$ and $`f_2(x,𝐳)L_2`$. In particular, if operator $`A:L_1L_2`$ is a function of the right translations generators $`\widehat{p}_\mu ^R`$, then one can’t map the function $`f_1(x,𝐳)`$ to the function $`f_2(x,𝐳)`$ by the group transformation, which leaves the interval square invariant. Therefore the physical equivalence of the states, that correspond to equivalent irreps in the decomposition of the scalar field $`f(x,𝐳)`$, is not evident at least.
Below we will consider a number of examples in various dimensions. In particular, in the framework of the representation theory of three-dimensional Euclidean group $`M(3)`$ irreps characterized by different spins (but with the same spin projection on the direction of propagation) are equivalent. There are no contradictions in the fact that in this case different particles are described by equivalent irreps since it is not possible to map corresponding wave functions one into another by the rotations or translations of the frame of references.
In some cases more general consideration may be based on the representation theory of an extended group. In the framework of the latter there are two possibilities: either irreps labelled by different eigenvalues of right generators of initial group are nonequivalent or some equivalent irreps of initial group are combined into one irrep. For example, in nonrelativistic theory spin becomes the characteristic of nonequivalent irreps after the extantion of $`M(3)`$ up to Galilei group. In 3+1 dimensions for $`m>0`$ the proper Poincaré group representations characterized by different chiralities are equivalent. If we going from the Lorentz group to the group $`SO(3,2)`$, then all characterized by spin $`s`$ states with different chiralities $`\lambda `$, $`\lambda =s,s+1,\mathrm{},s`$ are combined into one irrep.
The space of functions $`f(x,𝐳)`$ contains functions transforming under equivalent representations of the proper Poincaré group and is sufficiently wide to define discrete transformations, including space reflection, time reflection, and charge conjugation. These discrete transformations associated with automorphisms of the group also combine equivalent irreps of proper Poincaré group into one representation of the extended group. For example, in 3+1 dimensions space reflection combines two equivalent irreps of the proper group labelled by $`\lambda `$ and $`\lambda `$ into one irrep of the improper group.
Besides, as we will see below, the different types of RWE (finite-component and infinite-component equations) are also associated with equivalent representations in the decomposition of the left GRR.
Thus initially it is appropriate to consider all representations in the decomposition of the scalar field on the Poincaré group, including equivalent ones. In this sense we note the close analogy with the theory of nonrelativistic three-dimensional rotator . In the latter theory one considers functions on the rotation group $`SU(2)`$ and two sets of operators: angular momentum operators in an inertial laboratory (space-fixed) frame (left generators $`\widehat{J}_i^L`$) and angular momentum operators in a rotating (body-fixed) frame (right generators $`\widehat{J}_i^R`$). The classification of the rotator states is based on the use of the complete set of commuting operators which, apart from $`\widehat{𝐉}^2`$ and $`\widehat{J}_3^L`$, includes also $`\widehat{J}_3^R`$. Operator $`\widehat{J}_3^R`$ distinguishes equivalent representations in the decomposition of the left GGR of the rotation group and corresponds to the quantum number which does not depend on the choice of the laboratory frame. This quantum number plays a significant role in the theory of molecular spectra. In 3+1-dimensional case there exist two analogs of $`\widehat{J}_3^R`$, namely $`\widehat{B}_3^R=\widehat{S}_{03}^R`$ and $`\widehat{S}_3^R=\widehat{S}_{12}^R`$, which act in the space of functions on the Poincaré group. As we will see below, the first may be interpreted as a chirality operator, and the second allows to distinguish particles and antiparticles.
### G Quasiregular representations and spin description
The consideration of GRR of the Poincaré group ensures the possibility of consistent description of particles with arbitrary spin by means of scalar functions on $`^d\times \mathrm{Spin}(D,1)`$. At the same time, for description of spinning particles it is possible to use the spaces $`^d\times M`$, where $`M`$ is some homogeneous space of the Lorentz group (one or two-sheeted hyperboloid, cone, complex disk, projective space and so on); see, for example, and for 3+1 and 2+1-dimensional cases respectively. In some papers fields on homogeneous spaces are considered; in other papers such spaces are treated as phase spaces of some classic mechanics, and the latter are treated as models of spinning relativistic particles.
These spaces appear in the framework of the next group-theoretical scheme. Let us consider the left quasiregular representation of the Poincaré group
$$T(g)f(g_0H)=f(g^1g_0H),H\mathrm{Spin}(D,1).$$
(91)
$`H`$ is a subgroup of $`\mathrm{Spin}(D,1)`$, and since $`x`$ is invariant under right rotations (see (65))
$$g_0(X,Z),g_0H(X,ZH).$$
Therefore the relation (91) defines the representation of the Poincaré group in the space of functions $`f(x,zH)`$ on
$$^d\times (\mathrm{Spin}(D,1)/H).$$
(92)
In the decomposition of the representation in the space of functions on $`\mathrm{Spin}(D,1)/H`$ (or $`^d\times (\mathrm{Spin}(D,1)/H)`$) there is, generally speaking, only part of irreps of the Lorentz (or Poincaré) group. In particular, the case $`H\mathrm{Spin}(D,1)`$ corresponds to scalar field on Minkowski space. The classification and description of homogeneous spaces of 3+1 Poincaré and Lorentz groups one can find in .
Thus the consideration of quasiregular representations allows one to construct a number of spin models classified by subgroups $`H\mathrm{Spin}(D,1)`$. But the existence of nontrivial subgroup $`H`$ leads to rejection of the part of equivalent (with different characteristics with respect to the Lorentz subgroup) or, possibly, nonequivalent irreps of the Poincaré group.
### H Relativistic wave equations
The problem of RWE construction for particles with arbitrary spin in various dimensions is far from its completion and continues to attract significant attention. To describe massive particles of spin $`s`$ in four dimensions one usually employs the equations connected with the representations $`(\frac{s}{2}\frac{s}{2})`$ and $`(\frac{2s\pm 1}{4}\frac{2s1}{4})`$ of the Lorentz group (see, for example, ). These equations admit Lagrangian formulations , but for $`s>1`$ minimal electromagnetic coupling leads to a noncasual propagation . On the other hand, all known equations with casual solutions either have a redundant number of independent components (as equations for representations $`(s\mathrm{\hspace{0.17em}0})`$ and $`(0s)`$ have) or describe many masses and spins simultaneously, as Bhabha equations do. Besides the problem of interaction of higher spin fields, one may mention attempts to construct RWE with a completely positive energy spectrum and RWE for fractional spin fields .
With respect to mathematical methods used, it is possible divide all approaches to RWE construction in three groups.
The first approach, which follows Refs. , deals with equations for symmetric spin tensors. It allows one to describe fields with fixed mass and spin and also to construct RWE which admit Lagrangian formulation; however, as was mentioned above, for $`s>1`$ we face the problem of noncasual propagation.
The second approach, which follows Refs. , is devoted to studying RWE of the form $`(\alpha ^\mu \widehat{p}_\mu \varkappa )\psi (x)=0`$, and is based on the use of algebraic properties of $`\alpha `$-matrices. These equations admit Lagrangian formulation, however, for $`s>1`$ they describe a nonphysical spectrum of particles: a decreasing mass with increasing spin.
The third approach is connected to the use of some supplementary variables to describe spin degrees of freedom and initially was suggested for RWE with a mass spectrum (see ). It was used for constructing positive energy wave equations , equations describing gauge fields , and for anyon equations .
From the point of view of the approach which we developed above, the problem of constructing RWE looks like a selection of invariant subspaces in the space of functions on the group.
The classification of the scalar functions can be based on the use of the operators $`\widehat{C}_k`$ commuting with $`T_L(g)`$ (and correspondingly with all the left generators). For these operators, as a consequence of a relation $`\widehat{C}f(x,𝐳)=cf(x,𝐳)`$, one can obtain that $`\widehat{C}f^{}(x,𝐳)=cf^{}(x,𝐳)`$, where $`f^{}(x,𝐳)=T_L(g)f(x,𝐳)`$. Therefore, different eigenvalues $`c`$ correspond to subspaces, which are invariant with respect to action of $`T_L(g)`$. The invariant subspaces correspond to subrepresentations of the left GRR.
In addition to the Casimir operators, for the classification one may use the right generators since all the right generators commute with all the left generators. The right generators, as was mentioned, distinguish equivalent representations in the decomposition of the left GRR.
There is some freedom to choose the commuting operators which are functions of the right generators of the Poincaré group. We will use only functions of the generators of the right rotations (67), in particular, for the coordination with standard formulation of the theory.
Following the general scheme of harmonic analysis, for $`M(D,1)`$ one may consider the system consisting of $`d`$ equations
$$\widehat{C}_kf(x,𝐳)=c_kf(x,𝐳),$$
(93)
where $`\widehat{C}_k`$ are the Casimir operators of the Poincaré group and of the spin Lorentz subgroup. These operators constitute a subset of the complete set of commuting operators on the Poincare group. Just the system we will use for $`d=2+1`$ below.
On the other hand, there exist some additional requirements associated with the physical interpretation. In the first place, in massive case the system must be invariant under space reflection in order to describe states with definite parity. Secondly, it is often supposed that the system contains an equation of the first order in $`/t`$ (approach based on the first order equations advocated mainly in ). As a consequence of relativistic invariance, a linear in $`/t`$ equation can be either first order or infinite order in space derivatives (square-root Klein-Gordon equation ). The latter type of equations are naturally obtained in the theory of Markov processes for probability amplitudes .
The Casimir operators of the Poincaré group are the functions of the generators $`\widehat{p}_\mu `$ and $`\widehat{J}_{\mu \nu }`$. In odd dimensions there exists linear in $`\widehat{p}_\mu `$ Casimir operator, since the invariant tensor $`\epsilon ^{\mu \mathrm{}\nu }`$ has also odd number of indices. As we will see below, in 2+1 dimensions the system (93) is invariant under space reflections.
In even dimensions the invariant tensor $`\epsilon ^{\mu \mathrm{}\nu }`$ also has even number of indices, and therefore linear in $`\widehat{p}_\mu `$ Casimir operator does not exist. Besides, in even dimensions under space reflection irrep of proper Poincaré group is mapped onto equivalent irrep labelled by another eigenvalues of the Casimir operators of spin Lorentz subgroup. The linear combinations of basis elements of these two irreps form the bases of two labelled by intrinsic parity $`\eta =\pm 1`$ irreps of improper Poincaré group including space reflection.
Nevertheless, in even dimensions there exists operator $`\widehat{C}^{}=\widehat{p}_\mu \widehat{\mathrm{\Gamma }}^\mu `$, where $`\widehat{\mathrm{\Gamma }}^\mu =\widehat{\mathrm{\Gamma }}^\mu (𝐳,/𝐳)`$, commuting with all left generators and connecting the states which are interchanged under space reflections. In contrast to the Casimir operators this operator is not a function of generators of Poincaré group and does not commute with some right generators. A first order equation
$$\widehat{p}_\mu \widehat{\mathrm{\Gamma }}^\mu f(x,𝐳)=\varkappa f(x,𝐳)$$
(94)
interlocks, at least, two irreps of the group $`M(D,1)`$ characterized by different eigenvalues of the Casimir operator of spin Lorentz subgroup. Equations (93) and (94) have the same form; namely, invariant operator acts on the scalar function $`f(x,𝐳)`$ on the group $`M(D,1)`$. The addition of the operators $`\widehat{\mathrm{\Gamma }}^\mu `$ means in fact the extension of the Lorentz group up to more wide group (in particular, in four dimensions to the 3+2 de Sitter group $`SO(3,2)`$). Equation (94) replaces the equations of the system (93), which are not invariant under space reflection.
In the approach under consideration equations have the same form for all spins. The separation of the components with fixed spin and mass is realized by fixing eigenvalues of the Casimir operators of the Poincaré group (or the operator $`\widehat{p}_\mu \widehat{\mathrm{\Gamma }}^\mu `$). Fixing the representation of the Lorentz group, under which $`\varphi (𝐳)`$ transforms in the decomposition
$$f(x,𝐳)=\varphi ^n(𝐳)\psi _n(x),$$
one can obtain RWE in standard multicomponent form. This fixation is realized by the Casimir operator of spin Lorentz subgroup.
There are two types of equations to describe one and the same spin, one on functions $`f(x,𝐳)`$, where $`\varphi ^n(𝐳)`$ transforms under finite-dimensional nonunitary irrep of the Lorentz group, and another on functions $`f(x,𝐳)`$, where $`\varphi ^n(𝐳)`$ transforms under infinite-dimensional unitary irrep of the Lorentz group. In matrix representation these equations are written in the form of finite-component or infinite-component equations respectively. The latter type of equations (for example, Majorana equations ) is interesting because it gives the possibility to combine the relativistic invariance with probability interpretation. Desirability of this combination was emphasized in .
Let us briefly consider the possibility of existence of particles with fractional spin. The restrictions on the spin value arise in the representation theory of $`M(D)`$ and $`M(D,1)`$ if one restricts the consideration by (1) unitary, (2) finite-dimensional (with respect to the number of spin components) or (3) single-valued representations. (The latter means that the representation acts in the space of single-valued functions.) The restriction by single-valued functions (often supposed in mathematical papers related to a classification of representations) is omitted in some physical problems that allows to consider particles with fractional spin (anyons). Thus, we will also consider multi-valued representations of $`M(D)`$ and $`M(D,1)`$ in the space of the functions $`f(x,𝐳)`$ on the group. These representations correspond to single-valued representations of the universal covering group.
## III Two-dimensional case
### A Field on the group $`M(2)`$
In two-dimensional case the general formulas become simpler. Matrices $`U`$ (37) of $`SO(2)`$ subgroup depend on only one parameter, namely an angle $`\varphi `$, $`0\varphi 4\pi `$. Using the correspondence $`g_0(X,Z(\theta /2))`$, $`g(A,U(\varphi /2))`$ one may write the action of GRR:
$`T_L(g)f(x,\theta /2)=f(x^{},\theta /2\varphi /2),`$ (95)
$`x_1^{}=(x_1a_1)\mathrm{cos}\varphi +(x_2a_2)\mathrm{sin}\varphi ,x_2^{}=(x_2a_2)\mathrm{cos}\varphi (x_1a_1)\mathrm{sin}\varphi ,`$ (96)
$`T_R(g)f(x,\theta /2)=f(x^{\prime \prime },\theta /2+\varphi /2),`$ (97)
$`x_1^{\prime \prime }=x_1+a_1\mathrm{cos}\theta a_2\mathrm{sin}\theta ,x_2^{\prime \prime }=x_2+a_2\mathrm{cos}\theta +a_1\mathrm{sin}\theta .`$ (98)
Left and right generators, which correspond to parameters $`\theta `$ and $`\varphi `$, are given by
$`\widehat{p}_k=i_k,\widehat{J}=\widehat{L}+\widehat{S},`$ (99)
$`\widehat{p}_k^R=i\mathrm{\Lambda }_k^i_i,\widehat{J}^R=\widehat{S},`$ (100)
where
$$\widehat{L}=i(x_1_2x_2_1)=i\frac{}{\phi },\widehat{S}=i\frac{}{\theta },\mathrm{\Lambda }=\left(\begin{array}{cc}\mathrm{cos}\theta & \mathrm{sin}\theta \\ \mathrm{sin}\theta & \mathrm{cos}\theta \end{array}\right).$$
The functions on the group are ones on $`^2\times S^1`$, and invariant measure on the group is
$$d\mu (x,\theta )=(4\pi )^1dx_1dx_2d\theta ,\mathrm{}<x<+\mathrm{},0\theta <4\pi .$$
We will consider two complete sets of commuting operators: $`\widehat{p}_1`$, $`\widehat{p}_2`$, $`\widehat{S}`$ and $`\widehat{p}^2`$, $`\widehat{J}`$, $`\widehat{S}`$. The eigenfunctions of these operators are
$`x_1x_2\theta |p_1p_2s=(2\pi )^1\mathrm{exp}(ip_1x_1+ip_2x_2+is\theta ),`$ (101)
$`r\phi \theta |pjs=(2\pi )^{1/2}i^lJ_l(pr)\mathrm{exp}(il\phi )\mathrm{exp}(is\theta ),`$ (102)
where $`l=js`$ is orbital momentum, $`J_l(pr)`$ is the Bessel function. Irreps are labelled by eigenvalues $`p^2`$ of the Casimir operator $`\widehat{p}^2`$. For $`p0`$ the representation is irreducible, for $`p=0`$ it decomposes into one-dimensional irreps of spin subgroup $`U(1)`$, which are labelled by eigenvalues $`s`$ of the spin projection operator (or, simply speaking, spin operator) $`\widehat{S}`$.
At $`p0`$ the representations characterized by the spin $`s`$ and $`s^{}=s+n`$, where $`n`$ is integer number, are equivalent. Really, operator $`\widehat{S}`$ commutes with all left generators, but does not commute with the generators of right translations, which mix spin and space coordinates. Operators $`\widehat{p}_+^R=p_1^Rip_2^R`$ and $`\widehat{p}_{}^R=p_1^R+ip_2^R`$ are raising and lowering operators with respect to spin $`s`$
$$\widehat{p}_\pm ^R|p_1p_2s=(ip_1\pm p_2)|p_1p_2s\pm 1.$$
(103)
Right translations do not conserve both interval (distance) and spin $`s`$.
The functions (102) satisfy the relations of orthogonality and completeness
$`{\displaystyle pjs|r\phi \theta r\phi \theta |pjsr𝑑r𝑑\phi 𝑑\theta }={\displaystyle \frac{\delta (pp^{})}{p}}\delta _{jj^{}}\delta _{ss^{}},`$ (104)
$`{\displaystyle \underset{l,s}{}r\phi \theta |pjspjs|r\phi \theta dp}={\displaystyle \frac{\delta (rr^{})}{r}}\delta (\phi \phi ^{})\delta (\theta \theta ^{}).`$ (105)
It means that we have obtained the decomposition of left regular representation. Spin operator $`\widehat{S}`$ distinguishes equivalent irreps (except the case $`p=0`$, when irreps are labelled by its eigenvalues). The decomposition of the functions of $`\theta `$ on the eigenfunctions of $`\widehat{S}`$ corresponds to the Fourier series expansion of functions on a circle.
Thus the representations of $`M(2)`$ are single-valued for integer and half-integer $`s`$. The fractional values of $`s`$ correspond to multi-valued representations. Irreps are equivalent if are labelled by the same $`p0`$ and the difference $`ss^{}=n`$ is an integer number. For fixed $`p0`$ there are only two nonequivalent single-valued representations, which correspond to integer and half-integer spin. Nonequivalent multi-valued representations for fixed $`p0`$ are labelled by $`\stackrel{~}{s}[0,1)`$, $`\stackrel{~}{s}=s[s]`$.
### B Field on the group $`M(1,1)`$
Matrices $`U`$ (34) of $`SO(1,1)`$ subgroup, which is isomorphic to an additive group of real numbers, depend on a hyperbolic angle $`\varphi `$. Using the correspondence $`g_0(X,Z(\theta /2))`$, $`g(A,U(\varphi /2))`$, one may write the action of GRR:
$`T_L(g)f(x,\theta /2)=f(x^{},\theta /2\varphi /2),`$ (106)
$`x^0=(x^0a^0)\mathrm{cosh}\varphi +(x^1a^1)\mathrm{sinh}\varphi ,x^1=(x^1a^1)\mathrm{cosh}\varphi +(x^0a^0)\mathrm{sinh}\varphi ,`$ (107)
$`T_R(g)f(x,\theta /2)=f(x^{\prime \prime },\theta /2+\varphi /2),`$ (108)
$`x^{\prime \prime 0}=x^0+a^0\mathrm{cosh}\theta a^1\mathrm{sinh}\theta ,x^{\prime \prime 1}=x^1+a^1\mathrm{cosh}\theta a^0\mathrm{sinh}\theta .`$ (109)
The functions on the group are ones on $`^2\times `$, and invariant measure on the group can be written as
$$d\mu (x,\theta )=dx^0dx^1d\theta ,\mathrm{}<x,\theta <+\mathrm{}.$$
As above, we will consider two complete sets of commuting operators, $`\widehat{p}_1`$, $`\widehat{p}_2`$, $`\widehat{S}`$ and $`\widehat{p}^2`$, $`\widehat{J}`$, $`\widehat{S}`$, where $`\widehat{J}=\widehat{L}+\widehat{S}`$, $`\widehat{L}=i(x^0^0+x^1^1)`$, $`\widehat{S}=i/\theta `$. The eigenfunctions of the first set are
$$x^0x^1\theta |p_1p_2\lambda =(2\pi )^{3/2}\mathrm{exp}(ip_\mu x^\mu +i\lambda \theta ),$$
(110)
where $`\lambda `$ is an eigenvalue of the spin projection (chirality) operator $`\widehat{S}`$. The form of eigenfunctions of the second set depends on the type of irrep. There are four types of unitary irreps labelled by eigenvalue $`m^2`$ of operator $`\widehat{p}^2`$ .
1. $`m^2>0`$. Representations correspond to the particles of nonzero mass, the eigenfunctions of operators $`\widehat{p}^2`$, $`\widehat{J}`$, $`\widehat{S}`$ are
$$r\phi \theta |mj\lambda =(4\pi )^1i\mathrm{exp}(\pi l/2)H_{il}^{(2)}(\pm mr)\mathrm{exp}(il\phi )\mathrm{exp}(i\lambda \theta ),$$
(111)
where $`H_{il}^{(2)}(mr)`$ is Hankel function, $`r^2=(x^0)^2(x^1)^2`$, and $`\pm `$ corresponds to the sign of energy $`p_0`$.
2. $`m^2<0`$. Representations correspond to tachyons, which in $`d=1+1`$ are more similar to usual particles because of symmetry between space and time variables. The form of $`r\phi \theta |mj\lambda `$ coincides with (111), but $`m`$ is imaginary.
3. $`m=0`$, $`p_1=\pm p_0`$. Representations correspond to the massless particles. According to (64), for the action of finite transformations $`T_0(g)`$ on the functions $`f(p,\pm p,\theta /2)`$ one may obtain
$$T_0(g)f(p,\pm p,\theta /2)=e^{iap^{}}f(p^{},\pm p^{},\theta /2\varphi /2),p^{}=e^\varphi p.$$
Therefore the representation $`T_0(g)`$ is reducible and splits into four irreps differed by the signs of $`p_0`$ and $`p_1=\pm p_0`$, and reducible representation, which corresponds to $`m=p_0=0`$.
4. $`m=p_0=0`$. This representation decomposes into sum of one-dimensional irreps of the group $`SO(1,1)`$, which are labelled by eigenvalues of $`\widehat{S}`$.
There are no integer value restrictions for the spectrum of $`\widehat{S}`$, and chirality can be fractional, $`\mathrm{}<\lambda <+\mathrm{}`$. The decomposition of the functions $`f(x,\theta )`$ in therms of the eigenfunctions of $`\widehat{S}`$ corresponds to the Fourier integral expansion of functions on a line. The equivalence of the representations characterized by different $`\lambda `$ is related to the fact that, like in Euclidean case, operator $`\widehat{S}`$ does not commute with right translations.
One can convert vector indices into spinor indises and vice versa with the help of the formula (20). In the case under consideration matrices $`U`$ are real and symmeric, $`X^{}=UXU`$, or in component-wise form $`x^\nu \sigma _{\nu \alpha \alpha ^{}}=U_\alpha ^\beta \sigma _{\mu \beta \beta ^{}}x^\mu U_\alpha ^{}^\beta ^{}`$, and there exists one type of spinor indeces only. Denoting elements of the first column of matrix $`Z`$ transforming under spinor representation of $`SO(1,1)`$ by $`z_\alpha `$, $`z_1=\mathrm{cosh}(\theta /2)`$, $`z_2=\mathrm{sinh}(\theta /2)`$, we obtain for components of vector and antisymmetric tensor
$$q^\mu =\sigma ^{\mu \alpha \beta }z_\alpha z_\beta ,q^0=\mathrm{cosh}\theta ,q^1=\mathrm{sinh}\theta ,q^{01}=\sigma ^{01\alpha \beta }z_\alpha z_\beta =i.$$
(112)
There exist two invariant tensors $`\eta ^{\mu \nu }`$ and $`\epsilon ^{\mu \nu }`$, which can be used for raising of indices. This is related to the fact that vectors $`(x^0x^1)`$ and $`(x^1x^0)`$ have the same transformation rule, and one can construct invariant from two vectors by two different ways: $`\eta ^{\mu \nu }q_\mu q_\nu ^{}=\mathrm{cosh}(\theta \theta ^{})`$, $`\epsilon ^{\mu \nu }q_\mu q_\nu ^{}=\mathrm{sinh}(\theta \theta ^{})`$.
### C Relativistic wave equations in 1+1 dimensions
An irrep of the group $`M(1,1)`$ can be extract from GRR by fixing the sign of $`p_0`$ and eigenvalues of operators $`\widehat{p}^2`$, $`\widehat{S}`$,
$`\widehat{p}^2f(x,\theta )=m^2f(x,\theta ),`$ (113)
$`\widehat{S}f(x,\theta )=\lambda f(x,\theta ),`$ (114)
where chirality $`\lambda `$ distinguishes equivalent irreps labelled by identical eigenvalues $`m^2`$ of the Casimir operator $`\widehat{p}^2`$. Solutions of this system have the form $`f(x,\theta )=\psi (x)e^{i\lambda \theta }`$, where $`\widehat{p}^2\psi (x)=m^2\psi (x)`$.
According to (86), space reflection converts $`e^{i\lambda \theta }`$ to $`e^{i\lambda \theta }`$. Irreps of the improper Poincaré group are labelled by mass $`m`$, $`signp_0`$, intrinsic parity $`\eta =\pm 1`$, and spin $`s=|\lambda |`$ (as above, $`s`$ distinguishes equivalent irreps). In the rest frame it is easy to write down functions with mentioned characteristics:
$$e^{\pm imx^0}(e^{i\lambda \theta }\pm e^{i\lambda \theta }).$$
(115)
States with arbitrary momentum can be obtained from (115) by hyperbolic rotations and form the basis of unitary irrep of improper group. On the other hand, the problem arise to construct equations that unlike the system (113)-(114) are invariant under improper Poincaré group and have solutions with definite parity. These equations should combine states with chiralities $`\pm \lambda `$.
The general form of the linear in $`\widehat{p}^\mu `$ equations is
$$\widehat{p}_\mu \widehat{\mathrm{\Gamma }}^\mu f(x,\theta )=\varkappa f(x,\theta ),$$
(116)
where $`\widehat{\mathrm{\Gamma }}^\mu =\widehat{\mathrm{\Gamma }}^\mu (\theta ,/\theta )`$. For invariance of (116) under space reflection $`P`$ and hyperbolic rotations the operator $`\widehat{p}_\mu \widehat{\mathrm{\Gamma }}^\mu `$ must commute with $`P`$ and $`\widehat{J}=\widehat{L}+\widehat{S}`$, whence
$$\widehat{\mathrm{\Gamma }}^\mu \stackrel{P}{}(1)^{\delta _{1\mu }}\widehat{\mathrm{\Gamma }}^\mu ,[\widehat{\mathrm{\Gamma }}^0,\widehat{S}]=i\widehat{\mathrm{\Gamma }}^1,[\widehat{\mathrm{\Gamma }}^1,\widehat{S}]=i\widehat{\mathrm{\Gamma }}^0.$$
(117)
The operators
$$\widehat{\mathrm{\Gamma }}^0=s\mathrm{cosh}\theta \mathrm{sinh}\theta \frac{}{\theta },\widehat{\mathrm{\Gamma }}^1=s\mathrm{sinh}\theta \mathrm{cosh}\theta \frac{}{\theta },[\widehat{\mathrm{\Gamma }}^0,\widehat{\mathrm{\Gamma }}^1]=i\widehat{S}$$
(118)
obey these relations. One may construct the operators, which raise and lower chirality $`\lambda `$ by 1,
$$\widehat{\mathrm{\Gamma }}_+=\widehat{\mathrm{\Gamma }}^0\widehat{\mathrm{\Gamma }}^1=e^\theta (s+/\theta ),\widehat{\mathrm{\Gamma }}_{}=\widehat{\mathrm{\Gamma }}^0+\widehat{\mathrm{\Gamma }}^1=e^\theta (s/\theta ).$$
(119)
Operators $`\widehat{\mathrm{\Gamma }}^0`$, $`\widehat{\mathrm{\Gamma }}^1`$, and $`\widehat{\mathrm{\Gamma }}^2=i\widehat{S}=/\theta `$ obey the commutation relations of the generators of the $`SO(2,1)SU(1,1)`$ group:
$$[\widehat{\mathrm{\Gamma }}^a,\widehat{\mathrm{\Gamma }}^b]=ϵ^{abc}\widehat{\mathrm{\Gamma }}_c,\widehat{\mathrm{\Gamma }}_a=\eta _{ab}\widehat{\mathrm{\Gamma }}^b,\eta _{00}=\eta _{22}=\eta _{11}=1,\widehat{\mathrm{\Gamma }}_a\widehat{\mathrm{\Gamma }}^a=s(s+1).$$
Thus, if the symmetry with respect to the space reflection takes place, the condition of mass irreducibility (113) can be supplemented by equation (116) instead of (114). This means the passage to the new set of commuting operators, namely from $`\widehat{p}_\mu `$, $`\widehat{S}`$ to $`\widehat{p}_\mu `$, $`\widehat{p}_\mu \widehat{\mathrm{\Gamma }}^\mu `$. Let us consider the system
$`\widehat{p}^2f(x,\theta )=m^2f(x,\theta ),`$ (120)
$`\widehat{p}_\mu \widehat{\mathrm{\Gamma }}^\mu f(x,\theta )=msf(x,\theta ).`$ (121)
The operator $`\widehat{S}`$ does not commute with $`\widehat{p}_\mu \widehat{\mathrm{\Gamma }}^\mu `$, and the particle with nonzero mass described by equation (121) can’t be characterized by certain chirality. In the rest frame $`p_0=\pm m`$, and the functions $`f(x,\theta )=e^{\pm imx^0}\varphi (\theta )`$ should be eigenfunctions of operator $`\widehat{\mathrm{\Gamma }}^0`$ with eigenvalues $`\pm s`$. The equation
$$\widehat{\mathrm{\Gamma }}^0\varphi (\theta )=(s\mathrm{cosh}\theta (\mathrm{sinh}\theta )/\theta )\varphi (\theta )=\varkappa \varphi (\theta )$$
for $`\varkappa =\pm s`$ has solutions $`[\mathrm{cosh}(\theta /2)]^{2s}`$ and $`[\mathrm{sinh}(\theta /2)]^{2s}`$ respectively. Below we will consider two cases.
1. The solutions of the system (120)-(121) are sought in the space of polynomials of $`e^{\theta /2}`$ and $`e^{\theta /2}`$ that correspond to finite-dimensional nonunitary representations of $`SU(1,1)`$. Corresponding representations of $`SO(1,1)`$ subgroup are also nonunitary. For these representations generator $`\widehat{S}`$ is anti-Hermitian, and it is convenient to redefine chirality operator as $`i\widehat{S}`$. In the rest frame a general solution of the system (120)-(121) is
$$f(x,\theta )=C_1e^{imx^0}[\mathrm{cosh}(\theta /2)]^{2s}+C_2e^{imx^0}[\mathrm{sinh}(\theta /2)]^{2s},$$
(122)
where $`2s`$ is integer positive number. Therefore for an unique spin $`s`$ there are only two independent components (with positive and negative frequency). The space inversion takes $`\theta `$ to $`\theta `$, and in the rest frame solutions with different sign of $`p_0`$ and half-integer $`s`$ are characterized by opposite parity $`\eta `$. For integer $`s`$ all solutions are characterized by $`\eta =1`$. Plane wave solutions, which correspond to moving particle, can be obtained from (122) by a hyperbolic rotation by the angle $`2\varphi `$:
$$f_{m,s}(x,\theta )=C_1e^{ik_0x^0+ik_1x^1}\left(\mathrm{cosh}[(\theta +\varphi )/2]\right)^{2s}+C_2e^{ik_0x^0ik_1x^1}\left(\mathrm{sinh}[(\theta +\varphi )/2]\right)^{2s},$$
where $`k_0=m\mathrm{cosh}2\varphi `$, $`k_1=m\mathrm{sinh}2\varphi `$.
In the ultrarelativistic limit $`\varphi \pm \mathrm{}`$ we have two states with chirality $`\lambda =\pm s`$ respectively. Thus, if in the rest frame one may distinguish two components with positive and negative frequency, then in massless limit one may distinguish two components with positive and negative chirality.
Matrix form of the system (120)-(121) can be obtained by the decomposition of $`f(x,\theta )`$ over the basis $`e^{\lambda \theta /2}`$, $`\lambda =s,s+1,\mathrm{},s`$. There are $`2s+1`$ components $`\psi (x)`$ in this form, but only two of them are independent. Notice that representations of $`SO(1,1)`$ of the form $`e^{\lambda \theta }`$, are nonunitary for real $`\lambda `$ and integral over $`\theta `$ is divergent. One can redifine the norm of a state with the help of scalar product in the space of multicomponent functions $`\psi (x)`$, but this product is not positive definite.
For $`s=1/2`$, substituting the function $`f(x,\theta )=\psi _1(x)e^{\theta /2}+\psi _2(x)e^{\theta /2}`$ into equation (121), we obtain two-dimensional Dirac equation
$$\widehat{p}_\mu \gamma ^\mu \mathrm{\Psi }(x)=m\mathrm{\Psi }(x),\gamma ^0=\sigma _1,\gamma ^1=i\sigma _2,2\widehat{S}=\gamma ^3=\sigma _3.$$
(123)
where $`\mathrm{\Psi }(x)=(\psi _1(x)\psi _2(x))^T`$. Matrix $`\gamma ^3=\gamma ^0\gamma ^1`$ corresponds to chirality operator and satisfies the condition $`[\gamma ^3,\gamma ^\mu ]_+=0`$. On the other hand, this matrix corresponds to hyperbolic rotation, and similarly to 3+1 case one can write $`\gamma ^\mu \gamma ^\nu =\eta ^{\mu \nu }i\sigma ^{\mu \nu }`$, where $`\sigma ^{01}=i\gamma ^3`$. Invariant scalar product has the form $`|\psi _1(x)|^2|\psi _2(x)|^2`$.
For $`s=1`$, substituting the function $`f(x,\theta )=\psi _{11}(x)e^\theta +\psi _{12}(x)+\psi _{22}(x)e^\theta `$ into equation (121), we obtain
$`(\widehat{p}_\mu \mathrm{\Gamma }^\mu m)\mathrm{\Psi }(x)=0,`$ (124)
$`\mathrm{\Gamma }^0={\displaystyle \frac{1}{\sqrt{2}}}\left(\begin{array}{ccc}0& 1& 0\\ 1& 0& 1\\ 0& 1& 0\end{array}\right),\mathrm{\Gamma }^1={\displaystyle \frac{1}{\sqrt{2}}}\left(\begin{array}{ccc}0& 1& 0\\ 1& 0& 1\\ 0& 1& 0\end{array}\right),\widehat{S}=\left(\begin{array}{ccc}1& 0& 0\\ 0& 0& 0\\ 0& 0& 1\end{array}\right),`$ (134)
where $`\mathrm{\Psi }(x)=(\psi _{11}(x)\psi _{12}(x)/\sqrt{2}\psi _{22}(x))^T`$. Using (112) to convert spinor indices to vector ones: $`_0=\psi _{22}\psi _{11}`$, $`_1=\psi _{22}+\psi _{11}`$, $`F_{01}=F_{10}=i\psi _{12}`$, we obtain $`p_0_1p_1_0=imF_{01}`$, $`ip_0F_{10}=m_1`$, $`ip_1F_{10}=m_0`$. Thus one can rewrite 1+1 ”Duffin–Kemmer” equation (134) in the form, which is similar to Proca equations for in 3+1 dimensions (see (293),(298)),
$$_\mu _\nu _\nu _\mu =mF_{\mu \nu },^\nu F_{\mu \nu }=m_\mu .$$
(135)
As a consequence of (135) we obtain $`_\mu ^\mu =0`$, $`(\widehat{p}^2m^2)^\mu =0`$. But 1+1-dimensional case is distinctly different from 3+1-dimensional one because the component $`F_{01}=F_{10}`$ is characterized by zero chirality and thus the roles of $`F_{\mu \nu }`$ and $`_\nu `$ are interchanged.
In the massless case the system (135) splits into two independent equations for the components $`_\mu `$ and $`F_{\mu \nu }`$ respectively,
$`_\mu _\nu _\nu _\mu =0,`$ (137)
$`^\mu F_{\mu \nu }=0.`$
First equation has propagating solutions
$$_0=C_1e^{ip(x^0+x^1)}+C_2e^{ip(x^0x^1)},_1=C_1e^{ip(x^0+x^1)}C_2e^{ip(x^0x^1)}$$
obeying transversality condition $`_\mu ^\mu =0`$. Second equation (free two-dimensional Maxwell equation ) corresponds to the components with zero chirality and has trivial solution $`F_{\mu \nu }=const`$ only. Notice that for real field $`f^{}(x,\theta )=f(x,\theta )`$ components $`_\mu `$ and $`F_{\mu \nu }`$ also are real, and propagating solutions do not exist for $`m=0`$.
If for $`s=1/2`$ and for $`s=1`$ the first equation of the system (120)-(121) is the consequence of the second equation, then for $`s>1`$ there are the solutions of equation (121) with mass spectrum, $`m_i|s_i|=ms`$, $`s_i=s,s1,\mathrm{},s`$. For the extraction of characterized by certain mass $`m`$ and spin $`s`$ representations of improper Poincaré group it is necessary to use both equations of the system.
Notice that the chirality $`\lambda `$ of a particle described by (113)-(114) can be fractional, but the spin $`s`$ of a particle described by (120)-(121) can be only integer or half-integer for $`m0`$ and finite number of components $`\psi (x)`$.
Really, if $`2s`$ is not integer, then acting by the raising operator on the state with label $`\lambda =s`$, we not get into the state labelled by $`\lambda =s`$ and connected with initial state by the space reflection; moreover, the spectrum of $`\lambda `$ is not bounded above. On the other hand, it is possible to develop an alternative approach (in particular, for the massive particles with fractional spin) based on the using of infinite-dimensional unitary irreps of $`SO(2,1)`$.
2. Let us consider now the solutions of (120)-(121) in the space of the square-integrable functions of $`\theta `$. In the rest frame, as we have seen above, there are two types of the solutions. The solutions $`[\mathrm{sinh}(\theta /2)]^{2s}`$ are not square-integrable for any $`s`$ since corresponding integral is divergent either at zero or at infinity. The solutions $`[\mathrm{cosh}(\theta /2)]^{2s}`$ for $`s<0`$ are square-integrable:
$$_{\mathrm{}}^+\mathrm{}[\mathrm{cosh}(\theta /2)]^{4s}𝑑\theta =2B(1/2,\mathrm{\hspace{0.17em}2}s).$$
Therefore in the space of square-integrable functions equation (121) has only positive energy solutions. Solutions with $`p_0<0`$ correspond to the equation $`\widehat{p}_\mu \widehat{\mathrm{\Gamma }}^\mu f(x,\theta )=msf(x,\theta )`$. Normalized positive energy solutions of the system (120)-(121) for the particle with spin $`|s|`$ and momenta $`p_0=m\mathrm{cosh}\varphi `$, $`p_1=m\mathrm{sinh}\varphi `$ are
$$f(x,\theta )=(2\pi )^1\left(2B(1/2,\mathrm{\hspace{0.17em}2}s)\right)^{1/2}e^{ip_0x^0+ip_1x^1}\left(\mathrm{cosh}[(\theta +\varphi )/2]\right)^{2|s|}.$$
(138)
In contrast to the case $`d>2`$ solutions with distinct $`s`$ are nonorthogonal. The decomposition of the solutions (138) over the functions $`e^{i\lambda \theta }`$ (i.e. over $`SO(1,1)`$ unitary irreps) corresponds to the Fourier integral expansion. We will consider properties of the positive energy equations more detail in 2+1-dimensional case below.
## IV Three-dimensional case
### A Field on the group $`M(3)`$
The case of $`M(3)`$ group is characterized by many-dimensional spin space. On the other hand, the constructions allow the simple physical interpretation.
Using the operators $`\widehat{J}^i=\widehat{L}^i+\widehat{S}^i=(1/2)ϵ^{ijk}\widehat{J}_{jk}`$, it is possible to rewrite the commutation relations (62) in the more compact form
$$[\widehat{p}_i,\widehat{p}_k]=0,[\widehat{p}^i,\widehat{J}^j]=iϵ^{ijk}\widehat{p}_k,[\widehat{J}^i,\widehat{J}^j]=iϵ^{ijk}\widehat{J}_k.$$
(139)
The invariant measure on the group is given by the formulas
$`d\mu (x,𝐳)=Cd^3x\delta (|z_1|^2+|z_2|^21)d^2z_1d^2z_2={\displaystyle \frac{1}{16\pi ^2}}d^3x\mathrm{sin}\theta d\theta d\varphi d\psi .`$ (140)
$`\mathrm{}<x<+\mathrm{},0<\theta <\pi ,0<\varphi <2\pi ,2\pi <\psi <2\pi ,`$ (141)
where $`z_1=\mathrm{cos}\frac{\theta }{2}e^{i(\psi \varphi )/2}`$, $`z_2=i\mathrm{sin}\frac{\theta }{2}e^{i(\psi +\varphi )/2}`$ are the elements of the first column of matrix (71), $`z^2=z_1`$, $`z^1=z_2`$, and $`\theta ,\varphi ,\psi `$ are the Euler angles. The spin projection operators acting in the space of the functions on the group $`f(x,𝐳)`$ have the form
$`\widehat{S}_k={\displaystyle \frac{1}{2}}(z\sigma _k_z\stackrel{}{}z\stackrel{}{}\sigma _k_\stackrel{}{z}),z=(z^1z^2),_z=(/z^1/z^2)^T,`$ (142)
$`\widehat{S}_k^R={\displaystyle \frac{1}{2}}(\chi \stackrel{}{}\sigma _k_\chi \stackrel{}{}\chi \sigma _k_\stackrel{}{\chi }),\chi =(z^1\stackrel{}{}z^{\dot{2}}),_\chi =(/z^1/\stackrel{}{}z^{\dot{2}})^T.`$ (143)
In terms of Euler angles one can obtain
$$\widehat{S}_3=i/\varphi ,\widehat{S}_3^R=i/\psi .$$
(144)
The operator $`\widehat{𝐩}^2`$ and the operator of the spin projection on the direction of propagation $`\widehat{W}=\widehat{𝐩}\widehat{𝐉}=\widehat{𝐩}\widehat{𝐒}`$ are the Casimir operators. The eigenvalues $`S(S+1)`$ of the Casimir operator of rotation subgroup in $`z`$-space $`\widehat{𝐒}^2=\widehat{𝐒}_R^2`$ define spin $`S`$. complete sets of the commuting operators $`\{\widehat{p}_k,\widehat{W},\widehat{𝐒}^2,\widehat{S}_R^3\}`$, $`\{\widehat{𝐩}^2,\widehat{W},\widehat{𝐉}^2,\widehat{S}_3,\widehat{𝐒}^2,\widehat{S}_3^R\}`$ consist of six operators (two Casimir operators, two left generators, and two right generators). The Casimir operator $`\widehat{W}`$ does not commute with $`\widehat{L}_k`$ and $`\widehat{S}_k`$ separately but only with the generators $`\widehat{J}_k=\widehat{L}_k+\widehat{S}_k`$, therefore there are sets, which do not include $`\widehat{W}`$, for example, $`\{\widehat{𝐩}^2,\widehat{p}_3,\widehat{L}_3,\widehat{S}_3,\widehat{𝐒}^2,\widehat{S}_3^R\}`$ and $`\{\widehat{p}_\mu ,\widehat{S}_3,\widehat{𝐒}^2,\widehat{S}_3^R\}`$.
We will consider the first set since in this case eigenfunctions have the most simple form. This set includes two Casimir operators, the operator of spin square $`\widehat{𝐒}^2`$ and the generator $`\widehat{S}_3^R`$. The latter two generators commute with all left generators but do not commute with right generators and label equivalent representations in the decomposition of the left GRR.
According to (144), the eigenfunctions of $`\widehat{S}_3^R`$, $`\widehat{S}_3^R|\mathrm{}n=n|\mathrm{}n`$, have the form $`|\mathrm{}n=F(x,\theta ,\varphi )\mathrm{exp}(in\psi )`$ and are differed only by a phase factor. As a consequence of the commutation relations of generators $`\widehat{S}_k^R`$ the operators $`\widehat{S}_\pm ^R=\widehat{S}_1^R\pm i\widehat{S}_2^R`$ are the raising and lowering operators for the eigenfunctions of $`\widehat{S}_3^R`$
$$\widehat{S}_\pm ^R|\mathrm{}n=C(S,n)|\mathrm{}n\pm 1.$$
(145)
The intertwining operators $`\widehat{S}_\pm ^R`$ consist of the generators of right rotations, which conserve the interval square according to (59). Moreover, the right rotations do not act on $`x`$. But there are no transformations (rotations and translations) of the reference frame, which connect the representations with different $`n`$. Notice that the states labelled by $`n`$ and $`n`$ are interchanged under charge conjugation (88).
The operator $`\widehat{𝐒}^2`$ also labels equivalent representations of $`M(3)`$ group. This operator commutes with all generators except right translations, and therefore an intertwining operator is a function of the latter generators. Right translations change both the interval and spin. Therefore it is naturally to characterize free particle in three-dimensional Euclidean space not only by momentum and spin projection on the direction of propagation, but also by spin $`S`$.
There are two standard realizations of the representation spaces, which correspond to eigenvalues $`n=\pm 2S`$ and $`n=0`$ of the operator $`\widehat{S}_3^R`$.
The first realization is the space of analytic ($`n=2S`$) functions $`f(x,z)`$ or antianalytic ($`n=2S`$) functions $`f(x,\stackrel{}{}z)`$ of two complex variables $`z^1,z^2`$, $`|z^1|^2+|z^2|^2=1`$, i.e. the space of functions of two-component spinors. In particular, according to (143), for the space of analytic functions we have
$$\widehat{S}_k=\frac{1}{2}z\sigma _k_z,$$
(146)
$`\widehat{S}_3^R=(z^1/z^1+z^2/z^2`$), and $`\widehat{𝐒}^2=\widehat{S}_3^R(\widehat{S}_3^R1)`$. The eigenfunctions of the operator of spin square are polynomials of the power $`2S`$ in $`z^1`$, $`z^2`$. The charge conjugation transformation connects equivalent irreps labelled by $`n=\pm 2S`$ and the spaces of analytic and antianalytic function. This transformation reverses the direction of momentum and spin.
The second realization is the space of functions, which do not depend on the angle $`\psi `$, and corresponds to $`n=0`$. It is the space of functions of five real variables on the manifold
$$^3\times S^2,d\mu =(4\pi )^1d^3x\mathrm{sin}\theta d\theta d\varphi .$$
The point in the spin space (i.e. on the sphere $`S^2P^1SU(2)/U(1)`$) can be define by the spherical coordinates $`\theta ,\varphi `$, or by two complex variables $`z_1=\mathrm{cos}\frac{\theta }{2}e^{i\varphi /2},z_2=\mathrm{sin}\frac{\theta }{2}e^{i\varphi /2}`$ (in this case one may use (146) for the spin projection operators), or by one complex number $`z=z_1/z_2`$ (this case corresponds to the realization in terms of projective space $`P^1`$). In terms of variables $`\theta ,\varphi `$ the eigenfunctions of operators $`\widehat{S}`$, $`\widehat{S}_3`$ are $`P_S^s(\mathrm{cos}\theta )e^{is\varphi }`$, where $`P_S^s(\mathrm{cos}\theta )`$ are associated Legendre functions .
Let us consider eigenfunctions of the set of the operators $`\{\widehat{p}_\mu ,\widehat{W},\widehat{𝐒}^2\}`$ in the space of analytic functions of $`z^1`$, $`z^2`$:
$$\widehat{p}_\mu f(x,z)=p_\mu f(x,z),\widehat{𝐒}^2f(x,z)=S(S+1)f(x,z),\widehat{𝐩}\widehat{𝐒}f(x,z)=psf(x,z).$$
(147)
The eigenfunctions of $`\widehat{𝐒}^2`$ are polynomials of the power $`2S`$ in $`z`$ (the unitary irreps of $`SU(2)`$ are finite-dimensional, therefore spin $`S`$ and spin projection on the direction of propagation $`s`$ are integer or half-integer). Let $`p_\mu =(0,0,p)`$, then the normalized solutions of the system (147) are
$$|\mathrm{0\hspace{0.17em}0}pSs=(2\pi )^{3/2}\left(\frac{(2S)!}{(S+s)!(Ss)!}\right)^{1/2}(z^1)^{S+s}(z^2)^{Ss}e^{ix_3p}.$$
The states with arbitrary direction of vector $`p`$ may be obtain by the rotation $`P=UP_0U^{}`$, $`Z=UZ_0`$, $`P_0=p\sigma _3`$, $`Z_0=(z_1z_2)^T`$,
$$|p_1p_2p_3Ss=(2\pi )^{3/2}\left(\frac{(2S)!}{(S+s)!(Ss)!}\right)^{1/2}(z^1\stackrel{}{}u_1+z^2\stackrel{}{}u_2)^{S+s}(z^1u_2+z^2u_1)^{Ss}e^{ipx},$$
(148)
where $`u_1`$, $`u_2`$ are the elements of the first line of matrix $`U`$. Notice that it is sufficient to use only two angles for the parametrization of matrix $`U`$ since the initial state has a stationary subgroup $`U(1)`$.
For the rest particle $`\widehat{𝐩}^2=\widehat{𝐩}\widehat{𝐒}=0`$ and only in this case $`M(3)`$ irreps labelled by different $`S`$ are nonequivalent.
In general case functions corresponding to the particle of spin $`S`$ have the form
$`f_S(x,z)={\displaystyle \underset{n=0}{\overset{2S}{}}}\varphi ^n(z)\psi _n(x),\varphi ^n(z)=(C_{2S}^n)^{1/2}(z^1)^{Sn}(z^2)^n,`$ (149)
$`{\displaystyle }\stackrel{}{}f_S(x,z)f_S^{}^{}(x,z)d\mu (x,z)=\delta _{SS^{}}{\displaystyle \underset{n=0}{\overset{2S}{}}}\stackrel{}{}\psi _n(x)\psi _n^{}(x)d^{\mathrm{\hspace{0.17em}3}}x,`$ (150)
where $`C_n^{2S}`$ is the binomial coefficient, and $`d\mu (x,z)`$ is the invariant measure (140). The relation (149) gives the connection between the description by the functions $`f(x,z)`$ and the standard description by the multicomponent functions $`\psi _n(x)`$. It is easy to see that the action of the operators $`\widehat{S}_k=\frac{1}{2}z\sigma _k_z`$ on the function (149) reduces to the multiplication of the column $`\psi (x)`$ by $`(2S+1)\times (2S+1)`$ matrices $`S_k`$ of $`SU(2)`$ generators in the representation $`T_S`$, $`\widehat{S}_kf(x,z)=\varphi (z)S_k\psi (x)`$. Matrices $`S_k`$ obey the commutation relations of spin projection operators, namely $`[S^i,S^j]=iϵ^{ijk}S_k`$.
In particular, the linear function of $`z^1`$,$`z^2`$ corresponds to spin $`S=1/2`$, and the action of the operators $`\widehat{S}_k`$ on $`\psi (x)`$ is reduced to the multiplication by $`\sigma `$-matrices, $`\widehat{S}_kf(x,z)=\varphi (z)\sigma _k\psi (x)`$.
As was mentioned above, the operator $`\widehat{𝐒}^2`$ is not a Casimir operator of $`M(3)`$ and labels equivalent representations of the group. This operator is the direct analog of the Lorentz spin operator in pseudoeuclidean case, and we will consider its properties in details.
1. Operator $`\widehat{𝐒}^2`$ is composed of right generators commuting with all left generators and therefore is not changed under the coordinate transformation (left transformations of the Euclidean group). The right transformations do not change the spin projection $`s`$ on the direction of propagation but change both spin $`S`$ and interval (distance).
2. Operator $`\widehat{𝐒}^2`$ does not depend on $`x`$ and commutes with operators $`x_k,\widehat{p}_k,\widehat{S}_k`$, therefore in the presence of interactions is conserved for any Hamiltonian $`\widehat{H}=\widehat{H}(x_k,\widehat{p}_k,\widehat{S}_k)`$.
3. The eigenvalues of $`\widehat{𝐒}^2`$ label irreps of the rotation subgroup in the spin space and define the possible values of the spin projection $`s`$, which can arise under the interactions.
Notice that in the representation theory of Galilei group (symmetry group of nonrelativistic mechanics, which includes $`M(3)`$ and ensures more general description) irreps labelled by different eigenvalues of $`\widehat{𝐒}^2`$ are not equivalent. The classification of irreps of Galilei group can be based on the use of two invariant equations. The Schrödinger equation fixes the mass $`m`$, and the second equation fixes the eigenvalue of spin operator $`\widehat{𝐒}^2`$ .
### B Field on the group $`M(2,1)`$ and fractional spin
Using the operators $`\widehat{J}^\rho =\widehat{L}^\rho +\widehat{S}^\rho =(1/2)ϵ^{\rho \mu \nu }\widehat{J}_{\mu \nu }`$, it is possible to rewrite the commutation relations (62) in the next form:
$$[\widehat{p}_\mu ,\widehat{p}_\nu ]=0,[\widehat{p}^\mu ,\widehat{J}^\nu ]=iϵ^{\mu \nu \eta }\widehat{p}_\eta ,[\widehat{J}^\mu ,\widehat{J}^\nu ]=iϵ^{\mu \nu \eta }\widehat{J}_\eta .$$
(151)
The invariant measure on the group is given by the formulas
$`d\mu (x,\underset{¯}{z})=d\mu (z)d^3x=Cd^3x\delta (|z_1|^2|z_2|^21)d^2z_1d^2z_2={\displaystyle \frac{1}{8\pi ^2}}d^3x\mathrm{sinh}\theta d\theta d\varphi d\psi .`$ (152)
$`\mathrm{}<x<+\mathrm{},0<\theta <\mathrm{},0<\varphi <2\pi ,2\pi <\psi <2\pi ,`$ (153)
where $`z_1=\mathrm{cosh}\frac{\theta }{2}e^{i(\psi \varphi )/2}`$, $`z_2=\mathrm{sinh}\frac{\theta }{2}e^{i(\psi +\varphi )/2}`$ are the elements of the first column of matrix $`Z`$ (71), and $`\theta ,\varphi ,\psi `$ are the analogs of Euler angles, $`z^2=z_1`$, $`z^1=z_2`$. The spin projection operators acting in the space of the functions on the group $`f(x,𝐳)`$ have the form
$`\widehat{S}^\mu ={\displaystyle \frac{1}{2}}(z\gamma ^\mu _z\stackrel{}{}z\stackrel{}{}\gamma ^\mu _\stackrel{}{z}),z=(z^1z^2),_z=(/z^1/z^2)^T,`$ (154)
$`\widehat{S}_R^\mu ={\displaystyle \frac{1}{2}}(\chi \stackrel{}{}\gamma ^\mu _\chi \stackrel{}{}\chi \gamma ^\mu _\stackrel{}{\chi }),\chi =(z^1\stackrel{}{}z^{\dot{2}}),_\chi =(/z^1/\stackrel{}{}z^{\dot{2}})^T,`$ (155)
where $`\gamma ^\mu `$ are three-dimensional $`\gamma `$-matrices,
$$\gamma ^\mu =(\sigma _3,i\sigma _2,i\sigma _1),\gamma ^\mu \gamma ^\nu =\eta ^{\mu \nu }i\epsilon ^{\mu \nu \rho }\gamma _\rho .$$
(156)
Notice that nonequivalent set of $`\gamma `$-matrices, $`\gamma ^\mu \gamma ^\nu =\eta ^{\mu \nu }+i\epsilon ^{\mu \nu \rho }\gamma _\rho `$, is used in some papers. In terms of the Euler angles one may obtain $`\widehat{S}^0=i/\varphi `$, $`\widehat{S}_R^0=i/\psi `$. The sets of commuting operators are the same as in Euclidean case.
In consequence of the identity $`\sigma _1\stackrel{}{}U\sigma _1=U`$ one may show that matrix $`\sigma _1`$ is the invariant symmetric tensor converting dotted and undotted indices,
$$\stackrel{}{}z_\alpha =(\sigma _1)_\alpha ^{\dot{\alpha }}\stackrel{}{}z_{\dot{\alpha }}.$$
(157)
According to (72), the invariant tensor $`\sigma _{\mu \alpha \dot{\alpha }}`$ connects vector index and two spinor indices of different types. On the other hand, using the identity mentioned above one may rewrite (72) in the form $`x^\nu (\sigma _\mu \sigma _1)=x^\mu U(\sigma _\mu \sigma _1)U^T`$. Thus the invariant tensor, which we denote as
$$\stackrel{ˇ}{\sigma }_{\mu \alpha \beta }=(\sigma _\mu \sigma _1)_{\alpha \beta },\stackrel{ˇ}{\sigma }_{\mu \alpha \beta }=\stackrel{ˇ}{\sigma }_{\mu \beta \alpha },$$
(158)
connects vector index and two spinor indices of one type. Thus, one can write the generators $`\widehat{S}^\mu `$ in the form $`\widehat{S}^\mu =\frac{1}{2}\stackrel{ˇ}{\sigma }_{\alpha \beta }^\mu (z^\alpha ^\beta +\stackrel{}{}z^\alpha \stackrel{}{}^\beta )`$. An analog of $`\sigma ^{\mu \nu }`$-matrices in 2+1 dimensions is $`(\sigma ^{\mu \nu })_{\alpha \beta }=\epsilon ^{\mu \nu \lambda }\stackrel{ˇ}{\sigma }_{\lambda \alpha \beta }`$. Raising one of the spinor indices of $`\stackrel{ˇ}{\sigma }_{\mu \alpha \beta }`$, one may obtain two sets of three-dimensional $`\gamma `$-matrices differed only by signs of $`\gamma ^0`$ and $`\gamma ^2`$.
Similarly to the Euclidean case, there are two standard realizations of the representation spaces. These realizations correspond to eigenvalues $`n=\pm 2S`$ and $`n=0`$ of the operator $`\widehat{S}_3^R`$.
The first realization is the space of analytic ($`n=2S`$) functions $`f(x,z)`$ or antianalytic ($`n=2S`$) functions $`f(x,\stackrel{}{}z)`$ of two complex variables $`z^1,z^2`$, $`|z^2|^2|z^1|^2=1`$, i.e. the space of functions of two-component spinors. The eigenfunctions of $`\widehat{S}_\mu \widehat{S}^\mu `$ are homogeneous functions of degree $`2S`$ in $`z`$. According to (143), we have $`\widehat{S}_R^0=(z^1/z^1+z^2/z^2)`$ for the space of analytic functions and $`\widehat{S}_R^0=\stackrel{}{}z^{\dot{1}}/\stackrel{}{}z^{\dot{1}}+\stackrel{}{}z^{\dot{2}}/\stackrel{}{}z^{\dot{2}}`$ for the space of antianalytic functions. The eigenfunctions of $`\widehat{S}_\mu \widehat{S}^\mu `$ in these spaces are also eigenfunctions of $`\widehat{S}_R^0`$ with eigenvalues $`n=2S`$ respectively.
The second realization is the space of eigenfunctions of $`\widehat{S}_R^0`$ with zero eigenvalue. It is the space of functions of five real parameters on the manifold
$$^3\times D^1,d\mu =(2\pi )^1d^3x\mathrm{sinh}\theta d\theta d\varphi ,$$
where $`D^1SU(1,1)/U(1)`$ is a complex disk. This functions do not depend on the angle $`\psi `$.
Remember some facts from the representation theory of $`SU(1,1)`$. For finite-dimensional nonunitary irreps $`T_S^0`$ of 2+1 Lorentz group $`SU(1,1)SO(2,1)`$ spin projection $`s`$ (the eigenvalue of $`\widehat{S}^0`$) can be only integer or half-integer, $`s=S,\mathrm{},S`$, where $`S0`$. However, in 2+1 dimensions Lorentz group has not compact non-Abelian subgroup. Therefore there are infinite-dimensional unitary representations corresponding to fractional $`S`$. These representations are multi-valued representations of $`SU(1,1)`$. For single-valued representations of $`SU(1,1)`$ the spin projection $`s`$ can be only integer or half-integer (for $`SO(2,1)`$ only integer).
The representations of discrete seria correspond to $`S<1/2`$. Irreps of the positive discrete series $`T_S^+`$ are bounded by lowest weight $`s=S`$, irreps of the negative discrete series $`T_S^{}`$ are bounded by highest weight $`s=S`$, irreps of the principal series correspond to $`S=1/2+i\lambda `$ and can be bounded by highest (lowest) weight only for $`S=1/2`$. For other irreps of principal series the spectrum of $`s`$ is not bounded. Supplementary series correspond to $`1/2<S<0`$ and are characterized by nonlocal scalar product.
A visual picture for weight diagrams of all seria on the plain $`S,s`$ one can find in .
Thus there are only two possibilities for description of a particle with fractional spin by means of unitary irreps of $`SU(1,1)`$ with local scalar product. The first corresponds to the discrete or principal seria irreps bounded by lowest (highest) weight, $`|s||S|1/2`$. The second corresponds to the principal series irreps which is not bounded.
Unitary irreps of discrete seria are used for the description of anyons . In corresponding unitary infinite-component representations of $`M(2,1)`$ were constructed in the space of functions of $`x^\mu `$ and complex variable $`z=z^1/z^2`$, i.e. on the coset space $`M(2,1)/U(1)`$. It was shown that RWE associated with irreps of the discrete seria have solutions only with definite sign of energy. Thus mentioned RWE are analogs of Majorana equations in 3+1 dimensions; this aspect was considered more detail in . Various formulations of the higher spin theory based on finite-component representations were considered, in particular, in .
### C Relativistic wave equations in 2+1 dimensions
Let us fix eigenvalues of the Casimir operators of the Poincaré group and of spin Lorentz subgroup:
$`\widehat{p}^2f(x,𝐳)=m^2f(x,𝐳),`$ (159)
$`\widehat{p}_\mu \widehat{S}^\mu f(x,𝐳)=Kf(x,𝐳),`$ (160)
$`\widehat{S}_\mu \widehat{S}^\mu f(x,𝐳)=S(S+1)f(x,𝐳).`$ (161)
Below we will call the operator $`\widehat{S}_\mu \widehat{S}^\mu `$ as operator of the Lorentz spin square.
Equations (159),(160) define some sub-representation of the left GRR of $`M(2,1)`$, which is characterized by mass $`m`$, Lorentz spin $`S`$, and by the eigenvalue $`K`$ of Lubanski-Pauli operator. At $`m=0`$ we suppose $`K=0`$, that is true for irreps with finite number of spinning degrees of freedom. The general cases for $`m=0`$ and for $`m`$ imaginary were discussed in .
Possible values of $`K`$ can be easily described in the massive case. Here we can use a rest frame, where $`\widehat{p}_\mu \widehat{S}^\mu =\widehat{S}^0msignp_0`$. Thus, $`K=sm=s^0m`$ for $`p_0>0`$ and $`K=sm=s^0m`$ for $`p_0<0`$, where $`s^0`$ is the eigenvalue of $`\widehat{S}^0`$. The spectrum of $`\widehat{S}^0`$ depends on the representation of the Lorentz group.
Variable $`s`$ labels irreps of the group $`M(2,1)`$ and can take both positive and negative values. Thus there exist the analogy with characterized by helicity massless particles in 3+1 dimensions. In both cases $`SO(2)`$ is the little group, and single-valued irreps of $`SO(2)`$ are labelled by integer number $`2s`$. (It is a particular case of the connection between the massive fields in $`d`$ dimensions and massless fields in $`d+1`$ dimensions, see ). Therefore we will call $`s`$ as helicity and $`|s|`$ as spin.
Corresponding to (86), space reflection reduces to the rotation at $`\pi `$ around axis $`x^0`$ and converts $`Z`$ to $`(Z^{})^1=\sigma _3Z\sigma _3`$, or $`z_1z_1`$, $`z_2z_2`$. Operators $`\widehat{p}^0`$, $`\widehat{S}^0`$ do not change. Thus, distinct from 3+1-dimensional case, space reflection leaves helicity unaltered.
Fixing $`S`$ in (161), we pass to the space of homogeneous functions of degree $`2S`$ in $`z_1,z_2`$. According to the sign of $`S`$, below we will consider two possible choices of $`SU(1,1)`$ irreps bounded either with both sides or with one side respectively.
Finite-dimensional nonunitary irreps $`T_S^0`$ of $`SU(1,1)`$ are labelled by positive integer or half-integer $`S`$. The basis in the representation space is formed by the polynomials of power $`2S`$ in $`z`$, see (A2). We denote corresponding representations of $`M(2,1)`$ as $`T_{m,s}^0`$.
Infinite-dimensional unitary irreps $`T_S^{}`$ ($`T_S^+`$) of $`SU(1,1)`$ are labelled by negative $`S<1/2`$ and are bounded by highest (lowest) weight. The basis in the representation space is formed by the quasipolynomials of power $`2S`$ in $`z`$, see (A4). We denote corresponding representations of $`M(2,1)`$ as $`T_{m,s}^{}`$ ($`T_{m,s}^+`$).
One may present a function $`f(x,z)`$ in the form
$$f(x,z)=\varphi (z)\psi (x),$$
(162)
where $`\varphi (z)`$ is a line composed of the basis elements $`\varphi _n(z)`$ of the corresponding $`SU(1,1)`$ representation, and $`\psi (x)`$ is a column composed of the coefficients in the decomposition over this basis. The action of the differential operators $`\widehat{S}^\mu `$ on a function $`f(x,z)`$ may be presented in terms of matrices
$$\widehat{S}^\mu f(x,z)=\varphi ^n(z)(S^\mu )_n^n^{}\psi _n^{}(x),$$
(163)
where $`S^\mu `$ are $`SU(1,1)`$ generators in the representation $`T_S`$ (see Appendix and also ). They obey the commutation relations of the $`SU(1,1)`$ group $`[S^\mu ,S^\nu ]=iϵ^{\mu \nu \eta }S_\eta `$.
For fixed $`S`$ in the matrix representation equations (159),(160) have the form
$`(\widehat{p}^2m^2)\psi (x)=0,`$ (164)
$`(\widehat{p}_\mu S^\mu sm)\psi (x)=0,`$ (165)
According to (165),
$`\psi ^{}(x)(S^\mu \stackrel{}{\widehat{p}_\mu }+sm)=0.`$It follows from the explicit expressions (A8) that for $`T_{m,s}^0`$ the relation $`S^\mu =\mathrm{\Gamma }S^\mu \mathrm{\Gamma }`$, where relations $`(\mathrm{\Gamma })_{nn^{}}=(1)^n\delta _{nn^{}}`$, $`\mathrm{\Gamma }^2=1`$ take place. For $`T_{m,s}^+`$ and $`T_{m,s}^{}`$ matrices $`S^\mu `$ are Hermitian, $`S^\mu =S^\mu `$, according to (A11). Let us introduce the notation
$`\overline{\psi }=\psi ^{}\mathrm{\Gamma }\mathrm{for}T_{m,s}^0,`$ (166)
$`\overline{\psi }=\psi ^{}\mathrm{for}T_{m,s}^+,T_{m,s}^{}.`$ (167)
The function $`\overline{\psi }(x)`$ obeys the equation
$$\overline{\psi }(x)(S^\mu \stackrel{}{\widehat{p}_\mu }+sm)=0.$$
(168)
As a consequence of the relations $`S^\mu =\mathrm{\Gamma }S^\mu \mathrm{\Gamma }`$ and $`(S^\mu )^{}=(1)^{\delta _0\mu }S^\mu `$ we obtain that for irrep $`T_{m,s}^0`$ finite transformation matrices obey the equation $`\mathrm{\Gamma }T^{}(g)\mathrm{\Gamma }=T^1(g)`$. Therefore $`\overline{\psi }(x)\psi (x)`$ is a scalar density, and one may define a scalar product in the space of columns
$$(\psi ^{}(x),\psi (x))=\overline{\psi }^{}(x)\psi (x)d^3x.$$
(169)
The scalar density is positive definite for $`T_{m,s}^+`$ and $`T_{m,s}^{}`$ in contrast to the case of $`T_{m,s}^0`$.
As a consequence of (165) and (168), the continuity equation holds
$$_\mu j^\mu =0,j^\mu =\overline{\psi }S^\mu \psi .$$
(170)
Together with the current vector $`j^\mu `$, by analogy with four-dimensional case , one can associate with the linear equation (165) the energy-momentum tensor $`T^{\mu \nu }`$ and the energy density $`W=T^{00}`$:
$$T^{\mu \nu }=\mathrm{Im}(S^\mu \frac{\psi }{x^\nu },\psi ),W=T^{00}=\mathrm{Im}(S^0\frac{\psi }{x^0},\psi ).$$
(171)
If matrix $`S^0`$ is diagonable, then the positiveness of $`W(x)`$ is equivalent to the requirement that
$$(S^0\psi ,S^0\psi )0$$
(172)
for all vectors $`\psi `$ . In particular, for $`T_{m,s}^+`$ and $`T_{m,s}^{}`$ the relation $`(S^0\psi ,S^0\psi )=\psi ^{}S^0S^0\psi 0`$ takes place, and energy density is positive definite.
There are two cases when equation (164) is the consequence of (165). Indeed, multiplying equation (160) by $`\widehat{p}_\mu S^\mu +ms`$, one gets
$$(\widehat{p}_\mu S^\mu +ms)(\widehat{p}_\nu S^\nu ms)\psi (x)=\left(\frac{1}{2}\widehat{p}_\mu \widehat{p}_\nu [S^\mu ,S^\nu ]_+m^2s^2\right)\psi (x)=0.$$
(173)
In the particular case $`S=1/2`$ we have $`s=\pm 1/2`$, $`S^\mu =\gamma ^\mu /2`$, and (173) is merely the Klein-Gordon equation (164). In general case the matrices $`S^\mu `$ are not $`\gamma `$-matrices in higher dimensions, and the squared equation (173) does not coincide with the Klein-Gordon equation (164). Using the rest frame, one may show that the equation (164) follows from (165) also in the case of vector representation $`S=1`$, $`s=\pm 1`$. In the other cases for the identification of the irrep of $`M(2,1)`$ it is necessary to use both equations of the system (164),(165). Notice that another approach to the description of fields with fixed spin and mass was suggested in ; this approach is based on the system of spinor linear equations.
It is naturally to connect spin value with the highest (lowest) weight of the irrep of Lorentz group, $`s=\pm S`$. This means that up to a sign ($`+`$ for $`p_0>0`$, $``$ for $`p_0<0`$) $`s`$ is equal to maximal or minimal eigenvalue of the operator $`\widehat{S}^0`$ in the representation $`T_S`$ of the Lorentz group. According to (159)-(161), in this case functions $`f(x,z)`$ obey the equations
$`\widehat{p}^2f(x,𝐳)=m^2f(x,𝐳),`$ (174)
$`\widehat{p}_\mu \widehat{S}^\mu f(x,𝐳)=msf(x,𝐳),s=\pm S,`$ (175)
$`\widehat{S}_\mu \widehat{S}^\mu f(x,𝐳)=S(S+1)f(x,𝐳).`$ (176)
In the framework of the system (174)-(176) there are two possibilities to describe one and the same spin:
1. Equations for $`f(x,𝐳)=\varphi (𝐳)\psi (x)`$, where $`\varphi (𝐳)`$ transform under finite-dimensional nonunitary irrep of the Lorentz group.
2. Equations for $`f(x,𝐳)=\varphi (𝐳)\psi (x)`$, where $`\varphi (𝐳)`$ transform under infinite-dimensional unitary irrep of the Lorentz group. These equations allow us to describe also particles with fractional spin (anyons).
(1) At first, consider the Poincaré group representations $`T_{m,s}^0`$ associated with finite-dimensional non-unitary irreps of $`SU(1,1)`$. In this case $`S`$ has to be positive, integer or half-integer. In the rest frame the solutions of the system (174)-(176) in the space of analytic functions (polynomials of power $`2S`$ in $`z^1`$, $`z^2`$) are
$`s=S>0:f(x,z)=C_1(z^1)^Se^{imx^0}+C_2(z^2)^Se^{imx^0},`$ (177)
$`s=S<0:f(x,z)=C_3(z^1)^Se^{imx^0}+C_4(z^2)^Se^{imx^0}.`$ (178)
For unique mass and spin there exist four independent components differed by signs of $`p_0`$ and $`s`$, which correspond to four irreps of $`M(2,1)`$. The separation by the sign of helicity $`s`$ has absolute character since these states are solutions of different equations. But the states with different sign of $`p_0`$ are solutions of one and the same equation. Hence, the energy spectrum of solutions is not bounded below or above.
In the space of antianalytic functions (polynomials of power $`2S`$ in $`\stackrel{}{}z^{\dot{1}}`$, $`\stackrel{}{}z^{\dot{2}}`$) solutions of the system (174)-(176) are
$`s=S>0:f(x,\stackrel{}{}z)=C_1(\stackrel{}{}z^{\dot{1}})^Se^{imx^0}+C_2(\stackrel{}{}z^{\dot{2}})^Se^{imx^0},`$ (179)
$`s=S<0:f(x,\stackrel{}{}z)=C_3(\stackrel{}{}z^{\dot{1}})^Se^{imx^0}+C_4(\stackrel{}{}z^{\dot{2}})^Se^{imx^0}.`$ (180)
These solutions are connected with previous case (177),(178) by charge conjugation (88) and therefore may be treated as the solutions describing antiparticles.
The wave function (177) corresponding to the helicity $`s=S`$ has the form $`C(z^2)^{2S}e^{ip_0x^0}`$, $`p_0=m`$, in the rest frame. Acting on it by finite transformations, we get a solution in the form of the plane wave, which is characterized by the momentum $`p`$:
$`P=UP_0U^{},P_0=mI,Z=UZ_0,Z_0=(z_1z_2)^T,`$ (181)
$`f(x,z)=(2\pi )^{3/2}\left(z^2u_1z^1u_2\right)^{2S}e^{ipx}.`$ (182)
The state with $`P_0=mI`$ has the stationary subgroup $`U(1)`$, and we can take elements $`u_1=\mathrm{cosh}\theta /2`$ and $`u_2=\mathrm{sinh}\theta /2e^{i\omega }`$ of the first line of matrix $`U`$, that depends on two parameters only. Thus $`p_0=E=m\mathrm{cosh}\theta ,p_1+ip_2=m\mathrm{sinh}\theta e^{i\omega }`$, and one can express the parameters $`u_1`$ and $`u_2`$ via the momentum $`p`$:
$$\left(\genfrac{}{}{0pt}{}{u_1}{u_2}\right)=\frac{1}{\sqrt{2m(E+m)}}\left(\genfrac{}{}{0pt}{}{E+m}{p_1+ip_2}\right).$$
(183)
$`2S+1`$ components $`\psi _n(x)`$ are the coefficients in the decomposition of the function (182) over the basis $`\varphi ^n(z)`$, $`f(x,z)=\varphi ^n(z)\psi _n(x),n=0,1,\mathrm{},2S`$:
$`\psi _n(x)`$ $`=`$ $`(2\pi )^{3/2}\left(C_{2S}^n\right)^{1/2}(u_1)^{2Sn}(u_2)^ne^{ipx}`$ (184)
$`=`$ $`(2\pi )^{3/2}\left(C_{2S}^n\right)^{1/2}{\displaystyle \frac{\left(E+m\right)^{2Sn}\left(p_1ip_2\right)^n}{\left(2m(E+m)\right)^S}}e^{ipx}.`$ (185)
In the particular case $`S=1/2`$ we get
$$\psi (x)=\frac{1}{\sqrt{2m(Em)}}\left(\genfrac{}{}{0pt}{}{p_2ip_1}{E+m}\right)e^{ipx}.$$
Considering the system (175)-(176) without the condition of the mass irreducibility (174), it is easy to see that the charge density $`j^0=\psi ^{}\mathrm{\Gamma }S^0\psi `$ is positive definite only for $`S=1/2`$, and the energy density $`T^{00}`$ is positive definite only for $`S=1`$. The scalar density $`\overline{\psi }\psi =\psi ^{}\mathrm{\Gamma }\psi `$ is not positive definite.
Let us show that for the particles with half-integer spin described by the system (174)-(176) the charge density $`j^0`$ (170) is positive definite. In the rest frame solutions of the system (174)-(176) have only two components (labelled by $`s_0=\pm S`$), which we denote as $`\psi _S(x)`$ and $`\psi _S(x)`$. For half-integer spin an inequality $`j^0=\psi ^{}\mathrm{\Gamma }S^0\psi =S(|\psi _S|^2+|\psi _S|^2)>0`$ holds. At $`S3/2`$ from the explicit form of matrices $`S^1`$ and $`S^2`$ (A8) one can obtain that in the rest frame $`j^1=j^2=0`$, therefore the square of the current vector $`(j^0)^2(j^1)^2(j^2)^2`$ is positive. Therefore $`j^0>0`$ for any plane wave.
Thus the charge density $`j^0`$ is positive definite for half-integer spin particles described by representations $`T_{m,s}^0`$ of $`M(2,1)`$. The scalar density and the energy density are proportional to $`\psi ^{}\mathrm{\Gamma }\psi =|\psi _S|^2|\psi _S|^2`$ in the rest frame and therefore are indefinite.
Let us consider now particles with integer spin. In the rest frame solutions of the system also have only two components: $`\psi _S(x)`$ and $`\psi _S(x)`$, $`(S^0\psi ,S^0\psi )=\psi ^{}\mathrm{\Gamma }S^0S^0\psi =S^2(|\psi _S|^2+|\psi _S|^2)>0`$. Thus the energy density is positive definite for integer spin particles described by representations $`T_{m,s}^0`$ of $`M(2,1)`$. The charge density is proportional to $`|\psi _S|^2|\psi _S|^2`$ in the rest frame and therefore is indefinite.
Consider two particular cases explicitly. For $`S=1/2`$ the decomposition (162) has the following form
$$f(x,z)=z^1\psi _1(x)+z^2\psi _2(x),\psi ^{}(x^{})=U^1\psi (x),\psi (x)=(\psi _1(x)\psi _2(x))^T.$$
(186)
Taking into account the relation $`U^1=\sigma _3U^{}\sigma _3`$, which is valued for the $`SU(1,1)`$ matrices, we get the transformation law for the line $`\overline{\psi }=\psi ^{}\sigma _3`$, $`\overline{\psi }^{}(x^{})=\overline{\psi }(x)U`$. The product $`\overline{\psi }(x)\psi (x)=|\psi _1(x)|^2|\psi _2(x)|^2`$ is the scalar density.
Thus, in the case under consideration, we have two equivalent descriptions. One in terms of functions (186) and another one in terms of lines $`\overline{\psi }(x)`$ or columns $`\psi (x)`$. One can find the action of the operators $`\widehat{S}^\mu `$ in the latter representation, and equation (165) can be rewritten in the form of $`2+1`$ Dirac equation
$$\widehat{S}^\mu \psi (x)=\frac{1}{2}\gamma ^\mu \psi (x),(\widehat{p}_\mu \gamma ^\mu m)\psi (x)=0,$$
(187)
where minus corresponds to $`s=1/2`$, plus corresponds to $`s=1/2`$, and $`\gamma ^\mu `$ are $`2\times 2`$ $`\gamma `$-matrices (156) in $`2+1`$ dimensions. The functions $`\psi =(\psi ^1\mathrm{\hspace{0.33em}0})^T`$ and $`\psi =(0\psi ^2)^T`$ are eigenvectors of the operator $`\widehat{S}^0`$ with the eigenvalues $`\pm 1/2`$.
Sometimes two equations for $`s=\pm 1/2`$ are written as one equation for the four-component reducible representation , $`(\widehat{p}_\mu \mathrm{\Gamma }^\mu m)\mathrm{\Psi }(x)=0`$, where $`\mathrm{\Gamma }^\mu =diag(\gamma ^\mu ,\gamma ^\mu )`$, that corresponds to the simultaneous consideration of particles with opposite helicities.
For $`S=1`$ the decomposition (162) has the following form
$$f(x,z)=\psi _{11}(x)(z^1)^2+\psi _{12}(x)z^1z^2+\psi _{22}(x)(z^2)^2,$$
(188)
where $`\psi (x)=(\psi _{11}(x)\psi _{12}(x)/\sqrt{2}\psi _{22}(x))^T`$ is subjected to equation (165) with the matrices
$$S^0=\left(\begin{array}{ccc}\hfill 1& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 1\end{array}\right),S^1=\frac{1}{\sqrt{2}}\left(\begin{array}{ccc}\hfill 0& \hfill 1& \hfill 0\\ \hfill 1& \hfill 0& \hfill 1\\ \hfill 0& \hfill 1& \hfill 0\end{array}\right),S^2=\frac{i}{\sqrt{2}}\left(\begin{array}{ccc}\hfill 0& \hfill 1& \hfill 0\\ \hfill 1& \hfill 0& \hfill 1\\ \hfill 0& \hfill 1& \hfill 0\end{array}\right).$$
(189)
If instead of the cyclic components $`\psi _{\alpha \beta }(x)`$ one introduces new (Cartesian) components $`_\mu =\stackrel{ˇ}{\sigma }_\mu ^{\alpha \beta }\psi _{\alpha \beta }(x)`$, where $`\stackrel{ˇ}{\sigma }_{\mu \alpha \beta }`$ is defined in (158), $`_0=2\psi ^{12}`$, $`_1=\psi ^{11}+\psi ^{22}`$, $`_2=i(\psi ^{22}\psi ^{11})`$, then equation (165) takes the form
$$_\mu \epsilon ^{\mu \nu \eta }_\eta sm^\nu =0.$$
(190)
A transversality condition follows from (190), $`_\mu ^\mu =0`$. One can see now that equations (190) are in fact field equations of the so called ”self-dual” free massive field theory . As remarked in , this theory is equivalent to the topologically massive gauge theory with the Chern-Simons term . Indeed, the transversality condition allows introducing gauge potentials $`A_\mu `$, namely a transverse vector can be written as a curl $`^\mu =\epsilon ^{\mu \nu \lambda }_\nu A_\lambda =\epsilon ^{\mu \nu \lambda }F_{\nu \lambda }/2`$, where $`F_{\nu \lambda }=_\nu A_\lambda _\lambda A_\nu `$ is the field strength. Thus, $`^\mu `$ appears to be dual field strength, which is a tree-component vector in 2+1 dimensions. Then (190) implies the following equations for $`F_{\mu \nu }`$
$$_\mu F^{\mu \nu }\frac{sm}{2}\epsilon ^{\nu \alpha \beta }F_{\alpha \beta }=0,$$
which are the field equations of the topologically massive gauge theory.
To describe neutral spin 1 particle coinciding with its antiparticle one may consider a function
$$f(x,𝐳)=\psi _{11}(x)z^1\stackrel{}{}z^1+\psi _{12}(x)(z^1\stackrel{}{}z^2+\stackrel{}{}z^1z^2)/2+\psi _{22}(x)z^2\stackrel{}{}z^2,$$
(191)
where we have used (157) for the convertation to undotted indices. The spin part of the function (191) depends not on three angles as in the case (188), but on two angles only. This function is an eigenfunction of operator $`\widehat{S}_R^3`$ with zero eigenvalue. Substituting (191) into (175), we again obtain equation (190).
(2) Consider now Poincaré group representations $`T_{m,s}^+`$ and $`T_{m,s}^{}`$ associated with unitary infinite-dimensional irreps of $`SU(1,1)`$ with highest (lowest) weight. In this case $`S`$ can be non-integer, $`S<1/2`$ (discrete series) or $`S=1/2`$ (principal series). Eigenvalues of $`\widehat{S}^0`$ can take only positive values for discrete positive series, $`s^0=S+n`$, and only negative valuses for negative one, $`s^0=Sn`$, where $`n=0,1,2,\mathrm{}`$.
Let us consider the energy spectrum of the system (174)-(176) at $`m0`$. According to the first equation $`p_0=\pm m`$. The second equation ensures the relation between spectra of the operators $`\widehat{p}_0`$ and $`\widehat{S}^0`$,
$$p_0s^0=ms.$$
(192)
For representations $`T_{m,s}^0`$, which correspond to finite-dimensional irreps $`T_S^0`$ of the Lorentz group, the value of $`s^0`$ can be both positive and negative. Therefore for any $`s`$ there exist both positive-frequency and negative-frequency solutions, and the representations $`T_{m,s}^0`$ splits into two irreps, characterized by $`signp_0=\pm 1`$.
For unitary $`SU(1,1)`$ irreps with highest (lowest) weight the spectrum of $`\widehat{S}^0`$ has definite sign. For $`T_S^+`$, which act in the space of analytic functions, the spectrum of operator $`\widehat{S}^0`$ is positive, and for $`T_S^{}`$, which act in the space of antianalytic functions, is negative. Therefore the sign of energy $`p_0`$ coincides with the sign of $`s`$ for $`T_S^+`$, and the signs of $`p_0`$ and $`s`$ are opposite for $`T_S^{}`$. Thus $`T_{m,s}^+`$ and $`T_{m,s}^{}`$ are irreps of $`M(2,1)`$.
As well as in the case of representations $`T_{m,s}^0`$, for unique mass and spin there are four states distinguished in signs of $`p_0`$ and $`s`$. In the rest frame there are two solutions of the system in the space of analytic functions:
$`p_0>0,s>0:f(x,z)=(2\pi )^{3/2}(z^2)^Se^{imx^0},`$ (193)
$`p_0<0,s<0:f(x,z)=(2\pi )^{3/2}(z^2)^Se^{imx^0}.`$ (194)
The solutions are connected by time reflection $`T^{}`$ (87). In the space of antianalytic functions there are also two solutions:
$`p_0>0,s<0:f(x,\stackrel{}{}z)=(2\pi )^{3/2}(\stackrel{}{}z^{\dot{2}})^Se^{imx^0}`$ (195)
$`p_0<0,s>0:f(x,\stackrel{}{}z)=(2\pi )^{3/2}(\stackrel{}{}z^{\dot{2}})^Se^{imx^0}.`$ (196)
These solutions are connected respectively with (193),(194) by Schwinger time reversal $`T_{sch}=CT^{}`$, which turns particles into antiparticles. Thus, there exist four equations distinguished in sign of $`s`$ and by the used functional space (irrep $`T_S^+`$ or $`T_S^{}`$ of the Lorentz group), and any equation has the solutions only with definite sign of $`p_0`$.
In contrast to the case of $`T_{m,s}^0`$, where the energy spectrum $`p_0`$ is not bounded both above and below, the energy spectrum has definite sign. In any inertial frame the spectrum is bounded below by $`p_0=m`$ for the solutions (193), (195) and above by $`p_0=m`$ for the solutions (194), (196).
For the unitary irreps of $`M(2,1)`$ under consideration, which correspond to the irreps of the discrete seria of the Lorentz group, integration of the functions (A4) in the invariant measure (152) gives
$`{\displaystyle }\stackrel{}{}f_{S_1}(x,𝐳)f_{S_2}^{}(x,𝐳)d\mu (x,𝐳)=\delta _{S_1S_2}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}\psi ^n(x)\psi ^n(x)d^3x},`$ (197)
$`{\displaystyle }\stackrel{}{}f_{S_1}(x,𝐳)f_{S_2}^{}(x,𝐳)d\mu (𝐳)=\delta _{S_1S_2}\psi ^{}(x)\psi ^{}(x),`$ (198)
In particular, the states (193)-(196) have the norm $`\delta _{SS^{}}\delta (pp^{})`$. For the principal series $`j=1/2+i\lambda `$, and $`\delta _{j_1j_2}`$ in (197) is changed by $`\delta (\lambda _1\lambda _2)`$. At the same time, the integral over the spin space diverges for the representations $`T_{m,s}^0`$, which correspond to finite-dimensional irreps of the Lorentz group.
Arbitrary plain wave solutions can be obtained by analogy with considered above case of $`T_{m,s}^0`$. For example, for the states (193) one can get the formula (185), where now $`C_{2S}^n`$ are the coefficients from (A4) and $`n=0,1,2,\mathrm{}`$. The power $`2S`$ is negative, and the decomposition $`f(x,z)=\varphi _n(z)\psi ^n(x)`$ contains infinite number of terms.
Let us summarize some properties of the unitary irreps under consideration. Irreps $`T_{m,s}^+`$ and $`T_{m,s}^{}`$ of the Poincaré group describe particles and antiparticles respectively. Charge density $`j^0=\psi ^{}S^0\psi `$ is positive definite for particles and negative definite for antiparticles. The energy density is positive definite in both cases since $`(S^0\psi ,S^0\psi )=\psi ^{}S^0S^0\psi >0`$. Besides, for the unitary irreps the scalar density $`\psi ^{}\psi `$ is also positive definite in contrast to the finite-dimensional case. The existence of positive definite scalar density ensures the possibility of probability amplitude interpretation of $`\psi (x)`$.
Thus in 2+1 dimensions the problem of the construction of positive-energy RWEs is solved by the system (174)-(176) for the infinite-dimensional unitary irreps $`T_{m,s}^+`$ (signs of $`p_0`$ and $`s`$ are the same) or $`T_{m,s}^{}`$ (signs of $`p_0`$ and $`s`$ are opposite) characterized by the mass $`m`$ and the helicity $`s`$. These irreps of the Poincaré group are associated with irreps $`T_S^+`$ and $`T_S^{}`$ of the Lorentz group with lowest (highest) weight. Charge conjugation, changing signs of $`p_0`$ and $`s^0`$, leaves the helicity $`s`$ invariant and turns $`T_{m,s}^+`$ into $`T_{m,s}^{}`$.
An interesting problem is to find an explicit form of the intertwining operator $`A`$ for the unitary irreps $`T_{m,s}^+`$, $`T_{m,s}^{}`$ and the representation $`T_{m,s}^0`$ labelled by the same mass $`m`$ and helicity $`s`$ but assotiated with finite-dimensional nonunitary irreps of the Lorentz group, $`AT_{m,s}^0=(T_{m,s}^+T_{m,s}^{})A`$. The intertwining operator is nonunitary and must be a function of the generators of right translations, since other generators commute with Lorentz spin square operator $`\widehat{S}_\mu \widehat{S}^\mu `$ and can’t change the representation of spin Lorentz subgroup.
Notice that the 2+1 Dirac equation arises also in the case of unitary infinite-dimensional irreps $`T_S^+`$ and $`T_S^{}`$ of the Lorentz group not as an equation for a true wave function, but as an equation for spin coherent states evolution. In this case the equation includes effective mass $`m_s=|\frac{s}{S}|m`$, $`s=S,S+1,\mathrm{}`$ .
Among the above considered RWE there exist ones which describe particles with fractional real spin. These equations are associated with unitary multi-valued irreps of the Lorentz group and can be used to describe anyons.
In spite of the fact that the number of independent polarization states for massive 2+1 particle is one, the vectors of the corresponding representation space of irreps $`T_{m,s}^+`$, $`T_{m,s}^{}`$ have infinite number of components in matrix representation. Thus, $`z`$-representation is more convenient in this case.
## V Four-dimensional case
### A Field on the group $`M(3,1)`$
The generators and the action of the left GRR on the functions $`f(x,𝐳)`$ are given by formulas (62), (69). For spin projection operators it is convenient to use three-dimensional vector notation $`\widehat{S}_k=\frac{1}{2}ϵ_{ijk}\widehat{S}^{ij}`$, $`\widehat{B}_k=\widehat{S}_{0k}`$. The explicit calculation gives
$`\widehat{S}_k={\displaystyle \frac{1}{2}}(z\sigma _k_z\stackrel{}{}z\stackrel{}{}\sigma _k_\stackrel{}{z})+\mathrm{},`$ (199)
$`\widehat{B}_k={\displaystyle \frac{i}{2}}(z\sigma _k_z+\stackrel{}{}z\stackrel{}{}\sigma _k_\stackrel{}{z})+\mathrm{},z=(z^1z^2),_z=(/z^1/z^2)^T;`$ (200)
$`\widehat{S}_k^R={\displaystyle \frac{1}{2}}(\chi \stackrel{}{}\sigma _k_\chi \stackrel{}{}\chi \sigma _k_\stackrel{}{\chi })+\mathrm{},`$ (201)
$`\widehat{B}_k^R={\displaystyle \frac{i}{2}}(\chi \stackrel{}{}\sigma _k_\chi +\stackrel{}{}\chi \sigma _k_\stackrel{}{\chi })+\mathrm{},\chi =(z^1\underset{¯}{z}^1),_\chi =(/z^1/\underset{¯}{z}^1)^T;`$ (202)
Dots in the formulas replace analogous expressions obtaining by the substitutions $`zz^{}=(\underset{¯}{z}^1\underset{¯}{z}^2)`$, $`\chi \chi ^{}=(z^2\underset{¯}{z}^2)`$.
Since $`detZ=1`$, then any of $`z_\alpha `$, $`\underset{¯}{z}_\alpha `$ can be expressed in terms of other three, for example $`\underset{¯}{z}_2=(1z_2\underset{¯}{z}_1)/z_1`$. Invariant measure on $`^4\times SL(2,C)`$ has the form
$$d\mu (x,𝐳)=(i/2)^3d^4xd^2z_1d^2z_2d^2\underset{¯}{z}_1|z_1|^2.$$
(203)
The functions on the Poincaré group depend on 10 parameters, and correspondingly there are 10 commuting operators (two Casimir operators, four left and four right generators).
Both the Poincaré group and the spin Lorentz subgroup have two Casimir operators:
$`\widehat{p}^2=\widehat{p}_\mu \widehat{p}^\mu ,\widehat{W}^2=\widehat{W}_\mu \widehat{W}^\mu ,\mathrm{where}\widehat{W}^\mu ={\displaystyle \frac{1}{2}}ϵ^{\mu \nu \rho \sigma }\widehat{p}_\nu \widehat{J}_{\rho \sigma }={\displaystyle \frac{1}{2}}ϵ^{\mu \nu \rho \sigma }\widehat{p}_\nu \widehat{S}_{\rho \sigma },`$ (204)
$`{\displaystyle \frac{1}{2}}\widehat{S}_{\mu \nu }\widehat{S}^{\mu \nu }={\displaystyle \frac{1}{2}}\widehat{S}_{\mu \nu }^R\widehat{S}_R^{\mu \nu }=\widehat{𝐒}^2\widehat{𝐁}^2,{\displaystyle \frac{1}{16}}ϵ^{\mu \nu \rho \sigma }\widehat{S}_{\mu \nu }\widehat{S}_{\rho \sigma }={\displaystyle \frac{1}{16}}ϵ^{\mu \nu \rho \sigma }\widehat{S}_{\mu \nu }^R\widehat{S}_{\rho \sigma }^R=\widehat{𝐒}\widehat{𝐁}.`$ (205)
Let us consider a set of ten commuting operators
$$\widehat{p}_\mu ,\widehat{W}^2,\widehat{𝐩}\widehat{𝐒},\widehat{𝐒}^2\widehat{𝐁}^2,\widehat{𝐒}\widehat{𝐁},\widehat{S}_3^R,\widehat{B}_3^R.$$
(206)
This set consists of operators of momenta, the Lubanski-Pauli operator $`\widehat{W}^2`$, the proportional to helicity operator $`\widehat{𝐩}\widehat{𝐉}=\widehat{𝐩}\widehat{𝐒}`$, and four operators, which are the functions of the right generators. This four operators commute with the left rotations and translations and allow one to distinguish equivalent irreps in the decomposition of GRR. In the rest frame $`\widehat{𝐩}\widehat{𝐒}=0`$, and the complete set of commuting operators can be obtained from (206) with the help of the replacement of $`\widehat{𝐩}\widehat{𝐒}`$ by $`\widehat{S}_3`$.
Functions $`f(x,𝐳)`$ on the group $`M(3,1)`$ are the functions of four real variables $`x^\mu `$ and four complex variables $`z_\alpha `$, $`\underset{¯}{z}_\alpha `$ with the constraint $`z_1\underset{¯}{z}_2z_2\underset{¯}{z}_1=1`$.
The space of functions on the Poincaré group contains the subspace of analytic functions $`f(x,z,\stackrel{}{}\underset{¯}{z})`$ of the elements of the Dirac $`z`$-spinor
$$Z_D=(z^\alpha ,\stackrel{}{}\underset{¯}{z}_{\dot{\alpha }}).$$
(207)
Charge conjugation means the transition to subspace of antianalytic functions (i.e. analytic functions of $`\underset{¯}{z}^\alpha ,\stackrel{}{}z_{\dot{\alpha }}`$).
According to (86), for the space inversion we have $`Z\stackrel{P}{}(Z^1)^{}`$ or
$$\left(\begin{array}{cc}z^1& \underset{¯}{z}^1\\ z^2& \underset{¯}{z}^2\end{array}\right)\stackrel{P}{}\left(\begin{array}{cc}\hfill \stackrel{}{}\underset{¯}{z}_{\dot{1}}& \hfill \stackrel{}{}z_{\dot{1}}\\ \hfill \stackrel{}{}\underset{¯}{z}_{\dot{2}}& \hfill \stackrel{}{}z_{\dot{2}}\end{array}\right),$$
(208)
This transformation reverses the sign of the boost operators $`\widehat{B}_k`$. It is easy to see that, in contrast to charge conjugation, space inversion conserves the analyticity (or antianalyticity) of functions of $`Z_D`$.
Similarly to three-dimensional case (see (145)), eigenfunctions of $`\widehat{S}_3^R`$ and $`\widehat{B}_3^R`$ differ only by a phase factor. Fixing eigenvalues of operators $`\widehat{S}_3^R`$ and $`\widehat{B}_3^R`$, one may pass to the space of functions of $`x^\mu `$ and elements of Majorana $`z`$-spinor
$$Z_M=(z^\alpha ,\stackrel{}{}z_{\dot{\alpha }}),$$
(209)
i.e. the space of functions of 8 real independent variables on the manifold
$$^4\times ^2,d\mu =d^4xd^2z_1d^2z_2.$$
(210)
Thus in this space we have 8 commuting operators (2 Casimir operators, 4 operators distinguish states inside the irrep, 2 operators distinguish equivalent irreps). Notice that physical argumentation of necessity to make use at least 8 variables in order to describe spinning particles contains in . The space reflection takes the functions of $`Z_M`$ to the functions of $`\underset{¯}{Z}_M=(\underset{¯}{z}^\alpha ,\stackrel{}{}\underset{¯}{z}_{\dot{\alpha }})`$; as was mentioned above, $`\underset{¯}{z}^\alpha `$ and $`z^\alpha `$ have the same transformation rule. The charge conjugation leave the space of functions of $`Z_M`$ invariant.
Below we will consider the massive case characterized by the symmetry with respect to space reflection and therefore the space of the analytic functions of Dirac $`z`$-spinor $`Z_D`$, unless otherwise stipulated. In this space the action of $`M(3,1)`$ is given by a formula
$`T(g)f(x,z,\underset{¯}{z})=f(g^1x,g^1z,g^1\underset{¯}{z}),`$ (211)
$`(g^1x)^\mu =(\mathrm{\Lambda }^1)_\nu ^\mu x^\nu ,(g^1z)^\alpha =U_\beta ^\alpha z^\beta ,(g^1\underset{¯}{z})_{\dot{\alpha }}=(\stackrel{}{}U^1)_{\dot{\alpha }}^{\dot{\beta }}\stackrel{}{}\underset{¯}{z}_{\dot{\beta }}.`$ (212)
Spin projection operators have the form
$$\widehat{S}_k=\frac{1}{2}(z\sigma _k_z\stackrel{}{}\underset{¯}{z}\stackrel{}{}\sigma _k_{\stackrel{}{\underset{¯}{z}}}),\widehat{B}_k=\frac{i}{2}(z\sigma _k_z+\stackrel{}{}\underset{¯}{z}\stackrel{}{}\sigma _k_{\stackrel{}{\underset{¯}{z}}}).$$
(213)
It is known that one can compose the combinations $`\widehat{M}_k`$, $`\widehat{\overline{M}}_k`$ :
$`\widehat{M}_k={\displaystyle \frac{1}{2}}(\widehat{S}_ki\widehat{B}_k)=z\sigma _k_z,\widehat{M}_+=z^1/z^2,\widehat{M}_{}=z^2/z^1,`$ (214)
$`\widehat{\overline{M}}_k={\displaystyle \frac{1}{2}}(\widehat{S}_k+i\widehat{B}_k)=\stackrel{}{}\underset{¯}{z}\stackrel{}{}\sigma _k_{\stackrel{}{\underset{¯}{z}}},\widehat{\overline{M}}_+=\stackrel{}{}\underset{¯}{z}^{\dot{1}}/\stackrel{}{}\underset{¯}{z}^{\dot{2}},\widehat{\overline{M}}_{}=\stackrel{}{}\underset{¯}{z}^{\dot{2}}/\stackrel{}{}\underset{¯}{z}^{\dot{1}},`$ (215)
such that $`[\widehat{M}_i,\widehat{\overline{M}}_k]=0`$. For unitary representations of the Lorentz group $`\widehat{S}_k^{}=\widehat{S}_k`$, $`\widehat{B}_k^{}=\widehat{B}_k`$, and these operators obey the relation $`\widehat{M}_k^{}=\widehat{\overline{M}}_k`$ (for finite-dimensional nonunitary irreps $`\widehat{S}_k^{}=\widehat{S}_k`$, $`\widehat{B}_k^{}=\widehat{B}_k`$ and $`\widehat{M}_k^{}=\widehat{\overline{M}}_k`$). Introducing spin operators with spinor indices
$$\widehat{M}_{\alpha \beta }=(\sigma _{\mu \nu })_{\alpha \beta }\widehat{S}^{\mu \nu },\widehat{\overline{M}}_{\dot{\alpha }\dot{\beta }}=(\overline{\sigma }_{\mu \nu })_{\dot{\alpha }\dot{\beta }}\widehat{S}^{\mu \nu },$$
(216)
where $`\sigma _{\mu \nu }`$ and $`\overline{\sigma }_{\mu \nu }`$ are defined in (B6), we obtain
$`\widehat{S}^{\mu \nu }={\displaystyle \frac{1}{2}}\left((\sigma ^{\mu \nu })^{\alpha \beta }\widehat{M}_{\alpha \beta }+(\overline{\sigma }^{\mu \nu })^{\dot{\alpha }\dot{\beta }}\widehat{\overline{M}}_{\dot{\alpha }\dot{\beta }}\right),`$ (217)
$`\widehat{M}_{\alpha \beta }\widehat{M}^{\alpha \beta }=2\widehat{𝐌}^2,\widehat{\overline{M}}_{\dot{\alpha }\dot{\beta }}\widehat{\overline{M}}^{\dot{\alpha }\dot{\beta }}=2\widehat{\overline{𝐌}}^2.`$ (218)
In the space of analytic functions of $`z,\stackrel{}{}\underset{¯}{z}`$ we have:
$$\widehat{M}_{\alpha \beta }=\frac{1}{2}(z_\alpha _\beta +z_\beta _\alpha ),\widehat{\overline{M}}_{\dot{\alpha }\dot{\beta }}=\frac{1}{2}(\stackrel{}{}\underset{¯}{z}_{\dot{\alpha }}\underset{¯}{}_{\dot{\beta }}+\stackrel{}{}\underset{¯}{z}_{\dot{\beta }}\underset{¯}{}_{\dot{\alpha }}).$$
(219)
Taking into account that operators $`\widehat{M}_k`$, $`\widehat{\overline{M}}_k`$ are subjected to commutation relations of $`su(2)`$ algebra, we obtain spectra of the Casimir operators of the Lorentz subgroup:
$`\widehat{𝐒}^2\widehat{𝐁}^2=2(\widehat{𝐌}^2+\widehat{\overline{𝐌}}^2)=2j_1(j_1+1)+2j_2(j_2+1)={\displaystyle \frac{1}{2}}(k^2\rho ^24),`$ (220)
$`\widehat{𝐒}\widehat{𝐁}=i(\widehat{𝐌}^2\widehat{\overline{𝐌}}^2)=i\left(j_1(j_1+1)j_2(j_2+1)\right)=k\rho ,`$ (221)
$`\mathrm{where}\rho =i(j_1+j_2+1),k=j_1j_2.`$ (222)
Thus irreps of the Lorentz group $`SL(2,C)`$ are labelled by the pair $`(j_1,j_2)`$. It is convenient to label unitary irreps by $`[k,\rho ]`$, where irreps $`[k,\rho ]`$ and $`[k,\rho ]`$ are equivalent .
Notice that the formulas (212)-(222) are also valid if, using substitution $`\stackrel{}{}\underset{¯}{z}_{\dot{\alpha }}\stackrel{}{}z_{\dot{\alpha }}`$, we consider the functions of elements of Majorana $`z`$-spinor $`Z_M`$ instead of $`Z_D`$.
The formulas of reduction on the compact $`SU(2)`$-subgroup have the form
$$T_{(j_1,j_2)}=\underset{j=|j_1j_2|}{\overset{j_1+j_2}{}}T_j,T_{[k,\rho ]}=\underset{j=k}{\overset{\mathrm{}}{}}T_j$$
(223)
for finite-dimensional nonunitary irreps and infinite-dimensional unitary irreps of $`SL(2,C)`$ respectively . Analogously with $`2+1`$ case, there are two types of the Poincaré group representations describing the same spin $`s`$. These types correspond to finite-dimensional and infinite-dimensional unitary representations of the Lorentz group. In particular, one may choose: (i) $`s=j_{\mathrm{max}}=j_1+j_2`$ for nonunitary finite-dimensional irreps $`(j_1,j_2)`$; (ii) $`s=j_{\mathrm{min}}=j_0=|j_1j_2|`$ for unitary infinite-dimensional irreps $`[j_0,\rho ]`$, where $`j_{\mathrm{max}}`$ and $`j_{\mathrm{min}}`$ are maximal and minimal $`j`$ in the decomposition (223) of an irrep of the Lorentz group over irreps $`T_j`$ of compact $`SU(2)`$ subgroup. Below we will study only the case of finite-dimensional representations of the Lorentz group.
Consider monomial basis
$$(z^1)^a(z^2)^b(\stackrel{}{}\underset{¯}{z}_1)^c(\stackrel{}{}\underset{¯}{z}_2)^d$$
in the space of functions $`\varphi (z,\stackrel{}{}\underset{¯}{z})`$. The values $`j_1=(a+b)/2`$ and $`j_2=(c+d)/2`$ are conserved under the action of generators (215). Therefore the space of irrep $`(j_1,j_2)`$ is the space of homogeneous analytic functions depending on two pairs of complex variables of power $`(2j_1,2j_2)`$. We denote these functions as $`\varphi _{j_1j_2}(z,\stackrel{}{}\underset{¯}{z})`$.
For finite-dimensional nonunitary irreps of $`SL(2,C)`$ $`a,b,c,d`$ are integer nonnegative, therefore $`j_1,j_2`$ are integer or half-integer nonnegative numbers. One can write functions $`f_s(x,z,\stackrel{}{}\underset{¯}{z})`$, which are polynomial of the power $`2s=2j_1+2j_2`$ in $`z,\stackrel{}{}\underset{¯}{z}`$, in the form
$$f_s(x,z,\stackrel{}{}\underset{¯}{z})=\underset{j_1+j_2=s}{}\underset{m_1,m_2}{}\psi _{j_1j_2}^{m_1m_2}(x)\phi _{j_1j_2}^{m_1m_2}(z,\stackrel{}{}\underset{¯}{z}),$$
(224)
where functions
$`\phi _{j_1j_2}^{m_1m_2}(z,\stackrel{}{}\underset{¯}{z})=N^{\frac{1}{2}}(z^1)^{j_1+m_1}(z^2)^{j_1m_1}(\stackrel{}{}\underset{¯}{z}_{\dot{1}})^{j_2+m_2}(\stackrel{}{}\underset{¯}{z}_{\dot{2}})^{j_2m_2},`$ (225)
$`N=(2s)![(j_1+m_1)!(j_1m_1)!(j_2+m_2)!(j_2m_2)!]^1,`$ (226)
form basis of the irrep of the Lorentz group. This basis corresponds to chiral representation (see Appendix B). On the other hand, one can write the decomposition of the same function in terms of symmetric multispinors $`\psi _{\alpha _1\mathrm{}\alpha _{2j_1}}^{\dot{\beta }_1\mathrm{}\dot{\beta }_{2j_2}}(x)=\psi _{\alpha _{(1}\mathrm{}\alpha _{2j_1)}}^{\dot{\beta }_{(1}\mathrm{}\dot{\beta }_{2j_2)}}(x)`$:
$$f_s(x,z,\stackrel{}{}\underset{¯}{z})=\underset{j_1+j_2=s}{}f_{j_1j_2}(x,z,\stackrel{}{}\underset{¯}{z}),f_{j_1j_2}(x,z,\stackrel{}{}\underset{¯}{z})=\psi _{\alpha _1\mathrm{}\alpha _{2j_1}}^{\dot{\beta }_1\mathrm{}\dot{\beta }_{2j_2}}(x)z^{\alpha _1}\mathrm{}z^{\alpha _{2j_1}}\stackrel{}{}\underset{¯}{z}_{\dot{\beta }_1}\mathrm{}\stackrel{}{}\underset{¯}{z}_{\dot{\beta }_{2j_2}}.$$
(227)
Notice that similar generating functions summed over all $`s`$ have been used in to describe massless fields. Comparing decompositions (224) and (227), we obtain the relation
$$N^{\frac{1}{2}}\psi _{j_1j_2}^{m_1m_2}(x)=\psi _{\underset{j_1+m_1}{\underset{}{1\mathrm{}\mathrm{\hspace{0.17em}1}}}\underset{j_1m_1}{\underset{}{2\mathrm{}\mathrm{\hspace{0.17em}2}}}}^{\stackrel{j_2+m_2}{\stackrel{}{\dot{1}\mathrm{}\dot{1}}}\stackrel{j_2m_2}{\stackrel{}{\dot{2}\mathrm{}\dot{2}}}}(x).$$
(228)
Using invariant tensor $`\sigma _{\dot{\alpha }\alpha }^\mu `$ and spinors $`z^\alpha `$, $`\stackrel{}{}\underset{¯}{z}_{\dot{\alpha }}`$, $`_\alpha =/z^\alpha `$, $`\underset{¯}{}^{\dot{\alpha }}=/\stackrel{}{}\underset{¯}{z}_{\dot{\alpha }}`$, it is possible to construct just four vectors:
$`\widehat{V}_{12}^\mu ={\displaystyle \frac{1}{2}}\overline{\sigma }^{\mu \dot{\alpha }\alpha }\stackrel{}{}\underset{¯}{z}_{\dot{\alpha }}_\alpha ,\widehat{V}_{21}^\mu ={\displaystyle \frac{1}{2}}\sigma _{}^{\mu }{}_{\alpha \dot{\alpha }}{}^{}z^\alpha \underset{¯}{}^{\dot{\alpha }},`$ (229)
$`\widehat{V}_{11}^\mu ={\displaystyle \frac{1}{2}}\sigma _{\alpha \dot{\alpha }}^\mu z^\alpha \stackrel{}{}\underset{¯}{z}^{\dot{\alpha }},\widehat{V}_{22}^\mu ={\displaystyle \frac{1}{2}}\sigma _{\alpha \dot{\alpha }}^\mu ^\alpha \underset{¯}{}^{\dot{\alpha }}.`$ (230)
These operators are not functions of generators of $`M(3,1)`$ and interlock irreps with different $`(j_1,j_2)`$, however, as we will see below, play an impotent role in the theory of RWE.
### B Relativistic wave equations, invariant under proper Poincaré group
Let us fix eigenvalues of the Casimir operators of the Poincaré group and of the Lorentz subgroup:
$`\widehat{p}^2f(x,𝐳)=m^2f(x,𝐳),`$ (231)
$`\widehat{W}^2f(x,𝐳)=s(s+1)m^2f(x,𝐳),`$ (232)
$`\widehat{𝐌}^2f(x,𝐳)=j_1(j_1+1)f(x,𝐳),`$ (233)
$`\widehat{\overline{𝐌}}^2f(x,𝐳)=j_2(j_2+1)f(x,𝐳),`$ (234)
Spectrum (232) of the operator $`\widehat{W}^2`$ corresponds to the consideration of massive spin $`s`$ particles and massless particles with discrete spin. (For tachyons and massless particles with continuous spin spectrum differ from (232), see .) As a consequence of two last equations (remind that we consider the space of analytic functions of $`z,\stackrel{}{}\underset{¯}{z}`$) we obtain that eigenvalues of the belonging to the complete set (206) operators $`\widehat{S}_3^R`$ and $`\widehat{B}_3^R`$ are also fixed,
$$\widehat{S}_3^Rf(x,z,\stackrel{}{}\underset{¯}{z})=(j_1+j_2)f(x,z,\stackrel{}{}\underset{¯}{z}),i\widehat{B}_3^Rf(x,z,\stackrel{}{}\underset{¯}{z})=(j_1j_2)f(x,z,\stackrel{}{}\underset{¯}{z}).$$
(235)
Equations (231)-(234) define reducible representation of the proper Poincaré group $`M(3,1)`$. This representation splits into two representations labelled by the sign of $`p_0`$, which are irreducible for $`m0`$.
Nonequivalent representations are distinguished by eigenvalues of the Casimir operators $`\widehat{p}^2`$, $`\widehat{W}^2`$ and by the sign of $`p_0`$ (see also ). The case of zero eigenvalues of the operators $`\widehat{p}^2`$ and $`\widehat{W}^2`$ is an exception. This case corresponds to massless particles with discrete spin, and nonequivalent irreps are labelled by the helicity and by the sign of $`p_0`$. On the other hand, one can introduce a chirality as $`\lambda =j_1j_2`$ (or as the difference of numbers of dotted and undotted indices). The explicit form of the chirality operator in the space of analytic functions of $`z,\stackrel{}{}\underset{¯}{z}`$ is given by the formula (see (B4))
$$\widehat{\mathrm{\Gamma }}^5=\frac{1}{2}\left(z^\alpha _\alpha \stackrel{}{}\underset{¯}{z}_{\dot{\alpha }}\underset{¯}{}^{\dot{\alpha }}\right).$$
(236)
In the massless case helicity is equal to chirality up to sign . In the massive case irreps of the proper Poincaré group, which are labelled by the same $`m,s,signp_0`$ but by different chiralities, are equivalent. Thus, for fixed mass $`m`$ and spin $`s=j_1+j_2`$ the system (231)-(234) has $`2s+1`$ solutions differed by $`\lambda =j_1j_2`$.
Using (217), we rewrite the Lubanski-Pauli vector (204) and the Casimir operator $`\widehat{W}^2`$ in the form
$`\widehat{W}^\mu ={\displaystyle \frac{1}{2}}ϵ^{\mu \nu \rho \sigma }\widehat{p}_\nu \widehat{S}_{\rho \sigma }={\displaystyle \frac{1}{2}}i\widehat{p}_\nu \left((\sigma ^{\mu \nu })_{\alpha \beta }\widehat{M}^{\alpha \beta }(\overline{\sigma }^{\mu \nu })_{\dot{\alpha }\dot{\beta }}\widehat{\overline{M}}^{\dot{\alpha }\dot{\beta }}\right),`$ (237)
$`\widehat{W}^2=\widehat{p}^2(\widehat{𝐌}^2+\widehat{\overline{𝐌}}^2){\displaystyle \frac{1}{2}}\widehat{p}_\mu \widehat{p}_\nu (\sigma ^{\mu \rho })_{\alpha \beta }(\overline{\sigma }_\rho ^\nu )_{\dot{\alpha }\dot{\beta }}\widehat{M}^{\alpha \beta }\widehat{\overline{M}}^{\dot{\alpha }\dot{\beta }}.`$ (238)
Taking into account the explicit form of spin operators (219) and symmetry of $`(\sigma ^{\mu \nu })_{\alpha \beta }`$ and $`(\overline{\sigma }^{\mu \nu })_{\dot{\alpha }\dot{\beta }}`$ with respect to permutation of spinor indices, we rewrite the last relation as
$$\widehat{W}^2=\widehat{p}^2(\widehat{𝐌}^2+\widehat{\overline{𝐌}}^2)2\widehat{p}_\mu \widehat{p}_\nu (\sigma ^{\mu \rho })_{\alpha \beta }(\overline{\sigma }_\rho ^\nu )_{\dot{\alpha }\dot{\beta }}z^\alpha ^\beta \stackrel{}{}\underset{¯}{z}^{\dot{\alpha }}\underset{¯}{}^{\dot{\beta }}.$$
Finally, using the identity
$$4(\sigma ^{\mu \rho })_{\alpha \beta }(\overline{\sigma }_\rho ^\nu )_{\dot{\alpha }\dot{\beta }}=\eta ^{\mu \nu }ϵ_{\alpha \beta }ϵ_{\dot{\alpha }\dot{\beta }}+\sigma _{\alpha \dot{\alpha }}^{}{}_{}{}^{\mu }\sigma _{\beta \dot{\beta }}^{}{}_{}{}^{\nu }+\sigma _{\alpha \dot{\alpha }}^{}{}_{}{}^{\nu }\sigma _{\beta \dot{\beta }}^{}{}_{}{}^{\mu }$$
and the condition of mass irreducibility (231), we obtain
$$\widehat{W}^2=m^2(j_1+j_2)(j_1+j_2+1)+4\widehat{p}_\mu \widehat{V}_{11}^\mu \widehat{p}_\nu \widehat{V}_{22}^\nu ,$$
(239)
where operators $`\widehat{V}_{11}^\mu `$ and $`\widehat{V}_{22}^\mu `$ are defined in (230). Therefore for $`s=j_1+j_2`$ the necessary and sufficient condition of spin irreducibility is
$$\widehat{p}_\mu \widehat{V}_{11}^\mu \widehat{p}_\nu \widehat{V}_{22}^\nu f(x,z,\stackrel{}{}\underset{¯}{z})=0.$$
(240)
For the representations $`(s\mathrm{\hspace{0.17em}0})`$ and $`(0s)`$ we have $`\widehat{V}_{22}^\mu f(x,z,\stackrel{}{}\underset{¯}{z})=0`$ and the condition (240) is fulfilled identically. In general case, observing that in momentum representation an action of operator $`\widehat{V}_{11}^\mu \widehat{p}_\mu `$ reduces to multiplication by the number $`p_\mu \sigma _{\alpha \dot{\alpha }}^\mu z^\alpha z^{\dot{\alpha }}`$, we come to the alternative conditions:
$`\widehat{p}_\mu \widehat{V}_{11}^\mu =0,`$ (241)
$`\widehat{p}_\nu \widehat{V}_{22}^\nu f(x,z,\stackrel{}{}\underset{¯}{z})=0.`$ (242)
The first condition connects the components of momentum $`p_\mu `$ and complex spin variables $`q_\mu =\sigma _{\mu \alpha \dot{\alpha }}z^\alpha \stackrel{}{}\underset{¯}{z}^{\dot{\alpha }}/2`$, $`q_\mu q^\mu =0`$. Thus we have the space of functions of two 4-vectors $`p_\mu `$, $`q_\mu `$, which are subject to the invariant constraints
$$p^2=m^2,p_\mu q^\mu =0,q^2=0.$$
(243)
According to (243), in the rest frame we get $`z^1\stackrel{}{}\underset{¯}{z}^1+z^2\stackrel{}{}\underset{¯}{z}^2=0`$. Similar approach to the constructing of wave functions describing the elementary particles was suggested by E.Wigner in , where discussion was restricted to particles of integer spin and real $`q_\mu `$ with constraints $`p^2=m^2`$, $`p_\mu q^\mu =0`$, $`q^2=1`$. Different generalizations of the approach were considered later in .
The second condition (242) does not affect spin variables and can be written in terms of $`\psi (x)`$. Really, for fixed $`j_1,j_2`$, using the decomposition (227) of $`f(x,z,\stackrel{}{}\underset{¯}{z})`$ in terms of multispinors and also the relation $`_\alpha \psi _{\alpha _1\alpha _2\mathrm{}\alpha _k}z^{\alpha _{(1}}z^{\alpha _2}\mathrm{}z^{\alpha _{k)}}=\delta _{\alpha }^{}{}_{}{}^{\alpha _1}\psi _{\alpha _1\alpha _2\mathrm{}\alpha _k}z^{\alpha _{(2}}\mathrm{}z^{\alpha _{k)}}`$, one can rewrite the system
$$(\widehat{p}^2m^2)f_{j_1j_2}(x,z)=0,\widehat{p}_\mu \sigma _{}^{\mu }{}_{\alpha \dot{\alpha }}{}^{}^\alpha \underset{¯}{}^{\dot{\alpha }}f_{j_1j_2}(x,z)=0$$
(244)
in the form
$`(\widehat{p}^2m^2)\psi _{\alpha _1\mathrm{}\alpha _k\dot{\alpha }_1\mathrm{}\dot{\alpha }_l}(x)=0,`$ (245)
$`^{\dot{\alpha }\alpha }\psi _{\alpha \alpha _1\mathrm{}\alpha _{k1}\dot{\alpha }\dot{\alpha }_1\mathrm{}\dot{\alpha }_{l1}}(x)=0,`$ (246)
where $`^{\dot{\alpha }\alpha }=_\mu \overline{\sigma }^{\mu \dot{\alpha }\alpha }`$, $`k=2j_1`$, $`l=2j_2`$. These equations describe a particle with unique mass $`m`$ and spin $`s=j_1+j_2`$. Subsidiary condition (246) is necessary to exclude components corresponding to other possible spins $`s`$, $`|j_1j_2|s<j_1+j_2`$, see (223).
On the other hand, in order to describe spin $`s`$ one may use representations $`(j_1j_2)`$, $`j_1+j_2s`$. In this case, according to (239), the condition (246) should be changed by new subsidiary condition
$$_{\beta \dot{\beta }}^{\dot{\alpha }\alpha }\psi _{\alpha \alpha _1\mathrm{}\alpha _{k1}\dot{\alpha }\dot{\alpha }_1\mathrm{}\dot{\alpha }_{l1}}(x)=m^2[(j_1+j_2)(j_1+j_2+1)s(s+1)]\psi _{\beta \alpha _1\mathrm{}\alpha _{k1}\dot{\beta }\dot{\alpha }_1\mathrm{}\dot{\alpha }_{l1}}(x).$$
(247)
Notice that an approach using this general subsidiary conditions was not considered earlier.
Passing on to vector indices, one can see that for integer spins and irreps $`(\frac{s}{2}\frac{s}{2})`$ equations (245), (246) take the form
$$(\widehat{p}^2m^2)\mathrm{\Phi }_{\mu _1\mu _2\mathrm{}\mu _s}(x)=0,^\mu \mathrm{\Phi }_{\mu \mu _2\mathrm{}\mu _s}(x)=0,\mathrm{\Phi }_{}^{\mu }{}_{\mu \mathrm{}\mu _s}{}^{}(x)=0,$$
(248)
where
$$\mathrm{\Phi }_{\mu _1\mu _2\mathrm{}\mu _s}(x)=(1)^s2^s\overline{\sigma }_{\mu _1}^{\dot{\alpha }_1\alpha _1}\mathrm{}\overline{\sigma }_{\mu _s}^{\dot{\alpha }_s\alpha _s}\psi _{\alpha _1\mathrm{}\alpha _s\dot{\alpha }_1\mathrm{}\dot{\alpha }_s}(x).$$
Just equations (248) known also as massive tensor field equations or Fierz–Pauli equations are used most often to describe integer spins.
For half-integer spins and irreps $`(\frac{2s\pm 1}{4}\frac{2s1}{4})`$ after passage to vector indices subsidiary conditions (246) take the form
$`^\mu \mathrm{\Psi }_{\mu \mu _2\mathrm{}\mu _n\alpha }(x)=0,\overline{\sigma }^{\mu \dot{\alpha }\alpha }\mathrm{\Psi }_{\mu \mu _2\mathrm{}\mu _n\alpha }(x)=0,\mathrm{\Psi }_{}^{\mu }{}_{\mu \mu _2\mathrm{}\mu _n\alpha }{}^{}(x)=0,`$ (249)
$`^\mu \mathrm{\Psi }_{\mu \mu _2\mathrm{}\mu _n\dot{\alpha }}(x)=0,\sigma _{}^{\mu }{}_{\alpha \dot{\alpha }}{}^{}\mathrm{\Psi }_{\mu \mu _2\mathrm{}\mu _n}^{}{}_{}{}^{\dot{\alpha }}(x)=0,\mathrm{\Psi }_{}^{\mu }{}_{\mu \mu _2\mathrm{}\mu _n\dot{\alpha }}{}^{}(x)=0,`$ (250)
where $`n=(2s1)/2`$.
### C Relativistic wave equations, invariant under improper Poincaré group
Improper Poincaré group includes continuous transformations of the proper group and space reflection operator (parity operator) $`\widehat{I}_P`$. According to (86),(208) this operator obeys the condition $`\widehat{I}_P^2=\widehat{1}`$ and commutation relations
$`[\widehat{I}_P,\widehat{p}_0]=[\widehat{I}_P,\widehat{p}^2]=[\widehat{I}_P,\widehat{W}^2]=[\widehat{I}_P,\widehat{S}_k]=[\widehat{I}_P,\widehat{S}_k^R]=0,`$ (251)
$`[\widehat{I}_P,\widehat{p}_k]_+=[\widehat{I}_P,\widehat{B}_k]_+=[\widehat{I}_P,\widehat{B}_k^R]_+=0.`$ (252)
States with definite total parity are defined as eigenfunctions of operator $`\widehat{I}_P`$:
$$\widehat{I}_Pf(x,𝐳)=\pm f(x,𝐳).$$
(253)
For $`m>0`$ irreps of the improper Poincaré group are labelled by an orbit defining the mass $`m`$ and the sign of $`p_0`$, and by irrep of the little group $`O(3)`$ defining spin $`s`$ and intrinsic parity . In the rest frame the intrinsic parity coincides with the total.
The Casimir operators of the Lorentz group do not commute with parity operator, $`[\widehat{I}_P,\widehat{𝐌}^2]=\widehat{\overline{𝐌}}^2,[\widehat{I}_P,\widehat{\overline{𝐌}}^2]=\widehat{𝐌}^2`$, and parity transformation combines two labelled by Lorentz indices $`(j_1,j_2)`$ and $`(j_2,j_1)`$ (by chiralities $`\pm \lambda `$) equivalent irreps of the proper Poincaré group into one representation of the improper group. The latter representation is reductible and splits into two irreps differed by intrinsic parity $`\eta =\pm 1`$. Thus we can’t make use the operators $`\widehat{𝐌}^2`$, $`\widehat{\overline{𝐌}}^2`$ to select invariant subspaces, and instead of the set of eight commuting operators
$$\widehat{p}_\mu ,\widehat{W}^2,\widehat{𝐩}\widehat{𝐒},\widehat{𝐌}^2,\widehat{\overline{𝐌}}^2$$
(254)
used above in order to construct the system (231)-(234) we should consider an another set. Notice that parity operator $`\widehat{I}_P`$ can’t be used directly for identification of invariant subspaces since according to (252) it does not commute with translations and boosts.
The simplest possibility is to consider a system
$`\widehat{p}^2f(x,z,\stackrel{}{}\underset{¯}{z})=m^2f(x,z,\stackrel{}{}\underset{¯}{z}),`$ (255)
$`\widehat{W}^2f(x,z,\stackrel{}{}\underset{¯}{z})=s(s+1)m^2f(x,z,\stackrel{}{}\underset{¯}{z}),`$ (256)
$`\widehat{S}_3^Rf(x,z,\stackrel{}{}\underset{¯}{z})=sf(x,z,\stackrel{}{}\underset{¯}{z}).`$ (257)
The last equation fixes the power $`2s=2(j_1+j_2)`$ of the polynomial in $`z,\stackrel{}{}\underset{¯}{z}`$, see (235). The first two equations are the conditions of mass and spin irreducibility. Therefore the system describes fixed mass and spin, but the Poincaré group representation defined by this system is reducible. This representation decomposes into $`2(2s+1)`$ irreps differed by the chirality $`\lambda =s,\mathrm{},s`$ and sign of $`p_0`$.
Supplementing the system (255)-(257) by the equation
$$i\widehat{B}_3^Rf(x,z,\stackrel{}{}\underset{¯}{z})=\pm (j_1j_2)f(x,z,\stackrel{}{}\underset{¯}{z}),$$
(258)
which change the sign under space reflection, it is possible to extract components corresponding to the representation $`(j_1,j_2)(j_2,j_1)`$. If we consider only the components labelled by $`(j_1,j_2)`$ and $`(j_2,j_1)`$, then for $`j_1j_2`$ mass and spin irreducibility conditions (255),(256) leave $`4(2s+1)`$ independent components corresponding to the direct sum of four improper Poincaré group irreps differed by signs of energy $`p_0`$ and intrinsic parity $`\eta `$. But states with definite intrinsic parity arise in such an approach only as linear combinations of the solutions of two systems (255)-(258) with different sign in the last equation (i.e. solutions with fixed chirality).
Let us investigate the possibility to construct the system of equations, which remains invariant under space reflection and has solutions with definite intrinsic parity. For this purpose it is necessary to consider equations, which combine labelled by different chiralities $`\lambda =j_1j_2`$ equivalent irreps of the proper Poincaré group. In the other words, it is necessary to consider supplementary operators, which define some extension of the Lorentz group. These operators, replacing $`\widehat{𝐌}^2`$ and $`\widehat{\overline{𝐌}}^2`$ in the set (254), must commute with all the left generators of the proper Poincaré group and with parity operator $`\widehat{I}_P`$. We suppose that one of these commuting operators is linear in $`\widehat{p}`$.
A general form of the invariant equations linear in $`\widehat{p}`$ is
$$\widehat{p}_\mu \widehat{V}^\mu f(x,𝐳)=\varkappa f(x,𝐳),$$
(259)
where $`\widehat{V}^\mu `$ is a transforming as four-vector function of $`𝐳`$ and $`/𝐳`$.
The introduced above vector operators $`V_{ik}^\mu `$ (229), (230) interlock irreps with different $`(j_1,j_2)`$. Operators $`\widehat{V}_{12}^\mu `$, $`\widehat{V}_{21}^\mu `$ conserve $`j_1+j_2`$, and operators $`\widehat{V}_{11}^\mu `$, $`\widehat{V}_{22}^\mu `$ conserve $`j_1j_2`$. Any of four connecting two scalar functions relations
$`\widehat{p}_\mu \widehat{V}_{12}^\mu f_{j_1,j_2}(x,z,\stackrel{}{}\underset{¯}{z})=\varkappa _{12}f_{j_1\frac{1}{2},j_2+\frac{1}{2}}(x,z,\stackrel{}{}\underset{¯}{z}),\widehat{p}_\mu \widehat{V}_{21}^\mu f_{j_1,j_2}(x,z,\stackrel{}{}\underset{¯}{z})=\varkappa _{21}f_{j_1+\frac{1}{2},j_2\frac{1}{2}}(x,z,\stackrel{}{}\underset{¯}{z}),`$ (260)
$`\widehat{p}_\mu \widehat{V}_{11}^\mu f_{j_1,j_2}(x,z,\stackrel{}{}\underset{¯}{z})=\varkappa _{11}f_{j_1+\frac{1}{2},j_2+\frac{1}{2}}(x,z,\stackrel{}{}\underset{¯}{z}),\widehat{p}_\mu \widehat{V}_{22}^\mu f_{j_1,j_2}(x,z,\stackrel{}{}\underset{¯}{z})=\varkappa _{22}f_{j_1\frac{1}{2},j_2\frac{1}{2}}(x,z,\stackrel{}{}\underset{¯}{z}),`$ (261)
one may consider as a RWE. Thus the operator $`\widehat{V}^\mu `$ in (259) is a linear combination of $`\widehat{V}_{ik}^\mu `$.
Let us consider finite-component equations invariant with respect to space reflection. This means:
1. The operator $`\widehat{p}_\mu \widehat{V}^\mu `$ is invariant under space reflection.
2. The equation has solutions $`f(x,z,\stackrel{}{}\underset{¯}{z})=\psi _n(x)\varphi _n(z,\stackrel{}{}\underset{¯}{z})`$, where functions $`\varphi _n(z)`$ carry a representation containing finite number of irreps $`(j_1,j_2)`$.
It is easy to see that at $`\varkappa _{11}0`$ operator $`\widehat{V}_{11}^\mu `$ can’t be contained in $`\widehat{V}^\mu `$. In this case at $`\varkappa _{22}0`$ one can separate from the system of equations for functions $`f_{j_1,j_2}(x,z,\stackrel{}{}\underset{¯}{z})`$, $`f(x,z,\stackrel{}{}\underset{¯}{z})=f_{j_1,j_2}(x,z,\stackrel{}{}\underset{¯}{z})`$ the independent equation for the characterized by maximal $`j_1+j_2`$ function, which does not contain $`\widehat{V}_{22}^\mu `$. (Besides, it is not necessary to use operators $`\widehat{V}_{11}^\mu `$ and $`\widehat{V}_{22}^\mu `$ since these operators leave $`j_1j_2`$ invariable and can’t connect irreps with different $`\lambda `$.)
Relating to operators $`\widehat{V}_{12}^\mu `$ and $`\widehat{V}_{21}^\mu `$, one can see that only the combination $`\widehat{p}_\mu \widehat{\mathrm{\Gamma }}^\mu `$,
$$\widehat{\mathrm{\Gamma }}^\mu =\widehat{V}_{12}^\mu +\widehat{V}_{21}^\mu =\frac{1}{2}\left(\overline{\sigma }^{\mu \dot{\alpha }\alpha }\stackrel{}{}\underset{¯}{z}_{\dot{\alpha }}_\alpha +\sigma _{}^{\mu }{}_{\alpha \dot{\alpha }}{}^{}z^\alpha \underset{¯}{}^{\dot{\alpha }}\right),$$
(262)
is invariant under space reflections. Operators $`\widehat{\mathrm{\Gamma }}^\mu `$ connect representation $`(j_1j_2)`$ with $`(j_1+1j_21)`$ and $`(j_11j_2+1)`$ and conserve $`j_1+j_2`$. These operators obey the commutation relations
$`[\widehat{S}^{\lambda \mu },\widehat{\mathrm{\Gamma }}^\nu ]=i(\eta ^{\mu \nu }\widehat{\mathrm{\Gamma }}^\lambda \eta ^{\lambda \nu }\widehat{\mathrm{\Gamma }}^\mu ),`$ (263)
$`[\widehat{\mathrm{\Gamma }}^\mu ,\widehat{\mathrm{\Gamma }}^\nu ]=i\widehat{S}^{\mu \nu },`$ (264)
which coincide with the commutation relations of matrices $`\gamma ^\mu /2`$. The explicit calculation shows that $`\widehat{\mathrm{\Gamma }}_\mu \widehat{\mathrm{\Gamma }}^\mu `$ depends on irrep of the Lorentz subgroup,
$$\widehat{\mathrm{\Gamma }}_\mu \widehat{\mathrm{\Gamma }}^\mu =2j_1+2j_2+4j_1j_2.$$
(265)
Supplementing generators of the Lorentz group by four operators
$$\widehat{S}^{4\mu }=\widehat{\mathrm{\Gamma }}^\mu ,\widehat{S}^{ab}=\widehat{S}^{ba},$$
(266)
we obtain
$$[\widehat{S}^{ab},\widehat{S}^{cd}]=i(\eta ^{bc}\widehat{S}^{ad}\eta ^{ac}\widehat{S}^{bd}\eta ^{bd}\widehat{S}^{ac}+\eta ^{ad}\widehat{S}^{bc}),\eta ^{44}=\eta ^{00}=1.$$
(267)
Thus operators $`\widehat{S}^{ab}`$, $`a,b=0,1,2,3,4`$, obey the commutation relations of the generators of the 3+2 de Sitter group $`SO(3,2)Sp(4,R)`$. This group has two fundamental irreps, namely four-dimensional spinor irrep $`T_{[10]}`$ (by matrices $`Sp(4,R)`$) and five-dimensional vector irrep $`T_{[01]}`$ (by matrices $`SO(3,2)`$).
Using (205), (222) and (265), we obtain for the second order Casimir operator of the group $`Sp(4,R)`$
$$\widehat{S}_{ab}\widehat{S}^{ab}f(x,z,\stackrel{}{}\underset{¯}{z})=4S(S+2)f(x,z,\stackrel{}{}\underset{¯}{z}),S=j_1+j_2.$$
Thus we deal with symmetric representations of $`Sp(4,R)`$, which we denote as $`T_{[2S\mathrm{\hspace{0.17em}0}]}`$ (see Appendix). These irreps can be obtained as a symmetric term in the decomposition of the direct product $`(T_{[\mathrm{1\hspace{0.17em}0}]})^{2S}`$. Irreps $`T_{[2S\mathrm{\hspace{0.17em}0}]}`$ characterized by dimensionality $`(2S+3)!/(6(2S)!)`$ combines all finite-dimensional irreps of the Lorentz group with $`j_1+j_2=S`$.
However, it is obvious that the equation
$$\widehat{p}_\mu \widehat{\mathrm{\Gamma }}^\mu f(x,z,\stackrel{}{}\underset{¯}{z})=\varkappa f(x,z,\stackrel{}{}\underset{¯}{z})$$
(268)
by itself does not fix spin $`s`$ and mass $`m`$, defined by (231) and (232), or the power $`j_1+j_2`$ of the $`f(x,z,\stackrel{}{}\underset{¯}{z})`$ in $`z,\stackrel{}{}\underset{¯}{z}`$. In the rest frame it is easy to see that even for fixed $`S=j_1+j_2`$ this equation fix only product $`ms=\varkappa `$, $`sS`$.
Let us consider the set of eight commuting operators
$$\widehat{p}_\mu ,\widehat{W}^2,\widehat{𝐩}\widehat{𝐒}(\mathrm{or}\widehat{S}_3\text{in the rest frame}),\widehat{p}_\mu \widehat{\mathrm{\Gamma }}^\mu ,\widehat{S}_{ab}\widehat{S}^{ab}$$
(269)
acting in the space of functions of eight variables $`x^\mu ,z^\alpha ,\stackrel{}{}\underset{¯}{z}_{\dot{\alpha }}`$. In compare with the set (254) we have replaced two right operators $`𝐌^2`$, $`\widehat{\overline{𝐌}}^2`$ by invariant under parity transformation operators $`\widehat{p}_\mu \widehat{\mathrm{\Gamma }}^\mu `$, $`\widehat{S}_{ab}\widehat{S}^{ab}`$. Notice that instead of $`\widehat{S}_{ab}\widehat{S}^{ab}`$ one can use operator $`\widehat{S}_3^R`$ with eigenvalues equal to the minus power of polynomial in $`z,\stackrel{}{}\underset{¯}{z}`$, see (257).
Invariant subspaces are labelled by eigenvalues of operators
$`\widehat{p}^2f(x,z,\stackrel{}{}\underset{¯}{z})=m^2f(x,z,\stackrel{}{}\underset{¯}{z}),`$ (270)
$`\widehat{W}^2f(x,z,\stackrel{}{}\underset{¯}{z})=m^2s(s+1)f(x,z,\stackrel{}{}\underset{¯}{z}),`$ (271)
$`\widehat{p}_\mu \widehat{\mathrm{\Gamma }}^\mu f(x,z,\stackrel{}{}\underset{¯}{z})=m\stackrel{~}{s}f(x,z,\stackrel{}{}\underset{¯}{z}),`$ (272)
$`\widehat{S}_{ab}\widehat{S}^{ab}f(x,z,\stackrel{}{}\underset{¯}{z})=4S(S+2)f(x,z,\stackrel{}{}\underset{¯}{z}).`$ (273)
Unlike equations (233),(234), which fix $`j_1`$ and $`j_2`$ separately, the last equation of the system fixes irrep $`T_{[2S\mathrm{\hspace{0.17em}0}]}`$ of the 3+2 de Sitter group and therefore the power $`2S=2j_1+2j_2`$ of the polynomial in $`z,\stackrel{}{}\underset{¯}{z}`$. Irreps of the Poincaré group characterized by spin $`sS`$ can be realized in the space of these polynomials.
In the rest frame
$`\widehat{p}_0^2f(x,z,\stackrel{}{}\underset{¯}{z})=m^2f(x,z,\stackrel{}{}\underset{¯}{z}),`$ (274)
$`\widehat{p}_0\widehat{\mathrm{\Gamma }}^0f(x,z,\stackrel{}{}\underset{¯}{z})=m\stackrel{~}{s}f(x,z,\stackrel{}{}\underset{¯}{z}),\widehat{\mathrm{\Gamma }}^0={\displaystyle \frac{1}{2}}(\sigma ^{0\dot{\alpha }\alpha }\stackrel{}{}\underset{¯}{z}_{\dot{\alpha }}_\alpha +\sigma _{\alpha \dot{\alpha }}^0z^\alpha \underset{¯}{}^{\dot{\alpha }}).`$ (275)
According to the first equation, $`p_0=\pm m`$ and correspondingly $`\stackrel{~}{s}`$ is a product of eigenvalue of operator $`\widehat{\mathrm{\Gamma }}^0`$ and $`signp_0`$. For $`p_0=m`$, any characterized by $`n_1n_2=2s`$ function is the solution of equation (275), where $`n_1`$ is the power of homogeneity in the variables $`(z^1+\stackrel{}{}\underset{¯}{z}_{\dot{1}})`$, $`(z^2+\stackrel{}{}\underset{¯}{z}_{\dot{2}})`$, and $`n_2`$ is the power of homogeneity in the variables $`(z^1\stackrel{}{}\underset{¯}{z}_{\dot{1}})`$, $`(z^2\stackrel{}{}\underset{¯}{z}_{\dot{2}})`$. Therefore, for $`p_0=m`$ any characterized by $`n_1n_2=2s`$ function is the solution of equation (275).
Let us show that the relation
$$|\stackrel{~}{s}|sS$$
(276)
takes place. Variables $`z^\alpha `$ and $`z_{\dot{\alpha }}`$ have the same transformation rule under space rotations. Thus, the pairs $`(z^1+\stackrel{}{}\underset{¯}{z}_{\dot{1}})`$, $`(z^2+\stackrel{}{}\underset{¯}{z}_{\dot{2}})`$ and $`(z^1\stackrel{}{}\underset{¯}{z}_{\dot{1}})`$, $`(z^2\stackrel{}{}\underset{¯}{z}_{\dot{2}})`$ are spinors of rotation group, but are characterized by opposite parity. The polynomials of power $`2j^{}`$ in the first pair of variables or $`2j^{\prime \prime }`$ in the second pair of variables transform under $`T_j^{}`$ or $`T_{j^{\prime \prime }}`$ of the rotation group. At fixed $`j^{}`$ and $`j^{\prime \prime }`$ the relation $`\stackrel{~}{s}=(j^{}j^{\prime \prime })signp_0`$ takes place, and the space of polynomials of the power $`2S=2j^{}+2j^{\prime \prime }`$ corresponds to direct product of the representations $`T_j^{}`$ and $`T_{j^{\prime \prime }}`$. This direct product decomposes into sum of irreps, labelled by $`s=|j^{}j^{\prime \prime }|,\mathrm{},j^{}+j^{\prime \prime }`$, and therefore spin $`s`$ runs from $`|\stackrel{~}{s}|`$ up to $`S`$.
In particular, for $`|\stackrel{~}{s}|=S`$ the spin irreducibility condition (271) is a consequence of other equations of the system, and spin is equal to one half of the polynomial power. Below we restrict our consideration by this case, which allows to describe spin $`s`$ by means of the irrep of the 3+2 de Sitter group with minimal possible dimensionality. Correspondingly, for $`\stackrel{~}{s}=S`$ we will consider the system
$`\widehat{p}^2f(x,𝐳)=m^2f(x,𝐳),`$ (277)
$`\widehat{p}_\mu \widehat{\mathrm{\Gamma }}^\mu f(x,𝐳)=msf(x,𝐳),`$ (278)
$`\widehat{S}_{ab}\widehat{S}^{ab}f(x,𝐳)=4s(s+2)f(x,𝐳).`$ (279)
In the rest frame the general solution in the set of polynomial of the power $`2s`$ in $`z,\stackrel{}{}\underset{¯}{z}`$ has the form
$$f_{m,s}(x,𝐳)=\underset{s_3=s}{\overset{s}{}}C_{s_3}e^{imx^0}(z^1+\stackrel{}{}\underset{¯}{z}_{\dot{1}})^{s+s_3}(z^2+\stackrel{}{}\underset{¯}{z}_{\dot{2}})^{ss_3}+C_{s_3}^{}e^{imx^0}(z^1\stackrel{}{}\underset{¯}{z}_{\dot{1}})^{s+s_3}(z^2\stackrel{}{}\underset{¯}{z}_{\dot{2}})^{ss_3},$$
(280)
where $`s_3`$ is the spin projection,
$$\widehat{S}_3f(x,𝐳)=s_3f(x,𝐳),\widehat{S}_3=\frac{1}{2}(z^1_1+\stackrel{}{}\underset{¯}{z}_{\dot{1}}\underset{¯}{}^{\dot{1}}z^2_2\stackrel{}{}\underset{¯}{z}_{\dot{2}}\underset{¯}{}^{\dot{2}}).$$
(281)
Thus for unique mass $`m`$ and spin $`s`$ there are $`2s+1`$ independent positive-frequency solutions and $`2s+1`$ independent negative-frequency solutions belonging to two irreps of improper Poincaré group. In the case $`\stackrel{~}{s}=S`$, which corresponds to the change of sign in the equation (278), general solution is obtained from (280) by the substitution $`(z^\alpha +\stackrel{}{}\underset{¯}{z}_{\dot{\alpha }})(z^\alpha \stackrel{}{}\underset{¯}{z}_{\dot{\alpha }})`$. It follows from (280) that for half-integer spins the sign of $`\stackrel{~}{s}`$ is the product of $`signp_0`$ and intrinsic parity. According to (280), in the rest frame for half-integer spin positive-frequency and negative-frequency states are characterized by opposite parity. One can show (see ) that for fixed mass $`m`$ and representation $`(\frac{1}{2}\mathrm{\hspace{0.17em}0})(0\frac{1}{2})`$ of the Lorentz group this condition is sufficient to derive the Dirac equation.
Only corresponding to spin 1/2 four-dimensional irrep of the 3+2 de Sitter group remains irreducible under the reduction on the improper Lorentz group. For spin one 10-dimensional irrep splits into 6+4 (antisymmetric tensor and four-vector), for spin 3/2 20-dimensional irrep splits into 8+12, and so on.
Consider plain wave solutions corresponding to a moving along $`x^3`$ particle. They can be obtained from the solutions in the rest frame (280) by means of the Lorentz transformation
$$P=UP_0U^{},\mathrm{where}P_0=\pm diag\{m,m\},U=diag\{e^a,e^a\}SL(2,C),$$
where the sign corresponds to the sign of $`p_0`$,
$$p_\mu =k_\mu \mathrm{sign}p_0,k_0=m\mathrm{cosh}2a,k_3=m\mathrm{sinh}2a,e^{\pm a}=\sqrt{(k_0\pm k_3)/m}.$$
(282)
Thus it follows that
$`f_{m,s,s_3}^{}(x,𝐳)=`$ $`C_1e^{ik_0x^0+k_3x^3}(z^1e^a+\stackrel{}{}\underset{¯}{z}_{\dot{1}}e^a)^{s+s_3}(z^2e^a+\stackrel{}{}\underset{¯}{z}_{\dot{2}}e^a)^{ss_3}+`$ (283)
$`C_2e^{ik_0x^0k_3x^3}(z^1e^a\stackrel{}{}\underset{¯}{z}_{\dot{1}}e^a)^{s+s_3}(z^2e^a\stackrel{}{}\underset{¯}{z}_{\dot{2}}e^a)^{ss_3}.`$ (284)
In the ultrarelativistic case for positive $`a`$ it is convenient to rewrite (284) in the form
$`f_{m,s,s_3}(x,𝐳)=\left(\frac{k_0+k_3}{m}\right)^s\times `$ (285)
$`((C_1e^{ik_0x^0+k_0x^3}+C_2(1)^{ss_3}e^{ik_0x^0k_0x^3})(z^1)^{s+s_3}(\stackrel{}{}\underset{¯}{z}_{\dot{2}})^{ss_3}+O\left(\frac{k_0k_3}{k_0+k_3}\right)^{\frac{1}{2}})`$ (286)
The main term in (286) corresponds to functions carrying irrep $`(\frac{s+\lambda }{2}\frac{s\lambda }{2})`$, $`\lambda =s_3`$, of the Lorentz group. The contribution of other irreps $`(\frac{s+\lambda ^{}}{2}\frac{s\lambda ^{}}{2})`$ are damped by a factor $`(\frac{k_0k_3}{k_0+k_3})^{|\lambda \lambda ^{}|}`$. Passing to the limit $`a+\mathrm{}`$ (or $`m0`$), we obtain the states with certain chirality $`\lambda =j_1j_2=s_3`$ (for $`a\mathrm{}`$ with chirality $`\lambda =j_1j_2=s_3`$ respectively). In particular, in the limit the states characterized by helicity $`s_3=\pm s`$ correspond to the representation $`(s\mathrm{\hspace{0.17em}0})(0s)`$ of the Lorentz group.
Taking into account that operators $`\widehat{V}_{21}^\mu `$ ($`\widehat{V}_{12}^\mu `$) lower (raise) chirality $`\lambda `$ by 1 and the decomposition
$$f_s(x,z,\stackrel{}{}\underset{¯}{z})=\underset{\lambda =s}{\overset{s}{}}f_{j_1j_2}(x,z,\stackrel{}{}\underset{¯}{z}),\mathrm{where}s=j_1+j_2,\lambda =j_1j_2,$$
(287)
one can write equation (278) in chiral representation in the form
$$\left(\begin{array}{c}\widehat{p}_\mu \widehat{V}_{21}^\mu f_{s\frac{1}{2},\frac{1}{2}}\\ \widehat{p}_\mu \widehat{V}_{12}^\mu f_{s,0}+\widehat{p}_\mu \widehat{V}_{21}^\mu f_{s1,1}\\ \mathrm{}\\ \widehat{p}_\mu \widehat{V}_{12}^\mu f_{\frac{1}{2},s\frac{1}{2}}\end{array}\right)=ms\left(\begin{array}{c}f_{s,0}\\ f_{s\frac{1}{2},\frac{1}{2}}\\ \mathrm{}\\ f_{0,s}\end{array}\right).$$
(288)
For $`m0`$ this equation binds $`1+[s]`$ irreps of the improper Lorentz group and allows one to express components corresponding to irrep $`(s\mathrm{\hspace{0.33em}0})`$ in terms of components corresponding to irrep $`(s\frac{1}{2}\frac{1}{2})`$ and so on. This, in turn, for $`s=1,\mathrm{\hspace{0.17em}3}/2,\mathrm{\hspace{0.17em}2}`$ allows one to pass from the first order equations for the reducible representation to second order equations for irrep of improper Poincaré group. For example, for $`s=1`$, excluding $`f_{1,0}`$ and $`f_{0,1}`$, we obtain
$$m^2f_{\frac{1}{2}\frac{1}{2}}(x,𝐳)=[\widehat{p}_\mu \widehat{V}_{12}^\mu ,\widehat{p}_\nu \widehat{V}_{21}^\nu ]_+f_{\frac{1}{2}\frac{1}{2}}(x,𝐳).$$
(289)
In general case one also can to pass from the system of first order equations (288) on the reducible representation to higher order equations for irrep, for example, to the equations of $`1+[s]`$ order on the components transforming under irreps $`(\frac{s}{2}\frac{s}{2})`$ or $`(\frac{2s+1}{4}\frac{2s1}{4})(\frac{2s1}{4}\frac{2s+1}{4})`$ for the cases of integer or half-integer spin respectively.
Let us consider some particular cases.
For $`s=j_1+j_2=1/2`$ we have
$$f_{\frac{1}{2}}(x,𝐳)=\chi _\alpha (x)z^\alpha +\stackrel{}{}\psi ^{\dot{\alpha }}(x)\stackrel{}{}\underset{¯}{z}_{\dot{\alpha }}=Z_D\mathrm{\Psi }_D(x),\mathrm{\Psi }_D(x)=\left(\genfrac{}{}{0pt}{}{\chi _\alpha (x)}{\stackrel{}{}\psi ^{\dot{\alpha }}(x)}\right),$$
(290)
where $`Z_D`$ is given by the formula (207). If we substitute (290) into equation (278) and compare the coefficients at $`z^\alpha `$ and at $`\stackrel{}{}\underset{¯}{z}_{\dot{\alpha }}`$ in the left and right side, then we obtain the Dirac equation
$$\widehat{p}_\mu \gamma ^\mu \mathrm{\Psi }_D(x)=m\mathrm{\Psi }_D(x),\gamma ^\mu =\left(\begin{array}{cc}0& \sigma ^\mu \\ \overline{\sigma }^\mu & 0\end{array}\right).$$
(291)
According to (208) for space inversion we obtain $`Z_D\mathrm{\Psi }_D(x)\stackrel{P}{}Z_D\mathrm{\Psi }_D^P(\overline{x})=Z_D\gamma ^0\mathrm{\Psi }_D(\overline{x})`$, where $`\overline{x}=(x^0,x^k)`$. The matrix $`\gamma ^5=diag\{\sigma ^0,\sigma ^0\}`$ corresponds to the chirality operator (236).
A complex conjugate function corresponds to charge conjugate state,
$$\stackrel{}{}f_{1/2}(x,𝐳)=\psi _\alpha (x)\underset{¯}{z}^\alpha \stackrel{}{}\chi ^{\dot{\alpha }}(x)\stackrel{}{}z_{\dot{\alpha }},$$
(the minus sign is from anticommutation of spinors, $`\psi _\alpha z^\alpha =z_\alpha \psi ^\alpha `$) or in the matrix form
$$Z_D\mathrm{\Psi }_D(x)\stackrel{C}{}\stackrel{}{}Z_D\stackrel{}{}\mathrm{\Psi }_D(x)=\underset{¯}{Z}_D\mathrm{\Psi }_D^c(x),\mathrm{\Psi }_D^c(x)=\left(\genfrac{}{}{0pt}{}{\psi _\alpha (x)}{\stackrel{}{}\chi ^{\dot{\alpha }}(x)}\right)=i\sigma ^2\left(\genfrac{}{}{0pt}{}{\psi ^\alpha (x)}{\stackrel{}{}\chi _{\dot{\alpha }}(x)}\right),$$
(292)
where $`\underset{¯}{Z}_D=(\underset{¯}{z}^\alpha ,\stackrel{}{}z_{\dot{\alpha }})`$ and $`Z_D`$ obey the same transformation law. Thus we get the different scalar functions to describe particles and antiparticles and hence two Dirac equations for both signs of charge respectively. That matches completely with the results of the article . It was shown there that in the course of a consistent quantization of a classical model of spinning particle namely such (charge symmetric) quantum mechanics appears. At the same time it is completely equivalent to the one-particle sector of the corresponding quantum field theory.
Real functions $`f_{1/2}(x,𝐳)=\stackrel{}{}f_{1/2}(x,𝐳)`$ describing Majorana particle depend on the elements of $`Z_M`$ (209), and correspondingly $`\psi ^\alpha (x)=\chi ^\alpha (x)=i\sigma ^2\chi _\alpha (x)`$. Space reflection maps this functions into the functions of $`\underset{¯}{Z}_M`$.
For $`s=j_1+j_2=1`$ we have
$$f_1(x,𝐳)=\chi _{\alpha \beta }(x)z^\alpha z^\beta +\varphi _\alpha ^{\dot{\beta }}(x)z^\alpha \stackrel{}{}\underset{¯}{z}_{\dot{\beta }}+\psi ^{\dot{\alpha }\dot{\beta }}(x)\stackrel{}{}\underset{¯}{z}_{\dot{\alpha }}\stackrel{}{}\underset{¯}{z}_{\dot{\beta }}=\mathrm{\Phi }_\mu (x)q^\mu +\frac{1}{2}F_{\mu \nu }(x)q^{\mu \nu },$$
(293)
where
$`q^\mu ={\displaystyle \frac{1}{2}}\sigma _{}^{\mu }{}_{\alpha \dot{\beta }}{}^{}z^\alpha \stackrel{}{}\underset{¯}{z}^{\dot{\beta }},q_\mu q^\mu =0,q_{\mu \nu }=q_{\nu \mu }={\displaystyle \frac{1}{2}}\left((\sigma _{\mu \nu })_{\alpha \beta }z^\alpha z^\beta +(\overline{\sigma }_{\mu \nu })_{\dot{\alpha }\dot{\beta }}\stackrel{}{}\underset{¯}{z}^{\dot{\alpha }}\stackrel{}{}\underset{¯}{z}^{\dot{\beta }}\right),`$ (294)
$`\mathrm{\Phi }_\mu (x)=\overline{\sigma }_\mu ^{\dot{\beta }\alpha }\varphi _{\alpha \dot{\beta }}(x),F_{\mu \nu }(x)=2\left((\sigma _{\mu \nu })_{\alpha \beta }\chi ^{\alpha \beta }(x)+(\overline{\sigma }_{\mu \nu })_{\dot{\alpha }\dot{\beta }}\psi ^{\dot{\alpha }\dot{\beta }}(x)\right).`$ (295)
Substituting (293) into equation (278), we obtain
$`m\psi ^{\dot{\alpha }\dot{\beta }}(x)={\displaystyle \frac{1}{2}}\widehat{p}_\mu \overline{\sigma }^{\mu \dot{\alpha }\gamma }\varphi _\gamma ^{\dot{\beta }}(x),m\chi _{\alpha \beta }(x)={\displaystyle \frac{1}{2}}\widehat{p}_\mu \sigma _{}^{\mu }{}_{\dot{\gamma }\alpha }{}^{}\varphi _\beta ^{\dot{\gamma }}(x),`$ (296)
$`m\varphi _\alpha ^{\dot{\beta }}(x)=\widehat{p}_\mu (\overline{\sigma }^{\mu \dot{\beta }\gamma }\chi _{\alpha \gamma }(x)+\sigma _{\alpha \dot{\alpha }}^\mu \psi ^{\dot{\alpha }\dot{\beta }}(x)),`$ (297)
$`mF_{\mu \nu }(x)=_\mu \mathrm{\Phi }_\nu (x)_\nu \mathrm{\Phi }_\mu (x),m\mathrm{\Phi }_\mu (x)=^\nu F_{\mu \nu }(x).`$ (298)
The Duffin–Kemmer equation in the form (297) or (298) is the equation for irrep $`T_{[20]}`$ of the 3+2 de Sitter group and thus for the reducible representation $`(\mathrm{1\hspace{0.17em}0})(\frac{1}{2}\frac{1}{2})(\mathrm{0\hspace{0.17em}1})`$ of the Lorentz group. This representation contains both four-vector $`\mathrm{\Phi }_\mu (x)`$ and antisymmetric tensor $`F_{\mu \nu }(x)`$, which correspond to chiralities $`\lambda =0`$ and $`\lambda =\pm 1`$. Excluding components $`F_{\mu \nu }(x)`$, we obtain second order system only for the components $`\mathrm{\Phi }_\mu (x)`$ transforming under irrep $`(\frac{1}{2}\frac{1}{2})`$ of the Lorentz group:
$$(\widehat{p}^2m^2)\mathrm{\Phi }_\mu (x)=0,\widehat{p}^\mu \mathrm{\Phi }_\mu (x)=0.$$
(299)
One can rewrite operator $`\widehat{p}_\mu \widehat{\mathrm{\Gamma }}^\mu `$ in terms of complex variables $`q^\mu `$ and $`q^{\mu \nu }`$,
$$\widehat{p}_\mu \widehat{\mathrm{\Gamma }}^\mu =i\widehat{p}_\mu (q^{\mu \nu }/q^\nu q_\nu /q_{\mu \nu }).$$
(300)
Such conversion to vector indices is possible under considering of any integer spin. Notice that two sets of real spin variables with vector indices can be obtained by the substitution of elements of $`Z_M`$ and $`\underset{¯}{Z}_M`$ instead of $`Z_D`$ into (294).
One may describe neutral spin 1 field, in particular, by real function of the elements of Majorana $`z`$-spinor, $`f_1(x,𝐳)=\stackrel{}{}f_1(x,𝐳)`$. However, the spaces of quadratic functions of Dirac $`z`$-spinor $`Z_D`$ and Majorana $`z`$-spinor $`Z_M`$ are noninvariant with respect to charge conjugation and space reflection respectively. To describe spin 1 neutral particle coinciding with its antiparticle, one may use bilinear functions of $`Z_D`$ and $`\underset{¯}{Z}_D`$.
For the cases $`s=1/2`$ and $`s=1`$ first equation of the system (277)-(279) (Klein-Gordon equation) is the consequence of other equations. For $`s>1`$ the solutions of (278) are characterized by spin and mass spectrum, $`s_i=\{s,s1,\mathrm{}1\}`$ or $`s_i=\{s,s1,\mathrm{}1/2\}`$, $`m_i=ms/s_i`$. Thus for higher spin fields the Klein-Gordon equation is independent condition, allowing to exclude from spin spectrum all spins, except maximal $`s=j_1+j_2`$.
The cases $`s=1/2`$ and $`s=1`$ are also the exceptions in sense of simplicity of labelling the components by spinor or vector indices. The number of indices of symmetric spin-tensors, necessary for labelling higher spin components, increases in spite of the fact that it is sufficient to use only three operators and therefore only three numbers for labelling of the states belonging to symmetric irreps of $`SO(3,2)`$.
In particular, for spin 3/2 particle there exist four kinds of components, namely $`\psi _{\alpha \beta \gamma }`$, $`\psi _{\dot{\alpha }\beta \gamma }`$, $`\psi _{\dot{\alpha }\dot{\beta }\gamma }`$, $`\psi _{\dot{\alpha }\dot{\beta }\dot{\gamma }}`$, corresponding to four possible values of the chirality. For the spin 2 particle the representation in terms of $`q^\mu `$ and $`q^{\mu \nu }`$ is also cumbersome,
$$f_2(x,q)=\mathrm{\Phi }_{\mu \nu }(x)q^\mu q^\nu +\frac{1}{2}F_{\mu \nu ,\rho }(x)q^{\mu \nu }q^\rho +\frac{1}{4}F_{\mu \nu ,\rho \sigma }(x)q^{\mu \nu }q^{\rho \sigma },$$
(301)
with the necessity to fix independent components by means of relations $`q_\mu q^\mu =0`$, $`q_{\mu \nu }q^\mu +q_{\mu \nu }q^\nu =0`$ and so on.
Thus, beginning from the spin 3/2, it is convenient to use the universal notation $`\psi _{j_1j_2}^{m_1m_2}(x)`$ associated with the decomposition (224) over monomial chiral basis (226) (see also (B12)-(B14)). Two indices $`j_1,j_2`$ label spin $`s=j_1+j_2`$ and chirality $`\lambda =j_1j_2`$ and two indices $`m_1,m_2`$ label independent components inside the irrep of the Lorentz group. This notation are also suitable for infinite-dimensional representations.
By analogy with 2+1 case, one can find plain wave solutions of the system (277)-(278) for the any spin $`s`$ in general form without using matrix representation. Corresponding to the particle moving along $`x^3`$ states are eigenstates of the operator $`\widehat{p}_i\widehat{S}^i`$ with eigenvalues $`|p|\sigma `$, where $`\sigma =s_3\mathrm{sign}p_3`$ is the helicity. These states have the form
$`f_{m,s,\sigma }(x,𝐳)=`$ $`{\displaystyle \underset{\sigma =s}{\overset{s}{}}}C_\sigma e^{ik_0x^0+k_3x^3}(z^1e^a+\stackrel{}{}\underset{¯}{z}_{\dot{1}}e^a)^{s+\sigma }(z^2e^a+\stackrel{}{}\underset{¯}{z}_{\dot{2}}e^a)^{s\sigma }+`$ (302)
$`{\displaystyle \underset{\sigma ^{}=s}{\overset{s}{}}}C_\sigma ^{}e^{ik_0x^0k_3x^3}(z^1e^a\stackrel{}{}\underset{¯}{z}_{\dot{1}}e^a)^{s\sigma }(z^2e^a\stackrel{}{}\underset{¯}{z}_{\dot{2}}e^a)^{s+\sigma },`$ (303)
where $`e^a`$ is given by (282). For the rest particle one can obtain the general solution characterized by the spin projection $`s^{}`$ on the direction $`𝐧`$ from (280) by the rotation $`z_\beta ^{}=U_\beta ^\alpha z_\alpha `$, $`USU(2)`$. For particle characterized by momentum direction $`𝐧`$ and helicity $`\sigma `$, starting from the state (303), one can obtain the solution by the analogous rotation.
Improper Poincaré group includes space reflection, which interchanges representations $`(j_1j_2)`$ and $`(j_2j_1)`$. Therefore we consider equations binding these representations (and corresponding components of the solutions of the system (277)-(279)) more detail.
In the case $`j_1=j_2`$ solutions of the system (277)-(279) are characterized by fixed spin $`s=j_1+j_2`$ and mass $`m`$. Thus the relations (231)-(234) are valid, and corresponding $`2(2s+1)`$ components obey the equations for massive tensor field (248).
In general case equations connecting the components transforming under irreps $`(j_1,j_2)`$ and $`(j_2,j_1)`$ of the Lorents subgroup have the form
$`(2j_2)!(\widehat{p}_\mu \widehat{V}_{12}^\mu )^{2|\lambda |}f_{j_1j_2}(x,z,\stackrel{}{}\underset{¯}{z})=(2j_1)!m^{2|\lambda |}f_{j_2j_1}(x,z,\stackrel{}{}\underset{¯}{z}),`$ (304)
$`(2j_1)!(\widehat{p}_\mu \widehat{V}_{21}^\mu )^{2|\lambda |}f_{j_2j_1}(x,z,\stackrel{}{}\underset{¯}{z})=(2j_2)!m^{2|\lambda |}f_{j_1j_2}(x,z,\stackrel{}{}\underset{¯}{z}),`$ (305)
where $`j_1>j_2`$, $`|\lambda |=j_1j_2`$. This set of equations are invariant under space reflection. Using the decomposition (287) and explicite form of the general solution (280) of the system (277)-(279) in the rest frame, one can prove the validity of (305) by direct calculation. Going over to spin-tensor notation, we get
$`p_\mu \overline{\sigma }^{\mu \dot{\alpha }_{2j_2+1}\alpha _{2j_2+1}}\mathrm{}p_\nu \overline{\sigma }^{\nu \dot{\alpha }_{2j_1}\alpha _{2j_1}}\psi _{\alpha _1\mathrm{}\alpha _{2j_1}}^{\dot{\alpha }_1\mathrm{}\dot{\alpha }_{2j_2}}(x)=m^{2|\lambda |}\psi _{\alpha _1\mathrm{}\alpha _{2j_2}}^{\dot{\alpha }_1\mathrm{}\dot{\alpha }_{2j_1}}(x),`$ (306)
$`p_\mu \sigma _{\alpha _{2j_2+1}\dot{\alpha }_{2j_2+1}}^\mu \mathrm{}p_\nu \sigma _{\alpha _{2j_1}\dot{\alpha }_{2j_1}}^\nu \psi _{\alpha _1\mathrm{}\alpha _{2j_2}}^{\dot{\alpha }_1\mathrm{}\dot{\alpha }_{2j_1}}(x)=m^{2|\lambda |}\psi _{\alpha _1\mathrm{}\alpha _{2j_1}}^{\dot{\alpha }_1\mathrm{}\dot{\alpha }_{2j_2}}(x).`$ (307)
Equations (307) are consequences of the system (277)-(279), but unlike this system in general case require some supplementary conditions to fix mass and spin.
Equations (307) are the first order equations only in the case $`|\lambda |=1/2`$, which corresponds to representations $`(\frac{2j\pm 1}{4}\frac{2j1}{4})`$, $`j=j_1+j_2`$, describing half-integer spins. In this case, going over to vector indices and supplementing the equations by subsidiary conditions (250) (which also are the consequences of the system (277)-(279) and exclude components with $`s<j_1+j_2`$), we obtain the Rarita–Schwinger equations
$$(\widehat{p}_\mu \gamma ^\mu m)\mathrm{\Psi }_{\mu _1\mu _2\mathrm{}\mu _n}(x)=0,\gamma ^\mu \mathrm{\Psi }_{\mu \mu _2\mathrm{}\mu _n}(x)=0,$$
(308)
where $`n=2s1`$ and $`\mathrm{\Psi }_{\mu _1\mathrm{}\mu _n}(x)`$ is a four-component column composed of $`\mathrm{\Psi }_{\mu _1\mathrm{}\mu _n\alpha }(x)`$ and $`\mathrm{\Psi }_{\mu _1\mathrm{}\mu _n}^{}{}_{}{}^{\dot{\alpha }}(x)`$. The conditions $`^\mu \mathrm{\Psi }_{\mu \mu _2\mathrm{}\mu _n}(x)=0`$ and $`\mathrm{\Psi }_{}^{\mu }{}_{\mu \mathrm{}\mu _n}{}^{}(x)=0`$ appear as consequences of these two equations .
A case $`|\lambda |=s`$ corresponding to representations $`(s\mathrm{\hspace{0.17em}0})`$ and $`(0s)`$ is preferred because of minimal number of components. In this case the equations (307) are $`2s`$ order Joos-Weinberg equations of so called $`2(2s+1)`$-component theory,
$`p_\mu \overline{\sigma }^{\mu \dot{\alpha }_1\alpha _1}\mathrm{}p_\nu \overline{\sigma }^{\nu \dot{\alpha }_{2s}\alpha _{2s}}\psi _{\alpha _1\mathrm{}\alpha _{2s}}(x)=m^{2s}\psi ^{\dot{\alpha }_1\mathrm{}\dot{\alpha }_{2s}}(x),`$ (309)
$`p_\mu \sigma _{\alpha _1\dot{\alpha }_1}^\mu \mathrm{}p_\nu \sigma _{\alpha _{2s}\dot{\alpha }_{2s}}^\nu \psi ^{\dot{\alpha }_1\mathrm{}\dot{\alpha }_{2s}}(x)=m^{2s}\psi _{\alpha _1\mathrm{}\alpha _{2s}}(x).`$ (310)
In the rest frame as a consequence we obtain $`p_0^{4s}=m^{4s}`$ and for $`s1`$ Joos-Weinberg equations have solutions with complex energy $`p_0`$, $`|p_0|=m`$. The existence of such solutions was pointed out also in .
### D Relativistic wave equations, invariant under improper Poincaré group. <br>Equations for several scalar functions
Above we have considered the linear equations for one scalar function on the group. The condition of invariance under space reflection leads us to the system (277)-(279) for particle with spin $`s=j_1+j_2`$ and mass $`m`$.
For the construction of invariant wave equations one may also use the operators $`\widehat{p}_\mu \widehat{V}_{ik}^\mu `$, which are not invariant under space reflections. Using several scalar functions $`f(x,𝐳)`$, it is possible to restore the invariance under space reflections.
In particular, equations (260) containing operators $`\widehat{V}_{12(k)}^\mu `$ and $`\widehat{V}_{21(k)}^\mu `$ interlock two scalar functions. Using the decomposition (227) in terms of spin-tensors, we obtain Dirac–Fierz–Pauli equations ,
$`\widehat{p}_\mu \overline{\sigma }^{\mu \dot{\alpha }\beta }\psi _{\beta \beta _1\mathrm{}\beta _n}^{\dot{\alpha }_1\mathrm{}\dot{\alpha }_l}(x)=\varkappa \psi _{\beta _1\mathrm{}\beta _n}^{\dot{\alpha }\dot{\alpha }_1\mathrm{}\dot{\alpha }_l}(x)`$ (311)
$`\widehat{p}_\mu \sigma _{}^{\mu }{}_{\beta \dot{\alpha }}{}^{}\psi _{\beta _1\mathrm{}\beta _n}^{\dot{\alpha }\dot{\alpha }_1\mathrm{}\dot{\alpha }_l}(x)=\varkappa \psi _{\beta \beta _1\mathrm{}\beta _n}^{\dot{\alpha }_1\mathrm{}\dot{\alpha }_l}(x).`$ (312)
These equations connect two functions transforming under irreps $`(\frac{n}{2}+\frac{1}{2}\frac{l}{2})`$ and $`(\frac{n}{2}\frac{l}{2}+\frac{1}{2})`$ of the Lorentz group and at $`n=l`$ mapping to one another under parity transformation.
Let us consider a system of equations of the form (260),(261), which connect several scalar functions with different $`j_1,j_2`$. The equations of this system interlock the representation $`(j_1,j_2)`$ with at least one of the representations $`(j_1\pm 1,j_21)`$, $`(j_1\pm 1,j_2\pm 1)`$. This allows one to identify this system with general Gel’fand–Yaglom equations
$$(\alpha ^\mu \widehat{p}_\mu \varkappa )\psi =0,[S^{\lambda \mu },\alpha ^\nu ]=i(\eta ^{\mu \nu }\alpha ^\lambda \eta ^{\lambda \nu }\alpha ^\mu ).$$
(313)
In the present approach the latter relation is a consequence of the commutation relations $`[\widehat{S}^{\lambda \mu },\widehat{V}_{ik}^\nu ]=i(\eta ^{\mu \nu }\widehat{V}_{ik}^\lambda \eta ^{\lambda \nu }\widehat{V}_{ik}^\mu )`$. This relation is necessary for Poincaré invariance of the equation .
Supplemented by commutation relations $`[\alpha ^\mu ,\alpha ^\nu ]=S^{\mu \nu }`$ finite-component equations of the form (313) are known as Bhabha equations , although for the first time they were systematically considered by Lubanski . These equations are classified according to the finite-dimensional irreps of the 3+2 de Sitter group $`SO(3,2)`$. Other possible commutation relations of matrices $`\alpha ^\mu `$ are discussed in .
Considered above equation (278) on a scalar function is the particular case of Bhabha equations. This case correspond to symmetric irreps $`T_{[2s\mathrm{\hspace{0.17em}0}]}`$ of the 3+2 de Sitter group. Generally speaking, the Bhabha equations are characterized by finite number of different $`m`$ and $`s`$. Therefore, this equations connect the fields transforming under nonequivalent irreps of the Poincaré group.
If equations include the operators $`\widehat{p}_\mu \widehat{V}_{11}^\mu `$ and $`\widehat{p}_\mu \widehat{V}_{22}^\mu `$, then either the equations describe at least two different spins $`s`$ or the condition $`s=j_1+j_2`$ connecting spin $`s`$ with a highest weight of irrep of Lorentz group is not valid.
Cite as an example the system interlocking irreps (00) and $`(\frac{1}{2}\frac{1}{2})`$ of the Lorentz group:
$$\widehat{p}_\mu \widehat{V}_{11}^\mu f_{00}(x,𝐳)=\varkappa _1f_{\frac{1}{2}\frac{1}{2}}(x,𝐳),\widehat{p}_\mu \widehat{V}_{22}^\mu f_{\frac{1}{2}\frac{1}{2}}(x,𝐳)=\varkappa _2f_{00}(x,𝐳)$$
(314)
where $`f_{00}(x,𝐳)=\psi (x)`$, $`f_{\frac{1}{2}\frac{1}{2}}(x,𝐳)=\psi _\alpha ^{\dot{\beta }}(x)z^\alpha \stackrel{}{}\underset{¯}{z}_{\dot{\beta }}`$; in component-wise form we have $`\widehat{p}_\mu \psi =2\varkappa _1\psi _\mu `$, $`\widehat{p}_\mu \psi ^\mu =\varkappa _2\psi `$. In the rest frame one may obtain $`\varkappa _2=2\varkappa _1=m`$. Thus the system (314) is equivalent to Duffin equation for scalar particles, which correspond to five-dimensional vector irrep $`T_{[01]}`$ of $`SO(3,2)`$ group.
### E Relativistic wave equations, invariant under improper Poincaré group. <br>Equations for particles with composite spin
Many-particle systems are described by the functions of the sets of variables $`x_{(i)},z_{(i)},\stackrel{}{}\underset{¯}{z}_{(i)}`$. But here we will consider not many-particle systems in usual sense, but some objects corresponding to functions $`f(x,z_{(1)},\stackrel{}{}\underset{¯}{z}_{(1)},\mathrm{},z_{(n)},\stackrel{}{}\underset{¯}{z}_{(n)})`$, (or, briefly, $`f(x,\{𝐳_{(i)}\})`$), i.e. to functions of one set of $`x`$ and several sets of $`𝐳`$. These objects one may interpret as particles with composite spin.
As an example we will consider the Ivanenko-Landau-Kähler (or Dirac-Kähler) equation . Let us write some linear in $`𝐳_{(1)}`$ and $`𝐳_{(2)}`$ scalar function $`f(x,𝐳_{(1)},𝐳_{(2)})`$ in the form
$$f(x,𝐳_{(1)},𝐳_{(2)})=Z_D^{(1)}\mathrm{\Psi }(x)(Z_D^{(2)})^{}=\underset{i,j=1}{\overset{4}{}}(Z_D^{(1)})_i\mathrm{\Psi }_{ij}(x)(\stackrel{}{}Z_D^{(2)})_j,$$
(315)
where $`Z_D=(z^1z^2\stackrel{}{}\underset{¯}{z}_{\dot{1}}\stackrel{}{}\underset{¯}{z}_{\dot{2}})`$, and $`\mathrm{\Psi }(x)`$ is a $`4\times 4`$ matrix with a transformation rule
$$\mathrm{\Psi }^{}(x^{})=\stackrel{ˇ}{U}\mathrm{\Psi }(x)(\stackrel{ˇ}{U})^{},\stackrel{ˇ}{U}=diag\{U,(U^1)^{}\},$$
in contrast to the transformation rule $`\mathrm{\Psi }_D^{}(x^{})=\stackrel{ˇ}{U}\mathrm{\Psi }_D(x)`$ of Dirac spinor (291). Let us impose the equation on the first (”left”) spin subsystem,
$$(\widehat{p}_\mu \widehat{\mathrm{\Gamma }}_{(1)}^\mu m/2)f(x,𝐳_{(1)},𝐳_{(2)})=0,$$
(316)
and do not impose any conditions on the second (”right”) spin subsystem. Writing (316) in component-wise form, one can obtain Ivanenko-Landau-Kähler equation in spinor matrix representation
$$(\widehat{p}_\mu \gamma ^\mu m)\mathrm{\Psi }(x)=0.$$
(317)
According to (317), 16 components $`\mathrm{\Psi }_{ij}(x)`$ obey Klein-Gordon equation, therefore mass is equal to $`m`$. Spin of both subsystems is equal two 1/2. The spin of system is indefinite, and there are both spin 0 and spin 1 components.
The consideration of this equation is associated mainly with the attempts to describe fermions by the antisymmetric tensor fields (see, for example, and also as a good introduction). The spin subsystems (”left-spin” and ”right-spin”) were considered in ).
Let us consider now linear symmetric functions of $`𝐳_{(1)},\mathrm{},𝐳_{(n+l)}`$:
$$f_{\frac{n}{2}\frac{l}{2}}(x,\{𝐳_{(i)}\})=\psi _{\beta _1,\mathrm{},\beta _n}^{\dot{\alpha }_1,\mathrm{},\dot{\alpha }_l}(x)z_{(1)}^{\beta _1}\mathrm{}z_{(n)}^{\beta _n}\stackrel{}{}\underset{¯}{z}_{(n+1)\dot{\alpha }_1}\mathrm{}\stackrel{}{}\underset{¯}{z}_{(n+l)\dot{\alpha }_l},$$
(318)
where symmetric spinors $`\psi _{\beta _1,\mathrm{},\beta _n}^{\dot{\alpha }_1,\mathrm{},\dot{\alpha }_l}(x)`$ transform under irreps $`(n/2,l/2)`$, and all permutations of $`1,\mathrm{},n+l`$ are summed over. As a consequence of symmetry of the multispinors with respect to index permutations, spin subsystems are indistinguishable, and this allow us to use functions of several sets of spin variables for describing usual particles.
So, one may obtain Dirac–Fierz–Pauli equations (312), acting by the operators $`\widehat{V}_{12(k)}^\mu `$ and $`\widehat{V}_{21(k)}^\mu `$ on the functions (318) corresponding to irreps $`(\frac{n}{2}+\frac{1}{2},\frac{l}{2})`$ and $`(\frac{n}{2},\frac{l}{2}+\frac{1}{2})`$ of the Lorentz group:
$`\widehat{V}_{12(k)}^\mu f_{\frac{n}{2}+\frac{1}{2},\frac{l}{2}}(x,\{𝐳_{(i)}\})=\varkappa f_{\frac{n}{2},\frac{l}{2}+\frac{1}{2}}(x,\{𝐳_{(i)}\}),`$ (319)
$`\widehat{V}_{21(k)}^\mu f_{\frac{n}{2},\frac{l}{2}+\frac{1}{2}}(x,\{𝐳_{(i)}\})=\varkappa f_{\frac{n}{2}+\frac{1}{2},\frac{l}{2}}(x,\{𝐳_{(i)}\}).`$ (320)
In general case a linear symmetric function of $`𝐳_{(k)}`$, $`k=1,\mathrm{},2j`$ has the form
$$f_j(x,\{𝐳_{(i)}\})=\underset{n,l;n+l=2j}{}\psi _{\beta _1\mathrm{}\beta _n}^{\dot{\alpha }_1\mathrm{}\dot{\alpha }_l}(x)z_{(1)}^{\beta _1}\mathrm{}z_{(n)}^{\beta _n}\stackrel{}{}\underset{¯}{z}_{(n+1)\dot{\alpha }_1}\mathrm{}\stackrel{}{}\underset{¯}{z}_{(n+l)\dot{\alpha }_l}.$$
(321)
Functions (321) correspond to symmetric part of the representation $`\left((\frac{1}{2}\mathrm{\hspace{0.17em}0})(0\frac{1}{2})\right)^{2j}`$. This symmetric part expand into direct sum of irreps $`(j_1j_2)`$, $`j_1+j_2=j`$. Impose on each spin subsystem the condition
$$(\widehat{p}_\mu \widehat{\mathrm{\Gamma }}_{(k)}^\mu m/2)f(x,𝐳_{(1)},\mathrm{},𝐳_{(2j)})=0,k=1,\mathrm{},2j.$$
(322)
Rewriting this equations in four-component form, we obtain Bargmann–Wigner equations
$$(\widehat{p}_\mu \gamma _{(k)}^\mu m)_{\alpha _k\beta _k}\psi _{\beta _1\mathrm{}\beta _k\mathrm{}\beta _{2j}}(x)=0.$$
(323)
As a consequence of (322), one may obtain equations for system as a whole:
$`(\widehat{p}^2m^2)f(x,𝐳_{(1)},\mathrm{},𝐳_{(2j)})=0,`$ (324)
$`(\widehat{p}_\mu \widehat{\mathrm{\Gamma }}^\mu ms)f(x,𝐳_{(1)},\mathrm{},𝐳_{(2j)})=0,\widehat{\mathrm{\Gamma }}^\mu ={\displaystyle \widehat{\mathrm{\Gamma }}_{(k)}^\mu },`$ (325)
which are analogous with equations (277)-(278) at the case $`s=j_1+j_2`$. Both Bargmann–Wigner equations and system (325) have $`2(2s+1)`$ independent solutions $`\psi (x)`$, and therefore this systems are equivalent.
### F Relativistic wave equations: comparative consideration
In the framework of the group-theoretical classification of the scalar fields $`f(x,𝐳)`$ on the Poincaré group we have obtained two types of equations describing unique spin and mass, namely equations for the eigenfunctions of the Casimir operator of the Lorentz spin subgroup ($`j_1`$ and $`j_2`$ are fixed, see (234)) and equations for the eigenfunctions of the Casimir operator of the $`SO(3,2)`$ group (the sum $`j_1+j_2`$ is fixed). Below we will consider comparative characteristics of these equations and also the case $`(j_1j_2)(j_2j_1)`$ corresponding to irreps of the improper Poincaré group but requiring two scalar functions for its formulation.
1. Equations for the functions corresponding to the fixed irrep $`(j_1j_2)`$ of the Lorentz group. Mass and spin irreducibility conditions leave $`2(2s+1)`$ independent components corresponding to two improper Poincaré group irreps differed by sign of $`p_0`$. For $`s=j_1+j_2`$ the equations in spin-tensor form constitute the system of the Klein-Gordon equation and the subsidiary condition (246), which eliminates components with other possible values of spin $`s`$ for fixed $`j_1,j_2`$, $`|j_1j_2|s<j_1+j_2`$. For $`sj_1+j_2`$ one should consider general subsidiary condition (247). An alternative to the use of subsidiary condition is the consideration of functions of momentum and spin variables with invariant constraints (243).
There exist two preferred cases. The first corresponds to the representations $`(\frac{s}{2}\frac{s}{2})`$ mapping onto themselves under space reflection and are most often used to describe integer spins. The second corresponds to the representations $`(s\mathrm{\hspace{0.17em}0})`$ and $`(0s)`$. In this case there is not necessity to impose subsidiary conditions since they are fulfilled identically.
2. Equations for the functions corresponding to the representations $`(j_1j_2)`$ and $`(j_2j_1)`$, $`j_1j_2`$, which are interchanged under space reflection. Unlike the considered above equations for fixed $`j_1,j_2`$, these equations in general case do not assume formulation as equations for one scalar function $`f(x,𝐳)`$. The conditions of mass and spin irreducibility leave $`4(2s+1)`$ independent components corresponding to four improper Poincaré group irreps differed by sign of $`p_0`$ and intrinsic parity $`\eta `$. To choose $`2(2s+1)`$ components corresponding to fixed sign of $`\eta `$ or $`p_0\eta `$ it is necessary to supplement these conditions by equations (305) connecting components corresponding to $`(j_1j_2)`$ and $`(j_2j_1)`$.
Equations (305) are first order equations only for the representations $`(j+\frac{1}{2}j)(jj+\frac{1}{2})`$. These representations and associated with them Rarita–Schwinger equations (308) are most often used to describe half-integer spins. However, just as in the case of representations $`(jj)`$, subsidiary conditions supplement the field equations, and the number of equations exceeds the number of field components. Therefore one has an overdetermined set of equations which, although consistent in the free-field case, for $`s>1`$ becomes self-contradictory with minimal electromagnetic coupling . In order to avoid inconsistency it is possible to give a Lagrangian formulation, introducing auxiliary fields , but this formulation leads to acasual propagation with minimal electromagnetic coupling .
For the case $`(s\mathrm{\hspace{0.17em}0})(0s)`$ one can construct 2(2s+1)-component theory, but corresponding Joos-Weinberg equations of $`2s`$ order (see (310)) for $`s1`$ have also solutions with complex energy.
The second order equation for representation $`(s\mathrm{\hspace{0.17em}0})(0s)`$, $`(\widehat{P}^2\frac{e}{2s}\widehat{S}^{\mu \nu }F_{\mu \nu }m^2)\psi (x)=0`$, for free particle possesses $`4(2s+1)`$ independent components differed by spin projection and by signs of $`p_0`$ and $`\eta `$. On the other hand, this equation describes unique mass and spin and is characterized by casual solutions. In particular, exact solutions in external constant uniform electromagnetic field are known . One may rewrite the above equation as first order equation with minimal coupling for representations $`(s\mathrm{\hspace{0.17em}0})(s\frac{1}{2}\frac{1}{2})(\frac{1}{2}s\frac{1}{2})(0s)`$. As noted in , this is the simplest class of describing unique mass and spin representations, which led to first order equations without subsidiary conditions.
3. Equations (277)-(278) for eigenfunctions of the Casimir operator (279) of $`SO(3,2)`$ group with eigenvalues $`4s(s+2)`$, $`s=j_1+j_2`$:
$$(\widehat{p}_\mu \widehat{\mathrm{\Gamma }}^\mu ms)f(x,z,\stackrel{}{}\underset{¯}{z})=0,(\widehat{p}^2m^2)f(x,z,\stackrel{}{}\underset{¯}{z})=0.$$
(326)
The condition of spin irreducibility (256) is a consequence of this system.
The first equation of the system is the Bhabha equation corresponding to symmetric irrep $`T_{[2s\mathrm{\hspace{0.33em}0}]}`$ of the group $`Sp(4,R)SO(3,2)`$. This equation represent a straightforward higher spins generalization of Dirac and spin 1 Duffin–Kemmer equations. Both Bhabha equations and the problem of minimal coupling for these equations were most detail considered in papers of Krajcik and Nieto (see ; it contains references to the six earlier papers). The theory is casual with minimal electromagnetic coupling , but in general case Bhabha equations describe multi-mass systems. Notice that the connection of the Rarita–Schwinger and Bargmann–Wigner equations with Bhabha equations was considered also in .
The solutions of the system (326) have the components transforming under $`2s+1`$ irreps $`(j_1,j_2)`$, $`j_1+j_2=s`$, of the Lorentz group. But the components corresponding to different chiralities $`\lambda =j_1j_2`$ are not independent. In contrast to left generators of the Poincaré group operators $`\widehat{\mathrm{\Gamma }}_\mu `$ do not commute with chirality operator (which is the right generator of the Poincaré group) and combine $`2s+1`$ representations of the Lorentz group into one irrep of the 3+2 de Sitter group $`SO(3,2)`$.
The current component $`j^0`$ is positive definite for half-integer spin particles and the energy density is positive definite for integer spin particles, see Appendix B.
In the rest frame equations (326) have $`2s+1`$ positive- and $`2s+1`$ negative-frequency solutions labelled by different spin projections, see (280), and half-integer spins solutions with opposite frequency are characterized by opposite parity. In the ultrarelativistic limit two solutions with opposite sign of $`p_0`$ correspond to any of $`2s+1`$ of possible values of chirality, see (286).
Thus the system (326) describes a particle with unique spin and mass, is invariant under parity transformation and possesses $`2(2s+1)`$ independent components.
Let us briefly consider the problem of equivalence of the different RWE. In the case of free fields, using the relation
$$[_\mu ,_\nu ]=0,$$
(327)
one can establish the equivalence of wide class of RWE.
So, as we have established above, in a free case the system (326) and the Bargmann–Wigner equations (322), which both describe a particle by means wave functions with components transforming under $`2s+1`$ irreps $`(j_1j_2)`$, $`j_1+j_2=s`$, of the Lorentz group, are equivalent. However, the formulation (326) is more general since unlike Bargmann–Wigner equations can be considered also in the case of infinite-dimensional unitary representations of the Lorentz group, as was done above with an analogous system in 2+1-dimensional case.
The considered above free equations for representations $`(j_1j_2)`$ or $`(j_1j_2)(j_2j_1)`$ can be obtained as a consequence of the Bargmann–Wigner equations or the system (326) by excluding of other components. In general case for $`m0`$ one may express all components in terms of the components corresponding to two chiralities $`\pm \lambda `$, where $`s\lambda s`$.
It is obvious that the coupling, which is minimal for one system, is not minimal for another ”equivalent” system if one uses the relation (327) to prove this equivalence in a free case. These equations will differ by the terms proportional to the commutator of covariant derivatives $`[D_\mu ,D_\nu ]=igF_{\mu \nu }`$.
Therefore, when an interaction is introduced, the system of equations can be found inconsistent if, taking account of (327), some equations are the consequences of another. In particular, spin 1 Bargmann–Wigner equations with minimal electromagnetic coupling are inconsistent (for the prove see, for example, ), but equivalent to them in a free case Duffin–Kemmer and Proca equations with minimal coupling are consistent and characterized by casual solutions .
Recently different approaches to introduce interactions for higher spin massive fields have been considered (see, in particular, ). Keeping in mind the present approach, we hope that new possibilities to describe interacting higher spin fields will arise.
## VI Equations for fixed spin and mass: general features
Consider now the general properties of the obtained equations describing a particle with unique mass $`m>0`$ and spin $`s`$ in two dimensions
$`\widehat{p}^2f(x,\theta )=m^2f(x,\theta ),`$ (328)
$`\widehat{p}_\mu \widehat{\mathrm{\Gamma }}^\mu f(x,\theta )=msf(x,\theta ),`$ (329)
in three dimensions
$`\widehat{p}^2f(x,𝐳)=m^2f(x,𝐳),`$ (330)
$`\widehat{p}_\mu \widehat{S}^\mu f(x,𝐳)=msf(x,𝐳),`$ (331)
$`\widehat{S}_\mu \widehat{S}^\mu f(x,𝐳)=S(S+1)f(x,𝐳),`$ (332)
in four dimensions
$`\widehat{p}^2f(x,𝐳)=m^2f(x,𝐳),`$ (333)
$`\widehat{p}_\mu \widehat{\mathrm{\Gamma }}^\mu f(x,𝐳)=msf(x,𝐳),`$ (334)
$`\widehat{S}_{ab}\widehat{S}^{ab}f(x,𝐳)=4S(S+2)f(x,𝐳).`$ (335)
In the latter case in addition we suppose $`s=\pm S`$ to avoid nontrivial spin and mass spectrum.
In all dimensions the first equation (condition of the mass irreducibility) is the eigenvalue equation for the Casimir operator of the Poincaré group. But the other equations, although seem similar, have different origin in even and odd dimensions. This is related to the different role of space inversion.
In 2+1 dimensions other equations (331)-(332) are eigenvalue equations for the Casimir operator of the Poincaré group and the spin Lorentz subgroup.
In even dimensions the Casimir operators of the Lorentz subgroup do not commute with the space inversion operator, and space inversion combines two labelled by chiralities $`\pm \lambda `$ equivalent representations of the proper Poincaré group into representation of the improper Poincaré group. If one rejects equations, which fix chirality (in 3+1 dimensions this corresponds to the transition to the system (255)-(257)), then in the rest frame it is easy to see that there is redundant number of independent components. Thus it is necessary to construct equation, binding the states with different chiralities, and correspondingly new set of commuting operators. It can be down by using supplementary operators $`\widehat{\mathrm{\Gamma }}^\mu `$, which extend Lorentz group $`SO(D,1)`$ up to $`SO(D,2)`$ group with the maximal compact subgroup $`SO(D)SO(2)`$. Operator $`\widehat{\mathrm{\Gamma }}^0`$ is the generator of compact $`SO(2)`$ subgroup.
Third equation of the system fixes the power $`2S`$ of homogeneity of the functions $`f(x,𝐳)`$ in $`𝐳`$ and therefore fixes irrep of the Lorentz group in 2+1 dimensions or of the 3+2 de Sitter group in 3+1 dimensions. (In 1+1 dimensions there exists analogous equation $`\widehat{\mathrm{\Gamma }}_a\widehat{\mathrm{\Gamma }}^af(x,\theta )=s(s+1)f(x,\theta )`$, but, in fact, this equation defines the structure of $`\widehat{\mathrm{\Gamma }}^\mu `$.)
Positive (half-)integer $`S=s`$ correspond to finite-dimensional nonunitary irreps of the Lorentz (or de Sitter) group, which realize in the space of the power $`2s`$ polynomials in $`𝐳`$.
Negative $`S=s`$ correspond to infinite-dimensional unitary irreps. The unitary property allows one to combine probability amplitude interpretation and relativistic invariance (the desirability of this combination was stressed by Dirac in ). Thus the equations under consideration allow two approaches to the description of the same spin by means of both finite-dimensional nonunitary and infinite-dimensional unitary irreps.
In 1+1 and 2+1 dimensions there is the possibility of the existence of particles with fractional spin since the groups $`SO(1,1)`$ and $`SO(2,1)`$ do not contain compact Abelian subgroup. However, the description of massive particles with fractional spin can be given only in terms of the infinite-dimensional irreps of the group $`SO(2,1)`$. This is another reason to consider infinite-dimensional irreps.
Fixing the irrep of the Lorentz (or de Sitter) group with the help of the third equation of the system, one can come to usual multicomponent matrix description by the separation of space and spin variables: $`f(x,𝐳)=\varphi _n(𝐳)\psi _n(x)`$, where $`\varphi _n(𝐳)`$ form the basis in the representation space of the Lorentz (or de Sitter) group. Thus, depending on the choice of the solution of the third equation, second equation in matrix representation is either finite-component equation or infinite-component equation of Majorana type.
For fundamental spinor irreps the action of differential operators $`2\widehat{S}^\mu `$ in 2+1 dimensions and $`2\widehat{\mathrm{\Gamma }}^\mu `$ in 1+1 and 3+1 dimensions in the space of functions $`f(x,𝐳)`$ on the Poincaré group can be rewritten in terms of action of corresponding $`\gamma `$-matrices on the functions $`\psi (x)`$.
Differential operators $`\widehat{\mathrm{\Gamma }}^\mu `$ and matrices $`\gamma ^\mu /2`$ obey the same commutation relations
$$[\widehat{\mathrm{\Gamma }}^\mu ,\widehat{\mathrm{\Gamma }}^\nu ]=i\widehat{S}^{\mu \nu },[\widehat{S}^\mu ,\widehat{S}^\nu ]=iϵ^{\mu \nu \rho }\widehat{S}_\rho .$$
In 3+1 dimensions operators $`\widehat{\mathrm{\Gamma }}^\mu `$ and $`\widehat{S}^{\mu \nu }`$ obey the commutation relations of generators of $`SO(3,2)`$ group, see (267).
Anticommutation relations for operators $`\widehat{S}^\mu `$ in 2+1 and $`\widehat{\mathrm{\Gamma }}^\mu `$ in 1+1 and 3+1 dimensions are analogous with the relations for $`\gamma `$-matrices,
$$[\widehat{S}^\mu ,\widehat{S}^\nu ]_+=\frac{1}{2}\eta ^{\mu \nu },[\widehat{\mathrm{\Gamma }}^\mu ,\widehat{\mathrm{\Gamma }}^\nu ]_+=\frac{1}{2}\eta ^{\mu \nu },$$
and are valid only for fundamental spinor irreps. This is group-theoretical property connected with the fact that for these irreps the double action of lowering or raising operators on any state gives zero as a result. (Notice that, besides the case of spinor irreps of orthogonal groups, anticommutation relations also take place for fundamental $`N`$-dimensional irreps of $`Sp(N)`$ and $`SU(N)`$ groups .)
For $`s=1/2`$ and $`s=1`$ the first equation of the system (condition of mass irreducibility) is a consequence of (331) or (334). In general case the second equation of the system describes multi-mass systems $`m_is_i=ms`$. Thus for $`s>1`$ it is necessary to consider both equations.
Consider some characteristics of the equations associated with finite-dimensional irreps of the Lorentz group. If we reject the first equation of the system (i.e. the condition of mass irreducibility), then for the second equation of the system the component $`j^0`$ of the current vector is positive definite only for $`s=1/2`$, and the energy density $`T^{00}`$ (see (171)) is positive definite only for $`s=1`$. (The case $`s=1`$ in 3+1 dimensions has considered in detail in ). However, for the system as a whole the component $`j^0`$ of the current vector is positive definite for any half-integer spin, and energy density is positive definite for any integer spin. Besides, in the rest frame half-integer spin solutions with opposite sign of $`p_0`$ are characterized by opposite parity.
For the case of infinite-component equations in 2+1 dimensions, the energy is positive definite for any spin, and $`j^0`$ is positive or negative definite in accordance with the sign of charge.
The consideration of the field on the Poincaré group also allows one to ensure essential progress in the problem of practical computations for multicomponent equations. As was noted in , the general investigation of Gel’fand–Yaglom equations ”revealed a number of interesting features, but … the use of such equations (or more accurately, systems of a large or infinite number of equations) for any practical computations is not possible”. In the present approach, due to the use of spin differential operators instead of finite or infinite-dimensional matrices, from the technical point of view there is no essential distinction in the consideration of the equations associated with various finite-dimensional and infinite-dimensional representations of the Lorentz group. Therefore the present approach is adequate to work with higher spins and positive energy wave equations. For example, the use of spin variables $`𝐳`$ has allowed us to obtain explicit compact form of general plane wave solutions for any spin (including fractional spin in 2+1 dimensions).
Notice that unlike the equations for particles with unique mass and spin, in general case RWE with mass and spin spectrum can either interlock several scalar functions $`f(x,𝐳)`$ (as general Gel’fand–Yaglom equations and, in particular, Bhabha equations) or describe objects with composite spin, which correspond to the functions $`f(x,𝐳_{(1)},\mathrm{}𝐳_{(n)})`$ of one set of space-time coordinates $`x`$ and several sets of spin coordinates $`𝐳`$ (as Ivanenko-Landau-Kähler or Dirac-Kähler equation).
## VII Conclusion
In this paper we have elaborated a general scheme of analysis of fields on the Poincaré group and have applied it in two, three and four-dimensional cases.
Considering the left GRR of the Poincaré group, we introduce the scalar field $`f(x,𝐳)`$ on the group, where $`x`$ are coordinates in Minkowski space and $`𝐳`$ are coordinates on the Lorentz group. The connection between the left GRR and the scalar field allows one to use the powerful mathematical method of harmonic analysis on a group, at the same time supporting the consideration by physical motivations.
The consideration of the functions $`f(x,𝐳)`$ guarantees the possibility to describe arbitrary spin particles because any irrep of a group is equivalent to some sub-representation of GRR. Thus we deal with an unique field containing all masses and spins. As a consequence, we have:
1. The explicit form of spin projection operators does not depend on the spin value. These operators are the differential operators with respect to $`𝐳`$.
2. For this scalar field and thus for arbitrary spin discrete transformations $`C,P,T`$ are defined as the automorphisms of the Poincaré group.
3. RWE arise under the classification of the functions on the Poincaré group by eigenvalues of invariant operators and have the same form for arbitrary spin.
The switch to the usual multicomponent description by functions $`\psi _n(x)`$ corresponds to a separation of the space-time and spin variables, $`f(x,𝐳)=\varphi _n(𝐳)\psi _n(x)`$, where $`\varphi _n(𝐳)`$ and $`\psi _n(x)`$ transform under contragradient representations of the Lorentz group. The use of the transformation rules of $`x,𝐳`$ under automorphisms enables us to deduce the transformation rules of $`\psi _n(x)`$ under $`C,P,T`$ without any consideration of the specific form of equations of motion.
It is shown that in even dimensions the consistent consideration of invariant with respect to space reflection RWE requires to use the generators of group $`SO(D,2)`$, which is an extension of the corresponding Lorentz group $`SO(D,1)`$.
The interpretation of the right generators belonging to the complete set of commuting operators on the Poincaré group is given. This interpretation is similar to Wigner and Casimir interpretation of right generators of the rotation group in the nonrelativistic theory (see ). Like in the nonrelativistic case, right generators define some quantum numbers, which do not depend on the choice of the laboratory frame. In particular, in 3+1-dimensional case three right generators of the Poincaré group define Lorentz characteristics $`j_1,j_2`$ and chirality, and fourth right generator distinguishes particles and antiparticles.
Using complete sets of the commuting operators on the group, we classify scalar functions $`f(x,𝐳)`$. As one of the results of this classification we reproduce essentially all known finite-component RWE. Moreover, such an approach allows one to consider some alternative possibilities, which have not been formulated before. In particular, in 3+1-dimensional case we write out general subsidiary conditions (247) corresponding to $`sj_1+j_2`$. On the other hand, instead of subsidiary conditions one may consider functions of momentum $`p`$ and spin variables $`𝐳`$ with invariant constraints (243). It is shown that the set of operators related to higher spin equations in 3+1 dimensions obeys commutation relations of $`so(3,3)`$ algebra, which coincide with the algebra of $`\gamma `$-matrices for spin $`1/2`$. But unlike the latter case the set of operators for higher spin equations is not closed with respect to anticommutation.
In the framework of the classification of scalar functions we get also positive energy wave equations allowing probability amplitude interpretation and associated with infinite-dimensional unitary representations of the Lorentz group. Along with the alternative description of integer or half-integer spin fields, just these equations ensure description of fractional spin fields in 1+1 and 2+1 dimensions.
The consideration of the scalar field on the Poincaré group have allowed us both to obtain new results and to reproduce the main results of RWE theory, which earlier were obtained by means of different reasons and methods. Thus a general approach to the construction of different types of RWE is established. Besides, one may consider this approach as an alternative method to construct a detail theory of the Poincaré group representations.
Notice that the approach under consideration can be directly applied to higher dimensional cases and, as we hope, can be generalized to other space-time symmetry groups, such as de Sitter and conformal groups.
## Acknowledgments
This work was partially supported by Brazilian Agencies CNPq (D.M.G.) and FAPESP (D.M.G., A.L.Sh.). The authors would like to thank I L Buchbinder, L A Shelepin, S N Solodukhin, I V Tyutin and A A Deriglazov for useful discussions.
## A Bases of 2+1 Lorentz group representations and $`S^\mu `$ matrices
Spin projection operators $`\widehat{S}^\mu `$ acting in the space of the functions $`f(x,𝐳)`$ of $`x=(x^\mu )`$ and two complex variables $`z^1=z_2,z^2=z_1`$, $`|z_1|^2|z_2|^2=|z^2||z^1|=1`$, have the form
$$\widehat{S}^\mu =\frac{1}{2}(z\gamma ^\mu _z\stackrel{}{}z\stackrel{}{}\gamma ^\mu _\stackrel{}{z}),z=(z^1z^2),_z=(/z^1/z^2)^T,$$
(A1)
where $`\gamma ^\mu =(\sigma _3,i\sigma _2,i\sigma _1)`$. For $`z=(z_1z_2)`$ the relation $`\widehat{S}^\mu =\frac{1}{2}(z\stackrel{}{}\gamma ^\mu _z\stackrel{}{}z\gamma ^\mu _\stackrel{}{z})`$ is valid.
The polynomials of the power $`2S`$ in $`z`$, which correspond to finite-dimensional irreps $`T_S^0`$ of 2+1 Lorentz group, can be written in the form
$$T_S^0:f_S(x,z)=\underset{n=0}{\overset{2S}{}}\varphi ^n(z)\psi _n(x),\varphi ^n(z)=\left(C_{2S}^n\right)^{1/2}(z^1)^{2Sn}(z^2)^n,s^0=Sn,$$
(A2)
where $`s^0`$ is eigenvalue of $`\widehat{S}^0`$, and $`C_{2S}^n`$ are binomial coefficients. The quasipolynomials of the power $`2S1`$, which correspond to infinite-dimensional unitary irreps $`T_S^\pm `$ of 2+1 Lorentz group, can be written in the form
$`T_S^+:f_S(x,z)={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\varphi ^n(z)\psi _n(x),\varphi ^n(z)=\left(C_{2S}^n\right)^{1/2}(z^2)^{2Sn}(z^1)^n,s^0=S+n,`$ (A3)
$`T_S^{}:f_S(x,\stackrel{}{}z)={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\varphi ^n(z)\psi _n(x),\varphi ^n(z)=\left(C_{2S}^n\right)^{1/2}(\stackrel{}{}z^{\dot{2}})^{2Sn}(\stackrel{}{}z^{\dot{1}})^n,s^0=Sn,`$ (A4)
$`C_{2S}^n=\left({\displaystyle \frac{(1)^n\mathrm{\Gamma }(n2S)}{n!\mathrm{\Gamma }(2S)}}\right)^{1/2}.`$ (A5)
There is a correspondence between the action of differential operators $`\widehat{S}^\mu `$ on the functions $`f(x,𝐳)=\varphi (𝐳)\psi (x)`$ and the multiplication of matrices $`\widehat{S}^\mu `$ by columns $`\psi (x)`$ composed of $`\psi _n(x)`$, $`\widehat{S}^\mu f(x,𝐳)=\varphi (𝐳)S^\mu \psi (x)`$. For the finite-dimensional representations $`T_S^0`$ we have $`(S^0)^{}=S^0`$, $`(S^k)^{}=S^k`$,
$`(S^0)_n^n^{}=\delta _{nn^{}}(Sn),n=0,1,\mathrm{},2S,`$ (A6)
$`(S^1)_n^n^{}={\displaystyle \frac{i}{2}}\left(\delta _{nn^{}+1}\sqrt{(2Sn+1)n}+\delta _{n+1n^{}}\sqrt{(2Sn)(n+1)}\right),`$ (A7)
$`(S^2)_n^n^{}={\displaystyle \frac{1}{2}}\left(\delta _{nn^{}+1}\sqrt{(2Sn+1)n}\delta _{n+1n^{}}\sqrt{(2Sn)(n+1)}\right).`$ (A8)
Matrices $`S^\mu `$ satisfy the condition $`(S^\mu )^{}=\mathrm{\Gamma }S^\mu \mathrm{\Gamma }`$, where $`\mathrm{\Gamma }`$ is a diagonal matrix, $`(\mathrm{\Gamma })_n^n^{}=(1)^n\delta _{nn^{}}`$. The substitution $`z\stackrel{}{}z`$ in (A2) changes only signs of $`S^0`$ and $`S^2`$. For representations $`T_S^+`$ of discrete positive series is valid $`(S^\mu )^{}=S^\mu `$,
$`(S^0)_n^n^{}=\delta _{nn^{}}(S+n),n=0,1,2,\mathrm{},`$ (A9)
$`(S^1)_n^n^{}={\displaystyle \frac{1}{2}}\left(\delta _{nn^{}+1}\sqrt{(n12S)n}+\delta _{n+1n^{}}\sqrt{(n2S)(n+1)}\right),`$ (A10)
$`(S^2)_n^n^{}={\displaystyle \frac{i}{2}}\left(\delta _{nn^{}+1}\sqrt{(n12S)n}\delta _{n+1n^{}}\sqrt{(n2S)(n+1)}\right).`$ (A11)
For $`T_S^{}`$ matrices $`S^1`$ have the same form, whereas $`S^0`$, $`S^2`$ change only their signs.
The case of representations of principal series, which is not bounded by the highest (lowest) weight, was considered in .
For the representations, which correspond to finite-dimensional irreps $`T_S^0`$, the decomposition (A2) can be written in terms of symmetric spin-tensors: $`\psi _{\alpha _1\mathrm{}\alpha _{2S}}(x)=\psi _{\alpha _{(1}\mathrm{}\alpha _{2S)}}(x)`$,
$$f_S(x,z)=\psi _{\alpha _1\mathrm{}\alpha _{2S}}(x)z^{\alpha _1}\mathrm{}z^{\alpha _{2S}}.$$
(A12)
Comparing the decompositions (A2) and (A12), we obtain the relation
$$(C_{2S}^n)^{1/2}\psi _n(x)=\psi _{\underset{2Sn}{\underset{}{1\mathrm{}\mathrm{\hspace{0.17em}1}}}\underset{n}{\underset{}{2\mathrm{}\mathrm{\hspace{0.17em}2}}}}(x).$$
(A13)
## B Bases of 3+2 de Sitter and 3+1 Lorentz groups representations and $`\mathrm{\Gamma }^\mu `$ matrices
Consider polynomials of elements of Dirac $`z`$-spinor $`Z_D=(z^\alpha ,\stackrel{}{}\underset{¯}{z}_{\dot{\alpha }})`$. Any polynomial of power $`2S`$ can be decomposed in the basis of $`(2S+3)!/(6(2S)!)`$ monomials
$$(z^1)^a(z^2)^b\stackrel{}{}\underset{¯}{z}_{\dot{1}}^c\stackrel{}{}\underset{¯}{z}_{\dot{2}}^d,a+b+c+d=2S.$$
One may write out 16 operators, which conserve the power of polynomial:
$`\widehat{S}^{\mu \nu }={\displaystyle \frac{1}{2}}((\sigma ^{\mu \nu })_\alpha ^\beta z^\alpha _\beta +(\overline{\sigma }^{\mu \nu })_{\dot{\beta }}^{\dot{\alpha }}\stackrel{}{}\underset{¯}{z}_{\dot{\alpha }}\underset{¯}{}^{\dot{\beta }})c.c.,`$ (B1)
$`\widehat{\mathrm{\Gamma }}^\mu =\widehat{V}_{12}^\mu +\widehat{V}_{21}^\mu c.c.={\displaystyle \frac{1}{2}}(\overline{\sigma }^{\mu \dot{\alpha }\alpha }\stackrel{}{}\underset{¯}{z}_{\dot{\alpha }}_\alpha +\sigma _{}^{\mu }{}_{\alpha \dot{\alpha }}{}^{}z^\alpha \underset{¯}{}^{\dot{\alpha }})c.c.,`$ (B2)
$`\underset{¯}{\overset{^}{\mathrm{\Gamma }}}^\mu =i(\widehat{V}_{12}^\mu \widehat{V}_{21}^\mu )+c.c.={\displaystyle \frac{i}{2}}(\overline{\sigma }^{\mu \dot{\alpha }\alpha }\stackrel{}{}\underset{¯}{z}_{\dot{\alpha }}_\alpha \sigma _{}^{\mu }{}_{\alpha \dot{\alpha }}{}^{}z^\alpha \underset{¯}{}^{\dot{\alpha }})+c.c.,`$ (B3)
$`\widehat{\mathrm{\Gamma }}^5={\displaystyle \frac{1}{2}}\left(z^\alpha _\alpha \stackrel{}{}\underset{¯}{z}_{\dot{\alpha }}\underset{¯}{}^{\dot{\alpha }}\right)+c.c.,`$ (B4)
$`\widehat{𝒯}=\widehat{S}_3^R={\displaystyle \frac{1}{2}}\left(z^\alpha _\alpha +\stackrel{}{}\underset{¯}{z}_{\dot{\alpha }}\underset{¯}{}^{\dot{\alpha }}\right)c.c.,`$ (B5)
where $`_\alpha =/z^\alpha `$, $`\underset{¯}{}^{\dot{\alpha }}=/\stackrel{}{}\underset{¯}{z}_{\dot{\alpha }}`$,
$$(\sigma ^{\mu \nu })_\alpha ^\beta =\frac{i}{4}(\sigma ^\mu \overline{\sigma }^\nu \sigma ^\nu \overline{\sigma }^\mu )_\alpha ^\beta ,(\overline{\sigma }^{\mu \nu })_{\dot{\beta }}^{\dot{\alpha }}=\frac{i}{4}(\overline{\sigma }^\mu \sigma ^\nu \overline{\sigma }^\nu \sigma ^\mu )_{\dot{\beta }}^{\dot{\alpha }},$$
(B6)
and $`c.c.`$ is complex conjugate term corresponding the action in the space of polynomials of the elements of $`\underset{¯}{Z}_D=(\underset{¯}{z}^\alpha ,\stackrel{}{}z_{\dot{\alpha }})`$. Operator $`\widehat{𝒯}`$ commutes with other 15 operators and defines $`(\pm )`$ power of the polynomials, for functions of $`Z_D`$ and $`\underset{¯}{Z}_D`$ respectively. Operators (B1)-(B4) obey commutation relations of $`so(3,3)sl(4,R)`$ algebra,
$`[\widehat{\mathrm{\Gamma }}^5,\widehat{S}^{\mu \nu }]=0,[\widehat{\mathrm{\Gamma }}^5,\widehat{\mathrm{\Gamma }}^\mu ]=i\underset{¯}{\overset{^}{\mathrm{\Gamma }}}^\mu ,[\widehat{\mathrm{\Gamma }}^5,\underset{¯}{\overset{^}{\mathrm{\Gamma }}}^\mu ]=i\widehat{\mathrm{\Gamma }}^\mu ,`$ (B7)
$`[\underset{¯}{\overset{^}{\mathrm{\Gamma }}}^\mu ,\underset{¯}{\overset{^}{\mathrm{\Gamma }}}^\nu ]=i\widehat{S}^{\mu \nu },[\widehat{S}^{\lambda \mu },\underset{¯}{\overset{^}{\mathrm{\Gamma }}}^\nu ]=i(\eta ^{\mu \nu }\underset{¯}{\overset{^}{\mathrm{\Gamma }}}^\lambda \eta ^{\lambda \nu }\underset{¯}{\overset{^}{\mathrm{\Gamma }}}^\mu ),[\widehat{\mathrm{\Gamma }}^\mu ,\underset{¯}{\overset{^}{\mathrm{\Gamma }}}^\nu ]=i\eta ^{\mu \nu }\widehat{\mathrm{\Gamma }}^5,`$ (B8)
see also (263),(264). Using the notations $`\widehat{S}^{4\mu }=\widehat{\mathrm{\Gamma }}^\mu `$, $`\widehat{S}^{5\mu }=\underset{¯}{\overset{^}{\mathrm{\Gamma }}}^\mu `$, $`\widehat{S}^{54}=\widehat{\mathrm{\Gamma }}^5`$, one can rewrite commutation relations in the form (267), where $`\eta _{55}=\eta _{44}=\eta _{00}=1`$, $`\eta _{11}=\eta _{22}=\eta _{33}=1`$. However, for unitary representations of the Poincaré group all the generators and, in particular $`\widehat{B}_3^R=i\widehat{\mathrm{\Gamma }}^5`$ (for functions of $`Z_D`$ and $`\underset{¯}{Z}_D`$ respectively), are Hermitian. Thus, setting $`\widehat{S}^{5\mu }=i\underset{¯}{\overset{^}{\mathrm{\Gamma }}}^\mu `$, $`\widehat{S}^{54}=i\widehat{\mathrm{\Gamma }}^5`$, for these representations it is natural to consider an algebra $`so(4,2)su(2,2)`$ of Hermitian operators.
Supplementing generators $`\widehat{S}^{\mu \nu }`$ of the Lorentz group by four operators $`\widehat{\mathrm{\Gamma }}^\mu `$ (or $`\underset{¯}{\overset{^}{\mathrm{\Gamma }}}^\mu `$), we obtain the algebra of the 3+2 de Sitter group $`SO(3,2)`$. Generators in finite-dimensional representations of $`SO(3,2)`$ obey the relations $`\widehat{\mathrm{\Gamma }}^0=\widehat{\mathrm{\Gamma }}^0`$, $`\widehat{\mathrm{\Gamma }}^k=\widehat{\mathrm{\Gamma }}^k`$.
The linear functions of $`z`$ $`f(x,z)=Z_D\mathrm{\Psi }_D(x)`$ correspond to four-dimensional bispinor representation. In the space of columns $`\mathrm{\Psi }_D(x)`$ the operators act as matrices
$$\widehat{S}^{\mu \nu }\sigma ^{\mu \nu }/2,\widehat{\mathrm{\Gamma }}^\mu \gamma ^\mu /2,\widehat{\mathrm{\Gamma }}^5\gamma ^5/2,\underset{¯}{\overset{^}{\mathrm{\Gamma }}}^\mu i\gamma ^\mu \gamma ^5/2,\widehat{𝒯}1/2.$$
(B9)
In accordance with general theory, Dirac matrices and spin 1 Duffin–Kemmer matrices obey commutation relations of $`so(3,3)`$ algebra .
Using (86)-(88), we get for the action of the discrete transformations on the operators (B1)-(B5):
$$\begin{array}{cccccc}& \widehat{S}^{\mu \nu }& \widehat{\mathrm{\Gamma }}^\mu & \underset{¯}{\overset{^}{\mathrm{\Gamma }}}^\mu & \hfill \widehat{\mathrm{\Gamma }}^5& \hfill \widehat{𝒯}\\ C& 1& 1& 1& \hfill 1& \hfill 1\\ P,T^{}& (1)^{\delta _{0\mu }+\delta _{0\nu }}& (1)^{\delta _{0\mu }}& (1)^{\delta _{0\mu }}& \hfill 1& \hfill 1\\ T_{sch}& (1)^{\delta _{0\mu }+\delta _{0\nu }}& (1)^{\delta _{0\mu }}& (1)^{\delta _{0\mu }}& \hfill 1& \hfill 1\end{array}$$
(B10)
It is possible to construct two linear in $`\widehat{p}^\mu `$ equations for the scalar functions $`f(x,𝐳)`$, which are invariant under proper Poincaré group
$$(\widehat{p}_\mu \widehat{\mathrm{\Gamma }}^\mu \varkappa )f(x,𝐳)=0,(\widehat{p}_\mu \underset{¯}{\overset{^}{\mathrm{\Gamma }}}^\mu \varkappa )f(x,𝐳)=0,$$
(B11)
but in accordance with (B10) only operator $`\widehat{p}_\mu \widehat{\mathrm{\Gamma }}^\mu `$ is invariant under space reflection, operator $`\widehat{p}_\mu \underset{¯}{\overset{^}{\mathrm{\Gamma }}}^\mu `$ changes the sign. Thus only the first equation is invariant under space reflection.
Operators $`\widehat{\mathrm{\Gamma }}^5`$ and $`\widehat{p}_\mu \widehat{\mathrm{\Gamma }}^\mu `$ commute with all the left generators of the Poincaré group but do not commute with each other, $`[\widehat{\mathrm{\Gamma }}^5,\widehat{p}_\mu \widehat{\mathrm{\Gamma }}^\mu ]=i\widehat{p}_\mu \underset{¯}{\overset{^}{\mathrm{\Gamma }}}^\mu `$. Therefore chirality of massive particle describing by the equation $`(p_\mu \widehat{\mathrm{\Gamma }}^\mu ms)f(x,𝐳)=0`$ is uncertain. Operator $`\widehat{𝒯}`$ commutes both with all left generators of the Poincaré group and with operators $`\widehat{\mathrm{\Gamma }}^\mu `$; therefore one may relate to this operator some conserved quantum number changing the sign under charge conjugation.
On the polynomials of four complex variables $`z^\alpha ,\stackrel{}{}\underset{¯}{z}_{\dot{\alpha }}`$ one can realize symmetric irreps $`T_{[2S\mathrm{\hspace{0.17em}0\hspace{0.17em}0}]}`$ of $`SL(4,R)SO(3,3)`$. These irreps are a symmetric part of $`2S`$-times direct product of fundamental four-dimensional irreps $`T_{[\mathrm{1\hspace{0.17em}0\hspace{0.17em}0}]}`$ and remain irreducible after the reduction on the subgroup $`SO(3,2)`$, $`T_{[2S\mathrm{\hspace{0.17em}0\hspace{0.17em}0}]}T_{[2S\mathrm{\hspace{0.17em}0}]}`$. Notice that here we use the notation different from : $`[2Sj]`$ corresponds to $`(j+SS)`$ in the notation of .
We will consider two bases of finite-dimensional irrep $`T_{[2s\mathrm{\hspace{0.17em}0}]}`$ of $`SO(3,2)`$, namely, bases consisting of eigenfunctions of operators $`\widehat{\mathrm{\Gamma }}^5`$ or $`\widehat{\mathrm{\Gamma }}^0`$. The first basis corresponds to chiral representation,
$$\phi _{j_1j_2}^{m_1m_2}(z,\stackrel{}{}\underset{¯}{z})=N^{1/2}(z^1)^{j_1+m_1}(z^2)^{j_1m_1}\stackrel{}{}\underset{¯}{z}_{\dot{1}}^{j_2+m_2}\stackrel{}{}\underset{¯}{z}_{\dot{2}}^{j_2+m_2},$$
(B12)
where $`s=j_1+j_2`$, $`\lambda =j_1j_2`$, $`m_1`$ and $`m_2`$ are eigenvalues of the operators $`\widehat{M_3}`$ and $`\widehat{N_3}`$, which are the linear combinations of $`\widehat{S_3}`$ and $`\widehat{B_3}`$, see (215), $`N=(2s)!/((j_1+m_1)!(j_1m_1)!(j_2+m_2)!(j_2m_2)!)`$. Consisting of eigenfunctions of $`\widehat{\mathrm{\Gamma }}^0`$ basis
$$\varphi _{k_1k_2}^{n_1n_2}(z,\stackrel{}{}\underset{¯}{z})=(N^{})^{1/2}(z^1+\stackrel{}{}\underset{¯}{z}_{\dot{1}})^{k_1+n_1}(z^2+\stackrel{}{}\underset{¯}{z}_{\dot{2}})^{k_1n_1}(z^1\stackrel{}{}\underset{¯}{z}_{\dot{1}})^{k_2+n_2}(z^2\stackrel{}{}\underset{¯}{z}_{\dot{2}})^{k_2n_2},$$
(B13)
where $`s=k_1+k_2`$, $`N^{}=(2s)!/((k_1+n_1)!(k_1n_1)!(k_2+n_2)!(k_2n_2)!)`$, for $`s=1/2`$ corresponds to Dirac representation. The functions (B13) are eigenfunctions of operators $`\widehat{\mathrm{\Gamma }}^0,\underset{¯}{\overset{^}{\mathrm{\Gamma }}}^3,\widehat{S}_3`$ with eigenvalues $`k_1k_2`$, $`i(n_1n_2)/2`$, $`(n_1+n_2)/2`$ respectively. For fixed $`s`$ we have
$$f_s(x,z,\stackrel{}{}\underset{¯}{z})=\underset{j_1+j_2=s}{}\underset{m_1,m_2}{}\psi _{j_1j_2}^{m_1m_2}(x)\phi _{j_1j_2}^{m_1m_2}(z,\stackrel{}{}\underset{¯}{z})=\underset{k_1+k_2=s}{}\underset{n_1,n_2}{}\psi _{k_1k_2}^{n_1n_2}(x)\varphi _{k_1k_2}^{n_1n_2}(z,\stackrel{}{}\underset{¯}{z}).$$
(B14)
Below we will use the basis (B13). According to (280) in the rest frame for a particle describing by the system (277)-(279) we have
$`f(x,z,\stackrel{}{}\underset{¯}{z})=\psi ^+(x)\varphi _{s,s^3}^+(z,\stackrel{}{}\underset{¯}{z})+\psi ^{}(x)\varphi _{s,s^3}^{}(z,\stackrel{}{}\underset{¯}{z})=C_1e^{imx^0}\varphi _{s,s^3}^+(z,\stackrel{}{}\underset{¯}{z})+C_2e^{imx^0}\varphi _{s,s^3}^{}(z,\stackrel{}{}\underset{¯}{z}),`$ (B15)
$`\varphi _{s,s^3}^+(z,\stackrel{}{}\underset{¯}{z})=(z^1+\stackrel{}{}\underset{¯}{z}_{\dot{1}})^{s+s^3}(z^2+\stackrel{}{}\underset{¯}{z}_{\dot{2}})^{ss^3},\varphi _{s,s^3}^{}(z,\stackrel{}{}\underset{¯}{z})=(z^1\stackrel{}{}\underset{¯}{z}_{\dot{1}})^{s+s^3}(z^2\stackrel{}{}\underset{¯}{z}_{\dot{2}})^{ss^3}.`$ (B16)
The equation $`(\widehat{p}_\mu \widehat{\mathrm{\Gamma }}^\mu sm)f(x,z,\stackrel{}{}\underset{¯}{z})=0`$ has the matrix form
$$(\widehat{p}_\mu \mathrm{\Gamma }^\mu sm)\psi (x)=0,$$
(B17)
where $`\psi (x)`$ is a column. It is convenient to enumerate the basis elements (B13) (and the elements of the column $`\psi (x)`$) in order of decrease of $`k_1k_2=s,s1,\mathrm{},s`$. Matrices $`\mathrm{\Gamma }^\mu `$ obey the relations $`\mathrm{\Gamma }^0=\mathrm{\Gamma }^0`$, $`\mathrm{\Gamma }^k=\mathrm{\Gamma }^k`$. Matrix $`\mathrm{\Gamma }^0`$ is diagonal and has the elements $`k_1k_2`$. Matrices $`\mathrm{\Gamma }^1`$ and $`\mathrm{\Gamma }^3`$ are skew-symmetric real, and $`\mathrm{\Gamma }^2`$ is symmetric imaginary. According to (B2) matrices $`\mathrm{\Gamma }^k`$ have nonzero elements only in blocks corresponding to the transitions $`(k_1,k_2)(k_1\pm 1/2,k_21/2)`$. Using this property, it is easy to see that diagonal matrix $`\mathrm{\Gamma }`$ with the elements $`(1)^{2k_2}`$ commutes with $`\mathrm{\Gamma }^0`$ and anticommutes with $`\mathrm{\Gamma }^k`$, $`\mathrm{\Gamma }^\mu =\mathrm{\Gamma }\mathrm{\Gamma }^\mu \mathrm{\Gamma }`$. This allows one to rewrite the Hermitian-conjugate equation $`\psi ^{}(\stackrel{}{\widehat{p}_\mu }\mathrm{\Gamma }^\mu +sm)=0`$ in the form
$$\overline{\psi }(x)(\stackrel{}{\widehat{p}_\mu }\mathrm{\Gamma }^\mu +sm)=0,\overline{\psi }=\psi ^{}\mathrm{\Gamma },$$
(B18)
and to define invariant scalar product in the space of columns as $`\overline{\psi }(x)\psi (x)d^3x`$. As a consequence of (B17) and (B18), the continuity equation holds
$$_\mu j^\mu =0,j^\mu =\overline{\psi }\mathrm{\Gamma }^\mu \psi .$$
Now the question concerning the positive definiteness of current vector component $`j^0`$ and energy density may be consider similarly to 2+1-dimensional case, see Section 3. For half-integer spin particles describing by the system (277)-(279) charge density $`j^0`$ is positive definite, since in the rest frame (see (B16)) $`j^0=\overline{\psi }\mathrm{\Gamma }\mathrm{\Gamma }^0\psi =s(|\psi ^+(x)|+|\psi ^{}(x)|)>0`$. Energy density (defined in terms of energy-momentum tensor (171)) and the scalar product $`\overline{\psi }\psi `$ are indefinite since in the rest frame they are proportional to $`|\psi ^+(x)||\psi ^{}(x)|`$. For integer spin particles energy density is positive definite, the scalar product and $`j^0`$ are indefinite.
Consider discrete transformations in terms of the columns $`\psi (x)`$. According to (208) under space reflection $`\varphi _{k_1k_2}^{n_1n_2}(z,\stackrel{}{}\underset{¯}{z})(1)^{2k_1}\varphi _{k_1k_2}^{n_1n_2}(z,\stackrel{}{}\underset{¯}{z})`$. Whence, taking into account $`f(x,z,\stackrel{}{}\underset{¯}{z})f(x^{},z^{})=\varphi (z,\stackrel{}{}\underset{¯}{z})\psi ^{}(x^{})`$, we get
$$\psi (x)\stackrel{P}{}(1)^{2s}\mathrm{\Gamma }\psi (\overline{x}),\text{where}\overline{x}=(x^0,x^k).$$
(B19)
According to (88) under charge conjugation $`\varphi _{k_1k_2}^{n_1n_2}(z,\stackrel{}{}\underset{¯}{z})\varphi _{k_1k_2}^{n_1n_2}(\stackrel{}{}z,\underset{¯}{z})`$. Taking into account that $`\varphi _{k_1k_2}^{n_1n_2}(\stackrel{}{}z,\underset{¯}{z})`$ and $`(1)^{s+n_1n_2}\varphi _{k_2k_1}^{n_2n_1}(z,\stackrel{}{}\underset{¯}{z})`$ have the same transformation rule, we get
$$\psi _{k_1k_2}^{n_1n_2}(x)\stackrel{C}{}(1)^{s+n_1n_2}\stackrel{}{}\psi _{k_2k_1}^{n_2n_1}(x).$$
(B20)
In particular, for $`s=1/2`$, using the relation $`f(x,z,\stackrel{}{}\underset{¯}{z})=Z_D\mathrm{\Psi }(x)`$, we get $`\mathrm{\Psi }(x)\stackrel{P}{}\gamma ^0\mathrm{\Psi }(\overline{x})`$, $`\mathrm{\Psi }(x)\stackrel{C}{}\mathrm{\Psi }^c(x)=C\overline{\mathrm{\Psi }}^T(x)`$, where $`C`$ is the matrix with elements $`i\sigma _2`$ on secondary diagonal, $`C=i\gamma ^2\gamma ^0`$. The transformation properties of the bilinears $`\overline{\psi }\mathrm{\Gamma }^\mu \psi `$, $`\overline{\psi }\mathrm{\Gamma }^5\psi `$, $`\overline{\psi }\underset{¯}{\mathrm{\Gamma }}^\mu \psi `$ under $`C,P,T`$ coincide with ones of the corresponding operators, see (B10).
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# The Stern-Gerlach interaction between a traveling particle and a time varying magnetic field
## 1 Introduction
The Stern-Gerlach force acts on particles, carrying a magnetic moment, which cross inhomogeneous magnetic fields. In a reference frame where particles are at rest, the expression of this force is
$$\stackrel{}{f}_{SG}=U$$
(1)
where
$$U=\stackrel{}{\mu }\stackrel{}{B}$$
(2)
is the magnetic potential energy, and
$$\stackrel{}{\mu }=g\frac{e}{2m}\stackrel{}{S}$$
(3)
is the magnetic moment. Here $`e=\pm 1.602\times 10^{19}\mathrm{C}`$ is the elementary charge with $`+`$ for $`p,e^+`$ and $``$ for $`\overline{p},e^{}`$, making $`\stackrel{}{\mu }`$ and $`\stackrel{}{S}`$ either parallel or antiparallel, respectively. The rest mass, $`m`$, is $`1.67\times 10^{27}\mathrm{kg}`$ for $`p,\overline{p}`$ and $`9.11\times 10^{31}\mathrm{kg}`$ for $`e^\pm `$, and the relation between the gyromagnetic ratio $`g`$ and the anomaly $`a`$ is
$$a=\frac{g2}{2}=\{\begin{array}{cc}1.793(g=5.586)\mathrm{for}p,\overline{p}\hfill & \\ 1.160\times 10^3(g=2.002)fore^\pm \hfill & \end{array}$$
(4)
In the rest system, the quantum vector $`\stackrel{}{S}`$, named spin, has modulus $`|\stackrel{}{S}|=\sqrt{s(s+1)}\mathrm{}`$, and its component parallel to the magnetic field lines can take only the following values:
$$S_m=(s,s+1,\mathrm{}.,s1,s)\mathrm{},$$
(5)
where $`\mathrm{}=1.05\times 10^{34}\mathrm{Js}`$ the reduced Planck’s constant. Combining Eqs. (3) and (5) we obtain for a generic spin-$`\frac{1}{2}`$ fermion
$$\mu =|\stackrel{}{\mu }|=g\frac{|e|\mathrm{}}{4m}$$
(6)
or
$$\mu =\{\begin{array}{cc}1.41\times 10^{26}\mathrm{JT}^1\hfill & \\ 9.28\times 10^{24}\mathrm{JT}^1\hfill & \end{array}$$
(7)
Take note that the Bohr magneton is
$$\mu _B=2[\mu /g]_{\mathrm{electron}}=9.27\times 10^{24}\mathrm{JT}^1$$
(8)
Aiming to have the expression of the Stern-Gerlach force in the laboratory frame, we have first to carry out the Lorentz transformation of the electric and magnetic field from the laboratory frame, where we are at rest, to the center-of-mass frame, where particles are at rest and we can correctly evaluate such a force. Then this force must be boosted back to the laboratory frame. All of these rather cumbersome operations will be discussed in the next Section.
## 2 Lorentz Boost of a Force
In order to accomplish the sequence of Lorentz boosts more easily, we choose a Cartesian 4-dimensional Minkowski metric $`(x_1,x_2,x_3,x_4)=(x,y,z,ict)`$, where $`i=\sqrt{1}`$. Therefore, the back-and-forth Lorentz transformations between laboratory frame and particle’s rest frame (usually labeled with a prime) are the following:
$$\left(\begin{array}{c}x^{}\\ y^{}\\ z^{}\\ ict^{}\end{array}\right)=M\left(\begin{array}{c}x\\ y\\ z\\ ict\end{array}\right)=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& \gamma & i\beta \gamma \\ 0& 0& i\beta \gamma & \gamma \end{array}\right)\left(\begin{array}{c}x\\ y\\ z\\ ict\end{array}\right)\{\begin{array}{cc}x^{}=x\hfill & \\ y^{}=y\hfill & \\ z^{}=\gamma (z\beta ct)\hfill & \\ t^{}=\gamma \left(t\frac{\beta }{c}z\right)\hfill & \end{array}$$
(9)
$$\left\{\beta =|\stackrel{}{\beta }|=\frac{|\stackrel{}{v}|}{c},\gamma =\frac{1}{\sqrt{1\beta ^2}}\right\}$$
and
$$\left(\begin{array}{c}x\\ y\\ z\\ ict\end{array}\right)=M^1\left(\begin{array}{c}x^{}\\ y^{}\\ z^{}\\ ict^{}\end{array}\right)=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& \gamma & i\beta \gamma \\ 0& 0& i\beta \gamma & \gamma \end{array}\right)\left(\begin{array}{c}x^{}\\ y^{}\\ z^{}\\ ict^{}\end{array}\right)\{\begin{array}{cc}x=x^{}\hfill & \\ y=y^{}\hfill & \\ z=\gamma (z^{}+\beta ct^{})\hfill & \\ t=\gamma \left(t^{}+\frac{\beta }{c}z^{}\right)\hfill & \end{array}$$
(10)
Moreover, combining both eqs. (9) and (10), we obtain the following expressions for the partial derivatives:
$$\frac{}{x^{}}=\frac{}{x},\frac{}{y^{}}=\frac{}{y}$$
(11)
$$\frac{}{z^{}}=\gamma \left(\frac{}{z}+\frac{\beta }{c}\frac{}{t}\right)$$
(12)
The 4-vector formalism is still applied for undergoing the Lorentz transformation of a force. First of all, let us define as 4-velocity the quantity
$$u_\mu =\frac{dx_\mu }{d\tau }$$
(13)
where
$$d\tau =\frac{ds}{c}=\frac{dt}{\gamma }$$
(14)
is the differential of the proper time. We define the 4-momentum as the product of the rest mass $`m`$ times the 4-velocity, i.e.
$$P_\mu =mu_\mu =(\stackrel{}{p},i\gamma mc)$$
(15)
The 4-force is the derivative of the 4-momentum (15) with respect to the proper time, that is
$$F_\mu =\frac{dP_\mu }{d\tau }=(\gamma \frac{d\stackrel{}{p}}{dt},i\frac{\gamma }{c}\frac{d(\gamma mc^2)}{dt})=(\gamma \stackrel{}{f},i\frac{\gamma }{c}\frac{dE_{\mathrm{tot}}}{dt})$$
(16)
where $`\stackrel{}{f}`$ is the ordinary force. In the c.m. system eq. (16) reduces to
$$F_\mu ^{}=(\stackrel{}{f}^{},0)$$
(17)
since $`\gamma ^{}=1`$ and $`E_{\mathrm{tot}}^{}=mc^2`$ is a constant. Bearing in mind the last step of the whole procedure, i.e. the boost of any force from rest to laboratory frame, we have to use the relation
$$F_\mu =M^1F_\mu ^{}=\left(\begin{array}{c}\gamma f_x\\ \gamma f_y\\ \gamma f_z\\ F_4\end{array}\right)=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& \gamma & i\beta \gamma \\ 0& 0& i\beta \gamma & \gamma \end{array}\right)\left(\begin{array}{c}f_x^{}\\ f_y^{}\\ f_z^{}\\ 0\end{array}\right)=\left(\begin{array}{c}f_x^{}\\ f_y^{}\\ \gamma f_z^{}\\ i\beta \gamma f_z^{}\end{array}\right)$$
(18)
or
$$\stackrel{}{f}_{}=\frac{1}{\gamma }\stackrel{}{f}_{}^{}$$
(19)
$$\stackrel{}{f}_{}=\stackrel{}{f}_{}^{}(f_z=f_z^{})$$
(20)
## 3 Stern-Gerlach Force
The Stern-Gerlach force, as described by eq. (1), must be evaluated in the particle rest frame where it takes the form
$$\stackrel{}{f}_{SG}^{}=^{}(\stackrel{}{\mu }^{}\stackrel{}{B}^{})=\frac{}{x^{}}(\stackrel{}{\mu }^{}\stackrel{}{B}^{})\widehat{x}+\frac{}{y^{}}(\stackrel{}{\mu }^{}\stackrel{}{B}^{})\widehat{y}+\frac{}{z^{}}(\stackrel{}{\mu }^{}\stackrel{}{B}^{})\widehat{z}$$
(21)
having defined the magnetic moment as $`\mu ^{}`$, rather than $`\mu ^{}`$, for opportune reasons. By applying the transformations (11), (19) and (20), the force (21) is boosted to the laboratory system becoming
$$\stackrel{}{f}_{SG}=\frac{1}{\gamma }\frac{}{x}(\stackrel{}{\mu }^{}\stackrel{}{B}^{})\widehat{x}+\frac{1}{\gamma }\frac{}{y}(\stackrel{}{\mu }^{}\stackrel{}{B}^{})\widehat{y}+\frac{}{z^{}}(\stackrel{}{\mu }^{}\stackrel{}{B}^{})\widehat{z}$$
(22)
Bearing in mind the Lorentz transformation of the fields $`\stackrel{}{E},\stackrel{}{B}`$ and $`\stackrel{}{E}^{},\stackrel{}{B}^{}`$
$$\stackrel{}{E}^{}=\gamma (\stackrel{}{E}+c\stackrel{}{\beta }\times \stackrel{}{B})\frac{\gamma ^2}{\gamma +1}\stackrel{}{\beta }(\stackrel{}{\beta }\stackrel{}{E})$$
(23)
$$\stackrel{}{B}^{}=\gamma \left(\stackrel{}{B}\frac{\stackrel{}{\beta }}{c}\times \stackrel{}{E}\right)\frac{\gamma ^2}{\gamma +1}\stackrel{}{\beta }(\stackrel{}{\beta }\stackrel{}{B})$$
(24)
the energy $`(\stackrel{}{\mu }^{}\stackrel{}{B}^{})=\mu _xB_x^{}+\mu _yB_y^{}+\mu _zB_z^{}`$ becomes
$$(\stackrel{}{\mu }^{}\stackrel{}{B}^{})=\gamma \mu _x^{}\left(B_x+\frac{\beta }{c}E_y\right)+\gamma \mu _y^{}\left(B_y\frac{\beta }{c}E_x\right)+\mu _z^{}B_z$$
(25)
If we introduce eq. (25) into eq. (22) and take into account eq. (12), we can finally obtain the Stern-Gerlach force components in the laboratory frame:
$$f_x=\mu _x^{}\left(\frac{B_x}{x}+\frac{\beta }{c}\frac{E_y}{x}\right)+\mu _y^{}\left(\frac{B_y}{x}\frac{\beta }{c}\frac{E_x}{x}\right)+\frac{1}{\gamma }\mu _z^{}\frac{B_z}{x}$$
(26)
$$f_y=\mu _x^{}\left(\frac{B_x}{y}+\frac{\beta }{c}\frac{E_y}{y}\right)+\mu _y^{}\left(\frac{B_y}{y}\frac{\beta }{c}\frac{E_x}{y}\right)+\frac{1}{\gamma }\mu _z^{}\frac{B_z}{y}$$
(27)
$$f_z=\mu _x^{}C_{zx}+\mu _y^{}C_{zy}+\mu _z^{}C_{zz}$$
(28)
with
$$C_{zx}=\gamma ^2\left[\left(\frac{B_x}{z}+\frac{\beta }{c}\frac{B_x}{t}\right)+\frac{\beta }{c}\left(\frac{E_y}{z}+\frac{\beta }{c}\frac{E_y}{t}\right)\right]$$
(29)
$$C_{zy}=\gamma ^2\left[\left(\frac{B_y}{z}+\frac{\beta }{c}\frac{B_y}{t}\right)\frac{\beta }{c}\left(\frac{E_x}{z}+\frac{\beta }{c}\frac{E_x}{t}\right)\right]$$
(30)
$$C_{zz}=\gamma \left(\frac{B_z}{z}+\frac{\beta }{c}\frac{B_z}{t}\right)$$
(31)
## 4 The Rectangular Cavity
In order to simplify our calculations without loosing the general physical meaning, we shall consider a rectangular resonator, as the one shown in Fig.1, which is characterized by the following field components:
$$B_x=\frac{B_0}{K_c^2}\left(\frac{m\pi }{a}\right)\left(\frac{p\pi }{d}\right)\mathrm{sin}\left(\frac{m\pi x}{a}\right)\mathrm{cos}\left(\frac{n\pi y}{b}\right)\mathrm{cos}\left(\frac{p\pi z}{d}\right)\mathrm{cos}\omega t$$
(32)
$$B_y=\frac{B_0}{K_c^2}\left(\frac{n\pi }{b}\right)\left(\frac{p\pi }{d}\right)\mathrm{cos}\left(\frac{m\pi x}{a}\right)\mathrm{sin}\left(\frac{n\pi y}{b}\right)\mathrm{cos}\left(\frac{p\pi z}{d}\right)\mathrm{cos}\omega t$$
(33)
$$B_z=B_0\mathrm{cos}\left(\frac{m\pi x}{a}\right)\mathrm{cos}\left(\frac{n\pi y}{b}\right)\mathrm{sin}\left(\frac{p\pi z}{d}\right)\mathrm{cos}\omega t$$
(34)
$$E_x=B_0\left(\frac{n\pi }{b}\right)\frac{\omega }{K_c^2}\mathrm{cos}\left(\frac{m\pi x}{a}\right)\mathrm{sin}\left(\frac{n\pi y}{b}\right)\mathrm{sin}\left(\frac{p\pi z}{d}\right)\mathrm{sin}\omega t$$
(35)
$$E_y=B_0\left(\frac{n\pi }{b}\right)\frac{\omega }{K_c^2}\mathrm{sin}\left(\frac{m\pi x}{a}\right)\mathrm{cos}\left(\frac{n\pi y}{b}\right)\mathrm{sin}\left(\frac{p\pi z}{d}\right)\mathrm{sin}\omega t$$
(36)
$$E_z=0(\mathrm{as}\mathrm{typical}\mathrm{for}\mathrm{a}\mathrm{TE}\mathrm{mode})$$
(37)
where $`B_0`$ is the amplitude of the $`B_z`$-component and
$$K_c=\sqrt{\left(\frac{m\pi }{a}\right)^2+\left(\frac{n\pi }{b}\right)^2}$$
(38)
$$\frac{\omega }{c}=K=\frac{2\pi }{\lambda }=\sqrt{\left(\frac{m\pi }{a}\right)^2+\left(\frac{n\pi }{b}\right)^2+\left(\frac{p\pi }{d}\right)^2}$$
(39)
The wave’s phase velocity is $`v_{\mathrm{ph}}=\beta _{\mathrm{ph}}c`$ where
$$\beta _{\mathrm{ph}}=\frac{K}{\sqrt{K^2K_c^2}}=\sqrt{1+\left(\frac{md}{pa}\right)^2+\left(\frac{nd}{pb}\right)^2}$$
(40)
We have to recall that the polarization of a beam, revolving in a ring whose guide field is $`\stackrel{}{B}_{\mathrm{ring}}`$, can be defined as
$$P=\frac{N_{}N_{}}{N_{}+N_{}}$$
(41)
where $`N_{}`$ = No. Particles Spin Up (e.g. parallel to $`\stackrel{}{B}_{\mathrm{ring}}`$)
$`N_{}`$ = No. Particles Spin Down (antiparallel to $`\stackrel{}{B}_{\mathrm{ring}}`$)
and $`P`$ indicates the macroscopic average over the particle distribution in the beam, which is equivalent to the quantum mechanical expectation value found by means of the quantum statistical matrix. Obviously, an unpolarized beam has $`P=0`$ or $`N_{}`$ = $`N_{}`$.
A quick comparison among the SG-force components, given by the set of equations (26)-(31), suggests that $`f_z`$ will dominate at high energy, since it contains terms proportional to $`\gamma ^2`$, whereas the transverse components have terms independent of $`\gamma `$, not to mention the $`\gamma ^1`$ terms.
The most appropriate choice of the spin orientation seems to be the one parallel to $`\widehat{y}`$ i.e. to $`\stackrel{}{B}_{\mathrm{ring}}`$, i.e. the force component is the one given by eq. (28) with the insertion of eq. (30). This means that particles undergoing energy gain (or loss) don’t need any spin rotation while entering and leaving the rf cavity, beyond the advantage of having to deal with a force component proportional to $`\gamma ^2`$. Choosing the simplest $`\mathrm{TE}_{011}`$ mode, the quantities (38), (39) and (40) reduce to
$$k_c=\frac{\pi }{b}$$
(42)
$$\omega =c\sqrt{\left(\frac{\pi }{b}\right)^2+\left(\frac{\pi }{d}\right)^2}$$
(43)
$$\beta _{\mathrm{ph}}=\sqrt{1+\left(\frac{d}{b}\right)^2}$$
(44)
Setting $`x=\frac{a}{2}`$ and $`y=\frac{b}{2}`$ the field components along the beam axis become
$$B_x=B_z=0$$
(45)
$$B_y=B_0\frac{b}{d}\mathrm{cos}\left(\frac{\pi z}{d}\right)\mathrm{cos}\omega t$$
(46)
$$E_x=\omega B_0\frac{b}{\pi }\mathrm{sin}\left(\frac{\pi z}{d}\right)\mathrm{sin}\omega t$$
(47)
$$E_y=E_z=0$$
(48)
therefore the force component $`f_z`$ can be written as
$$f_z=\mu ^{}\gamma ^2B_0b\left\{\frac{1}{\pi }\left[\left(\frac{\pi }{d}\right)^2+\left(\frac{\beta \omega }{c}\right)^2\right]\mathrm{sin}\left(\frac{\pi z}{d}\right)\mathrm{cos}\omega t+\frac{2}{d}\left(\frac{\beta \omega }{c}\right)\mathrm{cos}\left(\frac{\pi z}{d}\right)\mathrm{sin}\omega t\right\}$$
(49)
For completeness, we shall also analyze the possibility of using a spin orientation parallel to $`\widehat{z}`$, i.e. to the motion direction, even though this option requires a system of spin rotators and looses a factor of $`\gamma `$ in the force component.
## 5 Involved Energy
The energy gained, or lost, by a particle with a magnetic moment after having crossed a rf cavity can be evaluated by integrating the Stern-Gerlach force (22) over the cavity length, namely:
$$\mathrm{\Delta }U=_0^d𝑑U=_0^d\stackrel{}{f}𝑑\stackrel{}{r}=_0^df_z𝑑z=_0^d\mu ^{}C_{zy}𝑑z$$
(50)
Bearing in mind eq. (49) and carrying out the trivial substitution $`\omega t=\frac{\omega z}{\beta c}`$, the integral (50) becomes
$$\mathrm{\Delta }U=\mu ^{}\gamma ^2B_0b\left\{\frac{1}{\pi }\left[\left(\frac{\pi }{d}\right)^2+\left(\frac{\beta \omega }{c}\right)^2\right]I_1+\frac{2}{d}\left(\frac{\beta \omega }{c}\right)I_2\right\}$$
with
$$I_1=_0^d\mathrm{sin}\left(\frac{\pi z}{d}\right)\mathrm{cos}\left(\frac{\omega z}{\beta c}\right)𝑑z=\frac{\frac{\pi }{d}}{\left(\frac{\pi }{d}\right)^2\left(\frac{\omega }{\beta c}\right)^2}\left[1+\mathrm{cos}\left(\frac{\omega d}{\beta c}\right)\right]$$
$$I_2=_0^d\mathrm{cos}\left(\frac{\pi z}{d}\right)\mathrm{sin}\left(\frac{\omega z}{\beta c}\right)𝑑z=\frac{\frac{\omega }{\beta c}}{\left(\frac{\pi }{d}\right)^2\left(\frac{\omega }{\beta c}\right)^2}\left[1+\mathrm{cos}\left(\frac{\omega d}{\beta c}\right)\right]$$
or
$$\mathrm{\Delta }U=\mu ^{}\gamma ^2B_0\frac{b}{d}\frac{\left(\frac{\pi }{d}\right)^2+\left(\frac{\beta \omega }{c}\right)^22\left(\frac{\omega }{c}\right)^2}{\left(\frac{\pi }{d}\right)^2\left(\frac{\omega }{\beta c}\right)^2}\left[1+\mathrm{cos}\left(\frac{\omega d}{\beta c}\right)\right]$$
(51)
Taking into account the stationary wave conditions (eqs. 43 and 44) pertaining to the $`\mathrm{TE}_{011}`$ mode, the length of the cavity can be expressed as
$$d=\frac{1}{2}\beta _{\mathrm{ph}}\lambda $$
(52)
which allows us to write eq. (51) as
$$\mathrm{\Delta }U=\gamma ^2\beta ^2\mu ^{}B_0\frac{b}{d}\frac{1+\beta _{\mathrm{ph}}^2(\beta ^22)}{\beta ^2\beta _{\mathrm{ph}}^2}\left(1+\mathrm{cos}\frac{\beta _{\mathrm{ph}}}{\beta }\pi \right)$$
(53)
In the ultrarelativistic limit ($`\gamma 1`$ and $`\beta 1`$),
$$\mathrm{\Delta }U\mu ^{}B_0\frac{b}{d}\gamma ^2(1+\mathrm{cos}\beta _{\mathrm{ph}}\pi )=2\mu ^{}B_0\frac{b}{d}\gamma ^2(\beta _{\mathrm{ph}}=\text{even integer})$$
(54)
As hinted before, let us evaluate the work-energy integral when the particle enters into the cavity with its spin parallel to $`\widehat{z}`$. In this example we must choose the mode $`\mathrm{TE}_{021}`$ as the lowest one; then we have from eqs. (34) and (31) respectively
$$B_z=B_0\mathrm{sin}\left(\frac{\pi z}{d}\right)\mathrm{cos}\omega t$$
(55)
$$f_z=\mu ^{}C_{zz}=\mu ^{}\gamma B_0\left[\frac{\pi }{d}\mathrm{cos}\left(\frac{\pi z}{d}\right)\mathrm{cos}\omega t\left(\frac{\beta \omega }{c}\right)\mathrm{sin}\left(\frac{\pi z}{d}\right)\mathrm{cos}\omega t\right]$$
(56)
and proceeding as above we obtain
$$\mathrm{\Delta }U=\mu ^{}B_0\gamma \frac{\pi }{d}\frac{\frac{\omega }{\beta c}\frac{\beta c}{\omega }}{\left(\frac{\pi }{d}\right)^2\left(\frac{\omega }{\beta c}\right)^2}\mathrm{sin}\left(\frac{\omega d}{\beta c}\right)$$
(57)
and
$$\mathrm{\Delta }U=\frac{\mu ^{}B_0}{\gamma }\frac{\beta _{\mathrm{ph}}\beta }{\beta _{\mathrm{ph}}^2\beta ^2}\mathrm{sin}\left(\frac{\beta _{\mathrm{ph}}}{\beta }\pi \right)$$
(58)
or ultrarelativistically
$$\mathrm{\Delta }U\frac{\mu ^{}B_0}{\gamma }\frac{\beta _{\mathrm{ph}}}{\beta _{\mathrm{ph}}^21}\mathrm{sin}\beta _{\mathrm{ph}}\pi ,\mathrm{\Delta }U_{\mathrm{max}}1.62\frac{\mu ^{}B_0}{\gamma }(\text{when }\beta _{\mathrm{ph}}1.13)$$
(59)
confirming a result already achieved.
Before making up our mind, we need to compare the energy gain/loss due to the Stern-Gerlach interaction with the same quantity caused by the electric field. To this aim, we emphasize that
$$dU_E=\stackrel{}{f}_Ed\stackrel{}{r}=eE_xdx$$
(60)
as can be easily understood looking at eqs. (47) and (48). Since the carrier particle travels from 0 to $`d`$ along the $`z`$-axis, the only integral which makes sense is the following:
$$\mathrm{\Delta }U_E=_0^deE_x𝑑x=_0^deE_x\frac{dx}{dz}𝑑z=_0^deE_xx^{}𝑑z$$
(61)
or
$$\mathrm{\Delta }U_E=x^{}e\omega B_0\frac{b}{\pi }_0^d\mathrm{sin}\left(\frac{\pi z}{d}\right)\mathrm{sin}\left(\frac{\omega z}{\beta c}\right)𝑑z=x^{}e\omega B_0\frac{b}{d}\frac{\mathrm{sin}\left(\frac{\omega d}{\beta c}\right)}{\left(\frac{\pi }{d}\right)^2\left(\frac{\omega }{\beta c}\right)^2}$$
or
$$\mathrm{\Delta }U_E=\left[e\omega B_0\frac{bd}{\pi ^2}\frac{\beta ^2}{\beta _{\mathrm{ph}}^2\beta ^2}\mathrm{sin}\frac{\beta _{\mathrm{ph}}}{\beta }\pi \right]x^{}=\kappa x^{}$$
(62)
having proceeded as before.
We recall that the Stern-Gerlach interaction in the realm of particle accelerators has been proposed either for separating in energy particles with opposite spin states, the well known spin-splitter concept, or for settling an absolute polarimeter .
As far as the spin-splitter is concerned, we quickly recall that spin up particles receive (or loose) that amount of energy given by eq. (54) at each rf cavity crossing, and this will take place all over the time required. Simultaneously, spin down particles behave exactly in the opposite way, i.e. they loose (or gain) the same amount of energy turn after turn. The actual most important issue is that the energy exchanges sum up coherently. More quantitatively, we may indicate as the final energy separation after $`N`$ revolutions:
$$\mathrm{\Delta }_{}=\{\mathrm{\Delta }_{}(\mathrm{\Delta }_{})\}=4\frac{b}{d}N\mu ^{}B_0\gamma ^24N\mu ^{}B_0\gamma ^2$$
(63)
Instead, the adding up of the energy contribution (62) due to the electric field is
$$(\mathrm{\Delta }U_E)_{tot}=\mathrm{\Delta }U_E=\kappa x^{}=0$$
(64)
since $`x^{}`$ changes continuously its sign with a periodicity related to the period of the betatron oscillations.
The result (63), together with the demonstration (64), would seem to provide very good news for the spin-splitter method!
As far as the polarimeter is concerned, we have to bear in mind that we are interested in the instantaneous interaction between magnetic moment and the rf fields: therefore the zero-averaging due to the incoherence of the betatron oscillations would not help us. Notwithstanding, if we set $`\beta _{\mathrm{ph}}`$ equal to an integer in eq. (62), we have for U.R. particles:
$$\mathrm{\Delta }U_E=\frac{x^{}e\omega B_0bd}{\pi ^2(\beta _{\mathrm{ph}}^21)}\mathrm{sin}\left(\beta _{\mathrm{ph}}\pi +\frac{\beta _{\mathrm{ph}}\pi }{2\gamma ^2}\right)\pm \frac{x^{}bd}{2\pi }\frac{\beta _{\mathrm{ph}}}{\beta _{\mathrm{ph}}^21}\frac{e\omega B_0}{\gamma ^2}$$
(65)
Then this $`1/\gamma ^2`$ dependence of the spurious signal, compared to the $`\gamma ^2`$ dependence of the signal (54) to be measured, sounds interesting for the feasibility of this kind of polarimeter; however, one must realize that if $`\beta _{\mathrm{ph}}`$ is not exactly an integer, then eq. (65) would become
$$\mathrm{\Delta }U_E\pm \frac{x^{}bd}{2\pi }\frac{e\omega B_0}{\beta _{\mathrm{ph}}^21}\left(ϵ+\frac{\beta _{\mathrm{ph}}}{\gamma ^2}\right)$$
(66)
where $`ϵ`$ is the error in $`\beta _{\mathrm{ph}}`$.
## 6 A Few Numerical Examples
The spin-splitter principle requires a repetitive crossing of $`N_{\mathrm{cav}}`$ cavities distributed along the ring, each of them resonating in the TE mode. After each revolution, the particle experiences a variation, or kick, of its energy or of its momentum spread
$$\zeta =\frac{\delta p}{p}=\frac{1}{\beta ^2}\frac{\delta E}{E}\frac{N_{\mathrm{cav}}\mathrm{\Delta }U}{E}\frac{2\sqrt{3}}{3}N_{\mathrm{cav}}\frac{B_0}{B_{\mathrm{}}}\gamma $$
(67)
having made use of eq. (54), further simplified by reasonably setting $`\beta _{\mathrm{ph}}=2`$, and with
$$B_{\mathrm{}}=\frac{mc^2}{\mu ^{}}=\frac{1.503\times 10^{10}\mathrm{J}}{1.41\times 10^{26}\mathrm{JT}^1}10^{16}\mathrm{T}$$
(68)
for (anti)protons. From eq. (67) we may find as the number of turns needed for attaining a momentum separation equal to 2$`\left(\frac{\mathrm{\Delta }p}{p}\right)`$
$$N_{\mathrm{SS}}=\frac{\left(\frac{\mathrm{\Delta }p}{p}\right)}{\zeta }=\frac{\sqrt{3}}{2N_{\mathrm{cav}}\gamma }\frac{B_{\mathrm{}}}{B_0}\left(\frac{\mathrm{\Delta }p}{p}\right)$$
(69)
Multiplying $`N_{\mathrm{SS}}`$ by the revolution period $`\tau _{\mathrm{rev}}`$ we obtain
$$\mathrm{\Delta }t=N_{\mathrm{SS}}\tau _{\mathrm{rev}}$$
(70)
as the actual time spent in this operation. For the sake of having some data, we consider RHIC and HERA whose essential parameters are shown in Table I together with what can be found by making use of eqs. (69) and (70) where $`B_00.1T`$ and $`N_{\mathrm{cav}}=200`$ are chosen as realistic values.
Table I: RHIC and HERA parameters
| | RHIC | HERA |
| --- | --- | --- |
| E(GeV) | 250 | 820 |
| $`\gamma `$ | 266.5 | 874.2 |
| $`\tau _{\mathrm{rev}}(\mu \mathrm{s})`$ | 12.8 | 21.1 |
| $`\frac{\mathrm{\Delta }p}{p}`$ | $`4.1\times 10^3`$ | $`5\times 10^5`$ |
| $`N_{\mathrm{SS}}`$ | $`6.67\times 10^9`$ | $`2.48\times 10^7`$ |
| $`\mathrm{\Delta }t`$ | $`8.52\times 10^4s\mathrm{\hspace{0.17em}23.7}h`$ | 523 s |
In the example of the polarimeter we have to pick up a signal generated at each cavity crossing. Therefore, making use of eq. (54) we have for a bunch train made up of $`N`$ particles the total energy transfer
$$\mathrm{\Delta }U2NP\mu ^{}B_0\frac{b}{d}\gamma ^2$$
(71)
where $`P`$ is the beam polarization slightly modified with respect the definition (41)
$$P=\frac{N_{}N_{}}{N_{}+N_{}}$$
(72)
The average power transferred will be
$$W=\frac{\mathrm{\Delta }U}{\tau _{\mathrm{rev}}}$$
(73)
If we operate our cavity as a parametric converter , with an initially empty level, we have for the power transferred to this empty level
$$W_2=\frac{\omega _{\mathrm{rf}}}{\omega _{\mathrm{rev}}}W=\frac{\nu _{\mathrm{rf}}}{\nu _{\mathrm{rev}}}W$$
(74)
where $`\nu _{\mathrm{rf}}`$ is the working frequency of the resonant cavity (typically in the GHz range), and $`\nu _{\mathrm{rev}}`$ is the revolution frequency. Putting all together we have
$$W_22P\frac{\nu _{\mathrm{rf}}}{\nu _{\mathrm{rev}}}\mu ^{}B_0\frac{b}{d}\gamma ^2$$
(75)
A feasibility test of the polarimeter principle has been proposed and studied to be carried out in the 500 MeV electron ring of MIT- Bates, whose main characteristics are
Table II: MIT-Bates parameters
| $`\tau _{\mathrm{rev}}`$ | 634 nsec |
| --- | --- |
| $`\nu _{\mathrm{rev}}`$ | 1.576 MHz |
| $`N_{electrons}`$ | $`3.6\times 10^8225=8.1\times 10^{10}`$ |
| $`\gamma `$ | $`10^3`$ |
| $`b/d`$ | $`\sqrt{3}/3`$ |
| $`B_0`$ | $`0.1`$ T |
| $`\nu _{\mathrm{rf}}/\nu _{\mathrm{rev}}`$ | $`10^3`$ |
| $`\mu ^{}`$ | $`9.27\times 10^{24}\mathrm{JT}^1`$ |
and, since polarized electrons can be injected into this ring but precessing on a horizontal plane, the $`\mathrm{TE}_{101}`$ mode is more appropriate than the $`\mathrm{TE}_{011}`$ as we shall have to use $`B_x`$ rather than $`B_y`$: a choice that does not make any substantial difference! From the above data we obtain
$$W_2137P\mathrm{watts}$$
(76)
Paradoxically, even for an almost unpolarized beam with $`N_{}N_{}=1`$ and, as a consequence of eq. (72), with $`P1.23\times 10^{11}`$, we should obtain $`W_21.7\mathrm{nW}`$, which can be easily measured.
As a last check, let us compare the energy exchanges ($`\stackrel{}{\mu }\stackrel{}{B}`$) and ($`e\stackrel{}{E}`$). Taking into account eqs. (52), (54) and (65), and setting $`x^{}1\mathrm{mrad}`$, $`\beta _{\mathrm{ph}}=2`$ and $`\lambda =10\mathrm{cm}`$, we have for the Bates-MIT ring:
$$r=\frac{\mathrm{\Delta }U_E}{\mathrm{\Delta }U}=\frac{x^{}}{8}\frac{\beta _{\mathrm{ph}}^3}{\beta _{\mathrm{ph}}^21}\frac{\lambda ec}{\mu ^{}}\frac{1}{\gamma ^4}=1.72\times 10^4$$
(77)
i.e. the spurious signal, depending upon the electric interaction between $`e`$ and $`\stackrel{}{E}`$, is absolutely negligible with respect the measurable signal generated by the magnetic interaction.
## 7 Conclusions
There is not too much to add to what has been found in the previous Sections, aside from performing more accurate calculations and numerical simulations. The Stern-Gerlach interaction seems very promising either for attaining the self polarization of a $`p(\overline{p})`$ beam or for realizing an absolute polarimeter.
In the first example the problem raised by the rf filamentation still holds on, although some tricks can be conceived: the extreme one could be the implementation of a triangular waveform in the TM cavity which bunches the beam.
The second example requires nothing but to implement that experimental test at the Bates-MIT electron ring.
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# Mass Loss On the Horizontal Branch: an Application to NGC 6791
## 1. Introduction
Understanding the origin and evolution of hot horizontal branch (HB) stars would have a significant impact on several areas of astronomy. For example, it is important in the study of globular cluster HB morphology, which constrains the formation history of the Galaxy (Searle & Zinn 1978). The global picture of HB morphology appears to be well understood. HB morphology is most affected by the metallicity of the cluster: more metal rich clusters exhibit redder HBs (“the first parameter” of HB morphology). However, observations in some globular clusters show HBs which exhibit color distributions not expected for their metallicities. To explain this deviation, a “second parameter” is invoked.
Age has been the most successful global candidate for the ”second parameter” (Searle & Zinn 1978; Lee, Demarque, & Zinn 1994; for recent review, see Stetson, VandenBerg, & Bolte 1996; Sarajedini, Chaboyer, & Demarque 1997). However, recent observations show an increasing number of HB morphologies with anomalous features, e.g. bimodality in the color (temperature) distribution, extended blue tails, and gaps (e.g., Sosin et al. 1997; Rich et al. 1997; Catelan et al. 1998). Most of these anomalies are found near the blue end of the HB. Within the framework of the canonical theory of single star evolution, they appear to be inconsistent with a scenario where metallicity and age are the only two parameters that govern the HB morphology.
NGC 6791 is an extreme example that exhibits deviations from the first and second parameter phenomena. This open cluster is believed to be extremely metal-rich, \[Fe/H\] $`=0.2`$ – 0.5 (Friel & Janes 1993; Kaluzny & Rucinski 1995; Peterson & Green 1998), which would suggest a predominantly red HB clump, following the first parameter phenomenon. Yet, nearly a third of its helium burning population (HB stars) are extremely blue (Liebert et al. 1994; Green et al. 1997) with a large temperature separation from the red clump. Dynamical effects do not seem to provide a promising explanation for this anomaly, because stellar densities are small in open clusters.
Figure 1 shows the observed color-magnitude diagram (CMD) and a synthetic HB model based on canonical stellar models. The observed CMD of NGC 6791, which is from Kaluzny & Rucinski (1995), shows approximately 23 red clump stars and 9 hot stars. These hot stars have also been identified by ultraviolet (UV) observations using the Ultraviolet Imaging Telescope (Landsman et al. 1998). The synthetic HB (diamonds) is based on the arbitrarily chosen chemical composition, i.e. \[Fe/H\]$`=0.33`$ and $`Y=0.31`$ (equivalent to $`\mathrm{\Delta }Y/\mathrm{\Delta }Z=1.5`$). For illustrative purposes, a 9 Gyr isochrone from the main sequence (MS) to the red giant branch (RGB), is shown on the same plot. For $`E(BV)=0.125`$ and $`mM=13.35`$, the match is reasonable. It is beyond the scope of this paper to derive the age of this cluster or any other parameters by matching the CMD with isochrones.
For the synthetic HB construction, we have assumed a Gaussian mass dispersion parameter $`\sigma _{HB}=0.04M_{}`$, and a Reimers mass loss parameter $`\eta =1.0`$ (Lamers & Cassinelli 1999; Willson, Bowen, & Struck 1996). Readers are referred to Yi, Demarque, & Oemler (1997b) for details of the synthetic HB construction. We assume that there are approximately 32 HB stars in NGC 6791, and we have added some random scatter in both $`M_V`$ and $`BV`$. The model reproduces, in accordance with the first parameter phenomenon, the prominent red clump seen in the data, but it lacks hot stars.
NGC 6791 is the oldest Galactic open cluster known to-date at 6.5-9 Gyr (Demarque, Green, & Guenther 1992; Phelps et al. 1994; Garvanich et al. 1994; Chaboyer, Green, & Liebert 1999). The most recent estimate from Chaboyer et al. (1999) suggests $`8.0\pm 0.5`$ Gyr for the age of this cluster. However, even at the largest age estimate so far (9 Gyr), a synthetic HB based on the conventional input parameters would not reproduce anything even close to the HB morphology of NGC 6791. By increasing its age by a couple of Gyrs, we could produce some hot HB stars, but their numbers would still be much smaller than observed. Thus, the age effect alone seems unsuccessful in explaining the HB color distribution.
A relatively large value of the dispersion ($`\sigma =0.06M_{}`$, in comparison to the conventional value 0.03 – 0.04) in the Gaussian mass distribution on the HB would result in a larger fraction of hot HB stars. This may match the blue-to-red HB star ratio found in NGC 6791 better, but the strong color-bimodality in NGC 6791 would not be reproduced.
One or more local parameters must be at work in this cluster. As candidate mechanisms that produce hot HB stars, a few additional scenarios have been proposed lately. Sweigart (1997) has suggested that helium mixing on the red giant branch would cause further evolution beyond the tip, resulting in higher mass loss on the giant branch. As the color dispersion on the HB is the result of varying envelope masses, the enhanced mass loss results in a bluer HB. D’Cruz et al. (1997) assumed an episode of extreme mass loss at the tip of the giant branch to explain blue HB stars.
Fig. 1 Observed CMD of NGC 6791 and a synthetic HB model. Note the presence of hot HB stars (enclosed in a box) only in the observed CMD.
In this letter, we explore yet another possible formation scenario for blue HB stars: mass loss during the core helium burning phase. Demarque & Eder (1985) first investigated this effect for low metallicity HB stars, and we extend on their work here. We will show that the inclusion of mass loss on the HB works in the direction of explaining the color distribution of the HB stars in NGC 6791.
## 2. Evolutionary Model Construction
We used the Yale stellar evolution code (YREC) for the construction of HB evolutionary models. The detail of physics are the same as in Guenther et al. (1992) with the addition of updated opacities. For log $`T4.00`$, OPAL opacities are used (Iglesias & Rogers 1996), and Alexander low temperature opacities are used for log $`T<4.00`$ (Alexander & Ferguson 1994). The ratio $`(\alpha )`$ of the mixing length to pressure scale height in the convection zone, was assumed to be 1.9, which is the value used in a standard solar model with the same input physics (Guenther et al. 1992). This parameter only affects the radii of the very coolest models, and has little impact on our conclusions. Although He diffusion and complex heavy element segregation processes are known to take place in the atmospheres of sdB and sdO stars (Michaud, Vauclair, & Vauclair 1983; Michaud et al. 1985; Unglaub & Bues 1998), their effect to the internal structure and radius calculation is negligible.
In addition, a mass-loss routine was included into YREC. As we propose no mechanism for mass loss, and any empirical mass loss rate, such as Reimers mass loss (which is based on red supergiant observations), does not apply here, the simplest approach was taken. The mass loss rate, which is supplied by the user, is assumed to be constant over time, and is used to determine the amount of mass loss for a given time step. This mass is then removed from the envelope, which is taken here to be the region of the hydrogen burning shell and outward, at each time step. The model is then allowed to evolve with its new mass. The energy needed to eject the mass against the gravitational field is also calculated and removed from the envelope (see also, Koopmann et al. 1994 and Demarque & Eder 1985).
While mass loss on the HB is an ad hoc assumption, it is possible to justify at least some mass loss. Mass loss is directly observed in the Sun at a rate of $`10^{14}M_{}yr^1`$. It has also been pointed out that mass loss should be enhanced in pulsating stars (Willson 1988). As RR Lyrae stars are pulsators on the horizontal branch, one should expect some mass loss in these stars. There is also both observational and theoretical evidence for pulsation in hot horizontal branch stars. Recent observations show that sdB stars pulsate as well (Kilkenny et al. 1998 and references therein). This pulsation was predicted theoretically by Charpinet et al. (1997a, 1997b). The effect of mass loss in HB stars while they evolve through the instability strip was studied by Koopmann et al. (1994). They found that the mass loss rate has an upper limit of $`10^9M_{}yr^1`$. A higher mass loss rate would break the good agreement between synthetic and observed HBs in the globular cluster M4. This provides an upper limit for the mass loss rate on the HB in the intermediate temperature (RR Lyrae) region.
HB stars hotter than RR Lyrae variables may be in conditions under which larger mass loss may take place. According to Abbott’s formula (1982), the mass loss rate for these stars for the solar composition should be around 2 – 4 $`10^{12}M_{}yr^1`$. We can think of two effects which should increase this number. The first is the fact that the sdB stars pulsate. Second is that the mass loss rate should scale with the metallicity (Abbott 1982; Willson et al. 1996; Lamers & Cassinelli 1999), so the assumption of higher mass loss rate in more metal rich stars, such as the blue HB stars in NGC 6791, could be justified. Although it would be very difficult to observe directly such low mass loss rates in HB stars, there is spectroscopic evidence which supports the presence of mass loss in sdB and sdO atmospheres. The abundance anomalies observed in the atmospheres of sdB and sdO stars, and their lifetimes, can only be explained when the interaction of both diffusion and mass loss are taken into account (Michaud et al. 1983, 1985; Unglaub & Bues 1998 for the sdO stars). More detailed spectroscopy and atmospheric modeling in the future will reveal the time dependence and intensity of these winds.
The mass loss rates used for the present study are in the range $`10^{10}\dot{M}(M_{}yr^1)10^9`$. Models have been constructed for the metallicity range of $`0.004Z0.1`$ and approximately for $`\mathrm{\Delta }Y/\mathrm{\Delta }Z=1.5`$ and 3.
The isochrones for the hydrogen burning phases that we have used in this paper are the Sukyoung Yi version of the Yale Isochrones (Yi, Demarque, & Oemler 1997b; available from http://shemesh.gsfc.nasa.gov). The input physics used in these isochrones are somewhat older than the ones used in our HB models, but they are qualitatively compatible to our HB models.
## 3. Results and Discussion
In order to study the possibility of forming blue HB stars through mass loss during the HB phase of stellar evolution, we have constructed evolutionary tracks for HB stars of varying metallicity and envelope mass. Figure 2 shows HB evolutionary tracks and 6 hot stars in NGC 6791 whose temperatures and luminosities have been measured by Landsman et al. (1998) based on their UV fluxes. The track furthest to the right is the evolutionary track without mass loss and each successive track to the left has a larger mass loss rate; $`10^{10}`$, $`10^{9.5}`$, $`10^9M_{}yr^1`$, respectively. The effect of mass loss shown in Figure 2 is representative of the effects of mass loss in the other models calculated. Several points can be made by examining Figure 2. First, the greater the mass loss rate assumed, the hotter the star becomes. This is not a surprising result. It is generally believed that the HB is a sequence of stars which have the same core mass and varying envelope masses. The smaller the envelope mass, the hotter the star appears, as a deeper and hotter region is exposed with a smaller envelope. It seems natural then that mass loss during the HB phase will produce hotter stars with smaller envelopes.
Fig. 2 Evolutionary tracks for stars starting on the zero-age HB for Z=0.04, $`\mathrm{\Delta }Y/\mathrm{\Delta }Z=3`$ and envelope mass of $`0.07M_{}`$. The tracks of $`\mathrm{\Delta }Y/\mathrm{\Delta }Z=1.5`$ show qualitatively the same phenomenon. The zero-age HB locus, which is insensitive to the adopted mass loss rate, is shown near the bottom of the plot. Tracks with increasing mass loss rates are shifted toward the hotter and fainter regions of the H-R diagram. The crosses indicate time intervals of $`10^7`$ years. Also shown are the hot stars in NGC 6791 whose temperatures and luminosities have been measured from their UV fluxes (Landsman et al. 1998).
At the same time, a larger mass loss rate also results in the stars becoming fainter than the stars without mass loss, particularly for the larger mass loss rates. As our mass loss scheme removes mass from the envelope, which is defined to be the hydrogen burning shell and outwards, the total amount of hydrogen burning is reduced. This results in the stars evolving with smaller luminosities. Figure 2 shows that the atmospheric properties of the hot stars in NGC 6791 (diamonds) are well matched by the theoretical models with mass loss if their envelope masses are approximately $`0.07M_{}`$. If there is no mass loss on the HB, their properties can be matched only if their envelope masses are as small as $`0.001M_{}`$ (the filled star symbol at log $`T_{\mathrm{eff}}4.4`$ and log L/L$`{}_{}{}^{}1.1`$). Whether HB stars can be born with such low envelope masses at all has been questioned by many groups: “the minimum envelope mass hypothesis” (see the discussion in §3.2 of Yi et al. 1997b). With mass loss on the HB, one can match the observed data without violating the minimum envelope mass hypothesis.
Fig. 3 Comparison of evolutionary tracks for a constant metallicity Z=0.04 and varying envelope masses. Mass loss has a more profound effect on low envelope mass stars.
It should be noted that the adoption of mass loss on the HB does not change the zero-age HB locus at all in our scenario, because we are only comparing the tracks with and without mass loss when they begin their HB evolution at the same location in the CMD.
The inclusion of mass loss on the HB also contributes to the wide color distribution in the HB of NGC 6791. Figure 3 shows evolutionary tracks for HB stars with varying envelope masses. Each panel is for a given envelope mass with four mass loss rates plotted for each. The composition for this particular set of models is $`Z=0.04`$ and $`\mathrm{\Delta }Y/\mathrm{\Delta }Z=3`$. Mass loss has little effect to the evolution of the HB stars with large envelope masses but significantly affects that of less massive stars. This mass dependency would generate a wider color distribution in the HB. It may even create a color bimodality (or multimodality), as seen in the hot ends of the HBs of several globular clusters, if the mass dependency is not a smooth function of mass. This may explain the lack of intermediate HB stars in NGC 6791. A more detailed modeling will show if this is a plausible explanation to the case of NGC 6791.
The effects of such mass loss in the HB phase have an impact on the study of ultraviolet spectral evolution of old stellar populations as well. Numerical studies by Demarque and Pinsonneault (1988) and by Horch, Demarque, & Pinsonneault (1992) suggested that in the advanced evolution of metal rich HB stars, stars with small envelopes evolve to become UV bright objects, the so-called slow blue phase (SBP) stars, rather than evolving into asymptotic giant branch stars (c.f. the AGB-manqué stars in Greggio & Renzini 1990). Yi, Demarque & Kim (1997a) elaborated on the theory and found a theoretical explanation to the SBP phenomenon and to its metallicity dependence, that is, the transitional mass (the maximum mass of a star that becomes a SBP star) increasing with metallicity. According to this theory, more metal-rich stars become UV bright SBP stars more easily.
With the additional mass loss on the HB that we consider in this paper, stars would become SBP stars even more easily. Moreover, if mass loss on the HB increases with increasing metallicity and increasing effective temperature, as various studies suggest, it would cause metal-rich stellar populations to develop UV bright stars more quickly. Because SBP stars are considered important UV sources in elliptical galaxies (e.g., Tantalo et al. 1996; Yi et al. 1997a; O’Connell 1999) and the UV flux of elliptical galaxies may constrain the ages of galaxies (Yi et al. 1999), this study has an impact on extragalactic astronomy. Landsman et al. (1998) already pointed out that a population like NGC 6791 would exhibit a UV upturn with the magnitude shown in giant elliptical galaxies. If mass loss on the HB indeed contributes to the production of hot stars in the metal-rich environments, our current age estimates of giant elliptical galaxies are likely to be systematically overestimated.
We should point out that there is no direct observational evidence for mass loss as substantial as we have assumed to occur on the HB. In fact, besides Koopmann et al.’s upper limit on the mass loss rate on the HB ($`10^9M_{}yr^1`$), the diffusion study of Michaud et al. (1985) set a more strict upper limit of $`10^{14}M_{}yr^1`$. They suggest that if mass loss rate is larger than this it would be difficult to reproduce the silicon underabundance observed in sdB stars. If this is true, the mass loss rates investigated in our study may be too large to justify, unless mass loss rate increases substantially in the metal-rich regime. Once again, we do not advocate any value or mechanism for the mass loss on the HB. Instead, we wanted to see what level of mass loss rate on the HB would be needed if it is the only process that is responsible for the production of hot HB stars. Based on a cursory inspection of our new HB tracks with mass loss, adoption of such low mass rates ($`<10^{14}M_{}yr^1`$) would not make any appreciable difference in the HB morphology.
Some of the blue HB stars have been observed to be spectroscopic binaries (Green et al. 1997). This brings up the possibility that the blue HB stars are formed through binary evolution. Some of the formation scenarios such as helium white dwarf mergers (Bailyn & Iben 1989, Iben 1990) can perhaps be ruled out based on the fact that the luminosities are in a small range (see Figure 2). This small range of the luminosities of the hot stars is also in conflict with the binary scenario (Landsman et al. 1998). Mengel, Norris, & Gross (1976) demonstrated that blue HB stars can form through binary evolution with mass transfer via Roche Lobe overflow. However, the conditions under which this occurs are probably too infrequently met, although further study, applicable to metal rich stellar populations, is needed. Furthermore, it is unclear whether all blue HB stars in NGC 6791 are in mass-exchanging binary systems. Therefore, a mechanism to form blue HB stars through single star evolution is still needed. Mass loss on the HB may be such a mechanism.
This work was supported in part by the Creative Research Initiative Program of the Korean Ministry of Science & Technology grant (S.Y.), and grant NAG5-8406 from NASA (P.D.).
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# Density fluctuations and entropy
## I Introduction
Entropy is one of the very important and challenging thermodynamic quantities in statistical mechanics; this is so because it depends on all the n-particle distribution functions. The problem is not the lack of exact expressions but to derive equations that are both accurate and manageable from a theoretical and numerical point of view. Amongst the exact expressions we cite the classical work of Nettleton and Green , and, more recently, , while approximate expressions can be found in . In we have shown how an infinite subset of terms (dependent exclusively on the one- and two- body distribution functions) can be analytically summed giving rise to the so-called ring approximation and, through a minimization of the free energy functional, the well known HNC approximation is obtained as an optimized superposition approximation. Later on, Bush et al derived sets of integral equations by analyzing several levels of approximation to the grand potential function and, in a recent and very interesting article, Puoskari extends the ring approximation to three particle functions showing how the HNC2 equations either of the Wertheim or the Baxter variety can be obtained. Baranyai and Evans showed that, even though the derivations are done in the grand canonical ensemble, the entropy equations are, in fact, ensemble invariants if local expressions are used for the entropy. In this way, the comparison with canonical ensemble numerical simulations is justified. They also analyze the convergence range of some needed integrals showing how this range increases at high densities. Wallace also worked with the same type of expression and, by analyzing the behaviour of density fluctuations, proposed an incompressible approximation for the entropy of dense fluids (which, as Wallace himself states, is decidedly wrong in the low density regime). He then concluded that the dilute gas and dense liquid regimes occupy different regions of the phase space. These conclusions have been criticized . Laird and Haymet have extended the ring approximation to mixtures and applied it to electrolytes. They found that the correct Debye-Huckel expression for the entropy in the low concentration limit is obtained when the ring approximation is included. They have also discussed the incompressible approximation in dense fluids, proposed an expression that differs from the one of Wallace, and applied it with good numerical results.
Summarizing, these works show that, at low densities the inclusion of the ring approximation gives a very accurate entropy equation and the incompressible approximation is hopelessly wrong, while, at high densities, the ring term overestimates the entropy and the incompressible approximation is reasonably accurate. It has been suggested that more than two bodies correlations must be incorporated in order to have an accurate expansion and, in this respect, Puoskari’s work is quite promising. It is the purpose of this article to elaborate on the compressibility related contribution to entropy and improve the entropy expansion when truncated to the pair distribution level. In Section II we discuss the conceptual structure of the entropy when written as a functional of the n-bodies distribution function. In Section III it is shown by functional differentiation that the compressibility related contribution only depends on the one body distribution function as well as on thermodynamic parameters and, in Section IV, by summing three subsets of terms, we shed some light on the nature of the incompressible approximation and show how the correct low density limit can also be obtained. In Section V we present our conclusions and propose a new entropy equation.
## II ENTROPY STRUCTURE
We can think on two different criteria for the analysis of the entropy. The first one can be called the functional criterium; is the one we take when interested in the entropy dependence on distribution functions, e.g., when a variational principle is formulated. The second point of view is the numerical one, i.e. the main goal is numerical accuracy. It is clear that it is not necessarily true that the same expression will fulfill satisfactorily both goals. In particular, the incompressible approximation (see eq. (12 )) is numerically correct at high densities, but has the wrong behaviour at low densities and also lacks a sound theoretical foundation.
When the entropy dependence on all the n-bodies distribution functions is explicitly written, we obtain an approximate expression with the following distinct structure contributions :
i) The ideal gas contribution
$$\frac{S^{(id)}}{k}=N\left[\frac{5}{2}\mathrm{ln}(\rho \lambda ^3)\right]$$
(1)
where $`\rho `$ is the number density and $`\lambda `$ the thermal wavelength.
ii) The ever present logarithmic contribution
$$\frac{S^{(\mathrm{ln})}}{k}=\underset{p1}{}d\{𝐩\}n_p(\{𝐩\})\omega _p(\{𝐩\})$$
(2)
where $`n_p(\{𝐩\})`$ is the p-particles distribution function and $`\omega _p(\{𝐩\})`$ the irreducible p-bodies contribution to the potential of average force. More specifically, the link with the more usual notation is
$$g_p(\{𝐩\})=\frac{n_p(\{𝐩\})}{_in_1(𝐢)}$$
(3)
$$e^{\omega _p(\{𝐩\})}=\frac{n_p_{\{𝐩\mathrm{𝟐}\}\{𝐩\}}n_{p2}_{\{𝐩\mathrm{𝟒}\}\{𝐩\}}n_{p4}\mathrm{}}{_{\{𝐩\mathrm{𝟏}\}\{𝐩\}}n_{p1}_{\{𝐩\mathrm{𝟑}\}\{𝐩\}}n_{p3}\mathrm{}}$$
(4)
As usual, we have that $`g_2=1+h_2`$ and, through the use of the generalized superposition approximation (GSA) we can write that, e.g.
$`g_3(\{\mathrm{𝟑}\})`$ $`=`$ $`\left[1+\mathrm{\Delta }_3(\{\mathrm{𝟑}\})\right]{\displaystyle \underset{\{\mathrm{𝟐}\}\{\mathrm{𝟑}\}}{}}g_2(\{\mathrm{𝟐}\})`$ (5)
$`g_4(\{\mathrm{𝟒}\})`$ $`=`$ $`\left[1+\mathrm{\Delta }_4(\{\mathrm{𝟒}\})\right]{\displaystyle \frac{_{\{\mathrm{𝟑}\}\{\mathrm{𝟒}\}}g_3(\{\mathrm{𝟑}\})}{_{\{\mathrm{𝟐}\}\{\mathrm{𝟒}\}}g_2(\{\mathrm{𝟐}\})}}=`$ (7)
$`\left[1+\mathrm{\Delta }_4(\{\mathrm{𝟒}\})\right]{\displaystyle \underset{\{\mathrm{𝟑}\}\{\mathrm{𝟒}\}}{}}\left[1+\mathrm{\Delta }_3(\{\mathrm{𝟑}\})\right]{\displaystyle \underset{\{\mathrm{𝟐}\}\{\mathrm{𝟒}\}}{}}g_2(\{\mathrm{𝟐}\})`$
which introduces the family of $`\mathrm{\Delta }_p`$ functions that, when different from zero, correct for the difference with the GSA. They can also be written as
$$\mathrm{\Delta }_p(\{𝐩\})=e^{\omega _p(\{𝐩\})}1$$
(8)
This $`S^{(\mathrm{ln})}`$ contribution is the one that, when functionally differentiated, gives rise to the $`\mathrm{ln}g`$ contribution in the integral equations
iii) The ring term, which in its simplest, two-body version, is
$$\frac{S^{(r)}}{k}=\frac{1}{2}\underset{p3}{}\frac{(1)^{p1}}{p}d\{𝐩\}\underset{i=1}{\overset{p}{}}n_1(𝐢)h_2(\mathrm{𝟏𝟐})h_2(\mathrm{𝟐𝟑})\mathrm{}h_2(\mathrm{𝐩𝟏})$$
(9)
and it can be summed in homogeneous systems . This term is the responsible of the contribution $`h_2c_2`$ in the integral equations. The three bodies version is derived in .
iv) The compressibility related contribution
$$\frac{S^{(c)}}{k}=\underset{p2}{}\frac{1}{p!}d\{𝐩\}\underset{i=1}{\overset{p}{}}n_1(𝐢)\mathrm{\Delta }_p(\{𝐩\})\mathrm{\Gamma }_p(\{n_p\})$$
(10)
$$\mathrm{\Gamma }_p(\{n_p\})=\frac{_{\{𝐩\mathrm{𝟏}\}\{𝐩\}}g_{p1}_{\{𝐩\mathrm{𝟑}\}\{𝐩\}}g_{p3}\mathrm{}}{_{\{𝐩\mathrm{𝟐}\}\{𝐩\}}g_{p2}_{\{𝐩\mathrm{𝟒}\}\{𝐩\}}g_{p4}\mathrm{}}$$
(11)
Its first term is essentially the compressibility ($`\mathrm{\Delta }_2h_2)`$ and the sequence of products in eq. (11) stops when reaching either $`g_3`$ or $`g_2`$ . As far as we know there are no previous studies of the whole series given in eq. (10); the compressibility approximation focuses on the first term of this series, which for a one component homogeneous system is
$$\frac{S_2^{(c)}}{k}=N\frac{\rho }{2}𝑑𝐫h_2(r)=\frac{N}{2}\left(1+\alpha _2\right)$$
(12)
As in the dense liquid limit is $`\alpha _21`$, the incompressible approximation considers $`\alpha _2=0`$ in the whole $`\rho T`$ space.
## III FUNCTIONAL DEPENDENCE
Here we prove that, in the thermodynamic limit, all the functional derivatives of the compressibility contribution with respect to the distribution functions can be summarized in the equation
$$\frac{\delta S^{(c)}/k}{\delta n_p(\{𝐩\})}=\delta _{1p}+(e^{<N>})$$
(13)
Therefore,
$$\frac{S^{(c)}}{k}=d\{\mathrm{𝟏}\}n_1(\{\mathrm{𝟏}\})+C(\rho ,T)+(e^{<N>})$$
(14)
C is an integration constant as far as the functional integration refers but, in fact, it depends on $`\rho ,T`$.
The derivation is straigthforward. The origin of the compressibility term $`S^{(c)}`$ is quite clear and eqs. (38-41) of Ref. are the equations to look at. Equation (38) is our eq. (12) and in eqs. (39-41) we see that each one of them has, amongst other terms, the integral $`n_1g_pd\{𝐩\}`$. When the GSA for $`g_p`$ is used (eq. (7)), the integral decomposes into a sum of two integrals $`n_1\mathrm{\Gamma }_p\left[1+\mathrm{\Delta }_p\right]d\{𝐩\}`$. The term without $`\mathrm{\Delta }_p`$ contributes to the ring approximation plus neglected terms (such as those shown in Ref. ) and the term with $`\mathrm{\Delta }_p`$ contributes to $`S^{(c)}`$.
As $`\mathrm{\Gamma }_p`$ can also be written as
$`\mathrm{\Gamma }_p(\{n_p\})={\displaystyle \frac{_{\{𝐩\mathrm{𝟏}\}\{𝐩\}}n_{p1}_{\{𝐩\mathrm{𝟑}\}\{𝐩\}}n_{p3}\mathrm{}}{_{\{𝐩\mathrm{𝟐}\}\{𝐩\}}n_{p2}_{\{𝐩\mathrm{𝟒}\}\{𝐩\}}n_{p4}\mathrm{}}}`$
using eqs. (4) and (8) we conclude that
$$\frac{S^{(c)}}{k}=\underset{p2}{}\frac{1}{p!}d\{𝐩\}\left[n_p\frac{_{\{𝐩\mathrm{𝟏}\}\{𝐩\}}n_{p1}_{\{𝐩\mathrm{𝟑}\}\{𝐩\}}n_{p3}\mathrm{}}{_{\{𝐩\mathrm{𝟐}\}\{𝐩\}}n_{p2}_{\{𝐩\mathrm{𝟒}\}\{𝐩\}}n_{p4}\mathrm{}}\right]$$
(15)
It is somewhat clear that each one of these integrals is related to p-bodies density fluctuations but a clearer explanation is to be found in the next section. This explains the origin of naming this contribution as compressibility related. Written in this way it is straightforward to show that the functional derivatives are
$$\frac{\delta S^{(c)}/k}{\delta n_1(𝐱)}=\underset{p1}{}(1)^p\frac{N^p}{p!}[1+\left(\frac{1}{N}\right)]=1+\left(e^N\right)$$
(16)
$$\frac{\delta S^{(c)}/k}{\delta n_s(\{𝐱_s\})}=\frac{1}{s!}\underset{p0}{}(1)^p\frac{N^p}{p!}[1+\left(\frac{1}{N}\right)]=\left(e^N\right)$$
(17)
and we arrive to eq. (13).
This result shows that, in the thermodynamic limit, the compressibility term does not contribute to any set of equations we may derive by functional diferentiation of a functional that includes the entropy; it only contributes to the constraint of fixed density. Therefore, if we are after a set of equations which are the consequence of a variational principle, we can rightly put the compressibility term aside from all the others, while if we are after a numerical approximation to the entropy we can assume, on physical grounds, absolute convergence of the whole entropy series and feel free to mix the compressibility term with all the others if numerical convergence is improved.
## IV SERIES SUMMATION
For the sake of this numerical goal we will cut the GSA to third order; in this way $`g_p`$ can be written in two equivalent forms
$$g_p(\{𝐩\})=\{\begin{array}{c}1+_{\{\mathrm{𝟐}\}\{𝐩\}}h_2(\{\mathrm{𝟐}\})+_{\{\mathrm{𝟑}\}\{𝐩\}}h_3(\{\mathrm{𝟑}\})+\mathrm{}h_p(\{𝐩\})\\ _{\{\mathrm{𝟑}\}\{𝐩\}}\left[1+\mathrm{\Delta }_3(\{\mathrm{𝟑}\})\right]_{\{\mathrm{𝟐}\}\{𝐩\}}\left[1+h_2(\{\mathrm{𝟐}\})\right]\end{array}$$
(18)
In eq. (15) for the compressibility contribution we will sum to infinite order three subsets of terms. These subsets are clearly identified in the $`p=3`$ summand of eq. (15), i.e.
$`{\displaystyle \frac{S_3^{(c)}}{k}}`$ $`=`$ $`{\displaystyle \frac{1}{3!}}{\displaystyle }d\{\mathrm{𝟑}\}{\displaystyle \underset{i=1}{\overset{3}{}}}n_1(𝐢)[g_3(\{\mathrm{𝟑}\}){\displaystyle \underset{\{\mathrm{𝟐}\}\{\mathrm{𝟑}\}}{}}[1+h_2(\{\mathrm{𝟐}\})]=`$ (20)
$`{\displaystyle \frac{1}{3!}}{\displaystyle d\{\mathrm{𝟑}\}\underset{i=1}{\overset{3}{}}n_1(𝐢)\left[h_3(\{\mathrm{𝟑}\})\underset{i=1}{\overset{3}{}}\underset{ki}{}h_2(\mathrm{𝐢𝐤})h_2(\mathrm{𝟏𝟐})h_2(\mathrm{𝟏𝟑})h_2(\mathrm{𝟐𝟑})\right]}`$
A) The first subset includes the contribution of the integrals $`h_pd\{𝐩\},p2`$. The series is
$$\frac{S_a^{(c)}}{k}=\underset{p=2}{\overset{\mathrm{}}{}}\frac{1}{p!}d\{𝐩\}\underset{i=1}{\overset{p}{}}n_1(𝐢)h_p(\{𝐩\})=\underset{p=2}{\overset{\mathrm{}}{}}\frac{C_p}{p!}$$
(21)
The moment-cumulant relation is
$$C_M(\{𝐌\})=h_M(\{𝐌\})\underset{i=1}{\overset{M}{}}n_1(𝐢)=\underset{k=1}{\overset{M}{}}\left\{k\{𝐦_i\}_𝐌\right\}(1)^{k1}(k1)!\underset{i=1}{\overset{k}{}}n_{m_i}(\{𝐦_i\})$$
(22)
Here, the partition of the coordinate set $`\{𝐌\}`$ in $`k`$ disjoint subsets $`\{𝐦_i\}_𝐌,1ik`$ is symbolized by $`\left\{k\{𝐦_i\}_𝐌\right\}`$ and therefore $`_{k=1}^M\left\{k\{𝐦_i\}_𝐌\right\}`$ indicates the sum over all the partitions in $`k`$ subsets and for each $`k`$ is $`1ik`$. In this way $`C_p`$ is related to the integrals $`n_{p_i}=𝑑p_in_{p_i}`$. On the other hand, $`n_{p_i}`$ can be expanded in terms of $`N^k`$
$$n_p=N(N1)\mathrm{}(Np+1)=\underset{k=1}{\overset{p}{}}s(p,k)N^k$$
(23)
for $`p1`$, where $`s(p,k)`$ are the Stirling numbers of first kind. One of its definitions is that $`(1)^{pk}s(p,k)`$ is the number of permutations of p elements which contain exactly k cycles. They satisfy the recurrence relation
$`s(p+1,k)=s(p,k1)ps(p,k),1kp`$
with starting values
$`s(p,0)=s(0,k)=\delta _{0n}`$
We also define the r-bodies density fluctuations $`\alpha _r`$ by
$$\alpha _r=\{\begin{array}{cc}1& r=1\\ \frac{\left(N<N>\right)^r}{<N>}& r>1\end{array}$$
(24)
The first few $`C_p`$ are then expressed in terms of the $`\alpha _r`$ and Stirling numbers as
$`C_2`$ $`=`$ $`N\left(1+\alpha _2\right)=N{\displaystyle \underset{i=1}{\overset{2}{}}}s(2,i)\alpha _i`$ (25)
$`C_3`$ $`=`$ $`N\left(23\alpha _2+\alpha _3\right)=N{\displaystyle \underset{i=1}{\overset{3}{}}}s(3,i)\alpha _i`$ (26)
$`C_4`$ $`=`$ $`N\left(6+11\alpha _26\alpha _3+\alpha _4\right)+3N^2\alpha _2^2=N{\displaystyle \underset{i=1}{\overset{4}{}}}s(4,i)\alpha _i+3N^2\alpha _2^2`$ (27)
The $`3N^2\alpha _2^2`$ term and similar ones from higher order $`C_p`$ will be included in the next partial sum. In order to sum to infinite order the contribution of each r-bodies density fluctuations we need the result
$`{\displaystyle \underset{t=k}{\overset{\mathrm{}}{}}}{\displaystyle \frac{s(t,k)}{k!}}x^k={\displaystyle \frac{\left[\mathrm{ln}(1+x)\right]^t}{t!}}`$
Therefore, the $`\alpha _r`$ contribution is
$$\mathrm{\Gamma }_{\alpha _r}=\{\begin{array}{cc}\mathrm{ln}21\hfill & r=1\hfill \\ \left(\mathrm{ln}2\right)^r/r!\hfill & r2\hfill \end{array}$$
(28)
and our first partial sum, which includes contributions to infinite order of all the $`\alpha _r`$, is
$$\frac{S_a^{(c)}}{k}=\underset{p=2}{\overset{\mathrm{}}{}}\frac{C_p}{p!}=N\left\{\mathrm{ln}21+\underset{p=2}{\overset{\mathrm{}}{}}\frac{\left(\mathrm{ln}2\right)^p}{p!}\alpha _p\right\}$$
(29)
Let us remark on some characteristics of this result: i) $`\mathrm{ln}21`$ is the contribution in the absence of density fluctuations and it is also its high density limit ($`\alpha _p1`$), ii) as $`\alpha _p1`$ when $`\rho 0`$, eq. (29) vanishes in the low density limit, iii) the series is rapidly convergent. Therefore, this sum goes in the right direction to improve on the incompressible approximation, both in its numerical results and in its theoretical foundation. This analysis makes also clear why this contribution is referred to as compressibility related. Lastly, the $`\alpha _p`$ are easily expressed as integrals of the correlation functions; the first ones are
$`\alpha _2`$ $`=`$ $`1+{\displaystyle \frac{1}{N}}{\displaystyle d\{\mathrm{𝟐}\}n_1(𝐢)h_2(\{\mathrm{𝟐}\})}`$
$`\alpha _3`$ $`=`$ $`2+3\alpha _2+{\displaystyle \frac{1}{N}}{\displaystyle d\{\mathrm{𝟑}\}n_1(𝐢)h_3(\{\mathrm{𝟑}\})}`$
B) When in the p-th term of eq. (15) the expansions given in eqs. (18) are inserted, each $`h_k,k<p`$ expanded in terms of $`h_2,\mathrm{\Delta }_3`$ and terms like the $`3N^2\alpha _2^2`$ which were left aside in the first partial series included, then all the unconnected (in the graph theory sense) terms cancel out and the first two sets of connected diagrams are those depicted in eq. (20). We first evaluate the sum of ”star ”products of $`h_2`$ bonds
$$\frac{S_b^{(c)}}{k}=\underset{p=3}{\overset{\mathrm{}}{}}\mathrm{\Psi }_p=\underset{p=3}{\overset{\mathrm{}}{}}\frac{1}{p!}\underset{i=1}{\overset{p}{}}d\{𝐢\}n_1(𝐢)\underset{ki}{}d\{𝐤\}n_1(𝐤)h_2(\mathrm{𝐢𝐤})$$
(30)
Each summand is easily evaluated as
$`\mathrm{\Psi }_p={\displaystyle \frac{N}{(p1)!}}\left(1+\alpha _2\right)^{p1}`$
and the second partial sum is
$$\frac{S_b^{(c)}}{k}=N\left[\mathrm{exp}(\alpha _21)\alpha _2\right]$$
(31)
Its low and high density limits are $`0`$ and $`e^1`$ respectively.
C) This series is a sum of rings very similar to eq. (9), its first term is given in eq. (20) and, as the symmetry number of p-rings is $`2p`$, it can be written as
$$\frac{S_c^{(c)}}{k}=\frac{1}{2}\underset{p3}{}\frac{1}{p}d\{𝐩\}\underset{i=1}{\overset{p}{}}n_1(𝐢)h_2(\mathrm{𝟏𝟐})h_2(\mathrm{𝟐𝟑})\mathrm{}h_2(\mathrm{𝐩𝟏})$$
(32)
which can be summed for homogeneous systems in the same way as the original ring approximation was
$$\frac{S_c^{(c)}}{k}=\frac{N}{2\rho }𝑑𝐤\left\{\mathrm{ln}\left(1\rho \stackrel{~}{h}_2(k)\right)\rho \stackrel{~}{h}_2(k)\frac{\left(\rho \stackrel{~}{h}_2(k)\right)^2}{2}\right\}$$
(33)
where $`\stackrel{~}{h}_2(k)=d𝐫h(r)\mathrm{exp}(2\pi i𝐤.𝐫)`$ is the Fourier transform of $`h_2(r)`$ and the integration is over the $`𝐤`$ space. This contribution can be added to the original ring approximation giving a renormalized ring approximation $`\stackrel{~}{S}^{(r)}`$, which is
$$\frac{\stackrel{~}{S}^{(r)}}{k}=\frac{1}{2}\underset{p2}{}\frac{1}{p}d\{\mathrm{𝟐}𝐩\}\underset{i=1}{\overset{2p}{}}n_1(𝐢)h_2(\mathrm{𝟏𝟐})h_2(\mathrm{𝟐𝟑})\mathrm{}h_2(\mathrm{𝟐}𝐩\mathrm{𝟏})$$
(34)
a sum over all even order rings and, for homogeneous systems, the result is
$$\frac{\stackrel{~}{S}^{(r)}}{k}=\frac{N}{2\rho }𝑑𝐤\left\{\mathrm{ln}\left(1(\rho \stackrel{~}{h}_2(k))^2\right)+\left(\rho \stackrel{~}{h}_2(k)\right)^2\right\}$$
(35)
## V CONCLUSIONS
Collecting together the different results obtained, i.e. eqs. (29), (31) and (35) with eqs. (1) and (2) we arrive to a new entropy equation which includes a partial summation of the compressibility related contribution, i.e.
$`{\displaystyle \frac{S}{k}}`$ $`=`$ $`N\left[{\displaystyle \frac{5}{2}}\mathrm{ln}(\rho \lambda ^3)\right]{\displaystyle \underset{p1}{}}{\displaystyle d\{𝐩\}n_p(\{𝐩\})\omega _p(\{𝐩\})}+`$ (38)
$`+{\displaystyle \frac{N}{2\rho }}{\displaystyle 𝑑𝐤\left\{\mathrm{ln}\left(1(\rho \stackrel{~}{h}_2(k))^2\right)+(\rho \stackrel{~}{h}_2(k))^2\right\}}+`$
$`+N\left\{\mathrm{ln}21+{\displaystyle \underset{p=2}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\left(\mathrm{ln}2\right)^p}{p!}}\alpha _p\right\}N\left[\mathrm{exp}(\alpha _21)\alpha _2\right]`$
If, for just a moment, we neglect the renormalization in the ring approximation, the contribution due to the r-bodies density fluctuations gives not only a theoretical understanding on the nature of the incompressible approximation, but also a description that is essentially correct in the low and high density limits ($`0`$ and $`\mathrm{ln}21e^1`$ respectively). As the ring approximation grows quite steeply when the density increases , its renormalization should have the right asymptotic behaviour. Let us also mention that these results extends trivially to mixtures (see, e.g. ) and, in this case, it is more convenient to work with the entropy per unite volume.
We have also shown that, in the thermodynamic limit, the compressibility term only depends (as a functional) on the one body distribution function. Therefore, this functional dependence is such that it only enters in the constant density constraint and, in this way, the conceptual structure of the equation for the entropy is significantly simplified. As our results apply to the full entropy functional, they are valuable to any functional minimization as, e.g., those in . It can also be mentioned that this theorem does not conflict with Laird-Haymet . They obtained the correct Debye-Huckel low density expression for the entropy by including the $`S^{(id)},S^{(\mathrm{ln})},`$ $`S^{(r)}`$ plus the compressibility related contribution of eq. 12). As this term and, in fact, all the sums we did, vanish when $`\rho 0`$ (including the one that renormalizes the ring approximation), there is no contradiction between ours and theirs results. Lastly, as this result does not depend on the potential, it is also valid for the associative Wertheim-Ornstein-Zernike equation as well as systems with directional forces .
This new functional provides a robust and systematic way to develop fully analytical theories of liquids , which will be examined in future work.
## VI ACKNOWLEDGMENTS
We acknowledge support from the National Science Foundation through grants CHE-95-13558, Epscor OSR-94-52893, by the DOE-EPSCoR grant DE-FCO2-91ER75674 and CONICET grant PIP 0859/98.
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# A Proof That Measured Data and Equations of Quantum Mechanics Can Be Linked Only by Guesswork
## 1. Introduction
This paper stems from earlier work and a proof presented here showing that inquiry in quantum physics continually presents a scientist with choices of equations and of instruments, unresolvable by calculation and measurement. Something else is demanded of the scientist, which may as well be called a guess.<sup>1</sup><sup>1</sup>1Other words are hypothesis, Ansatz, assumption, axiom, postulate, and sometimes principle. Challenged by the proof to look at its implications, we noticed that people in investigative endeavors inherit and frame choices open to guesswork, some of which show up clearly in the computers used in the endeavors.
Section 2 introduces the Classical Process-control Computer (CPC) with its special capacity to manipulate abstractions expressed as equations without contaminating them with its own physics.<sup>2</sup><sup>2</sup>2This capacity stems from regenerative amplifiers and clock-gated memory registers, two inventions used to make all computer hardware insensitive to manufacturing variations, so that, like the placing of a chess piece not quite in the center of a square, deviations in performance, within limits, do not matter. Its independence from its own physics distinguishes a classical computer from a quantum computer. A scientist can use a CPC not only to calculate with equations, but also to mediate the command of laboratory instruments via digital/analog (D/A) converters and to record experimental results returned from the instruments via analog/digital (A/D) converters. By noticing that both equations and instruments make contact in a CPC, we rewrite equations of quantum mechanics to explicitly include functions of CPC-commands to the instruments. This sets up the proof that the scientist’s choice in linking equations to instruments is unresolvable without guesswork to narrow the set of models. A lattice of sets of models is defined, two widely used guesses that narrow the set of models are noted, and the concept of statistical distance between probability distributions is applied to quantum-mechanical models.
Section 3 provides language for displaying and analyzing guess-demanding choices visible in CPC files. To this end Turing machines are introduced and adapted to formalize the definition of a CPC. This allows fragments of colored Petri-nets, opened to exogenous influences, to portray the programming and running of programs in a network of CPC’s operated by collaborating scientists. Many of these programs incorporate guesses. This general picture of process-control computation shows programs and other guesses as colors on tokens that a scientist enters on a Petri net that acts as a game board. Mechanisms for one scientist to judge programs (and hence guesses) made by another are sketched, leading to the first of many needs for concurrently operating CPC’s.
Section 4 describes some examples in which guessing, quantum mechanics and CPC structures must interact in the building of a quantum computer as a laboratory instrument specified by equations of quantum mechanics. We show the need for guesses to link equations and instruments brings with it a need to test the quantum computing instruments and to calibrate them, guided by test results and guesses. Quantum mechanics imposes a peculiar structure on this testing, related to the statistical distance between models. For the measure of precision conventionally used in quantum computing, the sample size needed for testing a quantum gate is shown to increase as the inverse square of the tolerated imprecision. While many questions are left open to future work, the example demonstrates a frame for analysis and experiment broader than any quantum model alone, a frame that includes the testing of the mathematical models by results of the use of instruments, and so distinguishes what the model says from what the instruments do, allowing provision for guesswork as an ingredient in advancing both models and instruments.
## 2. Quantum-mechanical models and their links to instruments
Proving the necessity of guesses demands language to describe the linking of numbers in mathematical models to numbers pertaining to laboratory instruments, starting with mathematical language to describe a scientist’s choosing one arrangement of laboratory instruments rather than another. We shall describe a situation in which a scientist chooses instruments by using a CPC keyboard to type strings of characters, much as Gödel, in mathematical logic, described equations as strings of characters. The scientist at the CPC keyboard writes and executes programs to command the operation of laboratory instruments and record their results. These programs make use of quantum-mechanical equations, which the scientist also writes into the CPC.
Quantum mechanics as a mathematical language expresses different measuring instruments by different operators, and thus has built into it a recognition that phenomena to be described cannot be independent of the instruments used to study them . Still, this dependence is emphasized more some times than others. Some modeling merely assumes that instruments can be found, without saying how, to implement various combinations of state vectors and operators. Such models appear in theories of quantum computing to relate the multiplication of unitary operators to the solving of problems of interest. To see the need for another kind of model, suppose a scientist has computer-controlled instruments (such as lasers) with the potential to implement a quantum computer, and faces the question of what commands the CPC should transmit to the instruments and when it should transmit them in order implement one or another quantum gate. Determining the commands and their timing to implement a quantum gate expressed as a unitary matrix $`U_j`$ takes a model that expresses the gates as unitary transformations in terms of commands that a process-control computer can transmit to the instruments. Curiously, models of this kind have not been much stressed in physics, and it is a merit of efforts to build quantum computers to make the importance of such models apparent.
### 2.1. Models and instruments make contact in a CPC
Part of a scientist’s control of instruments can work through the use of a process-control computer that transmits commands to the (computer-controlled) instruments and records results produced by them. We confine our analysis to this part, excluding from consideration here (but by no means denigrating) hand work beyond the reach of a process-control computer. We shall portray cases in which a scientist chooses arrangements of instruments, chooses models, and puts the two in contact, linking models to instruments, during a CPC session starting after the instruments have been set up and put under control of a CPC and ending before the scientist has to tinker with the instruments in ways unreachable by the CPC. Within the CPC, laboratory instruments and mathematical models make contact when:
1. a model resident in a CPC file is used to derive commands for the CPC to transmit to the instruments;
2. instrumental results collected by a CPC are used to narrow down a set of models. (We shall later see feedback as an example of this.)
Such contact does not spring from nothing, but is brought about by design and depends on choices made by a scientist, including choices of what set of models to start with, what model to choose for use by a CPC in generating commands, and what experiments to run. To picture the design and operation of contact between models and instruments, imagine eavesdropping on CPC’s used in various investigations. Commands sent to the instruments by the CPC and the results received from them, both numerical, are amenable to analysis, as is the scientist’s writing of equations, programs, calls for program execution, etc.; we also eavesdrop on displays produced by the CPC for the scientist.
Although the CPC puts instruments in contact with equations involving quantum superposition, the CPC itself is a classical machine, free of quantum superposition, for it needs no quantum behavior within itself, neither to manipulate equations of quantum mechanics nor to manage laboratory instruments. For example, the writing of an expression $`|0+|1`$ for a superposition of quantum states makes use of written characters that themselves exhibit no superposition. And any command to instruments is likewise a character string, including a command to rotate a polarizing filter by 45 degrees to implement the superposed state $`|0+|1`$. Similarly, results of the use of instruments interpreted as demonstrating superposition arrive as bit strings, themselves devoid of superposition.
The CPC is situated between a scientist to its left and laboratory instruments to its right, as shown in figure 1. Working at the CPC, a scientist is limited in action at any moment to the resolution of the choice presented by the CPC at that moment, a choice defined by the files stored in its memory and the state of its processor, and exemplified by a menu displayed by the CPC. Our analysis of the CPC cannot reach beyond its buffers: neither to its left into the scientist, nor to its right where, invisible to eavesdropping, reside digital-to-analog (D/A) and analog-to-digital (A/D) converters and beyond them the laboratory instruments.
### 2.2. Quantum-mechanical models that recognize com-mands sent to instruments
For equations of quantum mechanics to model effects of a scientist’s choices in arranging instruments, these choices must show up in the equations. To see how this can work, recall that quantum mechanics parses the functioning of instruments into state preparation, transformation, and measurement, three coordinated activities that generate outcomes, supposed visible in experimental results by some means unspecified. The three activities are described, respectively, by a state (as a unit vector representing a ray in a Hilbert space), a unitary operator, and a hermitian operator. The only way to make the scientist’s choices in arranging instruments show up in quantum-mechanical equations is to make the state vector $`|v`$, the transformation operator $`U`$, or the measurement operator $`M`$, or some combination of them, depend on how these choices are resolved.
A simple and yet, so far as we know, original way to analyze a scientist’s choice of arrangements of instruments is to suppose that during a CPC-mediated session the instruments are controlled by CPC-transmitted commands from a set $`B`$ of possible commands, where $`B`$ and $``$ is the set of all finite binary strings. We formulate a core set of quantum mechanical models that express the probability of an outcome of instruments in response to a command $`bB`$ sent to the instruments by the CPC, as follows. Let $`𝒱_B`$, $`𝒰_B`$, and $`_B`$ be the sets of all functions $`|v`$, $`U`$, and $`M`$, respectively, with
$`|v:B`$ $``$ $`,`$
$`U:B`$ $``$ $`\{\text{unitary operators on }\},\text{and}`$
$`M:B`$ $``$ $`\{\text{hermitian operators on }\}.`$
The core models exhibit discrete spectra for all $`M`$:
###### Property 1.
(2.1)
$$(bB)M(b)=\underset{j}{}m_j(b)M_j(b),$$
where $`m_j:B`$ (with $``$ denoting the real numbers) is the $`j`$-th eigenvalue of $`M`$, and $`M_j`$ is the projection onto the $`j`$-th eigenspace (so $`M_jM_k=\delta _{j,k}M_j`$).
Let $`\mathrm{Pr}(j|b)`$ denote the probability of obtaining the $`j`$-th outcome, given transmission by the CPC of a command $`b`$. Although not commonly seen in texts, this probability of an outcome given a command is the hinge pin for focusing on quantum mechanical modeling of uses of instruments. Quantum mechanics constrains the models to satisfy:
###### Property 2.
(2.2)
$$\mathrm{Pr}(j|b)=v(b)|U^{}(b)M_j(b)U(b)|v(b),$$
where the $``$ denotes the hermitian adjoint.
(Within this modeling scheme, the Schrödinger equation relates a model at a later time to a model at an earlier time by a certain transformation operator $`U`$, dependent on the situation.)
Any choice from the sets $``$, $`𝒱_B`$, $`𝒰_B`$, and $`_B`$ produces some quantum-mechanical model $`(|v,U,M)_B`$. Two models $`(|v,U,M)_B`$ and $`(|v^{},U^{},M^{})_B`$ generate the same probabilities $`\mathrm{Pr}(j|b)`$ if they are unitarily equivalent, meaning there exists a $`Q:B\{\text{unitary operators on }`$ $`\}`$ such that $`(bB)|v^{}(b)=Q(b)|v(b)`$, $`U^{}(b)=Q(b)U(b)Q^{}(b)`$ and $`M^{}(b)=Q(b)M(b)Q^{}(b)`$. For this reason, any model $`(|v,U,M)_B`$ can be reduced to $`(|v^{},\mathrm{𝟏},M)_B`$, where $`|v^{}=U|v`$ and $`M^{}=M`$ or, alternatively to $`(|v,\mathrm{𝟏},M^{})_B`$ where $`M^{}=U^{}MU`$.
More models are available in more general formulations. When we show that guesswork is necessary even to resolve choices among models of the core set, it follow that guesswork is necessary also to resolve the choices of among a larger set of models involving positive-operator-valued measures, superoperators, etc.
### 2.3. From results to quantum-mechanical outcomes
Before stating and proving the proposition that calculations and measurements cannot by themselves link models to outcomes obtained from instruments, we call to the reader’s attention that outcomes themselves, in the sense of quantum mechanics, are produced by instruments only with the help of interpretive guesswork.
###### Claim 1.
To speak of actual instruments in the language of quantum mechanics one needs to associate results of the use of the instruments, recorded in a CPC, with outcomes in the sense of quantum mechanics or with averages of outcomes.
Experimental results of the use of instruments become quantum-mechanical outcomes only by a scientist’s act of interpreting the results as outcomes. The interpretation involves judgment and guesswork, not only to sidestep imperfections in the instruments, but as a matter of principle, even for the limiting case of instruments supposed free of imperfections. For example, light detectors used in experiments described by models of quantum optics generate experimental results; typically, each of $`L`$ detectors reports to the CPC at each of a succession of $`K`$ time intervals a detection result, consisting of 0 (for no detection) or 1 (for detection), so a record contains $`LK`$ bits. Depending on judgments made about correlations from time interval to time interval and detector to detector, these $`LK`$ bits may constitute $`LK`$ quantum outcomes, or one quantum outcome, or some number in between. The number of outcomes in $`LK`$ bits is determined neither by the experimental results (which in this case are just these bits) nor by general principles of quantum mechanics; yet the parsing of results into outcomes must occur, at least provisionally, before any comparison between equations and measured outcomes can be made. Henceforth, when we speak of outcomes, we presuppose that this piece of guesswork has been accomplished and a decision made to define the parsing of results into outcomes.
### 2.4. Calculation and measurement by themselves cannot link quantum models to recorded outcomes
Could it be that the general properties 1 and 2 suffice to determine a model (up to unitary equivalence) if only one collects enough measured results interpreted as outcomes? The answer is: “no; unless some special properties restrict the model more tightly than the form established by properties 1 and 2 alone, one can always find many unitarily inequivalent models $`(|v,U,M)_B`$, all of which produce probabilities that match perfectly the relative frequencies of outcomes.”
To prove this we define some things to pose the issue more sharply. Let $`B`$ denote the set of commands used to generate some set of outcomes interpreted from measured results.<sup>3</sup><sup>3</sup>3Practically speaking, $`B`$ must be a finite set, but the proof holds also for $`B`$ denumerably infinite. For any $`bB`$, let $`N(b)`$ be the number of times that an outcome has been entered in the record for a run of the experiment for command $`b`$, and let $`J(b)`$ be the number of distinct outcomes for command $`b`$. For $`j=1,\mathrm{},J(b)`$, let $`\lambda _j(b)`$ be the $`j`$-th distinct outcome obtained for command $`b`$, and let $`n(j,b)`$ be the number of times this $`j`$-th distinct outcome $`\lambda _j`$ is recorded in response to command $`b`$. For all $`j>J(b)`$ let $`\mu _j(b)`$ be arbitrary real numbers, and for all $`j1`$ let $`\varphi (j,b)`$ be arbitrary real numbers.
###### Proposition 2.1.
Given any set of recorded outcomes associated with any set $`B`$ of commands, the set of models satisfying properties 1 and 2 contains many unitarily inequivalent models $`(|v,U,M)_B`$, each of which has a perfect fit with the set of outcomes, in the sense that
(2.3)
$$(b)(1jJ(b))\mathrm{Pr}(j|b)=n(j,b)/N(b).$$
###### Proof.
It is instructive to start with the special case in which for some $`bB`$, there exist two or more distinct values of $`j`$ for which $`n(j)>0`$. For this case let the set $`\{|j\}`$ be an orthonormal basis of a separable Hilbert space. Define a subset $`S`$ of models satisfying properties 1 and 2 as all models of the form $`(|v,U,M)_B`$, where
(2.4) $`|v(b)`$ $`\stackrel{\mathrm{def}}{=}`$ $`{\displaystyle \underset{j=1}{\overset{J(b)}{}}}[n(j,b)/N(b)]^{1/2}\mathrm{exp}(i\varphi (j,b)|j),`$
(2.5) $`U(b)`$ $`\stackrel{\mathrm{def}}{=}`$ $`\mathrm{𝟏},`$
(2.6) $`M(b)`$ $`\stackrel{\mathrm{def}}{=}`$ $`{\displaystyle \underset{j=1}{\overset{J(b)}{}}}\lambda _j(b)|jj|+{\displaystyle \underset{j=J(b)+1}{\overset{\mathrm{}}{}}}\mu _j(b)|jj|,`$
with $`\mu _j`$ and $`\varphi `$ arbitrary real-valued functions. By invoking property 2, one checks that any such model has the claimed perfect fit; yet the set contains many unitarily inequivalent models, which predict conflicting statistics for some possible quantum measurement.<sup>4</sup><sup>4</sup>4This happens e.g. for primed and unprimed models if for any $`b`$, $`n(1,b)0n(2,b)`$ and $`\varphi (1,b)\varphi (2,b)\varphi ^{}(1,b)\varphi ^{}(2,b)`$. This proves the special case.
For the general case, modify the definitions above to
(2.7) $`|v(b)`$ $`\stackrel{\mathrm{def}}{=}`$ $`{\displaystyle \underset{j=1}{\overset{J(b)}{}}}[n(j,b)/N(b)]^{1/2}|w_j,`$
(2.8) $`U(b)`$ $`\stackrel{\mathrm{def}}{=}`$ $`\mathrm{𝟏},`$
(2.9) $`M(b)`$ $`\stackrel{\mathrm{def}}{=}`$ $`{\displaystyle \underset{j=1}{\overset{J(b)}{}}}\lambda _j(b)P_j+{\displaystyle \underset{j=J(b)+1}{\overset{\mathrm{}}{}}}\mu _j(b)P_j,`$
where $`P_jP_k=\delta _{j,k}P_j`$, for all $`j`$ the projection $`P_j`$ has dimension greater than 1, and $`|w_j`$ ranges over all unit vectors of the eigenspace defined by $`P_j|w_j=|w_j`$. In particular, for any $`j`$, dim($`P_j`$) can be as large as one pleases. Then even if there is only one outcome that is ever recorded, there are still many unitarily inequivalent models that perfectly fit the data.∎
Proposition 2.1 implies that the density matrix, often supposed to be determined from measured data , is undetermined without assuming special properties shortly to be discussed; this follows by expressing the density matrix as $`|vv|`$ and noticing that the phases of the off-diagonal elements are undetermined. We leave to the future the demonstration of additional ambiguity in the link between any set of recorded outcomes and models expressed in the mathematical language of quantum mechanics.
### 2.5. Statistically significant differences between models
In practice, a scientist has little interest in a model chosen so that its probabilities exactly fit measured relative frequencies. Rather, the scientist wants a simpler model with some appealing structure that comes reasonably close to fitting. Quantum mechanics encourages this predilection, because on account of statistical variation in the sample mean, functions that perfectly fit outcomes on hand at one time are not apt to fit perfectly outcomes acquired subsequently. We show here that accepting statistics no way takes away from the proof that measurements and equations by themselves cannot link models to instruments.
One needs a criterion for the statistical significance of a difference between two quantum-mechanical models (or between a model and measured relative frequencies). Here we limit our attention to models $`\alpha `$ and $`\beta `$ which have a set $`B`$ of commands in common and for which the spectra of $`M_\alpha `$ and $`M_\beta `$ are the same. For a single command $`b`$, the question is whether the difference between the probability distributions $`\mathrm{Pr}_\alpha (|b)`$ and $`\mathrm{Pr}_\beta (|b)`$ is bigger than typical fluctuations expected in $`N(B)`$ trials. An answer is that two distributions are indistinguishable statistically in $`N(b)`$ trials unless
(2.10)
$$N(b)^{1/2}d(\mathrm{Pr}_\alpha (|b),\mathrm{Pr}_\beta (|b))>1,$$
where $`d`$ is the statistical distance defined by Wooters in Eq. (10) of . Furthermore, Wooters’s Eq. (12) shows for two models $`\alpha `$ and $`\beta `$ that differ only in the function $`|v`$,
(2.11)
$$d(\mathrm{Pr}_\alpha (|b),\mathrm{Pr}_\beta (|b))\mathrm{cos}^1|v_\alpha (b)|v_\beta (b)|.$$
To judge the significance of the difference between two models with respect to a set B of commands common to them, a scientist who chooses some weighting of different commands can define a weighted average of $`d(\mathrm{Pr}_\alpha (|b),\mathrm{Pr}_\beta (|b))`$ over all $`bB`$. The same holds if model $`\beta `$ is replaced by relative frequencies of outcomes interpreted from measured results.
It is noteworthy that the set of models statistically indistinguishable from a given model can be much larger than would be the case if the “$``$” of (2.11) were an equality, as follows.
###### Proposition 2.2.
For any set of outcomes, two models $`\alpha `$ and $`\beta `$ of the form $`(|v,\mathrm{𝟏},M)_B`$ can perfectly fit the relative frequencies of the outcomes (Proposition 2.1) and yet be mutually orthogonal in the sense that $`v_\alpha |v_\beta =0`$
###### Proof.
For any set of measured outcomes, there exists a perfectly fitting model $`\alpha `$ of the form in the proof of Proposition 2.1 for the general case, for which $`(j,b)dim(|w_j(b))>1)`$, and a corresponding perfectly fitting model $`\beta `$ such that $`(j,b)w_{\alpha ,j}(b)|w_{\beta ,j}(b)=0`$. For these two models, $`v_(b)\alpha (b)|v_\beta (b)=_j[n(j,b)/N(b)]w_{\alpha ,j}(b)|w_{\beta ,j}(b)=0`$ . ∎
Wooters extended the definition of statistical difference to unit vectors. While for any two unit vectors, there exist measurement operators that maximize the statistical distance between them, for any such operator there exist other vectors, mutually orthogonal, that have zero statistical distance relative to this operator. For this reason, among others, statistics still leaves the scientist needing something beyond calculation and measurement to determine a model, for the set of models closer than $`ϵ`$ in weighted statistical distance to certain measured results certainly includes all the models that exactly fit the data and, without special restrictions dependent on guesses, this set includes models that are mutually orthogonal. Models close to given measured data are not necessarily close to each other in the predictions they make.
### 2.6. Lattices of models
Properties 1 and 2 set up a big set of models $`(|v,U,M)_B`$, $`B`$, $`|v𝒱,U𝒰,M`$. Subsets of models of this set are a lattice under set intersection and union. Each command set $`B`$ establishes a smaller lattice of sets of models, and these lattices will play a part in the testing and calibrating of quantum computers, discussed in section 4, where a scientist encountering problems with a model chooses a set of possible alternatives, and then tries to narrow it. Often this narrowing is seen as choosing values of parameters within a form of model in order to obtain a best fit, say with a criterion of minimizing statistical distance between frequencies of outcomes interpreted from measured results and probabilities calculated from the model. One is free to think of the estimating of parameters in the language of a lattice of models as the using of measured results to select a model from a set of models.
From Proposition 2.1 that showed that the whole set of models defined by properties 1 and 2 is too big to permit measured results to select a model, we have:
###### Proposition 2.3.
For measured data to uniquely decide to within unitary equivalence which quantum-mechanical model of a set of models best fits experimental results interpreted as outcomes by a criterion of least statistical distance (or any other plausible criterion), the set of models must first be sufficiently narrowed, and this narrowing is underivable from the results and the basic properties 1 and 2 of quantum mechanics.
Something beyond measured results and calculations from equations is required to narrow a set of models so that measured results can select a model that is “best” by some criterion. Such an act of choosing undefined by calculation and results of observation is what we have called a guess.
### 2.7. Hidden guesswork in conventional quantum mechanical models
The proof casts in a clear light maneuvers conventionally made to narrow down the set of models. Sometimes a community of physicists is in mutual agreement about guesses deemed appropriate, and this agreement obscures from notice the fact that a guess is invoked. As an example of a widely invoked guess, most modeling in quantum physics supposes that the scientist can vary $`b`$ so as to vary $`U(b)`$ while holding $`v(b)`$ and $`M(b)`$ constant. Indeed, most models used in quantum physics are restricted to the subset of models having the special
###### Property 3.
The command $`b`$ is the concatenation of separate commands for the three types of operations, so that
(2.12)
$$b=b_vb_Ub_M,$$
where here the $``$ denotes concatenation of commands.
According to these models, one can vary any one of the three while holding the other two fixed. This specializes (2.1) to the more restrictive form:
(2.13)
$$\mathrm{Pr}(j|b)=v(b_v)|U^{}(b_U)M_j(b_M)U(b_U)|v(b_v).$$
An additional constraining guess characterizes models widely used in the analysis of quantum computers, a guess prompted by the desire to generate a unitary transformation as a product of other unitary transformations that serve as “elementary quantum gates.” For example, one may want to generate the unitary transformation $`U(b_{U,1})`$ $`U(b_{U,2})`$. To generate it one causes the CPC to transmit some $`b_U`$. For quantum computing to have an advantage over classical computing, the determination of this $`b_U`$ in terms of $`b_{U,1}`$ and $`b_{U,2}`$ must be of polynomial complexity . It is usually assumed that $`b_U`$ is the simplest possible function of $`b_{U,1}`$ and $`b_{U,2}`$, as follows.
Let $`B_UB`$ be a set of instrument-controlling commands, thought of as strings that can be concatenated. Suppose the function $`U`$ has the form $`U(b_1b_2)=U(b_2)U(b_1)`$ for all $`b_1b_2B_U`$ (note reversal of order). Then we say the function $`U`$ respects concatenation.
###### Property 4.
Quantum computation employs a subset of models in which $`U`$ respects concatenation.
###### Remark 2.1.
We present properties 1 through 4 not as properties of laboratory instruments, but as properties that a scientist can choose to demand of models. Whether the instruments act that way is another question. There are reasons, relaxation and other forms of decoherence among them, to expect limits to the precision with which instruments can behave in accord with properties 3 and 4. All four properties are used often enough to be conventions, in the sense that a convention is a guess endorsed by a community.
## 3. Petri nets to show choices open to guesswork
In orchestrating contact between mathematical models and laboratory instruments, scientists set up chains of cause and effect, expressed in computer programs with their “if-then” structure, not as static propositions but as designs for action. Such designs are implemented in experiments; an example is a feedback loop that adjusts the orientation of a filter according to a rule that tells what adjustment to make in immediate response to a result recorded by a light detector. On a more relaxed time scale, physicists make other connections by analyzing outcomes of one generation of experiment, using the equations of a model, to set up design instruments for a next generation. As remarked above, contact between equations and instruments depends on choices made by scientists, including choices of what set of models to start with, what model to choose for use by a CPC in generating commands, and what experiments to run. If these choices could be resolved by some combination of calculation and measurement, one could argue that they are irrelevant to physics. But the propositions of the preceding section show this is not the case, so the design and operation of contact between equations and instruments, with its ineradicable dependence on guesswork, cries out for attention as part and parcel of physics.
Although widespread in practice, the design of contact between equations and instruments is in its infancy as a topic for theoretical attention. A beginning can be seen in Benioff’s analysis of sequences of measurements (described quantum mechanically) in which subsequent measurements are functions of outcomes of preceding measurements . Called decision procedures, these involve classical feedback control equations to control instruments described quantum mechanically, in some cases with proved advantages . These efforts dealt with measurements occurring at a single location. Designs that put equations and instruments in contact over a network of cooperating investigators are wide open for future attention.
Logic in experiments, in feedback loops at many time scales, is logic in action. This is the logic of models that relate instrument commands to quantum vectors and operators. Here we adapt Petri nets to provide mathematical language by which to express and analyze designs for contact between equations and instruments, designs that include sequencing of effects, decision rules, and interactions among sequences of effects that scientists implement in their instruments. The nets will highlight choices resolvable only by resort to guesswork; they serve as a language with which one can express formally how guessing works in physics, case by case, within CPC-mediated investigations.
### 3.1. Requirements CPC’s
In order to adapt Petri nets to showing guess-demanding choices visible in CPC’s, we start by clarifying how a CPC differs from a Turing machine, on the way to adapting the Turing machine to process control and to use in a network of collaborating scientists. This lays the groundwork for introducing Petri nets.
#### 3.1.1. Timing in the execution of commands
The first thing that makes process-control computing special is timing. In the context of quantum-mechanical models, each unitary transformation maps states possible in one situation to states possible in another situation; for quantum computing this means mapping states possible at an earlier time to states possible at a later time. Thus a unitary transformation is implemented not all at once, but over a time duration. In practice, that duration depends on how the instruments implement the transformation. A written command $`b_U`$ acts as a musical score. Like sight reading at a piano, executing a program containing the command $`b_U`$ requires converting the character string $`b_U`$—the score—into precisely timed actions—the music. The piano keys, in this analogy, include the output buffers that control the amplitude, phase, frequency, and polarization of lasers of an ion-trap quantum computer or of radio-frequency transmitters for a nuclear-magnetic-resonance (NMR) quantum computer.
For this reason executing a command $`b_U`$ requires parsing it into pieces (signals) and implementing each signal at a time, the specification of which is contained in the string $`b_U`$. Either the CPC that executes a program in which $`b_U`$ is written parses the command into signals and transmits each signal at its appointed time, or the instruments receiving the command $`b_U`$, unparsed, contain programmable counters operating in conjunction with a clock that do this timed parsing. Such programmable counters themselves constitute a special-purpose CPC. So either the scientist’s CPC must execute commands by issuing an appropriately timed sequence of signals, or some other CPC attached to the instruments must do this. Either way, the capacity to execute programmed motion in step with a clock is a requirement for a CPC, distinct from and in addition to requirements to act as a Turing machine.
#### 3.1.2. Firewalls in a network of computers
Just as axioms set up branches of mathematics, guesses set up rules for the conduct of experiments and the interpretation of their results, rules often embedded in CPC’s. Collaborating scientists accept guesses from each other, at least provisionally, use these in experimenting and modeling; they evaluate some of them, sometimes refining or replacing them. This poses a problem for CPC-mediated inquiry, where guesses engender computer programs, for a scientist’s guess can reprogram a CPC, often for better but sometimes, by malice or accident, for worse. Scientists in a collaboration need to test each other’s programs and to limit the influence of any program, making the scope of influence of a CPC program a matter for negotiation among the collaborators.
An easy but narrow case is that of a computer running Gödel’s test for validity of a claimed derivation . To think about such testing, one models the computer by a Turing machine designed to start from a tape on which the claimed derivation is written and to halt leaving a “yes” or “no” on the tape, according to whether the claim is or is not valid. Such a Turing machine can be emulated by a universal Turing machine executing a testing program to check a passive (non-executed) file containing the claimed derivation.
Not just derivations, but also programs need to be tested with respect to what they do when they are executed. But what is to keep an executing program under test from infecting the program that tests it? Hardware walls of some kind are needed. By limiting our analysis to exclude remote login and insisting on computers that distinguish physically one interface from another, we can see a basic structure for testing programs and for limiting the reach of guesses of any one scientist in a network of CPC’s, based on operating two or more CPC’s concurrently with controlled interfaces between them, so the testing program and the program under test execute on separate CPC’s, with an interface controlled by the testing CPC. By virtue of concurrent operation of CPC’s with controlled interfaces, guesses made by collaborators can set up programs that frame choices open to guessing by any one scientist, and that test the performance of the scientist’s programs within that frame of choice, allowing freedom to a scientist to program one part of the investigation while insulating other parts. Hardware walls that limit the reach of one person’s guesses at any moment are one many motivations for stressing a network of concurrently operation CPC’s.
### 3.2. Turing machines and Petri nets
Here we provide language for displaying and analyzing guess-demanding choices visible in files of CPC’s used by collaborating scientists who on occasion reprogram those choices. As a model of a CPC, we assume that each CPC of a network is a Turing Machine adapted for Process-control (TMP), to be defined. Making sense of networks of TMP’s handling equations and controlling instruments calls for a descriptive capacity that allows for various viewpoints at various levels of detail. We introduce a specialized use of fragments of colored Petri-nets, opened to exogenous influences, to portray the programming and running of programs in a network of TMP’s operated by collaborating scientists.<sup>5</sup><sup>5</sup>5Our use of Petri nets is impressionistic and a more technical presentation will doubtless be rewarded by exposing issues here overlooked.
Different viewpoints and levels of detail are accommodated by morphisms in the category of nets. Isomorphisms between Petri nets trade net detail for color detail . These will be combined with coarsening maps that suppress detail, for example by mapping colored tokens to black tokens. We will show how the programming of a universal TMP (UTMP) portrayable as a single Petri net can produce any number of patterns of use of instruments and equations, portrayable by a host of different Petri nets. This general picture of process-control computation will show programs and other guesses as colors on tokens that a scientist enters on a game board defined by a fragment of a Petri net, and equations of quantum mechanics written as guesses by a scientist will be seen as colors on tokens that take part in directing and interpreting the use of laboratory instruments.
#### 3.2.1. Writing vs. executing a program
Computers rest on the writing of motionless characters on a page to describe something moving, a puzzle solved in music by writing notes on staves, to be read in step with a swinging pendulum that chops time into moments, so that written notes that portray a still picture for each moment direct the motion of the playing of a musical instruments . The logical machinery of a computer moves in response to triggering signals, “tick” and “tock”, synchronized to distinct phases of the swinging of a pendulum. Computer designers employ truth tables, each of which specifies the response of a clocked circuit at a tock to a stimulus present as an input at a preceding tick. A row of a truth-table can be drawn as a transition in a Finite State Machine (FSM). By coupling an FSM to a memory of unlimited capacity, one arrives at the theoretical concept of a Turing machine, various special cases of which perform various special tasks . And here is the crux of programming: because a state machine is describable by still writing—a table—a Turing machine can be designed to be universal. By coding into its memory the table that describes any given special Turing machine, one causes the universal Turing machine to emulate the given special Turing machine. So, apart from speed and memory requirements, the single universal Turing machine can be put to doing any of the things that any of the special Turing machines can do, making it potentially convenient, once adapted to process control, to designing and implementing contact between equations and instruments. (But demands for quick response require in some cases devices streamlined to a special task better modeled by a special Turing machine than by a universal one.) The next tasks are to adapt the Turing machines, special and universal, to process control, and after that to express them formally by use of colored Petri net fragments.
#### 3.2.2. Turing machine for process control (TMP)
To adapt a Turing machine as a model of a process-control computer, we leave the coupling of the FSM to the memory unchanged but add input and output buffers to the FSM. As for the FSM, at whatever level of detail of description one chooses, the control structure of a program (with its “if-then” statements) can be viewed as an FSM consisting of (classical) states drawn as circles, connected by directed arcs, with each arc labeled by an input $`I`$ that selects it and by an output $`O`$ ; a fragment of such a picture is shown in figure 2(a). An FSM serves as a game board on which a single token can be placed to mark the “current state.” Heading toward the hooking together of FSM’s to make a Petri net, we suppose that each arc in the FSM is punctuated by a tick event and a tock event, drawn as small boxes, enlarging the FSM into a special case of a condition-event Petri net fragment, as shown in figure 2(b). Once colors are introduced, states shown as dashed circles pointing into an event of the FSM from outside will become the means to express the entrance of guesses. These states are assumed to receive tokens put into them by scientists and instruments undescribed by events of the net. Similarly, dashed states pointed to by arcs from an event are assumed to have tokens taken from them by agents undescribed by events of the net. Figure 2(c) streamlines the picture to the form we shall use, in which more or less vertical arcs are understood to point downward, the dashed states are left undrawn, as are all states with one input and one output event. To emphasize the input and output arcs with their extra tokens, we often call this an FSM fragment to distinguish it from the FSM form of figure 2(a).
To define a Turing machine for Process-control (TMP), we adapt the FSM of a Turing machine to have for each of its states a cartesian product of states of a set of clocked internal registers and, in addition, input buffers and output buffers, which allow input/output transactions with a scientist, with laboratory instruments, and with other TMP’s.
#### 3.2.3. Colored tokens
By replacing the black tokens of an FSM fragment by a colored tokens and adjoining to each event a function that defines colors on output tokens in terms of colors on input tokens, any FSM fragment can be mapped one-to-one to the drastic form of figure 3, in which color changes substitute for most of the moves of black tokens on a bigger net. A “fork in the road” for black tokens, turns into a choice between red and green, so to speak, so the descriptive burden is taken up by the functions $`f_{\text{tick}}`$ and $`f_{\text{tock}}`$; $`f_{\text{tick}}`$ defines the color of a token placed on an internal state depending on a list of colors, one for each input, while $`f_{\text{tock}}`$ defines a list of output colors depending on the color of the token on the internal state. The vertical arc is to be read as directed downward, and the big circles at the top and bottom of a path signify that the path is wrapped around a cylinder, so the top is a continuation of the bottom, i.e. a loop. An FSM fragment in which the token carries a color will be called a colored FSM.
#### 3.2.4. Other mappings
Less drastic mappings are also possible. Any two states of a single FSM can be merged without breaking any arcs by augmenting the color rules in the events that feed them and the events fed by them. If a set of states connected to one another by events is mapped into a single state, the single state then connects to an event that loops back to it; this results in a place-transition Petri net, but not a condition-event net. We restrict the mappings dealt with here to ones that avoid pasting tick and tock events together, thereby avoiding self loops. Two events of an FSM that link the same pair of states can be merged by distinguishing external inputs and outputs by color instead of by place.
The mappings discussed so far are net isomorphisms: they map markings of one net bijectively to markings of the other and preserve the one-step reachability of one marking from another (by the firing of an event). Inverses of these bijections take more richly to less richly colored nets. Going in this direction depends on each state of a colored FSM having a set of possible colors associated with it ; then any colored transition corresponds one-to-one to a set of transitions obtained by partitioning sets of colors of input states, as illustrated in figure 4 for a two-in, two-out transition with color sets $`A`$, $`B`$, $`C`$, and $`D`$, each partitioned into “+” and “$``$” subsets. For this to make sense, it must be that an event which has tokens in all its inputs cannot fire unless the colors of the tokens comprise an element of the domain of its color function; we assume this firing rule.
One gets a coarser description by use of a surjective map that is not an isomorphism by dropping the color distinction and dropping the color functions from the transitions; this coarsening, however, preserves a one-to-one correspondence between the number of firings in one net and the number in another. All these maps are continuous in the net topology , and, as emphasized by Petri , nets form a category in which the morphisms are continuous maps, an idea that extends to nets with colored tokens .
#### 3.2.5. Disciplined coarsening of time
Some other kinds of continuous coarsening maps bundle up multiple event firings into a single firing; as when one describes e.g. “running a program” as a single event. This brings us to the first of several areas open to future work, for, more than other computing, process control benefits from well defined timing, and in particular from machine and software design that allows systematic, well controlled mappings that take a certain number of firings in an FSM to a single firing, so that one can think at a coarser level while still maintaining discipline in timing.
A striking example of the need to design programs that run in the same time for all inputs from some set $`I`$ occurs in quantum computing. For example, suppose that $`𝐔`$ is the universal unitary operator defined by Deutsch to operate on basis states of the form $`|s;𝐧;𝐦`$ where $`s`$ is the location of the scanned square, $`𝐧`$ is the state of the FSM-processor ($`𝐧=0`$ is the starting state and $`𝐧=1`$ the halt state) and $`𝐦`$ is the tape . For this to work in a computation that takes advantage of quantum superposition, one needs $`r[(xI)𝐔^r|0;0;x,0=|0;1;x,f(x)]`$; however, this is by no means implied for a program $`\pi _f`$ for which (as is usual in borrowed classical programs) one can assure only $`(xI)(r(x))[𝐔^{r(x)}|0;0;x,0=|0;1;x,f(x)]`$ . An interesting topic for future study is the complexity of converting various classes of programs with variable running time to programs running in a time independent of the input for some set of inputs.
#### 3.2.6. Cartoon of UTMP
Ignoring the laboratory instruments for the moment, by connecting input- and output-signals from a suitable FSM to a scientist and coupling the FSM to an unlimited memory, one gets a Universal Turing Machine (UTM) that provides for continual communication with a scientist, as shown in figure 5(a), in which boxes connected by a horizontal line are read as a single event. We cartoon the UTM in the condensed form of figure 5(b). By adding input- and output-signals from the FSM to laboratory instruments and to other FSM’s, one gets a Universal Turing Machine adapted for Process control (UTMP), as shown in figure 5(c); again almost all of the burden of description is in the color functions, here called $`T_1`$ and $`T_2`$ (for Turing) that define a finite state machine that operates a UTMP. We assume that at some level of description, the ticks and tocks of the UTMP slice time into moments not only for the UTMP but also for the scientist at a keyboard and the instruments on the laboratory bench; we assume that input tokens from the scientist and from the instruments arrive at the UTMP synchronized with the UTMP pendulum. If the scientist enters nothing at a given clock tick, then the token taken by the UTMP from the input buffer for the scientist carries the color “empty,” and if the instruments enter nothing, the token from the input buffer for the instruments carries the color “empty”; similarly the UTMP marks output tokens with the color “empty” if it writes nothing else on them.
#### 3.2.7. A scientist controls a UTMP
To see the structure imposed on physics by the UTMP, one must think as if the UTMP were delivered to a scientist in a bare condition: no installed software,<sup>6</sup><sup>6</sup>6The scientist can borrow software and install it, but is responsible for it. the FSM in a starting state, and the memory all blank. We assume that the function $`T_1`$ operating on empty input tokens, the starting state of the FSM, and a blank memory produces empty output tokens and makes no change in the FSM state or the memory or the memory location scanned. Finally, we invoke the universality of a UTM to assume that the functions $`T_1`$ and $`T_2`$ are fixed (by a manufacturer, so to speak) independent of whatever laboratory instruments need to be considered and independent of all action by the scientist. These assumptions imply
###### Proposition 3.1.
Whatever a UTMP does besides staying in its starting state and taking in and putting out empty tokens is in response to input tokens.
We invoke this proposition to view the scientist as precluded from defending questionable management of equations or instruments by saying “the computer did it.” If a CPC does something, it executes a program; we view the scientist as responsible for any program entered (as a colored token) into the UTMP and for running the program on any particular occasion.<sup>7</sup><sup>7</sup>7This rules out taking for granted the operating system, instrument-managing programs, a simulator, and whatever other programs come pre-installed in a commercially available CPC.
#### 3.2.8. Reprogramming always an option
We assume the UTMP is isomorphic to the net shown in figure 6, so that the scientist has a recurring choice of letting the UTMP run as programmed or of interrupting it to reprogram it. By programming a UTMP, a scientist can simulate an arbitrary special Turing machine. At will, the scientist can interrupt a program in execution to change to a program that simulates a different special Turing machine, corresponding to a different FSM and a different net. One can glimpse this in figure 4, where it is apparent that if the colors are limited to the sets $`A^+`$ and $`B^+`$, then six of the eight events are precluded from firing, and the net is in effect reduced to the fragment defined by the selected colors. In this way the part of the net that actually fires, corresponding to the event “Use existing program” of figure 6, is variable in how it acts and in the net by which one portrays it in more detail, according to the scientist’s actions in providing and running programs.
#### 3.2.9. Plug and play
To see how UTMP’s can be connected (as well as the detail of how the FSM of a TM or TMP is connected to the memory), we introduce a signal that is phased just opposite to an FSM: the signal takes an input at a tock event and issues an output at a tick event. Then FSM A can send a signal (which can convey a message as a token color) to B (which can be either another FSM or a memory), as shown in figure 7, provided the signal path is short enough compared to the clock rate of the machines. This use of a signal synchronizes A with B. For two-way communication, one adds a signal going the other way. If communication over a distance long compared to the clock period is called for, then a chain of communication over intermediating UTMP’s, is necessary, with the result that more firings of an event of A are required before a consequence of one firing can propagate to B and return as a property of a color on a token at a later firing of the A-event. The use of colored tokens sets up an area for future investigation of replacing the awkward definition of synchronic distance with a measure of synchronization that counts firings in circuits of color effects, without having to add artificial elements to a net.
### 3.3. Net fragments formalized
Portraying logic operating in CPC’s calls for fragments of Petri nets, not complete nets, to allow for guesses as token colors definable neither by results of experiments nor by calculations. From among the standard definitions of a Petri net, the one we use is $`(S,E,F)`$, where $`S`$ is a set of states, $`E`$ is a set of events, and $`FS\times EE\times S`$ is the flow relation. In order to make room for guesses from a scientist and results of instruments inexpressible in the logic defined by a net but essential to setting it up, the nets used are all net fragments, which we define as follows. A net fragment is a structure $`(S,S_I,S_O,E,F)`$ where $`S`$ is a set of states of CPC’s, and $`S_I`$ is a set of states of input signals (e.g. from A/D converters to a CPC input buffer), disjoint from $`S`$, allowing for input to the CPC from a scientist and laboratory instruments. $`S_O`$ is a set of states of output signals disjoint from both $`S`$ and $`S_I`$, allowing for output from the CPC; the flow relation is expanded so $`F[(SS_I)\times E][E\times (SS_O)]`$. States of $`S_I`$ are assumed to have tokens placed in them by some means beyond the net, and states of $`S_O`$ are assumed to have tokens removed from them by means beyond the net. Our pictures show stubs of arcs from states of $`S_I`$ to events and from events to states of $`S_O`$ while omitting the circles for these states. Associated with a net fragment is a “reduced net” obtained by omitting the states of $`S_I`$ and $`S_O`$ (and dropping the arc stubs); using this reduced net, one can explore issues of liveness and safety . The events of $`E`$ express computer logic and nothing else. As an example of a guess used in designing contact between equations and instruments, a mathematical model entered by a scientist as a colored token in an $`S_I`$ state can assert whatever rules the scientist chooses to relate tokens received from instruments in $`S_I`$ states to commands sent to them as colored tokens in $`S_O`$ states. In this way the net fragment expresses the difference between such a model, with its guesswork, as a color on a token and how the instruments actually behave by producing colored tokens on their own.
## 4. Net-based portrait of guesswork needed to test and calibrate a quantum computer
In section 2 choices of equations to link to instruments were shown inescapably open to guesswork, bidding to make guesswork part and parcel of physics. The availability of net fragments described in section 3 brings within physics the study of contacts between equations and instruments by making available to analysis relations of sequence, concurrency and choice expressed in these contacts and in the guess-dependent actions that set the contacts up. Here we turn from nets themselves to attention to an example problem in which a net illustrates an important structure needed to link equations to instruments. Besides the net explicitly shown in figure 8, the availability of nets provides a framework in which to view the main topic of this section, the problem of resolving a choice of commands by which a CPC manages a quantum computer. That framework can be used in the future to ask other questions, to do with: how do the necessities of quantum-mechanical models, classical process control, and guesswork interact; how are FSM’s as program structures affected by use of models that are quantum mechanical; how does the need for CPC’s to mediate between quantum-mechanical equations and instruments change our understanding of quantum mechanics?
Turning to the case at hand, some telling illustrations of guesswork needed to link models to instruments arise in quantum computing. To build a quantum computer, say to solve problems of factorizing and searching , a scientist must choose quantum-mechanical equations and laboratory instruments to work in harmony. Quantum computational models call for quantum gates that are unitary transformations, each a tensor product of an operator on a 1-bit or 2-bit subspace of the Hilbert space $``$ and identity operators for the other factors of the tensor product. Note that each permutation of a non-identity factor with an identity factor is a distinct gate, calling for a distinct command to the instruments that implement it. For this reason, the number of quantum gates for an $`n`$-bit quantum computer grows faster than $`n`$. Call this number $`G(n)`$ and let the set of gates be $`U_1,\mathrm{},U_G`$. The most commonly used models of quantum computers can be put in the form :
* Prepare a starting state independent of the input (e.g. the integer to be factorized).
* Transform the state by a product of quantum gates that depends on the input.
* Make a measurement independent of the input.
For an example, suppose the scientist assumes properties of models 1 through 4 and looks for the model that gives the least mean-square deviation between relative frequencies of outcomes and probabilities calculated by (2.13). To factorize an integer $`I`$, a classical computer program is converted to a product of $`K(I)`$ quantum gates, a number that rises faster than linearly with $`\mathrm{log}I`$. To obtain the effect of multiplying the gate transformations, the scientist must first have solved the model to determine the command $`b_{U,j}`$ for each gate $`U_j`$ occurring in the product. As in the portrait in section 3 of putting tokens into a net fragment, the scientist programs a CPC to transmit a command $`b_v`$ to prepare an initial state $`|v`$, commands $`b_{U,j}`$ for the gates needed, and a command $`b_M`$ for a measurement. This endeavor is known to exhibit the following four features:
1. The instruments are valuable as a quantum computer insofar as their results substitute for a more costly classical calculation defined by the model.
2. An inexpensive classical computation (e.g. with the CPC) tests whether outcomes interpreted from results correctly solve the problem.
3. Quantum indeterminacy imposes a positive probability that a result fails to provide a correct answer, so multiple tries with the instruments are the rule, and a wrong answer does not by itself imply a fault in the instruments.
4. The tolerable imprecision of instruments implementing the chosen model of a quantum gate diminishes as the inverse of the number of gates $`K(I)`$ in the sequence .
Because the number of gates required in the product rises with the size of the integer to be factorized, feature 4 implies that passing the test for smaller integers is no guarantee against failure of the instruments to factorize larger integers, unless the model or the instruments or both are refined. This requires, in turn, that a CPC intended for use on progressively larger integers be organized to switch between a mode of using the quantum computer and a mode of inquiring into its performance, e.g. so as to determine commands that make it behave more precisely in accord with the desired quantum gates. This calls for a program for the CPC that expands the events “Use existing program” of figure 6 to that of figure 8.
### 4.1. Navigating the lattice of models to get better commands
As an example of what goes on within the coarsely portrayed event “Calibrate,” suppose a scientist who uses a model $`\alpha `$ of the form $`(|v,U,M)_B`$ finds it works for small integers, but fails for bigger ones, which require more precise gates, which in turn requires calibrating (i.e. adjusting) the commands used to generate gates. This means giving up model $`\alpha `$ and choosing some alternative model $`\beta `$. A scientist does not choose a model all at once, but starts with some set of models and then narrows down on a smaller set, sometimes to a single model, a process open to guesswork at various stages. At one stage, the scientist may need to relax a constraint on models, leading to a bigger set of models from which to choose; at another stage the scientist may guess a new constraint, narrowing the set of models under consideration. By such a back and forth procedure, the scientist gives up $`U_\alpha `$ and arrives at a new function $`U_\beta `$ (and hence a new model) with the hope that solving this function for a command $`b_{U,j,\beta }`$ for gates $`U_j`$, $`j=1,2\mathrm{}`$, that will succeed for factorizing larger integers than did the commands obtained from $`U_\alpha `$. (This makes a need for models adapted to homing in on results, with some metric on B, so that a small change in the command $`b_U`$ results in a small change in e.g. $`U(b_U)`$; while properties 3 and 4 are a start, going beyond them is left to the future.)
To get a better model, the scientist guesses a set of models and hopes to find within it a model that better fits measured results interpreted as outcomes. If no model of the set adequately fits these outcomes, the scientist can first broaden the set of models and next try to guess a property that will narrow the set, not to the original model, but to one that fits better. The recognition of guesswork assures us that so long as progressively more ambitious goals of precision keep being introduced, there is no end to the need for adjusting both models and the laboratory instruments.
### 4.2. Sample sizes needed to choose between models ofgates
As discussed in section 2.5, the number of trials needed to statistically distinguish one model from another is bounded from below by the inverse square of a weighted statistical distance between the two models. Small numbers of experimental results can sometimes decide between distant models, but never between models that are close. In particular, distinguishing experimentally between two models for quantum gates can demand large samples:
###### Proposition 4.1.
Models $`\alpha `$ and $`\beta `$ that differ only in $`U`$, with spectral norm $`U_\alpha (b_U)U_\beta (b_U)=ϵ>0`$, are statistically indistinguishable for a command $`b`$ unless
(4.1)
$$N(b)ϵ^2.$$
###### Proof.
The models $`\alpha `$ and $`\beta `$ under the stated condition are unitarily equivalent to a pair of models that differ only in $`|v`$ with $`\mathrm{cos}|v_\alpha |v_\beta |`$ $`ϵ`$. The proposition then follows from (2.10) and (2.11).∎
We argue elsewhere that this is a serious and heretofore unappreciated challenge to bringing instruments into working order as quantum computers, made visible by attention to the need for guesswork in linking of laboratory instruments to equations of quantum mechanics .
## 5. Concluding remarks
Gödel proved that no one true structure could be generated by sitting in a room with blinds drawn, writing down axioms. Quantum mechanics tells us that with the blinds up and the world of physical measurement available, the situation remains much the same. Just as the openings for new axioms are uncloseable in mathematical logic, so in physics guesswork is part of the foundation.
The net formalism can be put to use both to address improving the contacts between equations and instruments, fostering advances in theory and in instrumentation, and, at a more abstract level, to pose problems pertaining to universal Turing machines adapted to process control. By formalizing commands to instruments, the techniques presented here extend the reach of set-based mathematics into the area of contact between equations and instruments, and open to study within physics of some of what physicist do in the course of doing physics. This extends a parallel beachhead established already in mathematics by Gödel’s study of what a mathematician does to prove a theorem and Turing’s analysis of a mathematician who makes a note by which to resume an interrupted computation.
## 6. Acknowledgment
We acknowledge Amr Fahmy for showing us our debt to Gödel’s proof of incompleteness in mathematical logic. We are indebted to Steffen Glaser, Raimund Marx, and Wolfgang Bermel for introducing us to the subtleties of laboratory work aimed at nuclear-magnetic-resonance quantum computers . To David Mumford we owe our introduction to quantum computing from the standpoint of pure mathematics. We are greatly indebted to Anatol W. Holt and to C. A. Petri for conversations years ago, in which each pointed in his own way to the still mysterious expressive potential of nets.
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# Effects of Photospheric Temperature Inhomogeneities on Lithium Abundance Determinations (2D)
## 1. Introduction
It would be an offense to the audience here to pretend to explain why it is important to determine accurately the abundance of <sup>6</sup>Li and <sup>7</sup>Li in the oldest stars. In this respect, we have nothing to add to the exposition by François Spite (this volume). We will right away mention the two major problems which may cast doubts on our real knowledge of the actual initial abundance of Li in the oldest stars, and consequently in the primordial matter. (i) While standard models of the internal structure of metal-poor dwarfs do not deplete <sup>7</sup>Li, more sophisticated models including rotationally induced mixing (Pinsonneault et al. 1992) have predicted that the measured abundance in the photosphere is 5 to 10 times less than the initial abundance representative of Big Bang material. (ii) On top of that, Kurucz (1995) claimed that the hot and cold convective structures produce large effects in metal-poor stellar photospheres, where the convection zone reaches the line formation layers. The claimed effect is an overionization of Li by a factor of 10, leading to an underestimation of the abundance of Li when derived from the resonance line of Li I ($`\lambda `$ 670.8 nm) in the usual way.
If these two statements are correct, the true abundance of Li in primordial matter is 50 to 100 times higher than the value derived from 1D, LTE models of halo subdwarfs so far. The first factor of 5 to 10 has been discussed in a previous paper by Ryan (this symposium), and shown to be likely much smaller, of the order of 1 to 1.4. We shall not come back to this point, which we consider as very well treated.
Before this symposium, a single paper (Asplund et al. 1999) has dealt with the question of the other factor of 10 claimed by Kurucz (1995), whose arguments were based on a simplified two-column model. In contrast, the work by Asplund et al. relies on realistic 3D hydrodynamical models, similar to the simulations of the solar granulation (Stein & Nordlund 1998), but with parameters appropriate for two metal-poor stars: HD 140283 and HD 84937, both subgiants. The computation of the lithium resonance line was made under the assumption of LTE, and the correction to be applied to the Li abundance derived from standard 1D models was found to be large, of the order of -0.2 to -0.35 dex. Note that these corrections have the opposite sign as Kurucz’s prediction! However, Kiselman (1997, 1998) had shown, in the solar case, that NLTE and LTE computations lead to significantly different values of equivalent widths of the Li I $`\lambda `$ 670.8 nm line over hot and cold structures (see Fig. 3 of his 1997 paper, top panel).
For this reason, we decided to undertake NLTE radiation hydrodynamics computations for the case of a metal-poor star, and we report here on the results of this investigation. In the next section we recall former work related to simulations of the solar granulation, a useful benchmark for checking the theory, but not directly applicable to metal-poor stars. In section 3 we describe the assumptions underlying the construction of the 2D RHD models used for the spatially resolved computation of the lithium resonance line. Section 4 gives the description of the NLTE treatment of the Li atom, and section 5 summarizes our results and compares them to those presented by M. Asplund (this symposium). Finally, our conclusions are listed in section 6.
## 2. Former work at solar metallicity
While there is only one paper dealing with multidimensional atmospheres for metal-poor stars (cited above), there are several studies for the solar case, aimed at understanding the variation of the continuous radiation intensity (granulation), and the behavior of spectral lines across the solar granulation pattern. Several of these works use snapshots from 3D simulations by Stein & Nordlund (1998), such as Kiselman (1997,1998) and Uitenbroek (1998). Gadun & Pavlenko (1997) use their own 2D simulations.
Of particular interest for us are the papers dealing with the combined effects of multidimensional structures and NLTE (Kiselman and Uitenbroek). It is clear from Kiselman (1997) that the NLTE behavior of the Li I $`\lambda `$ 670.8 nm resonance line is drastically different from its LTE behavior, in 2D as well as in 3D models. The result which is the most relevant for us is the difference of 30 per cent on the predicted mean equivalent width $`W`$ of the line, leading to a similar change in the derived lithium abundance (NLTE/LTE abundance correction +0.15 dex). But another interesting difference is the reverse behavior of the equivalent width $`W`$ as a function of surface continuum brightness $`I_c`$. While LTE computations result in a strongly positive slope in the $`W`$ versus $`I_c`$ diagram (with a large scatter around the mean relation), NLTE computations show a slightly negative slope and a much tighter (anti-)correlation between $`I_c`$ and $`W`$. This reflects the fact that the population of the Li I levels is much more controlled by the local temperature in LTE than in NLTE, where, for weak lines, the photoionization rates play the dominant role. So, even if the general conclusion of the above mentioned papers is that, in the solar case, the abundance determination of Li is not strongly affected by the combined effects of temperature fluctuations and NLTE, in comparison to what is obtained with classical 1D models having the same effective temperature and gravity, it appears unsafe to compute abundances from multidimensional stellar atmospheres based on the assumption of LTE line formation.
## 3. 2D radiation hydrodynamics models
Our LTE and NLTE computations have been performed on the basis of several snapshots from a 2D numerical simulation of convection in a stellar envelope having the same effective temperature and gravity as the Sun, but a 100 times smaller metallicity. Basically, the time dependent equations of hydrodynamics are solved for a compressible fluid, with an energy equation including 3 terms: turbulent and shock dissipation of kinetic energy, diffusive transport of heat, and radiative energy exchange. The main limitation of the code is the restriction of the flow to two spatial dimensions. Magnetic fields and rotation are ignored.
Apart from these simplifications, as much realistic physics as possible is included. The equation of state accounts for ionization of hydrogen and helium as well as H<sub>2</sub> molecule formation, opacities have been adapted from Kurucz’s ATLAS code and include line absorption. For the computation of the radiative energy balance, we employ a multi-dimensional, non-local, frequency-dependent radiative transfer scheme, actually solving the transfer equation along 26880 independent rays of various inclinations, using an efficient modified Feautrier method (Feautrier 1964). At the bottom boundary, inflowing matter has a given specific entropy, which is adjusted to produce the prescribed effective temperature of the atmosphere. Energy dissipation on small scales is roughly modeled by introducing a subgrid scale eddy viscosity, depending on the grid resolution and local velocity gradients in the usual way. Details of the employed hydrodynamics code can be found in Ludwig et al. (1994) and Freytag et al. (1996).
Fig. 1 shows a sample snapshot from our metal-poor Sun simulation. Note the complex velocity pattern and the occurrence of very strong temperature gradients. The relevant region for the formation of the Li I line is the $`\tau =0.1`$ contour line.
## 4. NLTE computation of the Li I resonance line
The computation of the Li I spectrum is greatly simplified by the fact that all lines of Li I are weak, Li being a trace element. So the radiation field in the line is, to first approximation, the same as the continuous radiation field, a single iteration being sufficient for taking care of the small perturbation of the monochromatic radiation field brought about by the line.
We have, as a first step, approximated the Li I atomic configuration by a five level atom, exactly as done by Uitenbroek (1998). This leaves six permitted bound-bound transitions, and five photoionization rates needing the computation of the continuous radiation field at frequencies above the threshold, until the contribution of the product of photoionization cross-section and mean intensity of the radiation becomes negligible. Because in the UV the contribution of lines to the opacity is important, we have used Kurucz’s Opacity Distribution Functions (ATLAS 9) for the relevant metallicity. This multiplies the computation time by 12, as each opacity bin is subdivided into 12 subintervals. So, for each snapshot, the transfer equations must be solved for about $`12\times 120`$ wavelengths, along 26880 different rays (note that these extensive computations are done only after the actual hydrodynamical simulation for a few selected snapshots). After this, it is possible to compute all the coefficients in the equations of statistical equilibrium (see Mihalas 1970, p. 144). Once the departure coefficients $`b_i`$ are evaluated, the line can be computed using the source function:
$$S=\frac{\kappa _c}{\kappa _l+\kappa _c}B_\nu +\frac{\kappa _l}{\kappa _l+\kappa _c}S_l$$
where:
$$S_l=S_{ij}=\frac{2h\nu _{ij}^3}{c^2}\frac{1}{(b_i/b_j)\mathrm{exp}(h\nu _{ij}/kT)1}$$
is the line source function. The departure coefficient for the lower level is $`b_i`$ and the one for the upper level $`b_j`$. The other notations are standard. Subscript “c” stands for continuum, and “l” for “line”. Note that in NLTE the expression of the partition function is modified and becomes:
$$U=\underset{i}{}b_ig_i\mathrm{exp}(\chi _i/kT)$$
where the $`g_i`$ and $`\chi _i`$ are the statistical weight and the excitation energy of level $`i`$, respectively .
## 5. Results and discussion
We have first tested our program on a Kurucz ATLAS 9 1D solar model to see whether the $`b_i`$ had the expected behavior, already computed by Carlsson et al. (1994). Fig. 2 shows the depth-dependence of the first 3 departure coefficients, applying to levels 2s, 2p and 3s, respectively.
Next, we have computed the equivalent width of the Li I resonance line for the Kurucz 1D solar model and for two 2D solar snapshots, still for the logarithmic Li abundance 2.2. Fig. 3 shows the variation of the equivalent width $`W`$ (for the intensity normal to the surface) of the Li I 670.8 nm line over the simulated granulation pattern, both as a function of the horizontal position $`x`$ (top) and as a function of the continuum intensity $`I_c`$ (bottom) for one particular snapshot. Note the wide variation of $`W`$ computed in LTE, compared to the much more limited excursion of $`W`$ computed in NLTE. The mean equivalent widths for the flux spectrum integrated over the full length of the sample are given in Table 1 for the two solar snapshots and for the 1D reference model having the same effective temperature, gravity and (solar) metallicity. In each case, the line is computed in LTE as well as in NLTE. We note that, in NLTE, the results for the 2D snapshots do not differ significantly from the 1D case.
Finally, we have carried out a similar procedure for our metal-poor stellar example, again computing the equivalent width of the Li I resonance line for a Kurucz 1D model and for five snapshots from our 2D hydrodynamical simulation; as before, a Li abundance of 2.2 was adopted. Fig. 4 shows the variation of the equivalent width $`W`$ of the Li I resonance line over the stellar granulation pattern for the typical snapshot displayed in Fig. 1. The difference between LTE and NLTE is even more pronounced than in the solar case, the NLTE correlation between $`W`$ and $`I_c`$ being much tighter and of opposite sign compared to LTE.
The mean equivalent widths derived from the horizontally averaged flux spectrum are listed in Table 2 for the 5 snapshots and for the Kurucz 1D reference model. In LTE, the 1D/2D difference is huge (granulation abundance correction $`0.45`$ dex). But remarkably, the 2D NLTE line strengths show very little dispersion and do not indicate any significant offset with respect to the 1D case. An obvious conclusion from these results is that the 2D LTE computations are way off, strongly underestimating the Li abundance. In NLTE, the error introduced by representing the inhomogeneous stellar atmosphere by a flux-constant 1D Kurucz model appears to be almost negligible.
The mechanism behind the spatial variation of the line strength is clearly identified on Figs 1 and 4: hot granules produce at the same time a steeper temperature gradient and a lower temperature in the line formation region. In LTE, the latter leads to an overpopulation of the lower level of the transition due to a shift of the ionization equilibrium towards neutral particles (Saha equation). Since both effects enhance the equivalent width of the line, the correlation between continuum intensity and LTE line strength is clearly positive (there is considerable dispersion around the mean $`W(I_c)`$ relation, however, due to the presence of inclined thermal inhomogeneities). This result is inverse of what is expected for a set of hot and cool radiative equilibrium atmospheres lined up side by side: here the line would weaken in the hot atmosphere, because the population of the lower level is the dominant factor. It is clear that the reasoning of Kurucz fails, essentially because the actual vertical stratification of hot and cold regions has little to do with the stratification of hot and cool radiative equilibrium atmospheres.
A positive $`W`$-$`I_c`$ correlation alone is not sufficient to explain the LTE result that the mean equivalent width is much larger in 2D than in 1D: as long as the fluctuations of the line strength with temperature remain in the linear domain, the mean equivalent width is not affected by the temperature fluctuations. But as the population of the ground level $`N_0`$ of Li I varies exponentially with T, a nonlinearity sets in: symmetrical temperature fluctuations produce asymmetrical population variations. Assume that Li is almost completely ionized and consider only the exponential factor of the Saha equation for simplicity. Then
$$N_0(T)=N_0(T_0)\mathrm{exp}(+\chi /kT)=N_0(T_0)\left\{1qs+qs^2+\frac{1}{2}q^2s^2+\mathrm{}\right\}$$
where $`\chi `$ is the ionization potential and we have defined $`q\chi /kT_0`$ and $`s(TT_0)/T_0`$, $`T_0`$ being the mean temperature. Taking the horizontal average assuming symmetrical temperature variations ($`s=0`$), we obtain $`NN_0(T_0)(1+0.5q^2s^2`$ ($`\sqrt{s^2}0.04`$; $`q12`$). So the mean of $`W`$ is biased towards larger equivalent widths in an inhomogeneous atmosphere. This explains part of the 1D/2D difference found in our numerical computations. The main contribution, however, is attributed to the lower mean photospheric temperature in the 2D model as a result of adiabatic cooling due to overshooting: $`T_0(2D)<T(1D)`$.
In NLTE, the local temperature plays little role. Rather, the radiation field is the dominant factor. Hence, photoionization overionizes the lower level over a hot granule, and the equivalent width becomes smaller. As a result, NLTE equivalent widths are anti-correlated with the continuum surface brightness. The variation of $`W`$ is much smaller than in LTE, because the angle-averaged radiation field depends only weakly on horizontal position at the height of line formation.
Finally, we would like to mention the work by M. Asplund, who has also presented his NLTE line formation calculations for Li I on this symposium. His results are based on the 3D snapshots that he had used previously for the LTE investigation of the subgiants HD 140283 and HD 84937. He reached, on this completely independent set of hydrodynamical models, and with a different NLTE code, the same conclusions as we did: NLTE line formation in multidimensional models is quite different from the LTE case, in the way that the NLTE multi-dimensional Li abundance is much closer to the abundance derived from classical flux-constant 1D models. A detailed comparison of our results shows that the remaining differences can be traced to the use of a Li I atom with 5 levels in our case, and with 20 levels in the case of Asplund et al. (this volume). Another difference is that the temperature inhomogeneities are somewhat enhanced in 2D models with respect of what occurs in 3D models, leading to correspondingly larger LTE granulation abundance corrections. But this is only a minor point. Certainly, the two groups agree that LTE Li abundance determinations relying on multidimensional hydrodynamical simulations of convection in metal-poor dwarfs are highly misleading.
## 6. Conclusions
1. The statement of Kurucz (1995) that abundances of lithium derived from standard 1D models of metal-poor stellar atmospheres is too small by a factor of 10 is not supported by actual multidimensional NLTE computations. Even the sign of the correction is doubtful, and the error is well below 0.1 dex, both according to our investigation and the one presented by M. Asplund on this symposium.
2. LTE abundance determinations based on inhomogeneous atmospheres are strongly discouraged. They produce large “granulation abundance corrections” due to non-linear effects in the direction opposite to Kurucz’ prediction, but the actual NLTE line formation mechanism couples the population of the atomic levels more closely to the mean radiation field than to the local temperature.
3. The combination of multidimensional models with NLTE line formation for the Li I $`\lambda `$ 670.8 nm resonance line leads to the same lithium abundance as that derived from NLTE analysis with flux-constant 1D models, abundance differences being less than 0.1 dex. However, this result must not be hastily generalized to other atoms with a different atomic structure.
4. An obvious future improvement is to extend the NLTE analysis to the continuum, which has be assumed here to be in LTE. We plan to do that in the near future. If it turns out that the H<sup>-</sup> ion is affected by NLTE, this will raise a new question: should such effects be included already in the radiation hydrodynamics code, which determines the amplitude of the thermodynamical fluctuations?
## References
Asplund, M., Nordlund, A, Trampedach, R., Stein, R.F. 1999, A&A, 346, L17
Carlsson, M., Rutten, R.J., Bruls, J.H.M.J., Shchukina, N.G. 1994, A&A, 288, 860
Freytag, B., Ludwig, H.-G., Steffen, M. 1996, A&A, 313, 497
Feautrier. P. 1964, C.R. Acad. Sci. Paris, 258, 3189
Gadun A.S., Pavlenko, Ya. V. 1997, A&A, 324, 281
Kiselman, D. 1997, ApJ, 489, L107
Kiselman, D. 1998, A&A, 333, 732
Kurucz, R.L. 1995, ApJ, 452, 102
Ludwig, H.-G., Jordan, S., Steffen, M. 1994, A&A, 284, 105
Mihalas, D. 1970 Stellar Atmospheres W. H. Freeman & Co.
Pinsonneault, M.H., Deliyannis, C. P., Demarque, P. 1992, ApJS, 78, 179
Stein, R.F., Nordlund, Å, 1998, ApJ, 499, 914
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# Statistical Model of Superconductivity in a 2D Binary Boson-Fermion Mixture
## 1 Introduction
Recent experiments indicate that composite bosons in ultra-cold clouds of most alkali atoms do indeed Bose-Einstein (BE) condense. Since Cooper pairs (CPs) of fermions (electrons or holes) in a many-fermion system form composite bosons in the sense of coupling to integer angular momentum, it is natural to consider the possible BE condensation of such pairs. The belief that some such condensate is central to superconductivity is more than 50 years old . High-$`T_c`$, as well as some organic, superconductors are quasi-two-dimensional (2D). Quasi-1D superconductors have also been found . BE condensation (BEC) is impossible in two or less space dimensions for usual or “ordinary” bosons (i.e., with a quadratic energy-momentum, or dispersion, relation). It is however still possible to have BEC in all dimensions $`d>1`$ for non-interacting bosons if they obey a linear dispersion relation —such as CPs moving in the Fermi sea. This possibility arises because the Hohenberg theorem , which prohibits BEC in 2D, relies on an $`f`$-sum rule based on the quadratic dispersion relation appropriate to bosons moving in a vacuum. Such a linear dispersion relation for the CPs in a binary boson-fermion mixture model was recently found to be consistent, without any adjustable parameters, with the anomalous linear (quadratic) temperature-dependence above $`T_c`$ in the resistivity of optimally-doped (overdoped) cuprates whether hole- or electron-doped. For the observed quadratic $`T`$-dependence in overdoped samples linear-dispersion CP charge carriers are essential.
Although extensive studies in the BCS-Bose “crossover” problem in superconductivity have spanned a period of over thirty years, we note that BEC is distinct from the standard (i.e., zero center-of-mass momentum CPs) BCS theory condensation where only that one bosonic state exists.
In this paper it is shown that in addition BEC is still possible in 2D even if the number of composite bosons (pairs of fermions) in a binary boson-fermion mixture is not fixed—as chemical/thermal equilibrium renders it coupling- and temperature-dependent—as long as the total number of fermions is fixed. This gives rise to an interesting statistical-mechanics problem irrespective of the particular mechanism for pair formation, and may have a vital application for superconductivity as well as for (neutral-atom) superfluidity such as in liquid <sup>3</sup>He , dilute mixtures of <sup>3</sup>He in <sup>4</sup>He , or in trapped Fermi gases . The statistical model dealt with here may be seamlessly linked to BCS theory, via the fermionic energy gap, when boson/unpaired-fermion interactions are included as, e.g., in Refs. and . However, in these two papers the quadratic CP dispersion relation has been assumed. The quadratic form has recently been shown to apply only in the zero-density or vacuum limit when the Fermi sea disappears.
In Section 2 we recall that at $`T=0`$ a 2D gas of fermions interacting via a constant pairing interaction in an annulus about the Fermi surface—viz., the BCS model interaction—the binding energy of a single pair near the Fermi surface (CP problem) decreases practically linearly with the center of mass momentum (CMM) of the pair for all values of the momentum below breakup, the breakup momentum typically being only about four orders of magnitude smaller than the Fermi momentum. In Section 3 we discuss why the interacting many-fermion system can be treated as a set of independent CPs (i.e., composite bosons with fermion number two) mixed int with pairable fermions which are not bound into pairs, i.e., unpaired fermions. In Section 4 the more realistic scenario is considered of the BEC of these pairs, incorporating pair breakup beyond a certain CMM. Although the number of pairs is not fixed but rather strongly coupling- and temperature-dependent, BEC is still possible in 2D. A simple binary boson- fermion statistical model is introduced by constructing the Helmholtz free energy for an ideal mixture of pairable but unpaired fermions plus paired fermions (both zero and nonzero CMM pairs), all in chemical and thermal equilibrium. The latter results through extrema of the free energy in: a) the pairable fermion occupation probabilities; b) the excited boson numbers (nonzero CMM CPs); and c) the ground boson number (zero CMM pairs). In Section 5 the coupling- and temperature-dependence of the boson number is derived. In Section 6 the critical BEC singularity temperature is obtained first by ignoring the unpaired fermions in a pure boson-gas model and then exactly for the boson-fermion binary mixture model from a $`T`$-dependent dispersion relation derived and calculated numerically, and results compared with empirical data. Finally, Section 7 gives conclusions.
## 2 Cooper-pair dispersion relation
Consider a 2D system of N fermions of mass $`m`$ confined in a square “pen” of area $`L^2`$ and interacting pairwise via the BCS model interaction
$$V_{𝐤,𝐤^{}}=\{\begin{array}{cc}V& \text{if}\mu (T)\mathrm{}\omega _D<ϵ_{k_1}(\mathrm{}^2k_1^2/2m),ϵ_{k_2}<\mu (T)+\mathrm{}\omega _D\hfill \\ 0& \text{otherwise},\hfill \end{array}$$
(1)
where $`𝐤\frac{1}{2}(𝐤_1𝐤_2)`$ is the relative wavevector of the two particles; $`V_{𝐤,𝐤^{}}`$ the 2D double Fourier integral of the underlying non-local interaction $`V(𝐫,𝐫^{})`$ in the relative coordinate $`𝐫=𝐫_1𝐫_2;\mu (T)`$ the ideal Fermi gas chemical potential which at $`T=0`$ becomes the Fermi energy $`E_F\mathrm{}^2k_F^2/2m`$ with $`k_F`$ the Fermi wavenumber; $`\mathrm{}\omega _D\mathrm{}^2k_D^2/2m`$ the width of the annulus about the Fermi circle in which the pairing interaction is nonzero, with $`\omega _D`$ being the Debye frequency. This model interaction mimics the net effect of an attractive electron-phonon interaction overwhelming the repulsive interfermion Coulomb repulsions, whenever $`V>0`$.
If $`\mathrm{}𝐊=\mathrm{}(𝐤_1+𝐤_2)`$ is the center-of-mass momentum (CMM) of a pair, let $`E_K`$ be its total energy (besides the CP rest-mass energy). The eigenvalue (CP ) equation for a pair of fermions at $`T=0`$ immersed in a background of $`N2`$ inert, spectator fermions within a (sharp) Fermi circular perimeter of radius $`k_F`$ is then
$$1=V\underset{𝐤}{}^{^{}}\frac{\theta (k_1k_F)\theta (k_2k_F)}{2ϵ_k(E_K\mathrm{}^2K{}_{}{}^{2}/4m)},$$
(2)
where again $`ϵ_k\mathrm{}^2k^2/2m`$, $`\theta (x)`$ is the Heaviside unit step function, and the prime on the summation sign denotes the conditions
$$k_1|𝐤+\frac{1}{2}𝐊|<(k_F^2+k_D^2)^{1/2}\text{and}k_2|𝐤\frac{1}{2}𝐊|<(k_F^2+k_D^2)^{1/2}$$
(3)
ensuring that the pair of fermions above the Fermi “surface” cease interacting beyond the annulus of energy thickness $`\mathrm{}\omega _D`$ in accordance with (1), thereby restricting the summation over $`𝐤`$ for a given fixed $`𝐊`$. Without these restrictions (2) would just be the Schrödinger equation in momentum space for the pair. Setting $`E_K2E_F\mathrm{\Delta }_K`$, the pair is bound if $`\mathrm{\Delta }_K>0`$, and (2) becomes an eigenvalue equation for the (positive) pair binding energy $`\mathrm{\Delta }_K`$. Our $`\mathrm{\Delta }_K`$ and $`\mathrm{\Delta }_0`$ should not be confused with the BCS energy gap $`\mathrm{\Delta }(T)`$.
Let $`\lambda g(E_F)V0`$ be a dimensionless coupling constant with $`g(E_F)`$ the electronic density-of-states (for each spin) at the Fermi surface in the normal (i.e., interactionless) state, which in 2D is constant
$$g(ϵ)=L^2m/2\pi \mathrm{}^2g.$$
(4)
The Cooper equation (2) for the unknown quantity $`\mathrm{\Delta }_K`$ is analyzed in Ref. . For zero CMM, $`K=0`$, it becomes a single elementary integral, with the familiar solution
$$\mathrm{\Delta }_0=\frac{2\mathrm{}\omega _D}{e^{2/\lambda }1}$$
(5)
valid for all coupling $`\lambda `$. For small $`K`$, it is not too difficult to extract the asymptotic result
$$\mathrm{\Delta }_K\underset{K0}{}\mathrm{\Delta }_0\frac{2}{\pi }\left[1+\frac{\mathrm{\Delta }_0}{2\mathrm{}\omega _D}(1+\sqrt{1+\nu })\right]\mathrm{}v_FK+O(K^2)\underset{\lambda 0}{}\mathrm{\Delta }_0\frac{2}{\pi }\mathrm{}v_FK+O(K^2)$$
(6)
where $`\nu \mathrm{\Theta }_D/T_F`$, and $`v_F`$ is the Fermi velocity defined through $`E_F\mathrm{}^2k_F^2/2m=\frac{1}{2}mv_F^2`$. For weak coupling, $`\lambda 0`$, this linear dispersion relation gives the 2D analog of the 3D result stated as far back as 1964 in Ref. , p. 33 (see also, Ref. , p. 336) but with the 2D coefficient $`2/\pi `$ of the last expression of (6) replaced by $`1/2`$.
## 3 Justification of boson formalism
These CP boson-like structures could be called “quasi-bosons” since their creation and annihilation operators are known not to obey the usual boson commutation relations , p. 38. However, they do obey the Bose-Einstein distribution since the energy $`E_K`$ of the CP is given only by the total CMM, $`K`$, but is independent of the relative momentum $`k`$. Thus, the possible energy states for the pair are $`E_K`$ as defined in (2). The number of pairs $`N_𝐊`$ that can occupy such a state can take on indefinite values since there exist also indefinitely many relative momenta, namely
$$N_𝐊\underset{𝐤}{}𝒩_{𝐤,𝐊}=0,1,2,\mathrm{}.$$
(7)
Here, $`𝒩_{𝐤,𝐊}=0,1`$ is the number of pairs characterized by both k and K, and is the same number as that characterized by definite $`𝐤_1`$ and $`𝐤_2`$, namely $`𝒩_{𝐤,𝐊}=n_{𝐤_1}n_{𝐤_2}=0,1`$ where $`n_{𝐤_i}=0,1`$ is the occupation number for a single fermion, these remarks all referring to singlet pairing. Much of all this has been known at least since 1958, albeit in somewhat different language.
This view of an actual Cooper pair should not be confused with, say, an Anderson phonon-like collective excitation (or modes) with weak-coupling dispersion relation—in 2D given by $`(1/\sqrt{2})\mathrm{}v_FK`$ in the long-wavelength limit, and which evolves into the plasmon when Coulomb repulsions between fermions are switched on. CPs here, like deuterons, carry fermion number two and as such are definite in number (although in the CP case this number is coupling- and temperature-dependent) and can thus undergo BEC. This is distinct from collective excitations which are indefinite in number. Park , e.g., distinguishes between “permanent” and “ephemeral” bosons, the latter sometimes being referred to as “quasiparticles” to distinguish from the former “particles”.
For $`N_B`$ ordinary bosons of mass $`m_B`$ and energy $`\epsilon _K=\mathrm{}^2K^2/2m_B`$ in any positive dimension, $`d>0`$, a temperature singularity $`T_c`$ appears in the number equation $`N_B=_𝐊[e^{(\epsilon _K\mu _B)/k_BT}1]^1`$ at vanishing bosonic chemical potential $`\mu _B\underset{}{<}0`$ when the number of $`𝐊=0`$ bosons just ceases to be negligible upon cooling. It is given by
$$T_c=\frac{2\pi \mathrm{}^2}{m_Bk_B}\left[\frac{n_B}{g_{d/2}(1)}\right]^{2/d}$$
(8)
with $`n_B`$ the boson particle density $`N_B/L^d`$, and $`g_{d/2}(z)`$ the usual Bose integrals
$$g_\sigma (z)\frac{1}{\mathrm{\Gamma }(\sigma )}_0^{\mathrm{}}𝑑x\frac{x^{\sigma 1}}{z^1e^x1}=\underset{l=1}{\overset{\mathrm{}}{}}\frac{z^l}{l^\sigma }\underset{z1}{}\zeta (\sigma ),$$
(9)
where $`\mathrm{\Gamma }(\sigma )`$ is the gamma function and $`\zeta (\sigma )`$ the Riemann zeta function of order $`\sigma `$. The last identification in (9) holds when $`\sigma >1`$ for which $`\zeta (\sigma )<\mathrm{}`$, while the series $`g_\sigma (1)`$ diverges for $`\sigma 1`$, thus giving $`T_c=0`$ for $`d2`$. For $`d=3`$ one has $`\zeta (3/2)2.612`$ so that (8) becomes the familiar formula $`T_c3.31\mathrm{}^2n_B^{2/3}/m_Bk_B`$ of “ordinary” BEC. On the other hand, for bosons with (positive) excitation energy $`\epsilon _K\mathrm{\Delta }_0\mathrm{\Delta }_K`$ given approximately by the linear term in (6) for all $`K`$, the singularity that lead to (8) now yields , for weak coupling,
$$T_c=\frac{a(d)\mathrm{}v_F}{k_B}\left[\frac{\pi ^{\frac{d+1}{2}}n_B}{\mathrm{\Gamma }(\frac{d+1}{2})g_d(1)}\right]^{1/d}$$
(10)
where $`a(d)=1,2/\pi `$ and $`1/2`$ for $`d=1,2`$ and $`3`$, respectively. Note that now $`T_c>0`$ for all $`d>1`$, which is precisely the dimensionality range of all known superconductors including the quasi-1D organo-metallic (Bechgaard) salts . This is not inconsistent with the Hohenberg theorem that there is no broken symmetry, i.e., long-range order, in a Bose fluid for $`d`$ = 1 or 2, since this is based on an $`f`$-sum rule for bosons with a quadratic dispersion relation. Indeed, both (8) and (10) are special cases of of the more general expression for any space dimensionality $`d>0`$ and any boson dispersion relation $`\epsilon _K=C_sK^s`$ with $`s`$ $`>0`$ and $`C_s`$ a constant, given by
$$T_c=\frac{C_s}{k_B}\left[\frac{s\mathrm{\Gamma }(d/2)(2\pi )^dn_B}{2\pi ^{d/2}\mathrm{\Gamma }(d/s)g_{d/s}(1)}\right]^{s/d}.$$
(11)
In what follows the number of bosons will be temperature-dependent and it is in conserving the fermion number that the singularity arises. As is the case for the pure boson gas, a linear rather than a quadratic dispersion relation will be needed to obtain BEC in 2D. This emerges in a statistical model of an ideal binary mixture of bosons (the CPs) and unpaired (both pairable and unpairable) fermions in chemical equilibrium , for which thermal pair-breaking into unpaired pairable fermions is explicitly allowed.
## 4 First-principles statistical model
Under interaction (1) at any $`T`$ the total number of fermions in 2D is $`N=L^2k_F^2/2\pi =N_1+N_2`$ and is just the number of non-interacting (i.e., unpairable) fermions $`N_1`$ plus the number of pairable ones $`N_2`$. The unpairable fermions obey the usual Fermi-Dirac distribution with fermionic chemical potential $`\mu `$. On the other hand, the $`N_2`$ pairable fermions are simply those in the interaction shell of energy width $`\mathrm{}\omega _D`$ so that
$$N_2=2_{\mu \mathrm{}\omega _D}^{\mu +\mathrm{}\omega _D}𝑑ϵ\frac{g(ϵ)}{e^{\beta (ϵ\mu )}+1}=2g\mathrm{}\omega _D,$$
(12)
since the density of electronic states (4) is constant and the remaining integral exact. At any interfermionic coupling and temperature these fermions form an ideal mixture of pairable but unpaired fermions plus CPs that are created near the single-fermion energy $`\mu (T)`$, with binding energy $`\mathrm{\Delta }_K(T)`$ $`0`$ and total energy
$$E_K(T)2\mu (T)\mathrm{\Delta }_K(T).$$
(13)
This is generalizes the $`T=0`$ equation $`E_K2E_F\mathrm{\Delta }_K`$ introduced below (3).
The Helmholtz free energy $`F=ETS`$, where $`E`$ is the internal energy and $`S`$ the entropy, for this binary “composite boson/pairable-but-unpaired-fermion system” at temperatures $`TT_c`$ is then
$`F_2`$ $`=`$ $`2{\displaystyle _{\mu \mathrm{}\omega _D}^{\mu +\mathrm{}\omega _D}}𝑑ϵg(ϵ)\left\{n_2(ϵ)ϵ+k_BT\left[n_2(ϵ)\mathrm{ln}n_2(ϵ)+\{1n_2(ϵ)\}\mathrm{ln}\{1n_2(ϵ)\}\right]\right\}`$ (14)
$`+[2\mu (T)\mathrm{\Delta }_0(T)]N_{B,0}(T)`$
$`+{\displaystyle \underset{K>0}{\overset{K_0}{}}}\{[2\mu (T)\mathrm{\Delta }_K(T)]N_{B,K}(T)`$
$`+k_BT[N_{B,K}(T)\mathrm{ln}N_{B,K}(T)\{1+N_{B,K}(T)\}\mathrm{ln}\{1+N_{B,K}(T)\}]\}.`$
The integral term is the contribution from the unpaired fermions and runs over all levels in the energy shell where the BCS model interaction is nonzero, $`n_2(ϵ)`$ being the average number of unpaired but pairable fermions with energy $`ϵ`$; the prefactor two comes from the spin. The second term gives the free energy of the bosons with CMM $`K=0`$ since their entropy is negligible in the thermodynamic limit; here $`N_{B,0}(T)`$ is the number of (bosonic) CPs with zero CMM at temperature $`T`$. The summation term represents the free energy of the bosons with nonzero CMM, while $`N_{B,K}(T)`$ is that with arbitrary nonzero CMM $`K,`$ and the cutoff $`K_0`$ is defined by $`\mathrm{\Delta }_{K_0}0`$. The free energy $`F_2`$ is to be minimized subject to the constraint that the total number of pairable fermions $`N_2`$ is conserved.
If $`N_{20}(T)`$ is the number of pairable but unpaired fermions, the relevant number equation for the pairable (i.e., active) fermions is then
$$N_2=N_{20}(T)+2[N_{B,0}(T)+N_{B,0<K<K_0}(T)]N_{20}(T)+2N_B(T),$$
(15)
where $`N_{B,0<K<K_0}(T)`$ denotes the total number of “excited” bosonic pairs (namely with CMM such that $`0<K<K_0`$), i.e., $`N_{B,0<K<K_0}(T)_{0<K<K_0}N_{B,K}(T)`$. Minimizing the free energy, subject to the constraint that (15) be a constant, is equivalent to minimizing the grand potential
$$\mathrm{\Omega }_2=F_2\mu _2N_2.$$
(16)
a) Minimizing $`\mathrm{\Omega }_2`$ with respect to the fermion occupation probabilities $`n_2(ϵ)`$ yields the Fermi-Dirac distribution with fermion chemical potential $`\mu _2`$, not $`\mu `$, namely
$$n_2(ϵ)=\frac{1}{e^{\beta (ϵ\mu _2)}+1};\beta (k_BT)^1.$$
(17)
Thus the total number of pairable (but unpaired) fermions then becomes
$$N_{20}(T)2_{\mu \mathrm{}\omega _D}^{\mu +\mathrm{}\omega _D}𝑑ϵg(ϵ)n_2(ϵ)=2_{\mu \mathrm{}\omega _D}^{\mu +\mathrm{}\omega _D}𝑑ϵ\frac{g(ϵ)}{e^{\beta (ϵ\mu _2)}+1},$$
(18)
and should be compared with (12) for $`N_2`$ which contains only $`\mu `$. Since in 2D $`g(ϵ)`$ is a constant (4), (18) becomes the exact expression
$$N_{20}(T)=\frac{2g}{\beta }\mathrm{ln}\left[\frac{1+e^{\beta (\mu \mu _2\mathrm{}\omega _D)}}{1+e^{\beta (\mu \mu _2+\mathrm{}\omega _D)}}\right].$$
(19)
b) Minimizing $`\mathrm{\Omega }_2`$ with respect to the excited boson numbers $`N_{B,K}(T)`$, $`K>0`$, yields the Bose-Einstein distribution summed over all $`0<K<K_0`$, namely
$$N_{B,0<K<K_0}(T)\underset{K>0}{\overset{K_0}{}}N_{B,K}(T)=\underset{K>0}{\overset{K_0}{}}[e^{\beta \{E_K(T)2\mu _2\}}1]^1.$$
(20)
The factor multiplying $`\beta `$ in (20) may be rewritten as $`\epsilon _K(T)\mu _B(T)`$, where $`\epsilon _K(T)\mathrm{\Delta }_0(T)\mathrm{\Delta }_K(T)0`$ is a (nonnegative) excitation energy as suggested by (6), while $`\mu _B(T)`$ turns out to be
$$\mu _B(T)=2\left[\mu _2(T)\mu (T)\right]+\mathrm{\Delta }_0(T).$$
(21)
This allows rewriting (20) in the more meaningful boson form
$$N_{B,0<K<K_0}(T)=\underset{K>0}{\overset{K_0}{}}[e^{\beta \{\epsilon _K(T)\mu _B(T)\}}1]^1$$
(22)
where $`\mu _B(T)`$ is clearly the bosonic chemical potential associated with the entire binary mixture.
c) Finally, minimizing $`\mathrm{\Omega }_2`$ with respect to the number of zero CMM (or, “ground state”) bosons $`N_{B,0}(T)`$ gives
$$2[\mu _2(T)\mu (T)]+\mathrm{\Delta }_0(T)=0(0TT_c),$$
(23)
valid only in the stated temperature range as $`N_{B,0}(T)`$ is negligible for all $`T>T_c`$. However, in view of (21) this implies that $`\mu _B(T)`$ $`=0`$ for all $`0TT_c`$—which is precisely the BEC condition for a pure boson gas, even though one now deals with a binary boson-fermion mixture.
## 5 Boson number
To determine $`N_B(T)`$ from (15) we need (19) which with (23) reduces to
$$N_{20}(T)=\frac{2g}{\beta }\mathrm{ln}\left[\frac{1+e^{\beta \{\mathrm{\Delta }_0(T)/2\mathrm{}\omega _D\}}}{1+e^{\beta \{\mathrm{\Delta }_0(T)/2+\mathrm{}\omega _D\}}}\right](0TT_c).$$
(24)
At $`T=0`$ two distinct coupling regimes emerge by inspecting (24): a) for $`\mathrm{\Delta }_0/2<\mathrm{}\omega _D`$ or, from (5) for $`\lambda 2/\mathrm{ln}22.89`$, we have that $`N_{20}(0)=2g(\mu )(\mathrm{}\omega _D\mathrm{\Delta }_0/2)`$; while b) for $`\mathrm{\Delta }_0/2>\mathrm{}\omega _D`$ (or $`\lambda 2.89`$) $`N_{20}(0)`$ is identically zero. Hence, the number of bosons $`N_B(0)`$ at $`T=0`$ from (15) is just $`N_B(0)=\frac{1}{2}[N_2N_{20}(0)]`$. Using (12) for $`N_2`$ the fractional number of pairable fermions that are actually paired at $`T=0`$, namely $`2N_B(0)/N_2=1N_{20}(0)/N_2`$, becomes simply
$$2N_B(0)/N_2=\{\begin{array}{cc}\mathrm{\Delta }_0/2\mathrm{}\omega _D=(e^{2/\lambda }1)^1\underset{\lambda 0}{}e^{2/\lambda }\hfill & (\mathrm{for}\lambda 2/\mathrm{ln}22.89)\hfill \\ 1\hfill & (\mathrm{for}\lambda 2/\mathrm{ln}22.89).\hfill \end{array}$$
This fraction is plotted against coupling $`\lambda `$ in Fig. 1. Since $`N_B(0)=\frac{1}{2}g\mathrm{\Delta }_0`$ for $`\lambda 2.89`$, only those fermions in an energy shell of width $`\mathrm{\Delta }_0/2`$ around the Fermi surface actually pair at $`T=0`$, while for $`\lambda 2.89`$ all pairable fermions actually pair up since then $`N_B(0)=g\mathrm{}\omega _D\frac{1}{2}N_2`$. This result contrasts sharply with the “heuristic model” , Eq. (22), where $`2N_B(0)/N_2`$ $`1`$ for all coupling, and is more in line with BCS theory which implies, in any $`d`$, a coupling-dependent fraction estimated (Ref. p. 128; see also ) to be $`[g(E_F)2\mathrm{\Delta }/2g(E_F)\mathrm{}\omega _D]^2=(\mathrm{\Delta }/\mathrm{}\omega _D)^2(\mathrm{sinh}1/\lambda )^2\underset{\lambda 0}{}4e^{2/\lambda }`$, where $`\mathrm{\Delta }\mathrm{}\omega _D/\mathrm{sinh}(1/\lambda )`$ (again, not to be confused with the CP binding energy $`\mathrm{\Delta }_0`$) is the $`T=0`$ BCS energy gap for the same BCS model interaction (1) used in this paper; this is graphed as the long-dashed curve in Fig. 1 and is seen to be much larger than (5) for fixed $`\lambda `$. The breakdown of BCS theory for BCS model interaction couplings larger than $`\lambda 1.13`$ is clear both because: a) the alluded fraction cannot exceed unity; and b) physically, if the fermionic energy gap $`\mathrm{\Delta }\mathrm{}\omega _D`$ no pairable fermions are available at all. This breakdown is indicated by the short-dashed curve in Fig. 1. (A strong-coupling many-body model differing from that of BCS theory but based on the BCS model interaction has been solved by Thouless ).
Also displayed in Fig. 1 are two finite-temperature results for $`2N_B(T)/N_2=1N_{20}(T)/N_2`$ which are obtainable from (24) for any $`T`$ provided one knows $`\mathrm{\Delta }_0(T)`$ for any $`T>0`$. For $`T>0`$, the $`\theta (k_1k_F)\theta (ϵ_{k_1}E_F)`$ in (2) becomes $`1n(\xi _{k_1})`$, where $`n(\xi _{k_1})(e^{\beta \xi _{k_1}}+1)^1`$ is the Fermi-Dirac distribution with $`\xi _{k_1}ϵ_{k_1}\mu (T)`$, with the ideal fermion gas chemical potential $`\mu (T)`$ in 2D being given exactly by
$$\mu (T)=\beta ^1\mathrm{ln}(e^{\beta E_F}1)\underset{T0}{}E_F.$$
(25)
Note that $`\mu (T)`$ decreases monotonically with temperature from its maximum value of $`E_F`$ but does not turn negative until $`T=T_F/\mathrm{ln}21.44T_F`$ so that the BCS model interaction (1), which requires $`\mu (T)`$ to be nonnegative, will not break down (i.e., become meaningless) over the entire range of temperatures relevant in this paper, see Fig. 4 below. Similar arguments hold for $`\theta (k_2k_F)`$. Since $`k_1=k_2`$ implies that $`\xi _{k_1}=\xi _{k_2}`$, (2) then leads to a simple generalization to finite-temperature of the $`K=0`$ CP equation, namely
$$1=\lambda _0^{\mathrm{}\omega _D}𝑑\xi (e^{\beta \xi }+1)^2[2\xi +\mathrm{\Delta }_0(T)]^1.$$
(26)
Its numerical solution for $`\mathrm{\Delta }_0(T)`$ is illustrated in Fig. 2 for $`\lambda gV=1/2`$ and $`\nu \mathrm{}\omega _D/E_F=0.05`$. Note that if one assumes a $`T`$ such that $`\mathrm{\Delta }_0(T)=0`$, the resulting integral in (26) diverges and the equation can only be satisfied for $`\lambda =0`$; thus, there is no temperature $`T`$ at which “depairing” will occur for any fixed $`\nu `$ and any nonzero $`\lambda `$.
## 6 Critical temperature
Neglecting the background unpaired fermions and modeling our system as a pure boson gas of CPs but with temperature-dependent number density $`n_B(T)`$, one converts the explicit $`T_c`$-formula (10) into an implicit one by allowing $`n_B`$ to be $`T`$-dependent. For $`d=2`$ (10) becomes, since $`g_2(1)\zeta (2)=\pi ^2/6`$,
$$T_c=\frac{4\sqrt{3}}{\pi ^{3/2}}\frac{\mathrm{}v_F}{k_B}\sqrt{n_B(T_c)}.$$
(27)
This requires $`n_B(T)N_B(T)/L^2`$ which in turn requires (24), along with $`\mathrm{\Delta }_0(T)`$ as determined from (26), and is given by the expression $`2N_B(T)/N_2=1N_{20}(T)/N_2`$. Solving (27) self-consistently with $`\lambda =1/2`$ gives the remarkably constant value $`T_c/T_F0.004`$ over the entire range of $`\nu \mathrm{}\omega _D/E_F`$ values $`0.030.07`$ typical of cuprate superconductors. On the other hand, the BCS formula $`T_c^{BCS}1.13\mathrm{\Theta }_De^{1/\lambda }`$ with $`\lambda =1/2`$ gives $`T_c/T_F`$ = 0.005, 0.008 and 0.011 for $`\nu `$ = 0.03, 0.05 and 0.07, respectively. Clearly, both sets of predictions are somewhat small compared with empirical cuprate values of $`T_c/T_F`$ that range from $`0.010.1`$.
To obtain the exact critical temperature without neglecting the background unpaired fermions, one needs the exact CP excitation energy dispersion relation $`\epsilon _K(T)\mathrm{\Delta }_0(T)\mathrm{\Delta }_K(T)`$ which is neither exactly linear in $`K`$ nor independent of $`T`$. To determine $`\mathrm{\Delta }_K(T)`$ we need a working equation that generalizes Ref. for $`T>0`$ via the new CP eigenvalue equation (26). Because of symmetry, see Fig. 3, one can restrict the angle $`\theta `$ to the interval $`(0,\pi /2)`$ where $`k_1k_2`$, i.e., to quadrant I. Recalling (13), in $`d`$-dimensions (2) becomes
$$1=V\left(\frac{L}{2\pi }\right)^d^{}𝑑𝐤\frac{\left[1n(\xi _{k_1})\right]\left[1n(\xi _{k_2})\right]}{\mathrm{}^2(k^2k_\mu ^2)/m+\mathrm{\Delta }_K(T)+\mathrm{}^2K^2/4m}.$$
(28)
Here $`k_\mu `$ is such that $`\mu \mathrm{}^2k_\mu ^2/2m`$ and becomes $`k_F`$ as $`T0`$, while $`k_D`$ is such that $`\mathrm{}\omega _D\mathrm{}^2k_D^2/2m`$. The prime on the integral sign now denotes the restrictions
$`k_2^2|𝐤\frac{1}{2}𝐊|^2`$ $`=`$ $`k^2kK\mathrm{cos}\theta +\frac{1}{4}K^2>k_\mu ^2,`$ (29)
$`k_1^2|𝐤+\frac{1}{2}𝐊|^2`$ $`=`$ $`k^2+kK\mathrm{cos}\theta +\frac{1}{4}K^2<k_\mu ^2+k_D^2.`$ (30)
In Fig. 3 the darkest shading corresponds to these (BCS model interaction) restrictions. The conditions (29) and (30) can be studied separately but must be satisfied simultaneously. If $`K<2\sqrt{k_\mu ^2k_D^2}`$, (29) and (30) are equivalent to
$$\begin{array}{ccccc}(k_\mu ^2\frac{1}{4}K^2\mathrm{sin}^2\theta )^{1/2}+\frac{1}{2}K\mathrm{cos}\theta & <& k& <& [(k_\mu ^2+k_D^2)\frac{1}{4}K^2\mathrm{sin}^2\theta ]^{1/2}\frac{1}{2}K\mathrm{cos}\theta .\end{array}$$
(31)
Note that for $`K>\sqrt{k_\mu ^2+k_D^2}\sqrt{k_\mu ^2k_D^2}`$ there exists a minimum value $`\theta _{\mathrm{min}}`$ of $`\theta `$ given by
$$\mathrm{cos}\theta _{\mathrm{min}}\frac{k_D^2}{K\sqrt{2(2k_\mu ^2+k_D^2)K^2}},$$
(32)
while $`\theta _{\mathrm{min}}`$ = 0 for $`K<\sqrt{k_\mu ^2+k_D^2}\sqrt{k_\mu ^2k_D^2}`$. We introduce the dimensionless variables
$$\kappa \frac{K}{2(k_F^2+k_D^2)^{1/2}}1,\xi \frac{k}{k_F},\stackrel{~}{\mathrm{\Delta }}_\kappa \frac{\mathrm{\Delta }_K}{E_F},\nu \frac{\mathrm{\Theta }_D}{T_F}\frac{k_D^2}{k_F^2},$$
(33)
with $`k_B\mathrm{\Theta }_D\mathrm{}\omega _D\mathrm{}^2k_D^2/2m`$ and$`k_BT_FE_F,`$ where $`k_B`$ is Boltzmann’s constant. Recall the $`d=2`$ constant expression (4) for $`g(ϵ)`$, the restrictions (31), and that for $`K0`$ and $`T>0`$ the step functions in (2) $`\theta (k_{1,2}k_F)\theta (|\frac{1}{2}𝐊\pm 𝐤|k_F)`$ become $`[\mathrm{exp}\{\beta [\mathrm{}^2(\frac{1}{2}𝐊\pm 𝐤)^2/2m\mu (T)]\}+1]^1`$—but with $`2ϵ_k`$ in (2) replaced by $`ϵ_{k_1}+ϵ_{k_2}`$, $`E_F`$ by $`\mu (T)`$ and $`\mathrm{\Delta }_K`$ by $`\mathrm{\Delta }_K(T)`$. One finally arrives at a working equation for the binding energy $`\mathrm{\Delta }_K(T)`$ that generalizes Eq. (18) of Ref. , namely
$`1`$ $`=`$ $`{\displaystyle \frac{4}{\pi }}\lambda {\displaystyle _{\theta _{min}}^{\pi /2}}𝑑\theta {\displaystyle _{\xi _{min}(\theta )}^{\xi _{max}(\theta )}}𝑑\xi \xi {\displaystyle \frac{[1+\mathrm{exp}\{\stackrel{~}{\beta }[\xi ^2+(1+\nu )\kappa ^2+2\sqrt{1+\nu }\kappa \xi \mathrm{cos}\theta 1]\}]^1}{2\xi ^2+2(1+\nu )\kappa ^22+\stackrel{~}{\mathrm{\Delta }}_\kappa (\stackrel{~}{T})}}`$ (34)
$`\times [1+\mathrm{exp}\{\stackrel{~}{\beta }[\xi ^2+(1+\nu )\kappa ^22\sqrt{1+\nu }\kappa \xi \mathrm{cos}\theta 1]\}]^1,`$
where $`\nu \mathrm{}\omega _D/\mu `$, $`\xi _{min}(\theta )\sqrt{1+\nu }\kappa \mathrm{cos}\theta +\sqrt{1(1+\nu )\kappa ^2\mathrm{sin}^2\theta }`$, $`\xi _{max}(\theta )`$
$`\sqrt{1+\nu }\kappa \mathrm{cos}\theta +\sqrt{(1+\nu )(1\kappa ^2\mathrm{sin}^2\theta )}`$ and
$$\theta _{min}=\{\begin{array}{cc}& 0\text{if}2\kappa <1\sqrt{(1\nu )/(1+\nu )},\hfill \\ & \mathrm{cos}^1(\nu /\{4\sqrt{1+\nu }\kappa \sqrt{1+\nu /2(1+\nu )\kappa ^2}\})\mathrm{otherwise}.\hfill \end{array}$$
In (34) we have introduced the more general dimensionless quantities $`\xi k/k_\mu `$, $`\stackrel{~}{\mathrm{\Delta }}_\kappa (\stackrel{~}{T})\mathrm{\Delta }_K(T)/\mu `$, where $`\stackrel{~}{T}k_BT/\mu `$ or $`\stackrel{~}{\beta }\mu \beta `$, and $`\kappa K/2\sqrt{k_\mu ^2+k_D^2}`$.
To obtain the critical temperature from the finite-temperature dispersion relation, besides solving (28) for $`\mathrm{\Delta }_K(T)`$, one needs (12), (15), (22) and (25). At $`T=T_c`$ both $`N_{B,0}(T_c)0`$ and $`\mu _B(T_c)0`$ so that one gets the implicit $`T_c`$-equation for the binary mixture gas
$$1=\frac{\stackrel{~}{T}_c}{\nu }\mathrm{ln}\left[\frac{1+e^{\{\stackrel{~}{\mathrm{\Delta }}_0(\stackrel{~}{T}_c)/2\nu \}/\stackrel{~}{T}_c}}{1+e^{\{\stackrel{~}{\mathrm{\Delta }}_0(\stackrel{~}{T}_c)/2+\nu \}/\stackrel{~}{T}_c}}\right]+\frac{8(1+\nu )}{\nu }_0^{\kappa _0(\stackrel{~}{T}_c)}𝑑\kappa \frac{\kappa }{e^{[\stackrel{~}{\mathrm{\Delta }}_0(\stackrel{~}{T}_c)\stackrel{~}{\mathrm{\Delta }}_\kappa (\stackrel{~}{T}_c)]/\stackrel{~}{T}_c}1}.$$
(35)
This must be solved numerically for the exact $`T_c`$ for each $`\lambda `$ and $`\nu `$ in conjunction with (26) for $`\stackrel{~}{\mathrm{\Delta }}_0(\stackrel{~}{T})`$ and (34) for both $`\stackrel{~}{\mathrm{\Delta }}_\kappa (\stackrel{~}{T})`$ and $`\kappa _0(\stackrel{~}{T}_c)`$. Results for $`\lambda =1/2`$ are shown in Table 1 and Fig. 4 for a range of $`\nu `$ values typical of cuprates. We have taken $`T_\mu /T_F1`$, a very good approximation up to the highest temperatures dealt with. For example, from Fig. 4 the highest $`T_c/T_F0.14`$ already gives $`T_\mu /T_F0.9999`$ from (25), while for smaller $`T_c/T_F`$ the values of $`T_\mu /T_F`$ are even closer to 1. The $`T_c`$ resulting from the exact dispersion relation for $`T=0`$ (dot-dashed curve) is somewhat higher than the exact result (full curve) but lower than that using the linear approximation for $`\mathrm{\Delta }_K(T)`$ (dotted curve). It is also clear that the effect of using the exact or linear (in $`K`$) cases dominates the effect of the dispersion relation $`T`$-dependence. For cuprates $`d2.03`$ has been suggested to be more realistic as it reflects inter-CuO-layer couplings but our results in that case would be very similar to those reported here for $`d=2`$.
Thus, for $`\nu =0.05`$ the exact $`T_c`$ is seen to be about 46% lower than the heuristic result found in Ref. , Eqs. (15) and (23). It is curious that all results depend very weakly on the $`T`$-dependence of the CP binding energy $`\mathrm{\Delta }_K(T)`$, in spite of its being substantial throughout the temperatures spanned in this paper, as seen in Fig. 2.
We defer study of the condensate fraction $`N_{B,0}(T)/N_B(T)`$ below $`T_c`$ and merely surmise that it may ultimately help explain the apparent absence in cuprates of the Hebel-Slichter peak of nuclear-spin (NMR) relaxation rates vs temperature for $`0TT_c`$. Such a peak, originally seen in aluminum, is perhaps the most stringent and qualitatively convincing experimental test of BCS theory (Refs. , p. 71 and , p. 79 ff). Besides cuprates, it is also absent in several quasi-1D Bechgaard and in several quasi-2D (ET) organic salt superconductors.
## 7 Conclusions
A simple statistical model treating CPs as non-interacting bosons in thermal and chemical equilibrium with unpaired fermions is proposed. The model gives rise to a boson number that is strongly coupling- and temperature-dependent. Since the CP dispersion relation is approximately linear, it exhibits a Bose-Einstein condensation of zero-CMM pairs at precisely two dimensions. Exact transition temperatures based upon the exact CP dispersion relation are in reasonable agreement with empirical cuprate data.
Needless to say, further corrections are yet to be included in the present simple binary mixture boson-fermion model, e.g., i) realistic Fermi surfaces, ii) Van Hove singularities or other means of accounting for periodic-crystalline effects, as well as iii) the all important $`d`$-wave interfermionic interaction, iv) the boson-fermion interaction and v) residual interbosonic interactions. As to the latter, also generally neglected in BCS theory, if the lowering of $`T_c`$ in liquid <sup>4</sup>He by about 29% with respect to the ideal Bose gas BEC $`T_c`$ is any guide, interbosonic interactions will also lower $`T_c`$ in a more realistic picture. As to the boson-fermion interaction, it is precisely this ingredient that enabled T.D. Lee and coworkers , and Tolmachev more generally, to link BCS and BEC through a relation stating that the BE condensate fraction is proportional to the (BCS-like) fermionic gap $`\mathrm{\Delta }(T)`$ squared.
## Acknowledgments
We thank A. Salazar for help with Figure 1. M.C., A.P. and A.R. are grateful for partial support from grant PB98-0124, and M.deLl. from grants PB92-1083 and SAB95-0312, both by DGICYT (Spain), and PAPIIT (Mexico) IN102198-9 as well as from CONACyT (Mexico) 27828-E. R.M.Q. and N.J.D. acknowledge the support of the FRD (South Africa). M.deLl. thanks D.M. Eagles, M. Fortes, O. Rojo, and A.A. Valladares for numerous discussions; V.V. Tolmachev for extensive correspondence; R. Escudero for calling his attention to Ref. ; and A.N. Kocharian for reading and commenting the manuscript.
Figure Captions
Figure 1. Fractional number of pairable fermions that are actually paired, at three different temperatures, vs. coupling $`\lambda `$ for the present first-principles model (5) and estimated for BCS theory at $`T=0`$ as explained below (5). The number of pairable fermions with the BCS model interaction used is just (12); all of them are actually paired at $`T=0`$ in the heuristic BEC model, Ref. Eq. (23).
Figure 2. Temperature dependence of $`K=0`$ CP binding energy $`\mathrm{\Delta }_0(T)`$ obtained numerically from (26) for $`\lambda =1/2`$ and $`\nu =0.05`$. Note that when $`T=\mathrm{}`$ (26) is analytical for $`\mathrm{\Delta }_0(\mathrm{})`$; the latter then turns out to be about $`10^8`$, so that the curve saturates from above to this value at $`T=\mathrm{}`$.
Figure 3. Cross-section of overlap “volume” in momentum space (darkest shading) where the tip of the relative wavevector $`𝐤`$ (for two fermions with wavevectors $`𝐤_\mathrm{𝟏}`$ and $`𝐤_\mathrm{𝟐}`$) must point for the attractive BCS model interaction (1) between them to be nonzero and form a Cooper pair of CMM magnitude $`\mathrm{}K`$.
Figure 4. Critical BEC temperature $`T_c`$ in units of $`T_F`$, resulting from (35) for $`\lambda =1/2`$ for varying $`\nu \mathrm{}\omega _D/\mu \mathrm{\Theta }_D/T_F`$: with no approximations (full curve); using $`\mathrm{\Delta }_K(T)`$ evaluated at $`T=0`$ (dot-dashed); using the linear-in-$`K`$ approximation for $`\mathrm{\Delta }_K(T)`$ (dotted). The dashed straight line is the BCS formula $`T_c1.13\mathrm{\Theta }_De^{1/2\lambda }`$ for $`\lambda =1/2`$. The very lowest full horizontal line is the solution of the implicit $`T_c`$-equation (27) for the pure boson gas for $`\nu =0.03,0.05`$ and $`0.07`$. Cuprate data are taken from Ref. .
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# 1 Introduction
## 1 Introduction
Non-trivial Kaluza-Klein sphere reductions of supergravity theories have been studied in a number of contexts. Long ago it was demonstrated that the linearised analysis of the zero-mode fluctuations of the $`S^7`$ reduction of $`D=11`$ supergravity can be extended to a fully non-linear and consistent embedding of maximal $`D=4`$ gauged $`SO(8)`$ supergravity in $`D=11`$. The fact that the truncation to the zero-mode sector is consistent, despite the non-linearities of the theory, is rather non-trivial since there appears to be no group-theoretic understanding of why it should work.
More recently, a similar, although more explicit, demonstration of the consistency of the $`S^4`$ reduction from $`D=11`$ to the maximal $`D=7`$ gauged $`SO(5)`$ supergravity was given . The consistent reduction of massive type IIA supergravity on a locally $`S^4`$ space, to give the $`N=2`$ gauged $`SU(2)`$ supergravity in $`D=6`$, has also been obtained . In addition, certain reductions to truncations of the maximal gauged supergravities in various dimensions have also been constructed. These have the advantage of being considerably simpler than the maximal theories, allowing the reduction Ansatz to be presented in a more explicit form. Cases that have been constructed include the $`N=2`$ gauged $`SU(2)`$ supergravity in $`D=7`$ ; the $`N=4`$ gauged $`SU(2)\times U(1)`$ supergravity in $`D=5`$ ; the $`N=4`$ gauged $`SO(4)`$ supergravity in $`D=4`$ ; and maximal abelian truncations in $`D=4`$, 5 and 7 . One can also consider non-supersymmetric truncations of the maximal gauged supergravities. In the consistent truncations of the $`D=7`$, $`D=5`$ and $`D=4`$ supergravities to the graviton plus scalar subsectors comprising only the diagonal scalars in the $`SL(N,\text{I}\mathrm{R})/SO(N)`$ scalar submanifolds were considered (with $`N=5`$, 6 and 8 respectively), and the consistent embeddings in $`D=11`$ and $`D=10`$ were constructed.
Although the full consistency of the $`S^7`$ and $`S^4`$ reductions of $`D=11`$ supergravity has essentially been demonstrated, no similar complete result exists for the $`S^5`$ reduction of type IIB supergravity. The field content of the full $`N=8`$ gauged supergravity consists of gravity; fifteen $`SO(6)`$ gauge fields; twelve 2-form gauge potentials in the $`6`$ and $`\overline{6}`$ representations of $`SO(6)`$; 42 scalars in the $`1+1+20^{}+10+\overline{10}`$ representations of $`SO(6)`$, and the fermionic superpartners. It is believed that this can arise from an $`S^5`$ reduction of type IIB supergravity; at the linearised level, the reduction Ansatz was given in . However, at the full non-linear level, the only complete demonstrations so far are for the consistent embedding of the maximal abelian $`U(1)^3`$ truncation , the $`N=4`$ gauged $`SU(2)\times U(1)`$ truncation , and the scalar truncation in . The full metric Ansatz was conjectured in .
In this paper, we obtain the consistent reduction Ansatz for a different truncation of the maximal gauged $`D=5`$ supergravity. One can consistently set the $`1+1+10+\overline{10}`$ of scalars to zero, at the same time as setting the $`6+\overline{6}`$ of 2-form potentials to zero. The bosons that remain, namely gravity, the $`15`$ Yang-Mills fields, and the $`20^{}`$ of scalars, come just from the metric plus the self-dual 5-form sector of the original type IIB theory. We shall obtain complete results for the consistent embedding of this subsector of the gauged $`N=8`$ theory, with all fifteen of the $`SO(6)`$ gauge fields $`A_{\left(1\right)}^{ij}`$ , and the twenty scalars that parameterise the full $`SL(6,\text{I}\mathrm{R})/SO(6)`$ submanifold of the complete scalar coset. These can be parameterised by a unimodular symmetric tensor $`T_{ij}`$.
Another way of expressing the truncation that we shall consider in this paper is as follows. The type IIB theory itself can be consistently truncated in $`D=10`$ so that just gravity and the self-dual 5-form remain. The fifteen Yang-Mills gauge fields and twenty scalars that we retain in our Ansatz are the full set of massless fields associated with Kaluza-Klein reduction from this ten-dimensional starting point. (The counting of massless fields is the same as the one that arises from a toroidal reduction from the same ten-dimensional starting point.) We shall see below that the self-duality condition on the 5-form plays an essential rôle in the consistency of the $`S^5`$ reduction.
This subsector of type IIB supergravity is particularly relevant for the AdS/CFT correspondence , because it is the metric and the self-dual 5-form that couple to the D3-brane.
We also address the more general question of the circumstances under which a theory admits a consistent sphere reduction. We show that in a $`D`$-dimensional theory of gravity coupled to an $`n`$-form field strength, a consistent $`S^n`$ reduction that retains the full set of $`SO(n+1)`$ Yang-Mills fields together with coupled massless scalars is possible only when $`D=11`$ and $`n=4`$ or $`7`$, or in $`D=10`$ with $`n=5`$. Furthermore, in $`D=11`$ the theory has to be the bosonic sector of eleven-dimensional supergravity, with the $`FFA`$ term, while in $`D=10`$ the 5-form must be self-dual. In all three cases the full set of massless scalars includes a subset $`T_{ij}`$ described by the coset $`SL(n+1,\text{I}\mathrm{R})/SO(n+1)`$. (Such a coset structure is absent for any other values of $`(D,n)`$, and so it would not be appropriate to look for a consistent reduction Ansatz with scalars $`T_{ij}`$ of $`SL(n+1,\text{I}\mathrm{R})/SO(n+1)`$ for generic $`(D,n)`$.) Of the three cases where such a consistent sphere reduction is possible, the ten-dimensional one is singled out as the only case where the consistent reduction includes only gravity plus the $`SO(n+1)`$ gauge fields $`A_{\left(1\right)}^{ij}`$ and the scalars $`T_{ij}`$. By contrast, for $`D=11`$ reduced on $`S^4`$ one must additionally retain a set of five 2-form potentials, while for $`D=11`$ reduced on $`S^7`$ one must instead additionally retain 35 pseudoscalars as well as the 35 scalars $`T_{ij}`$, in order to achieve consistent reductions.<sup>1</sup><sup>1</sup>1This is related to the fact that one can consistently truncate five-dimensional maximal gauged supergravity to the $`SO(6)`$ gauge fields plus scalars of the coset $`SL(6,\text{I}\mathrm{R})/SO(6)`$. By contrast, the analogous truncations cannot be performed in the maximal gauged supergravities in $`D=7`$ and $`D=4`$.
## 2 The $`SO(6)`$ reduction Ansatz on $`S^5`$
We parameterise the fields for this truncated theory as follows. The twenty scalars, which are in the $`20^{}`$ representation of $`SO(6)`$, are represented by the symmetric unimodular tensor $`T_{ij}`$, where $`i`$ is a 6 of $`SO(6)`$. The fifteen $`SO(6)`$ Yang-Mills gauge fields will be represented by the 1-form potentials $`A_{\left(1\right)}^{ij}`$, antisymmetric in $`i`$ and $`j`$. The inverse of the scalar matrix $`T_{ij}`$ is denoted by $`T_{ij}^1`$. In terms of these quantities, we find that the Kaluza-Klein reduction Ansatz is given by
$`d\widehat{s}_{10}^2`$ $`=`$ $`\mathrm{\Delta }^{1/2}ds_5^2+g^2\mathrm{\Delta }^{1/2}T_{ij}^1D\mu ^iD\mu ^j,`$ (1)
$`\widehat{H}_{\left(5\right)}`$ $`=`$ $`\widehat{G}_{\left(5\right)}+\widehat{}\widehat{G}_{\left(5\right)},`$ (2)
$`\widehat{G}_{\left(5\right)}`$ $`=`$ $`gUϵ_{\left(5\right)}+g^1(T_{ij}^1DT_{jk})(\mu ^kD\mu ^i)`$ (3)
$`\frac{1}{2}g^2T_{ik}^1T_j\mathrm{}^1F_{\left(2\right)}^{}{}_{}{}^{ij}D\mu ^kD\mu ^{\mathrm{}},`$
$`\widehat{}\widehat{G}_{\left(5\right)}`$ $`=`$ $`{\displaystyle \frac{1}{5!}}ϵ_{i_1\mathrm{}i_6}[g^4U\mathrm{\Delta }^2D\mu ^{i_1}\mathrm{}D\mu ^{i_5}\mu ^{i_6}`$ (4)
$`5g^4\mathrm{\Delta }^2D\mu ^{i_1}\mathrm{}D\mu ^{i_4}DT_{i_5j}T_{i_6k}\mu ^j\mu ^k`$
$`10g^3\mathrm{\Delta }^1F_{\left(2\right)}^{i_1i_2}D\mu ^{i_3}D\mu ^{i_4}D\mu ^{i_5}T_{i_6j}\mu ^j],`$
where
$`U2T_{ij}T_{jk}\mu ^i\mu ^k\mathrm{\Delta }T_{ii},\mathrm{\Delta }T_{ij}\mu ^i\mu ^j,`$
$`F_{\left(2\right)}^{ij}=dA_{\left(1\right)}^{ij}+gA_{\left(1\right)}^{ik}A_{\left(1\right)}^{kj},DT_{ij}dT_{ij}+gA_{\left(1\right)}^{ik}T_{kj}+gA_{\left(1\right)}^{jk}T_{ik},`$
$`\mu ^i\mu ^i=1,D\mu ^id\mu ^i+gA_{\left(1\right)}^{ij}\mu ^j,`$ (5)
and $`ϵ_{\left(5\right)}`$ is the volume form on the five-dimensional spacetime. Note that $`\widehat{}\widehat{G}_{\left(5\right)}`$ is derivable from the given expressions (1) and (3); we have presented it here because it is quite an involved computation. The coordinates $`\mu ^i`$, subject to the constraint $`\mu ^i\mu ^i=1`$, parameterise points in the internal 5-sphere. In obtaining the above Ansatz we have been guided by previous results in the literature, including the $`S^4`$ reduction Ansatz from $`D=11`$ that was constructed in .
It is consistent to truncate the fields of type IIB supergravity to the metric and self-dual 5-form $`\widehat{H}_{\left(5\right)}`$. The ten-dimensional equations for motion for these fields are then given by
$`\widehat{R}_{MN}`$ $`=`$ $`{\displaystyle \frac{1}{96}}\widehat{H}_{MPQRS}\widehat{H}_N{}_{}{}^{PQRS},`$
$`d\widehat{H}_{\left(5\right)}`$ $`=`$ $`0.`$ (6)
The Ansatz presented above satisfies these equations of motion if and only if the five-dimensional fields satisfy the equations
$`D(T_{ik}^1T_j\mathrm{}^1F_{\left(2\right)}^k\mathrm{})`$ $`=`$ $`2gT_{k[i}^1DT_{j]k}\frac{1}{8}ϵ_{ijk_1\mathrm{}k_4}F_{\left(2\right)}^{k_1k_2}F_{\left(2\right)}^{k_3k_4},`$
$`D(T_{ik}^1DT_{kj})`$ $`=`$ $`2g^2(2T_{ik}T_{jk}T_{ij}T_{kk})ϵ_{\left(5\right)}+T_{ik}^1T_\mathrm{}m^1F_{\left(2\right)}^\mathrm{}kF_{\left(2\right)}^{mj}`$
$`\frac{1}{6}\delta _{ij}\left[2g^2(2T_k\mathrm{}T_k\mathrm{}(T_{kk})^2)ϵ_{\left(5\right)}+T_{pk}^1T_\mathrm{}m^1F_{\left(2\right)}^\mathrm{}kF_{\left(2\right)}^{mp}\right],`$
together with the five-dimensional Einstein equation. These equations of motion can all be derived from the five-dimensional Lagrangian
$`_5`$ $`=`$ $`R\text{1}\mathrm{l}\frac{1}{4}T_{ij}^1DT_{jk}T_k\mathrm{}^1DT_\mathrm{}i\frac{1}{4}T_{ik}^1T_j\mathrm{}^1F_{\left(2\right)}^{ij}F_{\left(2\right)}^k\mathrm{}V\text{1}\mathrm{l}`$
$`\frac{1}{48}ϵ_{i_1\mathrm{}i_6}\left(F_{\left(2\right)}^{i_1i_2}F_{\left(2\right)}^{i_3i_4}A_{\left(1\right)}^{i_5i_6}gF_{\left(2\right)}^{i_1i_2}A_{\left(1\right)}^{i_3i_4}A_{\left(1\right)}^{i_5j}A_{\left(1\right)}^{ji_6}+\frac{2}{5}g^2A_{\left(1\right)}^{i_1i_2}A_{\left(1\right)}^{i_3j}A_{\left(1\right)}^{ji_4}A_{\left(1\right)}^{i_5k}A_{\left(1\right)}^{ki_6}\right),`$
where the potential $`V`$ is given by
$$V=\frac{1}{2}g^2\left(2T_{ij}T_{ij}(T_{ii})^2\right).$$
(9)
In (2) we have omitted the wedge symbols in the final topological term, to economise on space. The Lagrangian is in agreement with the one for five-dimensional gauged $`SO(6)`$ supergravity in .
To see this, we look first at the ten-dimensional equation $`d\widehat{H}_{\left(5\right)}=0`$, with the Ansatz given above. The terms involving structures of the form $`X_{\left(4\right)}^{ij}D\mu ^iD\mu ^j`$, where $`X_{\left(4\right)}^{ij}`$ represents a 4-form in the five-dimensional spacetime, give rise to the Yang-Mills equations above. Contributions of this kind come from $`d\widehat{G}_{\left(5\right)}`$ and also from the final term in $`d\widehat{}\widehat{G}_{\left(5\right)}`$. The terms involving structures of the form $`X_{\left(5\right)}^{ij}(\mu ^iD\mu ^j)`$, where $`X_{\left(5\right)}^{ij}`$ represents a 5-form in the five-dimensional spacetime, give rise to the scalar equations of motion in (2). Since $`\mu ^iD\mu ^i=\mu ^id\mu ^i=\frac{1}{2}d(\mu ^i\mu ^i)=0`$, there is a trace subtraction in the scalar equation, and we read off $`X_{\left(5\right)}^{ij}\frac{1}{6}\delta _{ij}X_{\left(5\right)}^{kk}=0`$ as the five-dimensional equation of motion. Contributions of this structure come only from $`d\widehat{G}_{\left(5\right)}`$. Finally, all other structures arising from calculating $`d\widehat{H}_5=0`$ vanish identically, without the use of any five-dimensional equations of motion. In deriving these results one needs to make extensive use of the Schoutens over-antisymmetrisation identity, $`ϵ_{[i_1\mathrm{}i_6}V_{i_7]}=0`$.
It is worth remarking that it is essential for the consistency of the reduction Ansatz that the 5-form $`\widehat{H}_{\left(5\right)}`$ should be self-dual. One cannot simply consider a reduction Ansatz for a ten-dimensional theory consisting of gravity plus a non-self-dual 5-form, whose Ansatz is given by $`\widehat{G}_{\left(5\right)}`$ in (3). Although the Bianchi identity $`d\widehat{G}_{\left(5\right)}=0`$ would give perfectly acceptable equations of motion for $`F_{\left(2\right)}^{ij}`$ and $`T_{ij}`$, the field equation $`d\widehat{}\widehat{G}_{\left(5\right)}=0`$ would produce the (unacceptable) constraint
$$ϵ_{ijk_1\mathrm{}k_4}F_{\left(2\right)}^{k_1k_2}F_{\left(2\right)}^{k_3k_4}=0.$$
(10)
It is only by combining $`\widehat{G}_{\left(5\right)}`$ and $`\widehat{}\widehat{G}_{\left(5\right)}`$ together into the self-dual field $`\widehat{H}_{\left(5\right)}`$ that a consistent five-dimensional result is obtained, with (10) now combining with terms from $`d\widehat{G}_{\left(5\right)}`$ to form part of the five-dimensional Yang-Mills equations given in (2).<sup>2</sup><sup>2</sup>2If one were to consider the reduction of a ten-dimensional theory with a non-self-dual 5-form there would be additional fields present in a complete massless truncation. These would comprise $`10=4+6`$ vector potentials and $`5=1+4`$ scalars. However the inclusion of these fields would still not achieve a consistent reduction Ansatz, since the current (10), which is in the $`15`$ of $`SO(6)`$, could still not acquire an interpretation as a source term for the additional fields. Thus the additional requirement of self-duality seems to be essential for consistency. (Of course anti-self-duality would be equally good.) It is interesting, therefore, that self-duality of the 5-form is apparently forced on us by the requirements of Kaluza-Klein consistency, if we try to “invent” a ten-dimensional theory that can be reduced on $`S^5`$. Thus once again we see that supersymmetry and Kaluza-Klein consistency for sphere reductions seem to go hand in hand.<sup>3</sup><sup>3</sup>3Of course the initial truncation of the type IIB theory to its gravity plus self-dual 5-form sector is itself a non-supersymmetric one, but the crucial point is that the consistency of the Kaluza-Klein $`S^5`$ reduction is singling out a starting point that is itself a subsector of a supersymmetric theory.
The Kaluza-Klein $`S^5`$ reduction that we have obtained here retains the full set of massless fields that can result from the reduction of gravity plus a self-dual 5-form in $`D=10`$. In other words, after the initial truncation of the type IIB theory in $`D=10`$, no further truncation of massless fields has been performed.
It should also be emphasised that it would be inconsistent to omit the fifteen $`SO(6)`$ gauge fields when considering the embedding of the twenty scalars $`T_{ij}`$. This can be seen from the Yang-Mills equations in (2), which have a source term $`gT_{k[i}^1DT_{j]k}`$ appearing on the right-hand side. This is a quite different situation from a toroidal reduction, where it is always consistent to truncate to the scalar sector, setting the gauge fields to zero. The new feature here in the sphere reduction is that the scalar fields are charged under the gauge group. This is a general feature of all the sphere reductions, to $`D=4`$, $`D=5`$ and $`D=7`$, and thus in all cases it is inconsistent to include the full set of scalar fields without including the gauge fields as well (see also ). (One can consistently truncate to the diagonal scalars in $`T_{ij}`$, setting all the gauge fields to zero, as in , since then $`T_{k[i}^1DT_{j]k}=0`$.)
Our testing of the consistency of the reduction Ansatz (1)-(4) has so far been restricted to checking the equations of motion for $`\widehat{H}_{\left(5\right)}`$. A full testing of the consistency of the ten-dimensional Einstein equations would be quite involved, and will be addressed in future work.<sup>4</sup><sup>4</sup>4The experience in all previous work on consistent reductions is that although the actual checking of the higher-dimensional Einstein equations is the most difficult part from a computational point of view, the consistency seems to be assured once it has been achieved for the equations of motion for the antisymmetric tensor fields, which itself is an extremely stringent requirement. In the next section we shall show that the reduction Ansatz that we have obtained in this paper reduces, with appropriate additional truncations, to results that have been obtained previously. Since the complete consistency was proven in these earlier results, including the ten-dimensional Einstein equations, this provides further supporting evidence for the complete consistency of the Ansatz that we have constructed here.
Finally in this section, we remark that the $`S^5`$ reduction that we have constructed here is also consistent if we include the dilaton $`\widehat{\varphi }`$ and axion $`\widehat{\chi }`$ of the type IIB theory. These simply reduce according to the Ansätze
$$\widehat{\varphi }=\varphi ,\widehat{\chi }=\chi ,$$
(11)
where the unhatted quantities denote the fields in five dimensions. In this reduction they do not appear in the previous Ansätze for the metric and self-dual 5-form, and in $`D=5`$ they just give rise to the additional $`SL(2,\text{I}\mathrm{R})`$-invariant Lagrangian
$$_{(\varphi ,\chi )}=\frac{1}{2}d\varphi d\varphi \frac{1}{2}e^{2\varphi }d\chi d\chi ,$$
(12)
which is added to (2). Note in particular that $`\varphi `$ and $`\chi `$ do not appear in the five-dimensional scalar potential $`V`$.
## 3 Truncations to previous results
We can consider three different truncations of the reduction scheme of the previous section, in order to make contact with previous results in the literature. The first of the three involves truncating the twenty scalars $`T_{ij}`$ of the $`SL(6,\text{I}\mathrm{R})/SO(6)`$ coset to the diagonal subset
$$T_{ij}=\mathrm{diag}(X_1,X_2,X_3,X_4,X_5,X_6),$$
(13)
where $`_iX_i=1`$, and setting all the fifteen gauge fields $`F_{\left(2\right)}^{ij}`$ to zero. This reduces to the embedding that was obtained in , for which the complete proof of consistency was constructed in .
The second possible truncation involves reducing the scalar sector still further, to
$$T_{ij}=\mathrm{diag}(\stackrel{~}{X}_1,\stackrel{~}{X}_1,\stackrel{~}{X}_2,\stackrel{~}{X}_2,\stackrel{~}{X}_3,\stackrel{~}{X}_3),$$
(14)
where $`_a\stackrel{~}{X}_a=1`$, but now retaining the three $`U(1)`$ gauge fields $`F_{\left(2\right)}^{12}`$, $`F_{\left(2\right)}^{34}`$ and $`F_{\left(2\right)}^{56}`$ of the maximal abelian $`U(1)^3`$ subgroup of $`SO(6)`$. (It is easy to see, by looking at the five-dimensional equations of motion in (2), that this is a consistent truncation.) The truncated theory is supersymmetric, and describes five-dimensional $`U(1)`$-gauged simple supergravity coupled to two $`U(1)`$ vector multiplets. The consistent embedding of this theory in type IIB supergravity was obtained in .
The third possible truncation involves retaining just a single scalar field $`X`$, by taking
$$T_{ij}=\mathrm{diag}(X,X,X,X,X^2,X^2),$$
(15)
at the same time retaining only the gauge fields of $`SU(2)\times U(1)`$. It is convenient now to take the $`SO(6)`$ indices $`i`$ to range over $`0i5`$. We take all the gauge potentials $`A_{\left(1\right)}^{ij}`$ to be zero except for the following:
$`A_{\left(1\right)}^{01}=A_{\left(1\right)}^{23}=\frac{1}{\sqrt{2}}A_{\left(1\right)}^1,A_{\left(1\right)}^{02}=A_{\left(1\right)}^{31}=\frac{1}{\sqrt{2}}A_{\left(1\right)}^2,A_{\left(1\right)}^{03}=A_{\left(1\right)}^{12}=\frac{1}{\sqrt{2}}A_{\left(1\right)}^3,`$
$`A_{\left(1\right)}^{45}=B_{\left(1\right)}.`$ (16)
We also parameterise the coordinates $`\mu ^i`$ of the internal 5-sphere as follows:
$`\mu ^0+\mathrm{i}\mu ^3=\mathrm{cos}\xi \mathrm{cos}\frac{1}{2}\theta e^{\mathrm{i}(\psi +\varphi )/2},\mu ^1+\mathrm{i}\mu ^2=\mathrm{cos}\xi \mathrm{sin}\frac{1}{2}\theta e^{\mathrm{i}(\psi \varphi )/2},`$
$`\mu ^4+\mathrm{i}\mu ^5=\mathrm{sin}\xi e^{\mathrm{i}\tau }.`$ (17)
Substituting (15), (16) and (17) into the metric Ansatz (1), we obtain
$`d\widehat{s}_{10}^2`$ $`=`$ $`\mathrm{\Delta }^{1/2}ds_5^2+g^2X\mathrm{\Delta }^{1/2}d\xi ^2+g^2\mathrm{\Delta }^{1/2}X^2s^2\left(d\tau gB_{\left(1\right)}\right)^2`$ (18)
$`+\frac{1}{4}g^2\mathrm{\Delta }^{1/2}X^1c^2{\displaystyle \underset{i}{}}(\sigma ^i\sqrt{2}gA_{\left(1\right)}^i)^2,`$
where $`c\mathrm{cos}\xi `$, $`s\mathrm{sin}\xi `$, $`\mathrm{\Delta }=X^2s^2+Xc^2`$, $`h^i\sigma ^i\sqrt{2}gA_{\left(1\right)}^i`$, and the $`\sigma _i`$ denote the three left-invariant 1-forms of $`SU(2)`$, given by $`\sigma _1+\mathrm{i}\sigma _2=e^{\mathrm{i}\psi }(d\theta +\mathrm{i}\mathrm{sin}\theta d\varphi )`$, $`\sigma _3=d\psi +\mathrm{cos}\theta d\varphi `$. This is the metric Ansatz obtained in for the embedding of five-dimensional $`N=4`$ gauged $`SU(2)\times U(1)`$ supergravity in $`D=10`$. The internal 5-sphere now has a geometrical interpretation as a foliation by $`S^3\times S^1`$, with $`\xi `$ parameterising the foliating surfaces.
Substituting (15), (16) and (17) into the Ansatz for $`\widehat{G}_{\left(5\right)}`$ given in (3), leads to
$`\widehat{G}_{\left(5\right)}`$ $`=`$ $`gU\epsilon _5{\displaystyle \frac{3sc}{g}}X^1dXd\xi +{\displaystyle \frac{c^2}{8\sqrt{2}g^2}}X^2F_{\left(2\right)}^ih^jh^k\epsilon _{ijk}`$ (19)
$`{\displaystyle \frac{sc}{2\sqrt{2}g^2}}X^2F_{\left(2\right)}^ih^id\xi {\displaystyle \frac{sc}{g^2}}X^4G_{\left(2\right)}d\xi (d\tau gB_{\left(1\right)}),`$
where $`U=2(X^2c^2+X^1s^2+X^1)`$, again in agreement with the Ansatz obtained in . As a consistency check, we can also verify that substituting (15), (16) and (17) into the Ansatz (4) for $`\widehat{}\widehat{G}_{\left(5\right)}`$ gives the same expression as the one obtained for $`\widehat{}\widehat{G}_{\left(5\right)}`$ in . Thus we have verified that the gauged $`SO(6)`$ embedding that we have obtained in this paper can be truncated to the $`SU(2)\times U(1)`$ embedding of the $`N=4`$ theory whose consistency was proven in . This again provides further supporting evidence for the consistency of the our gauged $`SO(6)`$ reduction Ansatz.
## 4 Consistency conditions for sphere reductions
One of the interesting outcomes from our analysis is that it is essential for the consistency of the 5-form reduction Ansatz that it be a self-dual 5-form, rather than an unconstrained one. It seems, therefore, that the requirement of consistency of the sphere reduction has singled out a ten-dimensional starting point that is itself embeddable in a supersymmetric theory.
This raises the more general question of what possible higher-dimensional theories might allow consistent Kaluza-Klein sphere reductions. All the known examples are associated with supersymmetric higher-dimensional theories, but one might wonder whether this was just a reflection of the fact that these are the cases that have received the most attention in the literature. However, the following argument seems to suggest that the supersymmetric cases may be the only ones that can allow consistent $`S^n`$ sphere reductions, in which all the massless fields (including the $`SO(n+1)`$ Yang-Mills fields) are retained.
Consider a $`D`$-dimensional theory of gravity plus an $`n`$-form field strength, with the Lagrangian
$$e^1_D=R\frac{1}{2n!}F_n^2,$$
(20)
where $`e=\sqrt{g}`$. If this were to give a $`(Dn)`$-dimensional theory with an $`SO(n+1)`$ gauge group, as a consistent reduction on $`S^n`$, it would be necessary that the ungauged $`(Dn)`$-dimensional theory obtained by reducing instead on the torus $`T^n`$ should have a global symmetry group $`G`$ that contains $`SO(n+1)`$ as a compact subgroup, since an $`SO(n+1)`$ factor in the denominator group would be gauged in the spherical reduction. A reduction on $`T^n`$ always produces a theory with a $`GL(n,\text{I}\mathrm{R})`$ global symmetry, which has $`SO(n)`$ as its maximal compact subgroup, and so this would be insufficient for allowing an $`SO(n+1)`$ gauging. In special cases the $`GL(n,\text{I}\mathrm{R})`$ global symmetry can be enhanced to $`SL(n+1,\text{I}\mathrm{R})`$, but this happens only if there is a “conspiracy” between axionic scalars coming from the metric and axions coming from the form-field $`F_n`$. For this conspiracy to occur, the strengths of the dilaton couplings to axions from these two sources must be the same.<sup>5</sup><sup>5</sup>5It is always the case that the counting of dilatons and axions would be consistent with the numerology required for an $`SL(n+1,\text{I}\mathrm{R})/SO(n+1)`$ coset structure, but this does not in general imply that the enhanced coset actually occurs. Specifically, if $`\stackrel{}{\varphi }`$ denotes the set of $`n`$ canonically-normalised dilatons that result from the $`T^n`$ reduction then all the dilaton/axion couplings should be of the form $`e^{\stackrel{}{a}_i\stackrel{}{\varphi }}(\chi _i)^2`$, where the constant vectors $`\stackrel{}{a}_i`$ satisfy $`(\stackrel{}{a}_i)^2=4`$ for each $`i`$. (In fact the full set of $`\stackrel{}{a}_i`$ vectors would constitute the positive-root vectors of $`SL(n+1,\text{I}\mathrm{R})`$ .) It was shown in that the strengths of dilaton couplings can be conveniently characterised in terms of the quantity $`\mathrm{\Delta }`$, related to the dilaton vector $`\stackrel{}{a}`$ by
$$\stackrel{}{a}^2=\mathrm{\Delta }\frac{2(n1)(Dn1)}{D2},$$
(21)
for an $`n`$-form field strength in $`D`$ dimensions, since the quantity $`\mathrm{\Delta }`$ is preserved under toroidal Kaluza-Klein reduction. Furthermore, it was shown that the Kaluza-Klein vectors and axions coming from a toroidal reduction of the metric always have $`\mathrm{\Delta }=4`$. It therefore follows that for the symmetry enhancement to $`SL(n+1,\text{I}\mathrm{R})`$ to take place, the dilaton-coupling $`\mathrm{\Delta }`$ for the original $`n`$-form field<sup>6</sup><sup>6</sup>6If there is no dilaton in the original higher-dimensional theory then the “dilaton coupling” $`\mathrm{\Delta }`$ is given by setting $`\stackrel{}{a}^2=0`$ in (21), and it is this value that is preserved under toroidal reduction. in (20) must also take the value $`\mathrm{\Delta }=4`$. An enumeration of all possible cases shows that for Lagrangians of the form (20) one gets $`\mathrm{\Delta }=4`$ only for
$$(D,n)=(11,4),(11,7),(10,5).$$
(22)
The above considerations seem to single out the three cases in (22) as the only ones where an $`S^n`$ reduction of a Lagrangian of the form (20) could consistently yield the gauge fields of $`SO(n+1)`$ and the scalars $`T_{ij}`$ of the coset $`SL(n+1,\text{I}\mathrm{R})/SO(n+1)`$.<sup>7</sup><sup>7</sup>7Consistent sphere reductions with scalars $`T_{ij}`$ for more general $`(D,n)`$ values for the Lagrangian (20) have been suggested in . We expect that a more complete analysis of the consistency of the reduction would exclude such possibilities. In particular, scalars $`T_{ij}`$ parametrising the coset $`SL(n+1,\text{I}\mathrm{R})/SO(n+1)`$ only occur in the three cases listed in (22). As we have seen in section 2 of this paper, it can turn out that additional structure is also required in order to achieve consistency, namely the self-duality of the 5-form in the case of $`(D,n)=(10,5)`$. By the same token one can expect that a consistent reduction would only be possible in the cases $`(D,n)=(11,4)`$ and $`(11,7)`$ if additional structure is also present in the eleven-dimensional Lagrangian.
For example, if we consider the $`S^7`$ reduction from $`D=11`$, then the total set of 70 spin-0 fields decompose as a $`35_v`$ of scalars and a $`35_c`$ of pseudoscalars. The 35 scalars in $`T_{ij}`$ correspond to the $`35_v`$. Turning on these forces all the 28 gauge fields of $`SO(8)`$ to be excited, and in turn these excite the remaining $`35_c`$ of pseudoscalars. In order for an $`S^7`$ reduction that retains all these fields to be consistent, it is necessary to include an $`F_{\left(4\right)}F_{\left(4\right)}A_{\left(3\right)}`$ term in the original Lagrangian (20), with precisely the coefficient dictated by $`D=11`$ supersymmetry. This point also emphasises that in the $`S^7`$ reduction one cannot consider just the 35 scalars $`T_{ij}`$ in isolation; the full set of bosonic fields of $`N=8`$ gauged supergravity (and not merely the 28 gauge fields) must be included if the full set of $`T_{ij}`$ scalars are present.
A similar situation arises with the $`S^4`$ reduction from $`D=11`$. Including the full set of 14 scalars $`T_{ij}`$ in a consistent reduction will force the complete set of massless fields to be non-vanishing, including not only the ten Yang-Mills fields of $`SO(5)`$ but also the five 3-form field strengths coming from the antisymmetric tensor. The consistency of the reduction is then only possible if the $`FFA`$ term of $`D=11`$ supergravity is included.
Another possibility for obtaining further examples of consistent sphere reductions is to include a dilaton in the higher-dimensional Lagrangian, whose coupling to the $`n`$-form field strength is arranged to have $`\mathrm{\Delta }=4`$:
$$e^1_D=R\frac{1}{2}(\varphi )^2\frac{1}{2n!}e^{a\varphi }F_n^2,$$
(23)
with $`a^2=42(n1)(Dn1)/(D2)`$. (Lagrangians of this type without the restriction on the value of the constant $`a`$ have also been discussed in in the context of sphere reductions.) In this case there will no longer be an AdS$`{}_{Dn}{}^{}\times S^n`$ vacuum solution, but rather a warped product of a domain wall and $`S^n`$. The possibilities for achieving $`\mathrm{\Delta }=4`$ couplings are in fact rather limited, given that $`a`$ should be a real number. First of all, for an $`n`$-form with $`4nD4`$ it can be seen that we must have $`D11`$. For example if $`n=4`$, then we must have $`D11`$, while for $`n=5`$ we must have $`D10`$. (It is only necessary to consider forms with $`nD/2`$ since Hodge duality maps those with $`n>D/2`$ into this range.) Interestingly, for $`n=3`$ one can achieve $`\mathrm{\Delta }=4`$ in an arbitrary dimension $`D`$; the Lagrangian (23) then corresponds to the low-energy effective theory of the $`D`$-dimensional bosonic string. We can expect that the full 3-sphere reduction of this Lagrangian (keeping the complete $`SO(4)`$ gauge fields of its isometry group, not merely the $`SU(2)`$ subset of left-invariant fields) will be consistent. Similarly, we can expect that it should be consistent to reduce the theory with $`n=3`$ on a $`(D3)`$-sphere. Before the gauging, the scalar coset in $`D=3`$ is $`SO(D2,D2)/(SO(D2)\times SO(D2))`$. One of the $`SO(D2)`$ denominator group factors can be gauged, and we obtain the scalar coset $`SL(D2,\text{I}\mathrm{R})/SO(D2)`$ together with the additional gauge fields of $`SO(D2)`$, and a singlet scalar (which is the original dilaton of the $`D`$-dimensional theory).
A consistent sphere reduction, albeit of a slightly different kind, has in fact been obtained in an example where there is a dilaton in the higher-dimensional theory that couples to the form-field. In it was shown that a reduction of the massive type IIA supergravity on an internal 4-dimensional space that is locally $`S^4`$ gives rise to the $`N=2`$ $`SU(2)`$ gauged supergravity in $`D=6`$.
## 5 Conclusions and further comments
In this paper we have constructed a consistent Kaluza-Klein reduction Ansatz for embedding the subset of the fields of five-dimensional $`N=8`$ gauged supergravity, comprising gravity, the fifteen $`SO(6)`$ gauge fields, and the twenty scalars of the $`SL(6,\text{I}\mathrm{R})/SO(6)`$ submanifold of the full scalar manifold, into type IIB supergravity in $`D=10`$. This embedding can equivalently be viewed as a complete reduction Ansatz (with no truncation of massless fields) for the ten-dimensional theory comprising just gravity plus a self-dual 5-form, which itself is a consistent truncation of type IIB supergravity.
A crucial point in the analysis is that in the gauged five-dimensional supergravity one cannot consistently set the Yang-Mills fields to zero, while retaining the full set of scalar fields, unlike the situation in ungauged supergravity. In the context of the truncation that we consider in this paper, where we retain the twenty scalars $`T_{ij}`$, we cannot ignore the fifteen $`SO(6)`$ gauge fields, since the scalars act as sources for them.<sup>8</sup><sup>8</sup>8This also implies that solutions built using these scalars will in general require non-vanishing Yang-Mills fields. We saw that the self-duality of the 5-form field of the type IIB theory plays an essential rôle in the consistency of the reduction.
More generally, we showed that if one starts from a theory comprising gravity and an $`n`$-form field strength, then a consistent reduction on $`S^n`$ that retains the scalars $`T_{ij}`$ of the coset $`SL(n+1,\text{I}\mathrm{R})/SO(n+1)`$ will also have to include at least the gauge fields of $`SO(n+1)`$, and furthermore will only be possible for the $`S^4`$ and $`S^7`$ reductions of $`D=11`$, and the $`S^5`$ reduction of $`D=10`$. In fact the consistency will in addition require that the theories in $`D=11`$ and $`D=10`$ have the additional structures associated with $`D=11`$ and type IIB supergravity, namely the $`FFA`$ term in $`D=11`$, and the self-duality of the 5-form in $`D=10`$. In $`D=7`$ the scalars and Yang-Mills fields must be supplemented by the five 2-form potentials, while in $`D=4`$ they must be supplemented by the 35 pseudo-scalars, in order to achieve consistency. On the other hand in $`D=5`$ no additional fields beyond the scalars $`T_{ij}`$ and the gauge fields $`A_{\left(1\right)}^{ij}`$ are required for consistency. In fact in all three cases this is related to the presence of the terms of the form (10), bilinear in Yang-Mills fields. In $`D=7`$ this term acts as a source for the five 3-form fields; in $`D=4`$ it acts as a source for the 35 pseudo-scalars; but in $`D=5`$ it acts as a “source” for the Yang-Mills fields themselves. This special feature of the $`D=5`$ gauged supergravity may have implications in the dual four-dimensional $`N=4`$ super Yang-Mills theory.
We also discussed the more general possibilities that might arise if one includes a dilaton in the higher-dimensional theory. The possibilities for further examples of consistent sphere reductions seem to be rather limited, as discussed in section 4.
In a full $`S^5`$ reduction of type IIB supergravity there will be additional fields coming from the reduction of the NS-NS and R-R 2-form potentials, and from the dilaton and axion. A complete analysis of the $`S^5`$ reduction can therefore be expected to be extremely complicated. In particular, for example, the $`10`$ and $`\overline{10}`$ of pseudo-scalars lead to a considerably more complicated metric reduction Ansatz. The Ansatz for a subset of the fields that included one scalar and one pseudoscalar was derived in , and in . However even in that case, the construction of the Ansatz for the antisymmetric tensor fields is rather involved, and has not yet been pushed to completion.
## Acknowledgement
C.N.P. is grateful to the University of Pennsylvania for hospitality during the course of this work.
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# G328.4+0.2: A large and luminous Crab-like supernova remnant
## 1. Introduction
When a massive star ends its life in a supernova explosion, the outer layers of the star are expelled at high velocity, and interact with the ambient interstellar medium (ISM) to produce a supernova remnant (SNR). “Shell” SNRs are usually characterized by radio synchrotron emission with a limb-brightened morphology, roughly centered on the site of the supernova explosion. Some SNRs are classified as “composite”, indicating that as well as a shell they have an additional central component, characterized by a filled-center morphology, a flat spectral index ($`0.3\alpha 0`$; $`S_\nu \nu ^\alpha `$) and significant linear polarization. This extra component is usually interpreted as a synchrotron nebula, powered by a pulsar also formed in the supernova explosion (Milne et al. (1979); Weiler & Panagia (1980); Reynolds & Chevalier (1984)). Thus even when a pulsar itself is not detected (as is the case much more often than not), the mere existence of one of these synchrotron nebulae tells us that an energetic pulsar must be located within.
The best known example of a pulsar-powered nebula is the Crab Nebula. However, it has long been apparent that the Crab is very different from most other SNRs, in that it consists of a synchrotron nebula and associated pulsar, but has no surrounding shell corresponding to the supernova blast wave. It is now recognized that about 5% of all SNRs similarly lack shells; these sources are generally referred to as “Crab-like” SNRs, or “plerions” (Weiler & Panagia (1978)). It is still an open question as to whether Crab-like SNRs have surrounding shells which are simply not detectable (Chevalier (1977); Sankrit & Hester (1997)), or whether they are fundamentally different from other SNRs in that they have no shell, invisible or otherwise (Nomoto (1987); Wallace, Landecker, & Taylor (1997)).
Crab-like SNRs are thus important to identify and study, both for understanding the nature and evolution of SNRs, and because they are an unambiguous indication of the presence of a young pulsar. The source G328.4+0.2 (MSH 15–57; Mills, Slee, & Hill (1961)) has been been the subject of several radio studies (Shaver & Goss (1970); Caswell et al. (1980); Whiteoak & Green (1996)), which have indicated a featureless filled-center appearance and a comparatively flat spectral index. These properties have resulted in G328.4+0.2 being proposed as a member of the Crab-like class of SNRs, and indeed is the only Galactic Crab-like SNR easily accessible to Southern hemisphere observers. H i measurements with a two-element interferometer have shown H i absorption in this direction to a distance in excess of 20 kpc (Caswell et al. (1975)). This large distance would make the SNR extraordinary luminous, even brighter than the Crab Nebula itself.
However, previous images of this source have been of only intermediate resolution (the best observations to date had only six beams across the SNR), while the lack of imaging capability in the H i measurements of Caswell et al. (1975) mean that the absorption they saw could have been coming from another bright source in the vicinity, rather than from the SNR itself. We therefore present new radio observations of G328.4+0.2 both in polarimetric continuum and in the H i line, aimed at confirming the Crab-like nature of this source, verifying the large distance claimed for it, and studying it at much higher resolution and sensitivity than in previous work. Our observations and analysis are described in §2, and our results are presented in §3. In §4 we confirm G328.4+0.2 as a Crab-like SNR, and infer the properties of the central pulsar presumed to be powering it. A complementary study of this source in X-rays is reported in a companion paper by Hughes, Slane & Plucinsky (2000).
## 2. Observations and Data Reduction
Our observations of G328.4+0.2 were made with the Australia Telescope Compact Array (ATCA; Frater, Brooks, & Whiteoak (1992)), a 6 km east-west synthesis array located near Narrabri, NSW, Australia, on dates as listed in Table 1. Continuum observations were made at 1.4 and 4.5 GHz; each band had a width of 128 MHz, divided into 32 channels. Simultaneous with the 1.4 GHz data, observations were made in the H i line, centered on 1421 MHz and using 512 channels across a 4 MHz bandwidth. All four Stokes parameters (XX, YY, XY and YX) were recorded in continuum observations; only XX and YY were obtained in the H i line. All observations consisted of a single pointing centered on G328.4+0.2. Flux density calibration was carried out using observations of PKS B1934–638 (with assumed flux densities of 14.9 and 6.3 Jy at 1.4 and 4.5 GHz respectively), while antenna gains and polarization were calibrated using regular observations of MRC B1456–367 at 1.4 GHz and PMN J1603–4904 at 4.5 GHz.
Data were reduced in the MIRIAD package. After flagging and calibration were carried out, total intensity images were formed at each of 1.4 and 4.5 GHz using uniform weighting, multi-frequency synthesis and maximum entropy deconvolution. The resulting images were then corrected for the mean primary beam response, with resolutions and sensitivities given in Table 1.
Images in linear polarization were formed by creating separate images in Stokes $`Q`$ and $`U`$, deconvolving each of these separately using the CLEAN algorithm, and then combining with appropriate debiasing to form $`L=(Q^2+U^2)^{1/2}`$. At 1.4 GHz, significant Faraday rotation across the band can depolarize the emission, and images of linear polarization were produced on a channel by channel basis and then averaged to form a single map. At 4.5 GHz, differential Faraday effects are negligible and the entire observing bandwidth was used to form a single pair of Stokes $`Q`$ and $`U`$ images. The resulting maps of linear polarization were clipped wherever Stokes $`I`$, $`Q`$ or $`U`$ fell below 3$`\sigma `$.
For the H i data, the continuum contribution was subtracted in the u-v plane (van Langevelde & Cotton (1990)), and a cube then formed using uniform weighting and discarding all baselines longer than 7 k$`\lambda `$. The cube contained planes between –200 km s<sup>-1</sup> and +200 km s<sup>-1</sup> (LSR) at intervals of 4 km s<sup>-1</sup>. The peak flux density in the image was sufficiently low that no sidelobes were apparent, and no attempt was made to deconvolve. The cube was then weighted by the corresponding 1.4 GHz continuum image. Absorption spectra against SNR G328.4+0.2 and the nearby H ii region G328.30+0.43 were then generated by integrating over an appropriate spatial region and renormalizing appropriately to give units of fractional absorption.
The rms noise, $`\sigma `$, in each absorption spectrum was estimated from the fluctuations in line-free channels. As in previous papers (e.g. Gaensler, Manchester, & Green (1998)), we adopt $`6\sigma `$ as a threshold for significance, to take into account the increase in system temperature of the receivers due to the brightness temperature of H i in the Galactic Plane.
## 3. Results
### 3.1. Continuum
Images of G328.4+0.2 at 1.4 and 4.5 GHz are shown in Fig 1, and properties of the source at each frequency are given in Table 1. The source is approximately circular, with a bar-like feature running east-west at the peak of the emission. A plateau of fainter emission surrounds this bar; this plateau is significantly more pronounced to the east and to the south.
G328.4+0.2 is of sufficiently small diameter that its entire flux density is recovered by the telescope at both observing frequencies, and the corresponding flux densities (after a background correction has been applied) are given in Table 1. These flux densities, together with the 0.8 GHz measurement of Whiteoak & Green (1996), imply a spectral index for the source $`\alpha =0.12\pm 0.03`$.
Examination of the whole 1.4 GHz field (not shown here) shows there to be no outer shell of emission surrounding G328.4+0.2, out to a radius of 15 arcmin and down to a 1$`\sigma `$ sensitivity of 0.8 mJy beam<sup>-1</sup>, corresponding to a surface brightness limit at 1 GHz (assuming a typical shell spectral index $`\alpha =0.5`$) of $`\mathrm{\Sigma }=1.6\times 10^{21}`$ W m<sup>-2</sup> Hz<sup>-1</sup> sr<sup>-1</sup>.
Discarding the shortest baselines to filter out extended emission, we can put $`5\sigma `$ limits at 1.4 GHz and 4.5 GHz of 7 and 0.3 mJy respectively on the flux density of any point source within G328.4+0.2.
### 3.2. Polarization
At 1.4 GHz, the only linearly polarized emission seen from G328.4+0.2 is in a couple of clumps seen at the position of the peak in total intensity. These clumps have a fractional polarized intensity of $``$1%. When imaged on a channel by channel basis, the variation of the position angle of this polarization, $`\varphi `$, shows a $`\varphi \lambda ^2`$ dependence across the 1.4 GHz band. This allows us to determine a rotation measure (RM) for this emission (cf. Gaensler, Manchester, & Green (1998)) of $`900\pm 100`$ rad m<sup>-2</sup>, where the uncertainty is dominated by spatial variations in RM rather than by uncertainty in the measurement process (for which typical uncertainties are $`\pm 30`$ rad m<sup>-2</sup>).
At 4.5 GHz, significantly more linear polarization is seen from G328.4+0.2, with a spatial distribution shown in Fig 2. Although the polarization appears to be reduced around the edges of the source, this is most likely an artifact of the clipping of polarized emission below $`3\sigma `$ in Stokes $`Q`$ and $`U`$; the fractional polarized intensity shows no decrease towards the edges. The level of fractional polarization ranges between 5% and 50% with a mean of 20%.
An RM of 1000 rad m<sup>-2</sup> corresponds to $`\mathrm{\Delta }\varphi 10^{}`$ across the entire 4.5 GHz observing bandwidth, which is smaller than the uncertainties in position angles of individual frequency channels. Thus no rotation measures could be extracted from the 4.5 GHz data.
### 3.3. H i absorption
H i absorption towards G328.4+0.2 needs to be compared to an H i emission spectrum in a similar part of the sky. But because most of the power in H i emission is on scales larger than were sampled by our observations, such a spectrum is difficult to extract from our ATCA data. Instead, we use an H i emission spectrum taken from the single-dish survey of Kerr et al. (1986), in the direction $`l=328\stackrel{}{\mathrm{.}}5`$, $`b=0\stackrel{}{\mathrm{.}}25`$. This emission profile is shown in the upper panel of Fig 3. Two ATCA absorption profiles are also shown in Fig 3: towards G328.4+0.2 and towards the H ii region G328.30+0.43, 15 arcmin to the north-west. The latter emits in various maser lines, and is known to have a systemic velocity of approximately –95 km s<sup>-1</sup> (Braz & Epchtein (1983); Cohen, Masheder, & Caswell (1995)).
At negative velocities, the two spectra are very similar, and match well the H i emission seen along similar lines-of-sight. However, a significant difference is seen at positive velocities, where H i at a velocity of +28 km s<sup>-1</sup> is clearly producing significant absorption against G328.4+0.2, but only a weak feature is seen for G328.30+0.43. Examination of the image plane corresponding to this velocity clearly shows the outline of all of G328.4+0.2 in absorption, leaving no doubt that this absorption feature is genuine. However, for the H ii region G328.30+0.43 we see no real match in absorption to the source’s morphology at this velocity, as expected given that it is known to be at the tangent point. This weak absorption probably corresponds to fluctuations in the H i emission at this velocity, appearing negative because of spatial filtering by the interferometer.
## 4. Discussion
### 4.1. Distance
The H i spectrum shown in Fig 3 demonstrates clearly that the systemic velocity of G328.4+0.2 is at least +28 km s<sup>-1</sup>. H i emission is seen at +28 km s<sup>-1</sup> not just in the direction of G328.4+0.2, but towards several other sources in the vicinity (Caswell et al. (1975)). Thus it is unlikely that the emission and absorption seen at this velocity correspond to some much closer cloud of gas which deviates dramatically from the Galactic rotation curve. While an H i emission feature at +58 km s<sup>-1</sup> shows no corresponding absorption, its low brightness temperature means that we can probably not use this result to derive an upper limit on the source’s radial velocity.
To calculate a lower limit on the distance to G328.4+0.2, we adopt standard IAU parameters for the Sun’s orbital velocity ($`\mathrm{\Theta }_0=220`$ km s<sup>-1</sup>) and distance from the Galactic Centre ($`R_0=8.5`$ kpc) (Kerr & Lynden-Bell (1986)), use the best fitting model for Galactic rotation of Fich, Blitz & Stark (1989), and assume uncertainties in systemic velocities of $`\pm `$7 km s<sup>-1</sup>. For a systemic velocity of +28 km s<sup>-1</sup>, it then follows that G328.4+0.2 must be at a distance of at least $`17.4\pm 0.9`$ kpc. In further discussion, we assume a distance $`d=17d_{17}`$ kpc; the diameter of G328.4+0.2 is then $`25d_{17}`$ pc.
### 4.2. Morphology
G328.4+0.2 is an extended Galactic source with a filled-center morphology, significant linear polarization across most of its extent, and a spectral index much flatter than for a typical shell-type supernova remnant. Down to our surface brightness limit, there is no evidence that G328.4+0.2 is surrounded by a larger, limb-brightened shell. This limit is sufficient to detect the shell of a typical young SNR (see Frail et al. (1995)), although not sensitive enough to detect the shell component of the composite SNR G322.5–0.1 (Whiteoak (1992)). Thus while the limits on a shell around G328.4+0.2 are not as stringent as around some other such sources, from the available data we class it as a Crab-like SNR.
The mere presence of a Crab-like or composite SNR is usually taken to indicate that there is an energetic pulsar within, even though in most such sources no pulsar has yet been detected. While indeed no pulsar has been seen within G328.4+0.2 in either radio waves (Manchester, D’Amico, & Tuohy (1985); Kaspi et al. (1996)) or in X-rays (Wilson (1986); Hughes, Slane, & Plucinsky (2000)), this can easily be accounted for by unfavorable beaming, insufficient sensitivity at the SNR’s large distance, and significant scattering, dispersion (at radio wavelengths) and absorption (in X-rays) along the long line of sight. We thus assume in further discussion that there is an unseen rotation-powered pulsar within G328.4+0.2, whose relativistic wind powers the radio emission from the remnant.
The bar running through the center of G328.4+0.2 is a feature similar to that seen in many other Crab-like SNRs. There are two interpretations of such a feature: one is that it corresponds to a trail of emitting particles left behind by a pulsar with a significant space velocity (e.g. Frail et al. (1996); Sun, Wang, & Chen (1999)), the other is that the bar is actually comprised of two opposed jet-like features, produced by a bipolar outflow from the pulsar (e.g. Fürst et al. 1988, 1989). High-resolution X-ray imaging is required to distinguish between these possibilities, since the much shorter synchrotron lifetime of X-ray emitting particles will result in X-ray emission being concentrated around the current point of injection into the SNR. So if the bar is a high velocity pulsar leaving a trail, we expect to see X-ray emission concentrated at one end of the bar, while if it is a pair of outflows, X-ray emission should peak near the center of the bar.
### 4.3. Physical Properties
We first determine the radio luminosity, $`L_R`$, of G328.4+0.2. Hughes et al. (2000) demonstrate that any break in the continuum spectrum of this source is at a frequency higher than 100 GHz. We thus assume that the radio spectrum of G328.4+0.2 is a single power-law of spectral index $`\alpha =0.12`$. Integrating between 100 MHz and 100 GHz, we then find a broad-band radio luminosity for G328.4+0.2 of $`L_R=3.3d_{17}^2\times 10^{35}`$ erg s<sup>-1</sup>.
In Table 2, we compare the size and radio luminosity of G328.4+0.2 to those of other Crab-like SNRs (as well as Crab-like cores of composite SNRs) with similar properties; it can be seen that G328.4+0.2 is one of the most radio-luminous and largest Crab-like SNRs yet discovered. Two sources closely resemble G328.4+0.2: N157B (Wang & Gotthelf (1998); Lazendic et al. (2000)), a Crab-like SNR in the Large Magellanic Cloud (LMC) in which a 16-ms X-ray pulsar has recently been discovered<sup>1</sup><sup>1</sup>1The pulsar has so far remained undetected in radio waves, but otherwise has properties typical of a radio pulsar (Crawford et al. (1998))., and G74.9+1.2 (Weiler & Shaver (1978); Wallace et al. (1997)), a Galactic SNR which as yet has had no pulsar detected in it. Below we will argue that these three SNRs constitute a particular sub-class of Crab-like SNR.
A pulsar’s spin-down luminosity is given by
$$\dot{E}4I\pi ^2\dot{P}/P^3,$$
(1)
where $`P`$ is the period of the pulsar and $`I=10^{45}`$ g cm<sup>-2</sup> is its moment of inertia. Defining $`ϵ=L_R/\dot{E}`$, we find that in the few cases where a pulsar has been detected within a Crab-like SNR, the radio luminosity of the SNR normally falls in the narrow range $`ϵ=(15)\times 10^4`$ (Frail & Scharringhausen (1997); Gaensler et al. (2000)). In particular, the Crab Nebula has $`ϵ=4\times 10^4`$, and we assume a similar value for G328.4+0.2. One then finds that its central pulsar has $`\dot{E}_{38}=8.3`$ (where $`\dot{E}=10^{38}\dot{E}_{38}`$ erg s<sup>-1</sup>), a value consistent with that derived by Hughes et al. (2000) from the X-ray properties of this source. This value of $`\dot{E}`$ ranks amongst the highest values seen for the radio pulsar population, and makes this source a good candidate for pulsed emission searches by future X-ray and $`\gamma `$-ray missions.
A pulsar drives a supersonic bubble into its surroundings, whose radius, $`R`$, is given by (Weaver et al. (1977); Arons (1983)):
$$R=0.82\dot{E}_{38}^{1/5}t_3^{3/5}n_0^{1/5}\mathrm{pc},$$
(2)
where $`t_3`$ kyr is the age of the pulsar and $`n_0`$ cm<sup>-3</sup> is the density of the ambient medium. It is usually assumed that Crab-like SNRs are propagating into a low density component of the ISM (Chevalier (1977)), for which we adopt a typical density $`n_0=0.003`$ cm<sup>-3</sup>. We can then determine an age for G328.4+0.2 of $`t_37d_{17}^{5/3}`$ kyr.
For a idealized braking index $`n=3`$ and an initial period $`P_iP`$, the age of the SNR is equal to its pulsar’s characteristic age, $`\tau _cP/2\dot{P}`$. We can then combine our inferred age with Equation (1) to determine parameters for the pulsar of $`P=11`$ ms and $`B3.2\times 10^{19}(P\dot{P})^{1/2}`$ G $`=6\times 10^{11}`$ G, where $`B`$ is the inferred surface magnetic field of the neutron star, assuming a dipole geometry. These inferred parameters, along with the results of similar calculations for G74.9+1.2 and N157B, are given in Table 2 – we find that all three pulsars are of much smaller period, lower magnetic field and larger age than the Crab Pulsar.
Various assumptions have gone into the calculations behind this last statement, and we now consider what effect each one has on the results. The assumption least likely to be true is that $`t=\tau _c`$; if we relax this requirement to include braking indices $`n3`$ and initial periods $`P_iP`$, we still find, for a wide range of parameters, that $`\tau _c0.7t`$ (see Fig 3 of Marshall et al. (1998)). Since $`P\tau _c^{1/2}`$ and $`B\tau _c^1`$, the consequent inferred spin periods and magnetic fields are still significantly smaller than for the Crab Pulsar. A second assumption is that $`ϵ=ϵ_{\mathrm{Crab}}`$; however the dependence on $`ϵ`$ is even weaker than on $`\tau _c`$ ($`Pϵ^{1/3}`$ and $`Bϵ^{1/6}`$), so that even invoking a value of $`ϵ`$ ten times larger than that for the Crab (which would be higher than for any pulsar observed) has little effect on the results. We have also assumed a value $`n=0.003`$ cm<sup>-3</sup>; a significantly lower density than this can indeed produce a more typical magnetic field ($`Bn^{2/3}`$), but still requires a fast period ($`Pn^{1/6}`$). Meanwhile, Equation (2) is only strictly true for constant $`\dot{E}`$. In reality, $`\dot{E}`$ will decrease as a function of time, so that the radius of the pulsar wind bubble will be larger than given by Equation (2) and will hence cause us to over-estimate the age of the system. While this effect is difficult to quantify, the ages given in Table 2 would have to be too large by a factor of 5–10 to produce periods as slow as that of the Crab Pulsar. Finally, we note that the calculations we have made for G328.4+0.2 all use the lower limit on the distance derived in Section 4.1; at greater distances, the inferred period and magnetic field are both even lower than those given in Table 2.
Thus even if not all of the assumptions which we have made are valid, one is still forced to conclude that the pulsars powering these SNRs are likely to be spinning rapidly and have low magnetic fields. In particular, the pulsar actually observed within N157B indeed fulfils this prediction, having observed properties quite similar to those we have inferred for it (see Table 2).
## 5. Conclusions
Our radio observations of G328.4+0.2 confirm it to be a Crab-like SNR at a distance of $`>`$17 kpc, making it the largest and most radio-luminous such object in our Galaxy. G328.4+0.2, together with the Galactic SNR G74.9+1.2 and N157B in the LMC, appear to form a small subset of Crab-like SNRs with both high radio luminosities and large diameters. The high luminosities of these remnants demand a high value of $`\dot{E}`$ for their central pulsars, while these remnants’ large extent, even assuming an energetic pulsar and a low ambient density, requires them to be significantly older than the Crab Nebula. Since $`\dot{E}\tau _cP^2`$ and $`\dot{E}\tau _c^2B^2`$, this combination of a high $`\dot{E}`$ and large age can only be produced by a pulsar which is spinning at least twice as fast as the Crab Pulsar, but which has a magnetic field $``$5 times weaker. Specifically, we infer G328.4+0.2 to be $``$7000 yrs old and powered by a pulsar with period $`P11`$ ms and dipole magnetic field $`B=6\times 10^{11}`$ G. Similar properties are predicted for G74.9+1.2 and for N157B, which in the latter case agree with those of the pulsar recently detected within this SNR.
Using the distribution of pulsar initial magnetic fields proposed by Stollman (1987), we can estimate that the birth-rate of pulsars with magnetic fields $`B6\times 10^{11}`$ G is a fraction 0.05 of that of pulsars with magnetic fields comparable to the Crab Pulsar. Thus of the $``$35 Crab-like and composite SNRs known in the Galaxy, we can expect $``$2 to be powered by such low-field pulsars, and G328.4+0.2 and G74.9+1.2 may thus represent the complete sample of such sources.
The radio beaming fraction for young pulsars is estimated to be in the range 50–70% (Frail & Moffett (1993); Brazier & Johnston (1999)). Thus the failure to detect radio pulsations from G328.4+0.2, G74.9+1.2 and N157B (Kaspi et al. (1996); Lorimer, Lyne, & Camilo (1998); Crawford et al. (1998)) is unlikely to be solely due to beaming, and more likely results from the large distance to these objects. With continuing improvements in sensitivity, searches for pulsations towards these sources should eventually be successful.
We thank Taisheng Ye for early assistance with this project, and Pat Slane and Jack Hughes for useful discussions on the manuscript. We also acknowledge some important suggestions made by the referee, Roger Chevalier. The Australia Telescope is funded by the Commonwealth of Australia for operation as a National Facility managed by CSIRO. B.M.G. acknowledges the support of NASA through Hubble Fellowship grant HF-01107.01-98A awarded by the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., for NASA under contract NAS 5–26555. J.R.D. acknowledges a Visitor’s Fellowship from the Netherlands Scientific Organization (NWO) during his stay at the NFRA.
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# Fidelity balance in quantum operations
## Abstract
I derive a tight bound between the quality of estimating the state of a single copy of a $`d`$-level system, and the degree the initial state has to be altered in course of this procedure. This result provides a complete analytical description of the quantum mechanical trade-off between the information gain and the quantum state disturbance expressed in terms of mean fidelities. I also discuss consequences of this bound for quantum teleportation using nonmaximally entangled states.
As a general rule, the more information is obtained from an operation on a quantum system, the more its state has to be altered. This heuristic statement was first exemplified by the Heisenberg microscope gedankenexperiment , where the spatial resolution of the apparatus was shown to scale inversely with the uncertainty of the momentum transfered during the observation. Presently, the disturbance caused by the information gain has become an issue of practical significance, as it underlies the security of quantum key distribution .
The balance between the information gain and the state disturbance attracts currently a lot of interest, particularly in the context of quantum cryptography . Information theory provides a selection of concepts to quantify both the information gain and the state disturbance. The choice of measures for these two effects is usually dictated by the relevance to a specific application. In most cases, however, derivation of the actual balance represents a highly nontrivial task, especially if one is tempted to resign from numerical means. The purpose of this Letter is to present a formulation of the information gain versus state disturbance trade-off which is completely solvable using elementary analytical techniques. This formulation is motivated by recent works on quantum state estimation , where the information obtained from the operation is converted into an estimate for the initial state of the system.
The problem considered in this Letter can be formulated as follows. Suppose we are given a single $`d`$-level particle in a completely unknown pure state $`|\psi `$. We want to make a guess about the quantum state of this particle, but at the same time we would like to alter the state as little as possible. One can associate two fidelities with such a procedure. The first one, which we will denote by $`F`$, describes how much the state after the operation resembles the original one. The second fidelity, denoted by $`G`$, characterizes the average quality of our guess. It is natural to expect that these two quantities cannot take simultaneously too large values. What is the actual quantitative bound between them?
Two extreme cases are well known: if nothing is done to the particle we have $`F=1`$, but then our guess about the state of the particle has to be random, which yields $`G=1/d`$. On the other hand, the optimal estimation strategy for a single copy yields $`G=2/(d+1)`$, but then the particle after the operation cannot provide any more information on the initial state; thus also $`F=2/(d+1)`$. I prove here that quantum mechanics imposes a general constraint between $`F`$ and $`G`$ in the form of the following inequality:
$$\sqrt{F\frac{1}{d+1}}\sqrt{G\frac{1}{d+1}}+\sqrt{(d1)\left(\frac{2}{d+1}G\right)}.$$
(1)
I also show that this inequality cannot be further improved, i.e. there exist quantum operations saturating the equality sign.
The most general strategy that can be applied to the particle has the form of a trace-preserving operation described by a set of operators $`\widehat{A}_r`$, where $`r=1,\mathrm{},N`$. These operators satisfy the completeness relation:
$$\underset{r=1}{\overset{N}{}}\widehat{A}_r^{}\widehat{A}_r=\widehat{𝟙}.$$
(2)
The classical information gained from this operation is given by the index $`r`$, which is subsequently used to estimate the initial state of the particle. The outcome $`r`$ of the operation performed on a state $`|\psi `$ is obtained with the probability $`\psi |\widehat{A}_r^{}\widehat{A}_r|\psi `$. This corresponds to the following conditional transformation of the quantum state:
$$|\psi \frac{\widehat{A}_r|\psi }{\sqrt{\psi |\widehat{A}_r^{}\widehat{A}_r|\psi }}.$$
(3)
We shall measure the resemblance of the transformed state to the original one using the squared modulus of the scalar product, equal $`|\psi |\widehat{A}_r|\psi |^2/\psi |\widehat{A}_r^{}\widehat{A}_r|\psi `$. Summation of this expression over $`r`$ with the weights $`\psi |\widehat{A}_r^{}\widehat{A}_r|\psi `$, and integration over all possible input states $`|\psi `$, yields the complete expression for the mean operation fidelity $`F`$:
$$F=\text{d}\psi \underset{r=1}{\overset{N}{}}|\psi |\widehat{A}_r|\psi |^2.$$
(4)
Here the integral $`\text{d}\psi `$ over the space of pure states is performed using the canonical measure invariant with respect to the group unitary transformations on the state vectors of the particle.
Given the outcome $`r`$ of the operation, we can make a guess $`|\psi _r`$ what the state originally was. The quality of this guess, assuming that the initial state was $`|\psi `$, can be quantified with the help of the overlap $`|\psi _r|\psi |^2`$. The mean estimation fidelity $`G`$ is given by the average of this expression over all outcomes $`r`$ with the probability distribution $`\psi |\widehat{A}_r^{}\widehat{A}_r|\psi `$, and by integration over states $`|\psi `$:
$$G=\text{d}\psi \underset{r=1}{\overset{N}{}}\psi |\widehat{A}_r^{}\widehat{A}_r|\psi |\psi _r|\psi |^2.$$
(5)
We will start derivation of the trade-off between the fidelities $`F`$ and $`G`$ by evaluating the integrals over $`|\psi `$. For this purpose, let us introduce in Eq. (4) two decompositions of unity in a certain orthonormal basis $`|i`$:
$`F`$ $`=`$ $`{\displaystyle \underset{r=1}{\overset{N}{}}}{\displaystyle \underset{i,j=0}{\overset{d1}{}}}\psi |ii|\widehat{A}_r^{}|\psi \psi |\widehat{A}_r|jj|\psi `$ (6)
$`=`$ $`{\displaystyle \underset{r=1}{\overset{N}{}}}{\displaystyle \underset{i,j=0}{\overset{d1}{}}}i|\widehat{A}_r^{}\widehat{M}_{ij}\widehat{A}_r|j`$ (7)
where by $`\widehat{M}_{ij}`$ we have denoted the following integrals of projectors on the states $`|\psi \psi |`$:
$$\widehat{M}_{ij}=\text{d}\psi \psi |ij|\psi |\psi \psi |=\frac{1}{d(d+1)}(\delta _{ij}\widehat{𝟙}+|ij|).$$
(8)
The second explicit form of the operators $`\widehat{M}_{ij}`$ has been derived in Ref. . This formula allows us to simplify the expression for the mean operation fidelity $`F`$ to the form:
$`F`$ $`=`$ $`{\displaystyle \frac{1}{d(d+1)}}\left({\displaystyle \underset{i=0}{\overset{d1}{}}}{\displaystyle \underset{r=1}{\overset{N}{}}}i|\widehat{A}_r^{}\widehat{A}_r|i+{\displaystyle \underset{r=1}{\overset{N}{}}}\left|{\displaystyle \underset{i=0}{\overset{d1}{}}}i|\widehat{A}_r|i\right|^2\right)`$ (9)
$`=`$ $`{\displaystyle \frac{1}{d(d+1)}}\left(d+{\displaystyle \underset{r=1}{\overset{N}{}}}|\text{Tr}\widehat{A}_r|^2\right)`$ (10)
Let us now consider the estimation fidelity $`G`$. The guess $`|\psi _r`$ can be represented as a result of a certain unitary transformation $`\widehat{U}_r`$ acting on a reference state, which we will take for concreteness to be $`|0`$:
$$|\psi _r=\widehat{U}_r|0$$
(11)
Using this representation, and changing the integration measure in Eq. (5) according to $`|\psi \widehat{U}_r|\psi `$, we can evaluate the integral over $`|\psi `$:
$`G`$ $`=`$ $`{\displaystyle \underset{r=1}{\overset{N}{}}}{\displaystyle \text{d}\psi |0|\psi |^2\psi |\widehat{U}_r^{}\widehat{A}_r^{}\widehat{A}_r\widehat{U}_r|\psi }`$ (12)
$`=`$ $`{\displaystyle \underset{r=1}{\overset{N}{}}}\text{Tr}(\widehat{U}_r^{}\widehat{A}_r^{}\widehat{A}_r\widehat{U}_r\widehat{M}_{00})`$ (13)
Inserting the explicit form of the operator $`\widehat{M}_{00}=(\widehat{𝟙}+|00|)/[d(d+1)]`$ yields:
$`G`$ $`=`$ $`{\displaystyle \frac{1}{d(d+1)}}\left({\displaystyle \underset{r=1}{\overset{N}{}}}\text{Tr}(\widehat{U}_r^{}\widehat{A}_r^{}\widehat{A}_r\widehat{U}_r)+{\displaystyle \underset{r=1}{\overset{N}{}}}0|\widehat{U}_r^{}\widehat{A}_r^{}\widehat{A}_r\widehat{U}_r|0\right)`$ (14)
$`=`$ $`{\displaystyle \frac{1}{d(d+1)}}\left(d+{\displaystyle \underset{r=1}{\overset{N}{}}}\psi _r|\widehat{A}_r^{}\widehat{A}_r|\psi _r\right)`$ (15)
This expression provides directly a recipe for optimal assignment of guesses $`|\psi _r`$ to outcomes of the operation: each of the components $`\psi _r|\widehat{A}_r^{}\widehat{A}_r|\psi _r`$ in the sum over $`r`$ is maximized if $`|\psi _r`$ is the eigenvector of $`\widehat{A}_r^{}\widehat{A}_r`$ corresponding to its maximum eigenvalue. Consequently, the maximum value of the mean estimation fidelity $`G`$ for a given operation $`\{\widehat{A}_r\}`$ can be written as:
$$G=\frac{1}{d(d+1)}\left(d+\underset{r=1}{\overset{N}{}}\widehat{A}_r^2\right)$$
(16)
where the operator norm is defined in the standard way:
$$\widehat{A}_r=\underset{\phi |\phi =1}{sup}\sqrt{\phi |\widehat{A}_r^{}\widehat{A}_r|\phi }.$$
(17)
In order to relate the fidelities $`F`$ and $`G`$ to each other, let us consider a polar decomposition of the operators $`\widehat{A}_r`$:
$$\widehat{A}_r=\widehat{V}_r\widehat{D}_r\widehat{W}_r$$
(18)
where $`\widehat{V}_r`$ and $`\widehat{W}_r`$ are unitary, and $`\widehat{D}_r`$ is a semi-positive definite diagonal matrix:
$$\widehat{D}_r=\underset{i=0}{\overset{d1}{}}\lambda _i^r|ii|,$$
(19)
with the diagonal elements put in a decreasing order: $`\lambda _0^r\mathrm{}\lambda _{d1}^r0`$. We will first show that only the diagonal matrices $`\widehat{D}_r`$ are relevant to the trade-off. Indeed, the modulus of the trace of the matrix $`\widehat{A}_r`$ appearing in Eq. (10) is bounded by:
$`|\text{Tr}\widehat{A}_r|`$ $`=`$ $`\left|{\displaystyle \underset{i=0}{\overset{d1}{}}}i|\widehat{W}_r\widehat{V}_r\widehat{D}_r|i\right|`$ (20)
$``$ $`{\displaystyle \underset{i=0}{\overset{d1}{}}}\lambda _i^r|i|\widehat{W}_r\widehat{V}_r|i|{\displaystyle \underset{i=0}{\overset{d1}{}}}\lambda _i^r`$ (21)
and moreover any quantum operation can be easily modified in such a way that the equality sign is reached. What needs to be done, is to follow the operation $`\{\widehat{A}_r\}`$ with an extra unitary transformation $`\widehat{W}_r^{}\widehat{V}_r^{}`$ depending on the outcome $`r`$. Let us note that this corresponds to the modification of the operation according to $`\widehat{A}_r\widehat{W}_r^{}\widehat{V}_r^{}\widehat{A}_r`$, which makes each element of the operation a semi-positive hermitian operator. As we are interested in the maximum value of $`F`$, we can further assume with no loss of generality that:
$$F=\frac{1}{d(d+1)}\left[d+\underset{r=1}{\overset{N}{}}\left(\underset{i=0}{\overset{d1}{}}\lambda _i^r\right)^2\right].$$
(22)
The expression for the estimation fidelity written in terms of $`\lambda _i^r`$ takes the form:
$$G=\frac{1}{d(d+1)}\left(d+\underset{r=1}{\overset{N}{}}(\lambda _0^r)^2\right).$$
(23)
In addition, the trace of the completeness condition given in Eq. (2) yields the following constraint on $`\lambda _i^r`$:
$$\underset{r=1}{\overset{N}{}}\underset{i=0}{\overset{d1}{}}(\lambda _i^r)^2=d.$$
(24)
To complete the proof of the inequality (1), it is convenient to introduce vector notation. Let us define $`d`$ real vectors $`𝐯_i=(\lambda _i^1,\mathrm{}\lambda _i^N)`$, where the index $`i`$ runs from $`0`$ to $`d1`$. Sums over $`r`$ appearing in Eqs. (22) and (23) can be written as:
$`f`$ $`=`$ $`{\displaystyle \underset{r=1}{\overset{N}{}}}\left({\displaystyle \underset{i=0}{\overset{d1}{}}}\lambda _i^r\right)^2={\displaystyle \underset{i,j=0}{\overset{d1}{}}}𝐯_i𝐯_j`$ (25)
$`g`$ $`=`$ $`{\displaystyle \underset{r=1}{\overset{N}{}}}(\lambda _0^r)^2=|𝐯_0|^2`$ (26)
where the dot denotes the scalar product, and $`||`$ is the standard quadratic norm. The completeness condition (24) for the operation $`\{\widehat{A}_r\}`$ written in the vector notation takes the form
$$\underset{i=0}{\overset{d1}{}}|𝐯_i|^2=d.$$
(27)
Let us now suppose that the vector $`𝐯_0`$ is fixed. The estimation fidelity is then given by $`G=(d+|𝐯_0|^2)/[d(d+1)]`$. What is the maximum operation fidelity $`F`$ that can be achieved with this constraint? The answer to this question is provided by an application of the Schwarz inequality to Eq. (25):
$$f\underset{i,j=0}{\overset{d1}{}}|𝐯_i||𝐯_j|=\left(\underset{i=0}{\overset{d1}{}}|𝐯_i|\right)^2=\left(\sqrt{g}+\underset{i=1}{\overset{d1}{}}|𝐯_i|\right)^2$$
(28)
We have excluded here from the sum over $`i`$ the norm of the vector $`𝐯_0`$ which is fixed and equal to $`\sqrt{g}`$. The sum of the norms of the remaining vectors can be estimated using the inequality between the arithmetic and quadratic means:
$$\frac{1}{d1}\underset{i=1}{\overset{d1}{}}|𝐯_i|\sqrt{\frac{1}{d1}\underset{i=1}{\overset{d1}{}}|𝐯_i|^2}=\sqrt{\frac{dg}{d1}},$$
(29)
where we have evaluated the sum $`_{i=1}^{d1}|𝐯_i|^2`$ using Eq. (27). Inserting this bound into Eq. (28) we finally obtain the inequality
$$f\left(\sqrt{g}+\sqrt{(d1)(dg)}\right)^2$$
(30)
which expressed in terms of the fidelities $`F`$ and $`G`$ takes the form of Eq. (1).
The necessary and sufficient conditions for a quantum operation to reach the equality sign can be most easily formulated in the vector notation. The Schwarz inequality (28) becomes equality if all the vectors $`𝐯_0,\mathrm{},𝐯_{d1}`$ are collinear. Furthermore, equation sign in Eq. (29) holds if and only if $`|𝐯_1|=\mathrm{}=|𝐯_{d1}|`$. It is straightforward to see that an exemplary operation satisfying these conditions for a given estimation fidelity $`G=(1+g/d)/(d+1)`$ is defined by:
$$\widehat{A}_r=\sqrt{\frac{g}{d}}|r1r1|+\sqrt{\frac{dg}{d(d1)}}(\widehat{𝟙}|r1r1|)$$
(31)
where the index $`r`$ runs from $`1`$ to $`d`$, and the projectors $`|r1r1|`$ are constructed using any orthonormal basis. This confirms the inequality (1) is indeed a tight one and cannot be further improved.
A simple transformation of Eq. (1) shows that the quantum mechanically allowed region for the fidelities $`F`$ and $`G`$ is bounded by a quadratic curve, which turns out to be a fragment of an ellipse given by the equation:
$$\begin{array}{ccc}(FF_0)^2+d^2(GG_0)^2\hfill & & \\ +2(d2)(FF_0)(GG_0)\hfill & =& \frac{d1}{(d+1)^2}\hfill \end{array}$$
(32)
with $`F_0=(d+2)/(2d+2)`$ and $`G_0=3/(2d+2)`$. The shape of the region for several values of $`d`$ is depicted in Fig. 1.
The balance between the operation and estimation fidelities derived in this Letter has interesting consequences in quantum teleportation based on nonmaximally entangled states. If two parties share a pure bipartite state of the Schmidt form $`|\text{tele}=_{k=0}^{d1}\mu _k|k|k`$, then the maximum teleportation fidelity attainable using this state is given by :
$$F_{\text{tele}}=\frac{1+\left(_{k=0}^{d1}\mu _k\right)^2}{d+1}.$$
(33)
Furthermore, for a nonmaximally entangled state the measurement performed during the teleportation protocol reveals some information on the teleported state. This information can be converted into an estimate for the initial state, whose maximum average fidelity has been shown to equal :
$$G_{\text{tele}}=\frac{1+\mu _0^2}{d+1}$$
(34)
where $`\mu _0`$ denotes the largest Schmidt coefficient for the state $`|\text{tele}`$. As the procedure of teleportation can be viewed as a special case of a quantum operation , the bound (1) applies as well to the pair of fidelities $`F_{\text{tele}}`$ and $`G_{\text{tele}}`$. Consequently, for a given teleportation fidelity $`F_{\text{tele}}`$, the maximum value of the estimation fidelity is achieved for the state $`|\text{tele}`$ satisfying the condition $`\mu _1=\mathrm{}=\mu _{d1}=\sqrt{(1\mu _0^2)/(d1)}`$. This condition defines a class of pure bipartite states which are optimal from the point of view of the trade-off between the teleportation fidelity and the estimation fidelity.
In conclusion, I have obtained a tight bound for the fidelities describing the quality of estimating the state of a single copy of a $`d`$-level particle, and the degree the initial state has to be changed during this operation. This result seems to be one of very few cases, when the trade-off between the information gain and the state disturbance can be derived in a closed analytical form.
I would like to acknowledge useful discussions with J. H. Eberly, C. A. Fuchs, N. Lütkenhaus, V. Vedral, and I. A. Walmsley. This research was partially supported by ARO–administered MURI grant No. DAAG-19-99-1-0125, NSF grant PHY-9415583, and KBN grant 2 P03B 089 16.
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# 1 Introduction
## 1 Introduction
Stochastic resonance(SR) is known as a phenomenon in which the presence of noise helps a nonlinear system in amplifying a weak (under barrier) signal . The features of systems exhibiting SR seem to be applicable to some natural systems such as sensory neurons, which are noisy and operate as threshold systems. Since SR produces an information-transmitting phenomenon, its positive role in the neuronal processes is to be revealed. In fact, a single neuron model described by the FitzHugh-Nagumo (FHN) equations exhibits SR behavior and this SR effect is found in the real sensory neurons located in the tail fan of crayfish . The FHN equations driven by white noise and an arbitrary aperiodic signal are also examined in the context of excitable systems with threshold dynamics . One can expect that the SR effect will be more pronounced in an ensemble of systems than in a single system. In view of a collective response of globally coupled bistable systems to periodic forcing, a neural network with dynamics of the Hopfield type is studied, under the assumption that white noise and the periodic signal are identical for all neurons .
Given the basic three ingredients, that is, a form of threshold, a source of noise and a weak input signal, SR can generally be observed in a large variety of systems. Considering the fact that deterministic chaos resembles the feature of noise and provides a source of fluctuation, we have a natural question whether SR-like behavior can be observed in deterministic dynamical systems in the absence of noise. Two different approaches to this problem are known. One way is to substitute the stochastic noise by a chaotic source. This situation in which the chaos is supplied as an additive noise, resembles the conventional setup for SR and yields SR-like enhancement as expected . The other approach is to use the intrinsic chaotic dynamics of a nonlinear map. No external source is necessary to provide the randomness. This method generates a sort of activated hopping process which is then synchronized by a weak periodic signal .
Until now, diverse types of chaos have been confirmed at several hierarchical levels in the real neural systems from single cells to cortical networks (e.g. ionic channels, spike trains from cells, EEG) . By producing chaos as effective noise spontaneously, biological systems may enhance their functions through signal amplification. This scenario is likely to have occurred even in the associative memory dynamics. The chaotic neural network model is known as a framework beyond the Hopfield neural network with only equilibrium point attractors. Its dynamic retrieving and learning features have been studied . In this paper, we set up a signal-driven scheme of the chaotic neural network with the coupling constants corresponding to certain information, and investigate the SR-like effects under its deterministic dynamics, comparing with the conventional case of the Hopfield network with stochastic noise.
## 2 Models
For N neurons connected by synaptic couplings $`w_{ij}`$ with $`w_{ii}=0`$, the system (signal-driven scheme) is described by
$`X_i(t+1)=f(h_i(t)+S_i(t)),`$ (1)
where $`X_i`$: output of neuron $`i(1X_i1)`$, $`w_{ij}`$: synaptic weight from neuron $`j`$ to neuron $`i`$, $`f`$: output function defined by $`f(y)=tanh(y/2\epsilon )`$ with the steepness parameter $`\epsilon `$, $`h_i`$: internal potential, $`S_i`$: contribution of external input signal. This form is a simple and possible incorporation of stimuli as the changes of neuronal activity. In the case of chaotic neural network (CNN) , the internal potential is given by
$`h_i(t)`$ $`=`$ $`\eta _i(t)+\zeta _i(t),`$ (2)
$`\eta _i(t)`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{N}{}}}w_{ij}{\displaystyle \underset{d=0}{\overset{t}{}}}k_f^dX_j(td),`$ (3)
$`\zeta _i(t)`$ $`=`$ $`\alpha {\displaystyle \underset{d=0}{\overset{t}{}}}k_r^dX_i(td)\theta _i,`$ (4)
where $`\theta _i`$: threshold of neuron $`i,k_f(k_r`$): decay factor for the feedback (refractoriness) $`(0k_f,k_r<1),\alpha `$: refractory scaling parameter. Owing to the exponentially decaying form of the past influence, the dynamics of $`\{\eta _i\}`$ and $`\{\zeta _i\}`$ can be described as follows :
$`\eta _i(t)`$ $`=`$ $`k_f\eta _i(t1)+{\displaystyle \underset{j=1}{\overset{N}{}}}w_{ij}X_j(t),`$ (5)
$`\zeta _i(t)`$ $`=`$ $`k_r\zeta _i(t1)\alpha X_i(t)\theta _i(1k_r).`$ (6)
When $`\alpha =k_f=k_r=0`$, the network corresponds to the conventional discrete-time Hopfield network (we call the Hopfield network point (HNP)):
$`X_i(t+1)=f({\displaystyle \underset{j=1}{\overset{N}{}}}w_{ij}X_j(t)\theta _i).`$ (7)
We also look into the case that stochastic fluctuations are attached to HNP in Eq.(7):
$`h_i(t)={\displaystyle \underset{j=1}{\overset{N}{}}}w_{ij}X_j(t)\theta _i+F_i(t),`$ (8)
where $`F_i(t)`$ is a neuron-independent Gaussian white noise defined by $`<F_i(t)>=0`$ and $`<F_i(t)F_j(t^{})>=D^2\delta _{t,t^{}}\delta _{i,j}`$, and $`D`$ is the noise intensity parameter.We call Eq.(8) a stochastic neural network(SNN).
The synaptic configuration $`\{w_{ij}\}`$ is determined by storing pattern information in the network as minima of the computational energy :
$`E={\displaystyle \frac{1}{2}}{\displaystyle \underset{ij}{}}w_{ij}X_iX_j`$ (9)
at HNP. This is done by using a local iterative learning rule for $`p`$ patterns $`\{\xi _i^\mu \}(\xi _1^\mu ,\mathrm{},\xi _N^\mu )`$, $`(\mu =1,\mathrm{},p`$; $`\xi _i^\mu =+1or1)`$ in the following form :
$`w_{ij}^{new}=w_{ij}^{old}+{\displaystyle \underset{\mu }{}}\delta w_{ij}^\mu `$ (10)
with
$`\delta w_{ij}^\mu ={\displaystyle \frac{1}{N}}\theta (1\gamma _i^\mu )\xi _i^\mu \xi _j^\mu ,`$ (11)
where $`\gamma _i^\mu \xi _i^\mu _{j=1}^Nw_{ij}\xi _j^\mu `$ and $`\theta (h)`$ is the unit step function.
## 3 Simulations and Results
To carry out computational experiments, we consider a network with $`N=156`$, $`\{\theta _i\}=0`$ and $`\epsilon =0.015`$(unless otherwise stated), and use non-orthogonal 20 random patterns $`R^1R^{20}`$ as a set of external signal: $`\{S_i\}=s\{\xi _i^\mu \}(\mu =1,\mathrm{},20)`$. $`s`$ is the strength factor of signal. $`\{\xi _i^\mu \}`$ ($`i=1,\mathrm{},N`$) is represented by $`12\times 13`$ binary data and a half of $`\xi _i^\mu `$’s has $`+1`$ and the other half has $`1`$. $`10`$ patterns $`R^1R^{10}`$ are stored with the above learning rule of Eq.(10) ($`p=10`$), and then the corresponding multi-stable landscape is made on the network dynamics.
As the temporal input signal $`\{S_i(t)\}`$, we take two kinds of signals. One is a random train composed of stored patterns $`R^1R^{10}`$ (e.g. $`[R^3R^7R^2\mathrm{}]`$) with a duration $`T_I`$ for every pattern and the other is a similar sequence of non-stored patterns $`R^{11}R^{20}`$. Figure $`1`$ shows typical examples of the results of temporal behaviors of the neural networks to input signal $`\{S_i(t)\}`$. An input random train of stored patterns ($`T_I=100`$) is given in Fig.1(a) by the relative overlap with pattern $`R^4`$:
$`m_I^4(t)={\displaystyle \frac{1}{N}}{\displaystyle \underset{i=1}{\overset{N}{}}}\stackrel{~}{S}_i(t)\xi _i^4,\stackrel{~}{S}_i=S_i/s.`$ (12)
Responses of CNN ($`k_f=0.1`$, $`k_r=0.7`$) and SNN to this signal with strength $`s=0.5`$ are shown in Fig.1(b) and (c) by
$`m_0^4(t)={\displaystyle \frac{1}{N}}{\displaystyle \underset{i=1}{\overset{N}{}}}X_i(t)\xi _i^4.`$ (13)
In CNN \[Fig.1(b)\], as we can see from three cases where $`\alpha `$ is $`0.125`$, $`0.375`$ and $`0.625`$, the performance of output response is largely affected by its refractory scale parameter $`\alpha `$. The response fits well the input signal when $`\alpha `$ is $`0.375`$. From Fig.1(c), similar behaviors are observed for SNN against the noise intensity parameter $`D(=0.25,0.75`$ and $`1.5)`$, except for the appearance of noise-driven fluctuations.
To evaluate the coherence between the signal {$`S_i(t)`$} and the response {$`X_i(t)`$}, we introduce the correlation coefficient $`r`$ and the discrimination efficiency $`n`$. These quantities are defined as
$`r={\displaystyle \frac{\overline{Dev(m_I^4)Dev(m_0^4)}}{[\overline{Dev(m_I^4)^2}]^{1/2}[\overline{Dev(m_0^4)^2}]^{1/2}}}`$ (14)
and
$`n=\overline{{\displaystyle \frac{1}{N}}{\displaystyle \underset{i=1}{\overset{N}{}}}\stackrel{~}{S}_i(t)X_i(t)},\stackrel{~}{S}_i=S_i/s,`$ (15)
where $`Dev(Y)`$ is the deviation $`Y\overline{Y}`$ and the overbar denotes an average over time. $`n`$ naively means the rate of information transfer from stimulus to response.
The numerical results of $`r`$ and $`n`$ plotted against refractory factor $`\alpha `$ are given in Fig.2, calculated using the temporal data at every different $`\alpha `$ values. The results for $`\{S_i(t)\}`$ of stored patterns $`R^1R^{10}`$ are in contrast with the results for $`\{S_i(t)\}`$ of non-stored patterns $`R^{11}R^{20}`$ in the range of $`\alpha `$ larger than about $`0.2`$ wherein the response is hopping out and into a well (minimum of the energy) corresponding to a stored pattern state. This movability causes an increased coherence with stored pattern stimuli. Under the same conditions of input signal, dependences of $`r`$ and $`n`$ on the noise intensity $`D`$ in SNN are examined as shown in Fig.3. Between stored and non-stored pattern signals, a quite difference like in CNN is also confirmed. Comparing the results for stored pattern signal, we can see that in CNN there appears a projected plateau for both $`r`$ and $`n`$, and in this flat region both of their values are kept very close to 1. Contrary to this, in SNN $`r`$ and $`n`$ have gradual variations with $`D`$ and $`n`$ degrades faster than $`r`$, which seems to be consistent with conventional SR. This comparison indicates that CNN can cause SR-like phenomena with high performance which cannot be attained in SNN.
In CNN, the stronger the signal strength $`s`$ is, this $`\alpha `$-range of the flat plateau becomes wider. Conversely, this range becomes narrower for weaker $`s`$. We illustrate this effect as coherent $`\alpha `$-range versus signal strength $`s`$ in Fig.4, together with the relationship of maximum Lyapunov exponent $`\lambda _1`$ to $`\alpha `$ when there exists no signal forcing. As $`s`$ decreases, the $`\alpha `$-range ($`r0.9`$) becomes narrow and at last disappears. Then the $`\alpha `$ value coincides with $`\alpha _{}`$ at the sudden rise point of the Lyapunov exponent $`\lambda _1`$. This fact tells that sensitive and flexible responses happen around the boundary between order and disorder, in other words, the edge of chaos . At $`\alpha <\alpha _{}`$, the refractory term ($`\alpha X_i`$) in Eq.(6) makes a well (minimum of the energy) shallow in effect and helps the input signal drive the network state. On the other hand, at $`\alpha >\alpha _{}`$ the chaotic attractor is driven by the input signal and stabilized to the corresponding network state.
Figure 5 shows return maps of the internal potential $`h_i`$ for a neuron ($`i=12`$) in CNN with no signal forcing ($`s=0`$) when $`\alpha >\alpha _{}`$. In the case of $`\alpha =58/128`$, the trajectories of $`h_i`$ are attracted into the region of square ($`0.50.5`$) such that controllable by the input signal $`s=0.5`$. When $`\alpha =96/128`$ (a deeper chaotic state), however, the trajectories are attracted outward so that chaos cannot be suppressed by the input signal $`s=0.5`$.
We have investigated all the above phenomena in other conditions of the parameters and the input signals, and have found similar results and overall tendency.
## 4 Conclusion
We have shown that the chaotic neural network can enhance weak subthreshold signals and have higher coherence abilities between stimulus and response than those attained by the conventional stochastic neural network model. The high coherent response is found to arise around the edge of chaos. This implies that some of SR phenomena may be realized by the inherent properties of deterministic nonlinear systems without any external noise. Analytical study to explain these results will be important in our future work. The coherent response concept is expected to be related to researches in cognitive neuroscience from dynamical system viewpoints .
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# Acknowledgements
## Acknowledgements
We would like to thank G. Arutyunov for valuable discussions. We would also like to thank K. Intriligator for his very illuminating remarks after the first version of the manuscript. R. M. and A. C. P. are supported by the Alexander von Humboldt Foundation.
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# Wilson loop via AdS/CFT duality11footnote 1Talk given by S. Förste.
## 1 Introduction
One important ingredient of the string dualities is the twofold description of D-branes as solitonic supergravity solutions and as submanifolds of spacetime where open strings may end. The second description leads to a gauge field theory on the world-volume of D-branes. Based on this general idea is Maldacena’s conjecture relating superconformal field theories to supergravity (or superstrings) living in a higher dimensional space with boundary. (For a review containing also a comprehensive list of references, see .) For example, $`𝒩=4`$ supersymmetric $`SU(N)`$ Yang Mills theory in four dimensions is dual to the type IIB string theory on $`adS_5\times S^5`$ space. The radii of the $`S^5`$ and the $`adS_5`$ are equal and related to the Yang-Mills coupling via $`R^2/\alpha ^{}=\sqrt{4\pi g_{YM}^2N}`$. (The string coupling $`g`$ is $`g=g_{YM}^2`$.) For supergravity to be a good effective description, we need the radius of curvature to be large and also the string coupling to be small. This means that we need $`g_s`$ small but $`g_sN`$ large. The latter is however the ’t Hooft coupling constant in the large $`N`$ field theory which is thus strongly coupled. Maldacena’s conjecture provides a new possibility to gain insight into strongly coupled gauge theory by studying weakly coupled string theory. As an application, Wilson loops have been computed in and . The string configuration for a quark-antiquark pair separated by a distance $`L`$, is a long string on $`adS_5\times S^5`$, the ends of which are restricted to the (four dimensional) boundary of $`adS_5`$, where they are at a distance $`L`$ apart. The expectation value of the Wilson loop is then given by the effective energy of the string. We will review this computation in the next section. In the third section leading corrections are discussed. As a further application of our techniques we briefly review in section four membrany corrections to the Wilson surface in M5 theory. We conclude with a short summary.
## 2 Review
The dual description of $`𝒩=4`$ super Yang-Mills theory is a type IIB string living in $`adS_5\times S^5`$. In particular the target space metric ($`G_{MN}`$) is
$`ds^2`$ $`=`$ $`R^2[U^2(\left(dx^0\right)^2+\left(dx^1\right)^2+\left(dx^2\right)^2`$ (1)
$`+\left(dx^3\right)^2)+{\displaystyle \frac{dU^2}{U^2}}+d\mathrm{\Omega }_5^2].`$
Compared to we have rescaled $`UR^2U`$, and also set $`\alpha ^{}=1`$. In addition there is a constant dilaton and $`N`$ units of RR-4-form flux through $`S^5`$. In order to compute the Wilson loop one minimizes the Nambu-Goto action <sup>2</sup><sup>2</sup>2After switching off the world-sheet fermions, the type IIB action reduces to the Nambu-Goto action.
$$S_{NG}=\frac{d^2\sigma }{2\pi }\sqrt{detG_{MN}_iX^M_jX^N}$$
(2)
with the boundary condition that the ends of the string are separated by a distance $`L`$ at $`U=\mathrm{}`$. We work in the static gauge ($`X^0=\tau `$, $`X^1=\sigma `$) and restrict to the case that all coordinates but $`U`$ are constant. The radial coordinate $`U`$ depends on $`\sigma `$. Then the (implicit) solution reads<sup>3</sup><sup>3</sup>3In the following we will just take the upper sign with the understanding that the square root stands for both branches.
$$_\sigma U=\pm \frac{U^2}{U_0^2}\sqrt{U^4U_0^4},$$
(3)
where
$$U_0=\frac{\left(2\pi \right)^{\frac{3}{2}}}{\mathrm{\Gamma }\left(\frac{1}{4}\right)^2L}$$
(4)
is the minimal $`U`$-value the string obtains. The energy of the quark-antiquark pair is the length of the geodesic (open string) connecting them. One finds (after subtracting an $`L`$ independent infinite contribution from the self energy of the heavy quarks)
$$E=\frac{4\pi \sqrt{g_{YM}^2N}}{\mathrm{\Gamma }\left(1/4\right)^4L}.$$
(5)
This strong coupling result differs from the perturbative field theory computation ($`g_{YM}^2N`$ small), which predicts a Coulomb law with $`Eg_{YM}^2N/L`$. In general the numerator is a function of $`g_{YM}^2N`$ which interpolates between these two forms, and hence there ought to be corrections to (5) which is the result of a classical supergravity computation. Since $`adS_5\times S^5`$ is an exact string background , the first correction comes from the fluctuations of the superstring ($`R^2/\alpha ^{}`$ correction). In this talk we will report on work in this direction. Corrections due to string fluctuations have been discussed in , and subsequently in . Ref. considered corrections to the field theoretical result.
## 3 Fluctuations
The quantum theory of type IIB string in this background is described by the action in . It is a Green-Schwarz type sigma model action on a target supercoset. The usual sigma model expansion in $`R^2/\alpha ^{}`$ results in a power series in $`1/\sqrt{Ng_{YM}^2}`$. Since conformal invariance prevents the appearance of a new scale these corrections are not expected to change the $`1/L`$ dependence of $`E`$ on dimensional grounds. However they can modify the ‘Coulomb charge’.
The leading order correction is obtained by expanding around the classical configuration (3) to second order in fluctuations. We parameterize the bosonic fluctuations by normal coordinates: $`\xi ^a`$ on $`adS_5\times S^5`$, (here $`a=0,\mathrm{},4;5,\mathrm{},9`$ are local (flat) Lorentz indices; $`\xi ^4`$ is in the $`U`$ direction). Using the normal coordinate expansion one ensures that the functional measure for the fluctuations is translation invariant. This takes care of potential problems due to the curved target space. Ref. has an extensive discussion on additional subtleties in the functional measures due to a curved world sheet.
At second order, the bosonic fluctuations in $`adS_5`$ and $`S^5`$ directions and the fermionic fluctuations decouple. Before writing their action, we define the combinations
$`\xi ^{}`$ $`=`$ $`{\displaystyle \frac{U_0^2}{U^2}}\xi ^1+{\displaystyle \frac{\sqrt{U^4U_0^4}}{U^2}}\xi ^4,`$
$`\xi ^{}`$ $`=`$ $`{\displaystyle \frac{\sqrt{U^4U_0^4}}{U^2}}\xi ^1+{\displaystyle \frac{U_0^2}{U^2}}\xi ^4,`$ (6)
which parameterize fluctuations along the longitudinal and perpendicular directions of the background string in the one-four plane. Now the $`adS_5`$ part of the action becomes
$`S_{adS}^{(2)}`$ $`=`$ $`{\displaystyle \frac{1}{4\pi }}{\displaystyle }d^2\sigma \sqrt{h}[h^{ij}\left({\displaystyle \underset{a=2,3,}{}}_i\xi ^a_j\xi ^a\right)`$ (7)
$`+`$ $`2\left(\xi ^2\right)^2+2\left(\xi ^3\right)^2+2{\displaystyle \frac{U^4U_0^4}{U^4}}\left(\xi ^{}\right)^2]`$
where $`h_{ij}`$ is (up to a factor of $`R^2`$) the classical induced world-sheet metric<sup>4</sup><sup>4</sup>4For our purpose it is more convenient to work with the induced metric rather than the standard (conformally) flat one on the world-sheet.
$$ds^2=U^2\left(d\sigma ^0\right)^2+\frac{U^6}{U_0^4}\left(d\sigma ^1\right)^2.$$
(8)
Observe that $`\xi ^0`$ and $`\xi ^{}`$ do not appear in (7) (total derivatives have been dropped). Hence a natural choice to fix world-sheet diffeomorphisms is
$$\xi ^0=\xi ^{}=0.$$
(9)
The action quadratic in fluctuations in $`S^5`$ directions comes out to be
$$S_{S^5}^{(2)}=\frac{1}{4\pi }d^2\sigma \sqrt{h}h^{ij}\underset{a=5}{\overset{9}{}}_i\xi ^a_j\xi ^a.$$
(10)
The fermionic part of the action is given by plugging in the background (4) into the action of and keeping terms quadratic in fermions. This action has local fermionic $`\kappa `$-symmetry which has to be fixed for the one-loop calculation. There is a proposal in the literature to this end. For our purpose, however, it turns out that the following choice is most convenient.<sup>5</sup><sup>5</sup>5We will comment on a different choice below. We will set (in the notation of )
$$\left(\gamma ^{}\right)_\beta ^\alpha \theta ^{1\beta \beta ^{}}=0,\left(\gamma ^+\right)_\beta ^\alpha \theta ^{2\beta \beta ^{}}=0$$
(11)
where $`\gamma ^\pm =\gamma ^0\pm \gamma ^{}`$ with $`\gamma ^{}=\frac{U_0^2}{U^2}\gamma ^1+\frac{\sqrt{U^4U_0^4}}{U^2}\gamma ^4`$ (Cf. (6)). With this choice the target space spinors ‘metamorphose’ into world-sheet spinors, and the action is found to simplify substantially. The corresponding equations of motion are most compactly written for the ‘two-component’ world-sheet spinors $`\left(\begin{array}{c}\theta ^1\\ \theta ^2\end{array}\right)`$:
$$\left[i\left(\rho ^m_m\right)_\beta ^\alpha \delta _\beta ^\alpha \rho ^3\right]\left(\begin{array}{c}\theta ^{1\beta \alpha ^{}}\\ \theta ^{2\beta \alpha ^{}}\end{array}\right)=0.$$
(12)
The notation needs explanation. Firstly, the covariant derivatives act as
$`\left(_\pm \theta ^I\right)^{\alpha \alpha ^{}}`$ $`=`$ (13)
$`\left[\delta _\beta ^\alpha \left(_\pm \pm {\displaystyle \frac{\omega _\pm }{2}}\right)+\left(A_\pm \right)_\beta ^\alpha \right]\theta ^{I\beta \alpha ^{}},`$
where the tangent space derivatives
$$_\pm =\frac{1}{U}_\tau \pm \frac{U_0^2}{U^3}_\sigma $$
(14)
are defined with the help of a (inverse) zweibein $`ϵ_m`$ of the metric (8), $`\omega _m^{01}=ϵ_m^\tau \omega _\tau ^{01}`$ being the corresponding spin connection. There is an additional gauge field
$$A_\pm =\pm \frac{U_0^2}{U^2}\gamma ^{14}$$
(15)
for local rotations in the tangent one-four-plane. Finally, the matrices
$$\rho ^+=\left(\begin{array}{cc}0& 0\\ \gamma ^0& 0\end{array}\right),\rho ^{}=\left(\begin{array}{cc}0& \gamma ^0\\ 0& 0\end{array}\right)$$
(16)
satisfy a two dimensional Clifford algebra, and $`\rho ^3=[\rho ^+,\rho ^{}]`$. The fermionic action is easily inferred from the equations of motion (12).
Collecting our results (7), (10) and (12) one can write a formal expression for the 1-loop contribution to the effective action as a ratio of determinants of two dimensional generalized Laplace operators. These determinants suffer from divergences and can be regularized by, say, the heat kernel technique. The quadratic divergences cancel, but naively a logarithmic divergence remains, which may be absorbed in the (infinite) mass of the external quarks. As unsatisfactory as it may be, it does not affect our result for the correction to the Coulomb charge. More recently the authors of have argued that this divergence should, as in flat space, actually cancel. As far as the corrections to the Coulomb charge are concerned the results of and ours are actually equivalent. In the following we demonstrate that the apparently different expressions for the fermionic operators in and are related to each other by a local Lorentz rotation<sup>6</sup><sup>6</sup>6From the sigma model point of view this is just a field redefinition with unit Jacobian.. To this end, define
$$\theta ^I=S\psi ^I,$$
(17)
where we suppressed target space spinor indices. For the matrix $`S`$ we choose the one given in Ref., Section 6.3,
$$S=\mathrm{cos}\frac{\alpha }{2}\mathrm{sin}\frac{\alpha }{2}\gamma ^{14},$$
(18)
where
$`\mathrm{cos}\alpha `$ $`=`$ $`{\displaystyle \frac{U_0^2}{U^2}},`$ (19)
$`\mathrm{sin}\alpha `$ $`=`$ $`{\displaystyle \frac{\sqrt{U^4U_0^4}}{U^2}},`$ (20)
implying that
$$_\sigma \alpha =2U.$$
(21)
The Dirac operator for $`\psi `$ is given by conjugation by $`S`$ from the one for $`\theta `$. Using the fact that $`S`$ commutes with the $`\rho _m`$ and using ((14), (15))
$$S^1A_\pm S=S^1_\pm S=\pm \frac{U_0^2}{U^2}\gamma ^{14},$$
(22)
we find that the Dirac operator acting on the fields $`\psi ^I`$ (17) takes again the form (12), but now the connection $`A_\pm `$ has been gauged away. It is also easy to show that for the redefined spinors $`\psi ^I`$ the kappa-fixing condition takes the form (11), but with
$$\gamma ^{}S^1\gamma ^{}S=\gamma ^1,$$
(23)
which is the same as (6.35) in . This shows that the results in are equivalent to ours.
## 4 Wilson surface in M5 theory
The Maldacena conjecture also applies to the case of M-theory living on $`adS_7\times S^4`$ being dual to the field theory on a stack of M5 branes. The metric of $`adS_7\times S^4`$ is
$$ds^2=l_p^2R^2\left[U^2\eta _{\mu \nu }dx^\mu dx^\nu +4\frac{dU^2}{U^2}+d\mathrm{\Omega }_4^2\right],$$
(24)
where we have rescaled the five-brane coordinates $`x^\mu `$ by $`R^{3/2}`$ as compared to, $`d\mathrm{\Omega }_4^2`$ is the metric on $`S^4`$. In the limit that the eleven dimensional Planck length $`l_p`$ goes to zero M-theory on $`adS_7\times S^4`$ is conjectured to be dual to the field theory on $`N`$ M5-branes, where the adS radius and the number of five-branes are related,
$$R=\left(\pi N\right)^{\frac{1}{3}}.$$
(25)
Higher curvature corrections will be small as long as $`N`$ is taken to be large. In analogy to the previously discussed Wilson loop one can study a Wilson surface in M5-theory. The set up is a membrane extending along the $`x^{0,1,2}`$ and the $`U`$ direction ending in two parallel lines separated by a distance $`L`$ at the boundary of $`adS_7`$. In the following we will recall this set-up (with slightly changed conventions) and thereafter study corrections due to fluctuations of the membrane. This will be a brief summary of the work presented in. The classical background corresponding to the Wilson surface is obtained by minimizing the world volume of the membrane
$$S=\frac{1}{2\pi }\sqrt{det\left(G_{MN}_aX^M_bX^N\right)}$$
(26)
with appropriate boundary conditions. The indices $`M,N`$ label the eleven target space coordinates and $`a,b`$ are world volume coordinates ($`\tau ,\sigma ,\varphi `$). By choosing the static gauge $`X^0=\tau `$, $`x^1=\sigma `$, $`X^2=\varphi `$ and assuming all other coordinates but $`U=U\left(\sigma \right)`$ to be constant one finds the solution,
$$_\sigma U=\pm \frac{U^2}{2U_0^3}\sqrt{U^6U_0^6}.$$
(27)
Requiring that the membrane ends in two parallel strings at distance $`L`$ determines the integration constant
$$U_0=\frac{2}{3L}B(\frac{2}{3},\frac{1}{2}),$$
(28)
where $`B`$ denotes Euler’s beta-function. Computing the vacuum energy density of this configuration one obtains (again after subtracting an $`L`$ independent infinite contribution to the self-energy of the strings) the potential between the two strings living on the M5-branes,
$$\epsilon _{pot}=\frac{2R^3}{27\pi }B(\frac{2}{3},\frac{1}{2})^3\frac{1}{L^2}.$$
(29)
This result is reliable for large $`N`$ where the geometry is not corrected and the classical approximation dominates. In it was argued that there are no corrections to the geometry due to finite $`N`$. Another potential source for corrections are fluctuations of the membrane around its classical background described above. Again we expand in normal coordinates and obtain the action quadratic in bosonic fluctuations. We trade the fluctuations in one- and six-direction<sup>7</sup><sup>7</sup>7The fluctuations are labeled by Lorentz indices; $`\xi ^6`$ is in the $`U`$ direction. for normal and parallel ones,
$`\xi ^{}`$ $`=`$ $`{\displaystyle \frac{U_0^3}{U^3}}\xi ^1+{\displaystyle \frac{\sqrt{U^6U_0^6}}{U^3}}\xi ^6`$ (30)
$`\xi ^{}`$ $`=`$ $`{\displaystyle \frac{\sqrt{U^6U_0^6}}{U^3}}\xi ^1+{\displaystyle \frac{U_0^3}{U^3}}\xi ^6.`$ (31)
The contribution from the $`adS`$ directions is
$`S_{adS}^{(2)}`$ $`=`$ $`{\displaystyle \frac{1}{4\pi }}{\displaystyle }d^3\sigma \sqrt{h}[h^{ij}{\displaystyle \underset{a=3}{\overset{5,}{}}}_i\xi ^a_j\xi ^a`$ (32)
$`+{\displaystyle \frac{3}{4}}{\displaystyle \underset{a=3}{\overset{5}{}}}\left(\xi ^a\right)^2+({\displaystyle \frac{9}{4}}R^{(3)})\left(\xi ^{}\right)^2]`$
where $`h_{ij}`$ is (up to a factor of $`R^2`$) the induced metric,
$$ds^2=U^2d\tau ^2+\frac{U^8}{U_0^6}d\sigma ^2+U^2d\varphi ^2,$$
(33)
and $`R^{(3)}`$ is the corresponding scalar curvature,
$$R^{(3)}=\frac{3}{2}\frac{U^6+U_0^6}{U^6}.$$
(34)
Again the three longitudinal directions $`0,2,`$ drop out of the action and we gauge them to zero. The bosonic fluctuations in $`S^4`$ direction have a simple action,
$$S_S^{(2)}=\frac{1}{4\pi }d^3\sigma \sqrt{h}h^{ij}\underset{a=7}{\overset{10}{}}_i\xi ^a_j\xi ^a.$$
(35)
In order to discuss the fermionic fluctuations we take the $`\kappa `$ symmetric action of. For our background the part quadratic in fermions consists out of terms containing $`\mathrm{\Gamma }^a`$ ($`a=0,\mathrm{},6`$), where $`\mathrm{\Gamma }^a`$ is an eleven dimensional gamma matrix. These can be written as $`\mathrm{\Gamma }^a=\gamma ^a\gamma ^5^{}`$ where the lower case gammas are gamma matrices in the tangent spaces of $`adS_7`$ and $`S^4`$, respectively. We split the 32-component spinors into two 16-component spinors ($`\theta ^1`$, $`\theta ^2`$) according to their eigenvalue with respect to $`\gamma ^5^{}`$. ($`\gamma ^5^{}\theta ^I=()^I\theta ^I`$.) A convenient kappa-fixing condition turns out be
$$\left(1(1)^I\gamma ^{02}\right)\theta ^I=0,$$
(36)
where now $`\gamma ^{}=\frac{U_0^3}{U^3}\gamma ^1+\frac{\sqrt{U^6U_0^6}}{U^3}\gamma ^6`$ . Imposing the kappa-fixing condition we find that the equations of motion for e.g. $`\theta ^1`$ are
$$\rho ^ae_a^i\left(_i+\frac{1}{4}\omega _i^{bc}\rho _{bc}+A_i\right)\theta ^1=\frac{3}{4}\theta ^1,$$
(37)
where $`e_a^i`$ and $`\omega _i^{bc}`$ are the vielbeine and spin-connections computed from (33) (for details see ). The matrices $`\rho `$ are
$$\rho ^0=\gamma ^0,\rho ^1=\gamma ^{02},\rho ^2=\gamma ^2$$
(38)
satisfying a 3d Clifford algebra. The field $`A_\sigma =\frac{3U}{4}\gamma ^{16}`$ is a background value for a gauge connection belonging to local rotations in the 1-6 plane. (For $`\theta ^2`$ one obtains the same result but with a minus sign in the definition of $`\rho ^1`$.) Note that the Dirac operator appearing in (37) is covariant from a world volume perspective. Collecting the results (32), (35) and (37) one can express the contribution to the energy density in terms of determinants of covariant operators. These can be analyzed using for example the heat-kernel method. In difference to the previously discussed string case one finds divergent contributions as well to the self-energy density as to the potential energy density. It would be interesting to investigate whether one can extend the arguments of to cancel those divergences. (Since the discussion in is quite heavily based on conformal invariance and 2d calculus this is not straightforward.) Finally, let us point out that also in the membrane case one can gauge away the connection appearing in (37). This can be achieved by a field redefinition $`\theta ^I=S\psi ^I`$ with
$$S=\mathrm{cos}\frac{\alpha }{2}\mathrm{sin}\frac{\alpha }{2}\gamma ^{16},$$
(39)
where
$`\mathrm{cos}\alpha `$ $`=`$ $`{\displaystyle \frac{U_0^3}{U^3}},`$
$`\mathrm{sin}\alpha `$ $`=`$ $`{\displaystyle \frac{\sqrt{U^6U_0^6}}{U^3}}.`$ (40)
The kappa-fixing condition is again (36) but with $`\gamma ^{}`$ replaced by $`\gamma ^1`$. This coincides with the kappa-fixing condition advertised in.
## 5 Summary
In this talk we presented techniques for computing stringy corrections to the Wilson loop in $`𝒩=4`$ supersymmetric Yang-Mills theory. The final result can be expressed in terms of determinants of 2d covariant operators. A result for the leading correction to the Coulomb charge in terms of a number is still missing (a rough estimate can be found in ). We commented on the relation between our results and those obtained in . In the end we reviewed the application of our techniques to the case of a Wilson surface in M5-theory. There the result is less satisfactory as divergences also affect the Coulomb charge.
Acknowledgments
This work was supported in part by GIF (German Israeli Foundation for Scientific Research) and the EC programs ERB-FMRX-CT-96-0045 and 96-0090. D.G. acknowledges the support of the Humboldt Foundation. S.F. and S.T. thank the organizers of the Paris meeting for creating a pleasant and stimulating atmosphere during the conference.
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# 1 Introduction
## 1 Introduction
One could say that the aim of noncommutative geometry is to generalize geometrical tools to “spaces whose coordinates fail to commute” , (see also , , ). One way to implement this program is to start with a given geometrical theory involving sets $`X`$ endowed with additional structures and formulate them algebraically by using suitable subalgebras of the algebra of complex valued functions over $`X`$. Then one extends parts of the theory to noncommutative algebras, which are thought of as functions over “noncommutative spaces”. Although much of the construction takes place at the algebraic level, it is necessary, in order to use the powerful machinery of functional analysis, to represent these algebras as operators on a Hilbert space. Accordingly, this can be seen from a physicist’s point of view, as analogous to quantum mechanics: one trades the commutative algebra of functions over phase space for a noncommutative algebra of operators acting on a Hilbert space. Most of the geometrical ideas of classical mathematics can be “quantized”. For example, topology can be formulated in terms of C-algebras, commutative C-algebras corresponding to locally compact spaces. Thus, noncommutative ones are referred to as “continuous functions over noncommutative locally compact spaces”. In the compact case (i.e. when the algebra is unital), one can further define noncommutative vector bundles as finitely generated projective modules over a given unital C-algebra which plays the role of functions over the base space. When this algebra is commutative, Serre-Swan’s theorem asserts that these modules correspond to module of sections of vector bundles. Furthermore, methods of differential topology are also available within the realm of cyclic and Hochschild cohomology and this leads, via the coupling of the former with K-theory, to quantities that are stable under deformation and that generalize topological invariants, like, for instance, winding numbers and topological charges.
Noncommutative geometry has already proved to be useful in understanding various physical phenomena, like the integral quantum Hall effect or the classical aspect of the Higgs sector of the standard model (see for a review). Recent developments and also indicate that it is helpful in string theory. These last developments involve Yang-Mills fields defined on noncommutative spaces that fit into a broad formalism for gauge fields in noncommutative geometry which allows one to define connections, their curvature or the associated Yang-Mills action while preserving most of their classical aspects. For instance, one can prove a topological bound for the Yang-Mills action in dimension 4 . Also, one can construct a Chern-Simons type theory and interpret its behavior under large gauge transformations as a coupling between cyclic cohomology and K-theory .
In this report, we will be interested in constructing analogues of two dimensional non-linear $`\sigma `$-models within the noncommutative world. Since these models usually exhibit a very rich and easily accessible geometrical structure, we expect their noncommutative counterparts to be an ideal playground for a probe into the interplay between noncommutative geometry and field theory. This we shall try to exemplify by means of three different models: a continuous analogue of the Ising model which admits instantonic solutions, the analogue of the principal chiral model together with its infinite number of conserved currents and the noncommutative Wess-Zumino-Witten model together with its modified conformal invariance.
All ideas will be presented in a rather sketchy form and we refer to (for fields with values in finite spaces) and (for $`S^1`$-valued fields) for a detailed account.
## 2 General Aspects
In ordinary field theory, non-linear $`\sigma `$-models (see for a review) are field theories whose configuration space consists of maps $`X`$ from a Riemannian manifold $`\mathrm{\Sigma }`$ with metric $`g`$, which we assume to be compact and orientable, to an other Riemannian manifold $``$ whose metric we denoted by $`G`$. In the physics literature, these manifolds are called source and target space respectively. By using local coordinates, the action functional is defined as
$$S[X]=\frac{1}{2\pi }_\mathrm{\Sigma }\sqrt{g}g^{\mu \nu }G_{ij}(X)_\mu X^i_\nu X^j,$$
(1)
where as usual $`g=detg_{\mu \nu }`$ and $`g^{\mu \nu }`$ is the inverse of $`g_{\mu \nu }`$. When $`\mathrm{\Sigma }`$ is two dimensional, the action $`S`$ is conformally invariant, since a rescaling of the metric $`gge^\sigma `$, with $`\sigma `$ being any map from $`\mathrm{\Sigma }`$ to $``$, leaves it invariant. Accordingly, the action only depends on the conformal class of the metric and may be rewritten using a complex structure on $`\mathrm{\Sigma }`$ as
$$S[X]=\frac{i}{\pi }_\mathrm{\Sigma }G_{ij}(X)X^i\overline{}X^j$$
(2)
where $`=_zdz`$ and $`\overline{}=_{\overline{z}}d\overline{z}`$, $`z`$ being a suitable local complex coordinate.
Different choices of the source and target spaces lead to different field theories, some of them playing a major role in physics. Especially interesting are their applications to statistical field theory, and (supplemented by some supersymmetries) they are the basic building blocks of superstring theories.
From the mathematical point of view, the stationary points of the action functional (1) are harmonic maps from $`\mathrm{\Sigma }`$ to $``$ and describe the extremal surfaces embedded in $``$. Thus, a noncommutative generalization of the action functional of the non-linear $`\sigma `$-model should yield, as stationary points, noncommutative analogues of harmonic maps.
To generalize such a construction to the noncommutative case, we must dualize the previous picture and reformulate it in terms of the $``$-algebras $`𝒜`$ and $``$ of complex valued smooth functions defined respectively on $`\mathrm{\Sigma }`$ and $``$. Embeddings $`X`$ of $`\mathrm{\Sigma }`$ into $``$ are in one to one correspondence with $``$-algebra morphisms $`\pi `$ from $``$ to $`𝒜`$, the correspondence being simply $`\pi _X(f)=fX`$. Since this makes perfectly sense in the noncommutative case, we define our configuration space, for fixed, not necessarily commutative algebras $`𝒜`$ and $``$, as the space of all $``$-algebra morphisms from $``$ to $`𝒜`$.
The construction of the action functional is more tricky since it involves noncommutative generalizations of the conformal and Riemannian geometries. Following an idea of Connes and , the former can be understood within the framework of positive Hochschild cohomology. Without entering into details, one observes that, in the commutative case, the trilinear map on $`𝒜^3`$ defined by
$$\varphi (f_0,f_1,f_2)=\frac{i}{\pi }_\mathrm{\Sigma }f_0f_1\overline{}f_2$$
(3)
is an extremal element of the space of positive Hochschild cocycles that belongs to the cohomology class of the cyclic cocycle $`\psi `$ defined by
$$\psi (f_0,f_1,f_2)=\frac{i}{2\pi }_\mathrm{\Sigma }f_0𝑑f_1df_2.$$
(4)
Again, we refer to for the general definitions and we simply notice that (3)-(4) still make sense for a general noncommutative algebra $`𝒜`$.
Roughly speaking, one can say that $`\psi `$ allows to integrate 2-forms in dimension 2,
$$\frac{i}{2\pi }a_0𝑑a_1𝑑a_2=\psi (a_0,a_1,a_2)$$
(5)
so that it is a metric independent object, whereas $`\varphi `$ defines a suitable scalar product
$$a_0da_1,b_0db_1=\varphi (b_0^{}a_0,a_1,b_1^{})$$
(6)
on the space of 1-forms and thus depends on the conformal class of the metric. Furthermore, this scalar product is positive and invariant with respect to the action of the unitary elements of $`𝒜`$ on 1-forms, and its relation to the cyclic cocylic $`\psi `$ allows to prove various inequalities involving topological quantities.
Having such a cocycle $`\varphi `$, it is natural to compose it with a morphism $`\pi :𝒜`$ in order to obtain a positive cocycle on $``$ defined by $`\varphi _\pi =\varphi (\pi \pi \pi )`$. Since our goal is to build an action functional, which assigns a number to any morphism $`\pi `$, we have to evaluate the previous cocycle on a suitably chosen element of $`^3`$. Such an element is provided by the noncommutative analogue of the metric on the target, which we take simply as a positive element $`G=_ib_0^i\delta b_1^i\delta b_2^i`$ of the space of universal 2-forms $`\mathrm{\Omega }^2()`$. Thus
$$S[\pi ]=\varphi _\pi (G)$$
(7)
is well defined and positive and we take it as a noncommutative analogue of the action functional of the non linear $`\sigma `$-model. Of course we consider $`\pi `$ as the dynamical variable (the embedding) whereas $`\varphi `$ (the conformal structure on the source) and $`G`$ (the metric on the target) are background structures that have been fixed.
As an alternative, one could consider that only the metric $`G`$ on the target is a background field, since the morphism $`\pi :𝒜`$ allows to define the induced metric $`\pi _{}G`$ on the source as
$$\pi _{}G=\underset{i}{}\pi (b_0^i)\delta \pi (b_1^i)\delta \pi (b_2^i),$$
(8)
which is obviously a positive universal 2-form on $`𝒜`$. To such an object one can associate, by means of a variational problem (see and ), a positive Hochschild cocycle that stands for the conformal class of the induced metric. As a result, the critical points of the corresponding $`\sigma `$-model describe “minimally embedded surfaces” in the noncommutative space associated with $``$.
A scrupulous reader may be puzzled by such a formal and sketchy construction. However, in what follows we will mainly work out examples involving the noncommutative torus and only consider a fixed $`\varphi `$ and fixed metrics on the two target space we will consider (the circle and the two point space). Accordingly, $`\varphi `$ and $`G`$ could be replaced by their explicit expressions. Nevertheless, we think that it may be useful to have a general setting. In particular, one easily reconstructs ordinary $`\sigma `$-models with suitable choices of $`\varphi `$ and $`G`$.
## 3 Two points as a target space
### 3.1 A General Construction
The simplest example of a target space one can think of is that of a finite space made of two points, like in the Ising model. Of course, any continuous map from a connected surface to a discrete space is constant and the resulting (commutative) theory would be trivial. However, this is no longer true if the source is a noncommutative space and one has, in general, lots of such maps (i.e. algebra morphisms).
Let us first notice that the algebra $`=^2`$ of functions over a two point space is the unital algebra generated by a hermitian projection $`e,e^2=e^{}=e`$. Thus, any $``$-algebra morphism $`\pi `$ from $``$ to $`𝒜`$ is given by a hermitian projection $`p=\pi (e)`$ in $`𝒜`$. Choosing the metric $`G=\delta e\delta e`$ on the space of two points, the action functional simply reads
$$S[p]=\varphi (1,p,p),$$
(9)
where $`\varphi `$ is a given Hochschild cocycle standing for the conformal structure. Of course, one could choose other metrics on the two points space, but $`G`$ is more interesting since it will lead to a topological bound for the corresponding action. We shall prove this fact for the noncommutative torus, but the procedure is general and only uses the idea of positivity in Hochschild cohomology.
### 3.2 The noncommutative two torus as a source space
For the sake of completeness we recall the very basic aspects of the noncommutative torus and refer the reader to for a thorough survey. The algebra $`𝒜_\theta `$ of smooth functions on the noncommutative torus is the unital $``$-algebra made of power series of the form
$$a=\underset{m,n^2}{}a_{mn}U_1^mU_2^n,$$
(10)
with $`a_{mn}`$ a complex-valued Schwarz function on $`^2`$ that is, the sequence of complex numbers $`\{a_{mn},(m,n)^2\}`$ decreases rapidly at ‘infinity’. The two unitary elements $`U_1,U_2`$ have commutation relations
$$U_2U_1=e^{2\pi i\theta }U_1U_2.$$
(11)
On $`𝒜_\theta `$ there is a unique normalized positive definite trace which we shall unusually denote by an integral symbol $`:𝒜_\theta `$ and which is given by
$$(\underset{(m,n)^2}{}a_{mn}U_1^mU_2^n)=:a_{00}.$$
(12)
This trace is invariant under the action of the commutative torus $`T^2`$ on $`𝒜_\theta `$ whose infinitesimal form is determined by two commuting derivations $`_1,_2`$ acting by
$$_\mu (U_\nu )=2\pi i\delta _\mu ^\nu U_\nu ,\mu ,\nu =1,2.$$
(13)
The invariance just being the statement that $`_\mu (a)=0,\mu =1,2`$ for any $`a𝒜_\theta `$.
All the previous properties, even if elementary, turn out to be important for our construction since they allow us to use all tools of elementary calculus on a commutative torus. However one must bear in mind that, in order to develop a more general setting, one should work only with the corresponding cyclic and Hochschild cocycles that we shall now describe.
The cyclic 2-cocycle associated to the integration of 2-forms is simply given by
$$\psi (a_0,a_1,a_2)=\frac{i}{2\pi }ϵ_{\mu \nu }a_0_\mu a_1_\nu a_2,$$
(14)
where $`ϵ_{\mu \nu }`$ is the standard antisymmetric tensor. Its normalization ensures that for any hermitian projector $`p𝒜_\theta `$, the quantity $`\psi (p,p,p)`$ is an integer: it is indeed the index of a Fredholm operator.
Working with the standard Euclidean metric on the torus, the Hochschild cocycle $`\varphi `$ is
$$\varphi (a_0,a_1,a_2)=\frac{2}{\pi }a_0a_1\overline{}a_2$$
(15)
where the complex derivations $`=1/2(_1i_2)`$ and $`\overline{}=1/2(_1+i_2)`$ are combination of the previous derivations. Note that we consider $``$ and $`\overline{}`$ as maps with values in $`𝒜_\theta `$ and not in the bimodule of 1-forms. A construction of the cocycle (15) as the conformal class of the Euclidean metric on the torus can be found in and . We remark that one can also work with a general constant metric whose conformal class is parametrized by a complex number $`\tau `$ that belongs to the upper half plane.
Accordingly, the action functional for our non-linear $`\sigma `$-model reads
$$\varphi (1,p,p)=\frac{1}{2\pi }_\mu p_\mu p=\frac{1}{\pi }p_\mu p_\mu p,$$
(16)
the contraction with the Euclidean metric being understood.
As a subset of a topological vector space, the space $`𝒫_\theta `$ of all hermitian projectors of $`𝒜_\theta `$ comes equipped with a natural topology (in fact it is an infinite dimensional manifold) and we are interested in the study of the critical points of the action in a given connected component of $`𝒫_\theta `$. By carefully taking into account the non linear structure of the space $`𝒫_\theta `$, we get the field equations
$$p\mathrm{\Delta }(p)\mathrm{\Delta }(p)p=0.$$
(17)
where $`\mathrm{\Delta }=_\mu _\mu `$ is the laplacian.
The previous equation is a non linear equation of the second order and it is rather difficult to explicit its solutions in closed form. Following a standard route, we shall show that the absolute minima of (3.2) in a given connected component of $`𝒫_\theta `$ actually fulfill a first order equation which is easily solved.
Given a projector $`p𝒫_\theta `$, there is a ‘topological charge’ (the first Chern number) defined by
$$Q(p)=:\frac{1}{2\pi i}p[_1(p)_2(p)_2(p)_1(p)].$$
(18)
As in four dimensional Yang-Mills theory, this topological quantity yields a bound for the action functional.
Due to positivity of the trace $``$ and its cyclic property, we have
$$\left[_\mu (p)p\pm iϵ_{\mu \alpha }_\alpha (p)p\right]^{}\left[_\mu (p)p\pm iϵ_{\mu \beta }_\beta (p)p\right]0,$$
(19)
from which we obtain the inequality
$$S(p)\pm 2Q(p).$$
(20)
The inequality (20), which gives a lower bound for the action, is the analogue of the one for ordinary $`\sigma `$-models . Also, it is a two dimensional analogue of the inequality that occurs in four dimensional Yang-Mills theory. A similar bound for a model on the fuzzy sphere has been obtained in .
From (19) it is clear that the equality in (20) occurs exactly when the projector $`p`$ satisfies the following self-duality or anti-self duality equations
$$\left[_\mu p\pm iϵ_{\mu \alpha }_\alpha p\right]p=0.$$
(21)
By using the derivations $`,\overline{}`$ , the self duality equation (21) reduce to
$$\overline{}(p)p=0,$$
(22)
while the anti-self duality one is
$$(p)p=0.$$
(23)
Simple manipulations show directly that each of the equations (22) and (23) implies the field equations (17), as it should be.
In the next section, we will partially solve these equations.
### 3.3 The instantons of charge $`1`$.
Before we proceed further, let us be more precise about the connected components of $`𝒫_\theta `$ . The latter are parametrized by two integers $`m`$ and $`n`$ such that $`m+n\theta >0`$. When $`\theta ]0,1[`$ is irrational, the corresponding projectors are exactly the projectors of trace $`m+n\theta `$ and the topological charge $`Q(p)`$ appearing in (18) is just $`n`$. We shall construct our solutions for $`m=0`$ and $`n=1`$ and postpone the general discussion to . Thus we have to find projectors that belongs to the previous homotopy class and satisfy the self-duality equation $`(\overline{}p)p=0`$ or, equivalently, $`pp=0`$.
Although these equations look very simple, they are far from being easy to solve because of their non linear nature. To reduce them to a linear problem, we shall introduce the following material and mimic the original construction of Rieffel but with the constraint arising from the self-duality equation.
The space $`=𝒮()`$ of Schwarz functions of one variable is made into a right module over $`𝒜_{1/\theta }`$ by defining
$`(\xi V_1)(s)=:\xi (s1/\theta ),`$
$`(\xi V_2)(s)=:e^{2\pi is}\xi (s),`$ (24)
for any $`\xi `$. It is easily checked that this defines an action on the right of the algebra generated by $`V_1`$ and $`V_2`$ and that the latter is isomorphic to $`𝒜_{1/\theta }`$.
Furthermore, $``$ admits also a left action of $`𝒜_\theta `$ given by
$`(U_1\xi )(s)=:\xi (s1),`$
$`(U_2\xi )(s)=:e^{2\pi is\theta }\xi (s).`$ (25)
and one easily proves that the latter commutes with the right action of $`𝒜_{1/\theta }`$. Besides, the elements of $`𝒜_\theta `$ acting on the left are exactly all linear operators from $``$ to itself that commute with the right action of $`𝒜_{1/\theta }`$, namely $`𝒜_\theta \mathrm{End}_{𝒜_{1/\theta }}()`$.
On the module $``$ there is also a $`𝒜_{1/\theta }`$-valued hermitian structure, namely a sesquilinear map (antilinear in the first variable) $`,:\times 𝒜_{1/\theta }`$ which is compatible with the right $`𝒜_{1/\theta }`$-module structure of $``$ (see for explicit formulae). As a consequence, if $`\xi `$ is such that $`\xi ,\xi `$ is an invertible element of $`𝒜_{1/\theta }`$, the endomorphism
$$p=|\xi \frac{1}{\xi ,\xi }\xi |$$
(26)
is a self-adjoint idempotent (that is a projector) in the algebra $`𝒜_\theta `$ (due to the identification $`𝒜_\theta \mathrm{End}_{𝒜_{1/\theta }}()`$). Here we are using a physicist’s notation for an element $`|\xi `$ and the dual element $`\xi |^{}`$ is defined by means of the hermitian structure as $`\xi |(\eta )=\xi ,\eta 𝒜_{1/\theta }`$ for any $`\eta `$.
In order to translate the self-duality equations for $`p`$ to equations for $`\xi `$, we need to introduce a connection on $``$. This is done by means of two covariant derivatives explicitly given by $`_1,_2:`$,
$$(_1\xi )(s)=:2\pi i\theta s\xi (s),_2\xi =:\frac{d\xi }{ds},$$
(27)
These two operators fulfill a Leibniz rule with respect to the right action
$$_\mu (\xi a)=(_\mu \xi )a+\xi (_\mu a),\mu =1,2.$$
(28)
for any $`\xi `$ and $`a𝒜_{1/\theta }`$ and $`b𝒜_\theta `$. Furthermore, they are compatible with the hermitian structure in the sense that
$$_\mu \xi ,\eta =_\mu \xi ,\eta +\xi ,_\mu \eta ,\mu =1,2,$$
(29)
for any $`\xi ,\eta `$.
By introducing the operator $`\overline{}=1/2(_1+i_2)`$, it is easy to show that $`p`$ satisfies the self-dual equations (22) if and only if there is an element $`\rho 𝒜_{1/\theta }`$ such that
$$\overline{}\xi =\xi \rho .$$
(30)
Thus, we manage to reduce the self-duality equation to a linear equation for $`\xi `$ that can be easily solved in some simple cases.
When $`\rho `$ is a constant element (i.e. it is proportional to the unit of $`𝒜_{1/\theta },\rho =\lambda \mathrm{\hspace{0.17em}1}`$, with $`\lambda `$), equation (30) reduces to the simple differential equation
$$\frac{d\xi }{dt}+(2\pi \theta t+2i\lambda )\xi =0$$
(31)
whose solutions are the gaussians
$$\xi _\lambda (t)=Ae^{\pi \theta t^22i\lambda t},$$
(32)
and $`A^{}`$ is an inessential normalization parameter.
We will show in that, at least for $`\theta `$ small enough, the norms $`\xi _\lambda ,\xi _\lambda `$ are invertible. Accordingly, the gaussians (32) provide a two (real) parameter family of solutions $`p_\lambda =|\xi _\lambda (\xi _\lambda ,\xi _\lambda )^1\xi _\lambda |`$ of the self-duality equations (22), and one can show that the freedom we have in $`\lambda `$ just corresponds to the action of the ordinary torus on $`𝒜_\theta `$ by translation. Thus, we interpret the solution as a two dimensional “instanton” in this simple “ noncommutative Ising model” (remember that the target is just made of two points) and the freedom in $`\lambda `$ in a sense corresponds to its location.
However, it is not obvious that different solutions of the self duality equations on $``$, yield different projectors. In fact, $`\xi `$ and $`\xi ^{^{}}`$ provide different projectors if and only if they belong to different orbits of the action of the group of invertible elements of $`𝒜_{1/\theta }`$ that acts on the right on $``$. Obviously, this action preserves the invertibility of $`\xi ,\xi `$ while the structure of the self-duality equation (30) is preserved provided $`\rho `$ is modified according to
$$\rho g^1\rho g+g^1\overline{}g.$$
(33)
In a more physical language, this means that in trading $`p`$ for $`\xi `$ we have introduced spurious gauge degrees of freedom that we must get rid of. In the case of the Gaussians $`\xi _\lambda `$, it is easy to show that $`\xi _\lambda `$ and $`\xi _\lambda ^{^{}}`$ are gauge equivalent if and only if $`\xi _\lambda ^{^{}}=\xi _\lambda U_1^{n_1}U_2^{n_2}`$ where $`n_1`$ and $`n_2`$ are integers. More generally, given a solution $`\xi `$ of the self-dual equation
$$\overline{}\xi \xi \rho =0$$
(34)
with $`\rho 𝒜_{1/\theta }`$, it is not clear that we can find a complex gauge transformation $`g`$ (i.e. an invertible element of $`𝒜_{1/\theta }`$) that allows to gauge transform $`\xi `$ into one of the gaussians. If this were the case, it would mean that we had indeed constructed all self-dual solutions belonging to the corresponding homotopy class. This problem is tantamount to solve the following equation in $`g`$ and $`\lambda `$,
$$\rho =\lambda +g\overline{}g^1.$$
(35)
Again, this can be done if $`\theta `$ is small enough . The corresponding idea is simple: we first notice that the problem is trivial when $`\theta =1/n`$, with $`n^{}`$ because $`𝒜_{1/\theta }`$ is commutative in this case. Indeed, the existence of the gauge transformation results from the Hodge decomposition of 1-forms. Then, we use the implicit function theorem in order to find how to deform the commutative solution, considered as functions of $`\theta `$ .
A few additional remarks are in order. Even if many of the methods we have used are similar to the ones used in the $`CP^N`$ model rather than to the ones pertaining to the Ising model, we refrain from calling these “noncommutative $`CP^N`$ models” since we want to emphasize the fact that our target space in made of two points and is not the manifold $`CP^N`$ (or more general grassmanian manifolds). But obviously the ordinary grassmanian models can also be considered as “noncommutative Ising models” with a source described by matrix valued functions over an ordinary Riemann surface.
It is also worth remarking that we have been working with the Euclidean metric, but all constructions are readily extended to constant metrics whose conformal class are parametrized by a complex number $`\tau `$ in the upper half-plane. Then, the corresponding moduli space turns out to be a complex torus.
## 4 An analogue of the principal chiral model
Apart from finite spaces, the simplest possible target spaces are circles. Ordinary two dimensional field theories compactified on a circle have been extensively studied (see for instance ) and they essentially behave like free fields (with minor deviations). As we shall show in our next example, this is not the case for noncommutative models, the interaction arising from the noncommutative nature of the source.
To proceed, let us first recall that the algebra of function over the circle $`S^1`$ is generated by a unitary element $`U`$. Thus, specifying a $``$-algebra morphism $`\pi `$ from the algebra of functions on the circle to a another $``$-algebra $`𝒜`$ is tantamount to select a unitary element $`g=\pi (U)`$ in $`𝒜`$. Accordingly, our configuration space is made of all unitary elements of $`𝒜`$. For the metric on the circle we shall take the most natural one, $`G=\delta U\delta U^{}`$, while for the target space we take the Euclidean noncommutative torus, extension to other constant metrics being straightforward. Then, the Hochschild cocycle is the one in (15) and our action functional simply reads $`S[g]=\varphi (1,g,g^1)`$, which reduces to
$$S[g]=\frac{1}{2\pi }_\mu g_\mu g^1,$$
(36)
the variables being unitary elements in the algebra of the noncommutative torus $`𝒜_\theta `$.
Our model is analogous to a principal chiral model, with values in a unitary group of matrices with which it shares lots of properties, apart from non-locality. For the time being we shall limit our study to the existence of infinitely many conserved currents.
From the action functional (36), one readily obtain the equations of motion by varying $`gg+\delta g`$. As in the commutative case, they are equivalent to a current conservation
$$_\mu \left(g^1_\mu g\right)=0,$$
(37)
which expresses the invariance under the global $`U(1)`$ symmetry.
To construct infinitely many such currents, we use a standard induction that relies on the Hodge decomposition of differential forms. Since the latter reduce to a simple problem in linear algebra on the noncommutative torus we shall use it without further discussion and write any differential form as a unique sum of a harmonic one (i.e. a constant one), an exact one and a coexact one.
Let us assume that we have constructed the conserved current $`J_\mu ^{(n)}`$ and let us build $`J_\mu ^{(n+1)}`$. Since $`J_\mu ^{(n)}`$ is conserved, the Hodge decomposition just tells that it is a sum of a constant form and a co-exact one. After an incorporation of the possible constant term into $`J_\mu ^{(n)}`$, one can find $`\chi 𝒜_\theta `$ such that
$$J_\mu ^{(n)}=ϵ_{\mu \nu }_\nu \chi ,$$
(38)
where $`ϵ_{\mu \nu }`$ is the standard antisymmetric tensor. Then, let us introduce the gauge field $`A_\mu =g^1_\mu g`$ and the covariant derivative $`D_\mu =_\mu +A_\mu `$. We define the next current as
$$J_\mu ^{(n+1)}=D_\mu \chi .$$
(39)
It is easy to check that it is conserved, owing to the easily verified commutation rules $`[_\mu ,D_\mu ]=0`$ and $`[D_\mu ,D_\nu ]=0`$. Starting with $`J_\mu ^{(1)}=g^1_\mu g`$, by repeating the construction we can construct an infinite number of non local conserved currents.
Of course the series of new currents would stop whenever there appears a constant current. With some more work one can show that this does not happen unless one starts with a trivial solution of the equations of motion which is a product of the generators. One could also object that non trivial solutions of the equations of motion may not exist. Again this is not the case, since one can take $`g=2p1`$, where $`p`$ is one of the instantonic solutions we constructed in the previous section.
All previous construction is very elementary and follows directly from the ordinary field theoretical construction of the currents. The only point we want to emphasize is that the latter still works in noncommutative geometry. A more thorough survey of our theory along the classical lines will be given in , including a generalization of unitons.
## 5 Addition of the Wess-Zumino term
Although the previous considerations are purely classical, the models can be quantized. This amounts to define and compute the partition function
$$𝒵=[𝒟g]e^{S[g]},$$
(40)
together with the correlation functions
$$gg\mathrm{}g=\frac{1}{𝒵}[𝒟g](gg\mathrm{}g)e^{S[g]},$$
(41)
as well as possible insertions of composite operators. Of course, none of these functional integrals are well defined and to give a precise meaning to them, one has to set up the renormalization procedure which yields some non trivial problems even in dimension two, these models being power counting renormalizable only in that dimension.
However, as far as the one-loop level is considered, this is easily achieved within the background field method. As its non-abelian cousins, our model exhibits a negative $`\beta `$ function so that one may say that it is asymptotically free. Accordingly, one can definitely exclude the possibility of having a free field theory.
We shall see that, after addition of the so called “Wess-Zumino term”, the model behaves almost like a free field (see, for instance, and for recent pedagogical reviews of the ordinary WZW model). Once again, we will only be sketchy because of lack of space, and refer to for a detailed account.
To construct the Wess-Zumino term, let us start with a given unitary element $`g`$ of $`𝒜_\theta `$. It is known from K-theory that there always exist a curve $`g_t,t[0,1]`$ in the group of unitary elements of $`𝒜_\theta `$ that fulfills $`g_1=g`$ and $`g_0=(U_1)^{n_1}(U_2)^{n_2}`$, where $`(n_1,n_2)`$ denotes the class of $`g`$ in K$`{}_{1}{}^{}(𝒜_\theta )`$. Therefore, we can define the Wess-Zumino term as
$$S_{WZ}[g]=\frac{ik}{4\pi }ϵ^{\mu \nu }_0^1𝑑tg_t^1\frac{dg_t}{dt}_\mu g_t^1g_t,$$
(42)
where $`k`$ is an a priori arbitrary real number. As in the classical case, this term can be expressed as the integral over a solid noncommutative torus, but the latter depends on the class of $`g`$ in K-theory, different classes yielding isomorphic and cobordant solid tori.
Although the model (42) depends on the curve $`g_t`$ and not only on $`g`$ one can show that, given any other curve $`\stackrel{~}{g}_t`$ connecting the same boundaries, the difference of the two Wess-Zumino terms can be expressed as an integral over a loop in the group of unitary elements of $`𝒜_\theta `$. Such a quantity may be easily identified with a coupling of a 3-cyclic cocycle with a unitary element of $`C^{\mathrm{}}(S^1)𝒜_\theta `$ and it can be shown to be proportional to an integer, as follows from a straightforward application of the index theorem (see for a very elementary treatment). It turns out that if $`k`$, the Wess-Zumino term is defined up to integral multiples of $`2i\pi `$.
Accordingly, we construct the Wess-Zumino-Witten action just by adding the previous term to the non-linear $`\sigma `$-model and we get,
$$S_{WZW}[g]=\frac{k}{8\pi }_\mu g\mu g^1+\frac{ik}{4\pi }ϵ^{\mu \nu }_0^1𝑑tg_t^1\frac{dg_t}{dt}_\mu g_t^1g_t,$$
(43)
for positive $`k`$.
By introducing the usual operators $``$ and $`\overline{}`$, algebraic manipulations involving integrations by part show that a Polyakov-Wiegman identity holds, namely
$$S_{WZ}[gh]=S_{WZ}[g]+S_{WZ}[h]+\frac{1}{4i\pi }g^1\overline{}ghh^1.$$
(44)
When $`h=1+g^1\delta g`$, this identity allows one to write the variation of $`S_{WZW}[g]`$ as
$$\delta S_{WZW}[g]=\frac{k}{2i\pi }g^1\delta g\left(g^1\overline{}g\right).$$
(45)
Then, the equations of motion can be written equivalently as
$$\overline{J}=0,\mathrm{or}\overline{}J=0,$$
(46)
with $`\overline{J}=g^1\overline{}g`$ and $`J=gg^1`$.
One readily sees that there are very few solutions of the previous equation, since any holomorphic function on the noncommutative torus, defined as an element of the algebra in the kernel of $`\overline{}`$, is constant.
In order to get non trivial solutions, we equip the torus with the Minkowski metric. Then the equations of motion are
$$_+\left(g^1_{}g\right)=0,$$
(47)
with $`_\pm =_1_2`$. Apart from products of the generators $`U_1`$ and $`U_2`$, we will show that the general solution of equations (47) can be factorized as
$$g=g_+g_{},$$
(48)
where $`g_\pm `$ are unitary elements of $`𝒜_\theta `$ satisfying the equations $`_{}g_\pm =0`$.
To proceed, let us first assume that $`g`$ belongs to the connected component of the identity. If this is not the case, we multiply it by a suitable product of the generators, given by the class of $`g`$ in K-theory (the result is still a solution of the equation of motion). Now, it follows from the equation (47) that $`J_{}=g^1_{}g`$ belongs to the kernel of $`_+`$ so that it can be expanded as a Laurent series in $`U_1U_2^1`$. Besides, since we are assuming that $`g`$ belongs to the connected component of the identity, the constant mode of the expansion vanishes, since it is invariant under deformation of $`g`$. Therefore, the primitive
$$_{}J_{}$$
(49)
is well defined (one simply has to divide the coefficient in front of any monomial by the corresponding non vanishing power). As a consequence, the solutions of the remaining equation
$$_{}g=gJ_{}$$
(50)
are easily expressed as
$$g=g_+e^_{}J_{}$$
(51)
with $`g_+`$ an arbitrary unitary element of the algebra of Laurent series in $`U_1U_2`$. Thus, $`g_\pm `$ can be expanded as
$$g_\pm =\underset{n}{}g_\pm ^{(n)}\left(U_1U_2^{\pm 1}\right)^n,$$
(52)
and both can be interpreted as maps from circles $`S_\pm ^1`$, which are the spaces of characters of the commutative algebras generated by $`U_1U_2^{\pm 1}`$, to $`U(1)`$. Note however that the coordinates on $`S_+^1`$ do not commute with the ones on $`S_{}^1`$.
Although this model almost looks like a free field theory, with commutative left and right movers, the standard parity symmetry that exchanges left and right has been broken and the theory, due to noncommutativity, always remembers that left movers must appear on the left. Alternatively, this may be understood as a lack of invariance of the Wess-Zumino term under the inversion $`gg^1`$, while the kinetic term obviously enjoys this symmetry. From the strict point of view of solving the equation of motion, this is the only remainder of the noncommutative nature of the source space.
Obviously, the space of solutions of the equations of motion is invariant under gauge symmetry (respective multiplication on the left and the right by left and right moving unitaries) and under conformal symmetry (reparametrisation of $`S^\pm `$). However, general conformal transformations are not symmetries of the noncommutative torus in the sense that they do not correspond to automorphisms of the algebra $`𝒜_\theta `$; it is only the translations that can be lifted to automorphisms. Therefore, there is no a priori satisfactory way to define conformal transformations of $`g`$ when it is not a solution of the equations of motion. Note that gauge transformation do not create any trouble since they correspond to inner automorphisms of the algebra.
Fortunately, one can construct analogues of conformal transformations that do leave the action invariant and reduce both on-shell and in the commutative case to ordinary conformal transformations. To proceed, let us introduce the (not necessarily conserved) left and right currents $`J_\pm `$, analogous of $`J`$ and $`\overline{J}`$. Furthermore, let us introduce infinitesimal multiplets $`ϵ_\pm =(ϵ_\pm ^{(1)},\mathrm{},ϵ_\pm ^{(n)})`$ of left and right moving elements of $`𝒜_\theta `$. then, we define the infinitesimal transformations $`\delta _{ϵ_\pm }(g)`$ as
$$\delta _ϵ_{}(g)=g\left(\underset{\mathrm{permutations}}{}ϵ_{}^{(i_1)}J_{}ϵ_{}^{(i_2)}J_{}\mathrm{}ϵ_{}^{(i_{n1})}J_{}ϵ_{}^{(i_n)}\right)$$
(53)
and
$$\delta _{ϵ_+}(g)=\left(\underset{\mathrm{permutations}}{}ϵ_+^{(i_1)}J_+ϵ_+^{(i_2)}J_+\mathrm{}ϵ_+^{(i_{n1})}J_+ϵ_+^{(i_n)}\right)g,$$
(54)
with the sums running over all permutations of the indices $`i_1,\mathrm{},i_n`$.
One readily sees that, by replacing $`\delta g`$ with $`\delta _{ϵ_\pm }(g)`$ in the variation of the Wess-Zumino-Witten action, one gets the integral of a total derivative. Thus the variation vanishes even if $`g`$ is not a solution of the equation of motion.
For the particular $`n=1`$ case, the previous transformations reduce to gauge transformations. On shell or in the commutative case, the $`n=2`$ transformations are just the conformal ones. The case $`n>2`$ is more exotic, since two such transformations acting on the left and on the right do not in general commute. Furthermore, it is not clear whether these transformations close off-shell or not. Probably the closure of this algebra requires more general transformations. For instance, $`[\delta _{ϵ_+},\delta _ϵ_{}]`$ is a new symmetry that one has to introduce into the algebra. In the same vein, one also introduces transformations involving the derivatives of the currents. All these transformations are of the form
$$\delta _{ϵ_{},ϵ_+}(g)=gK_{}(ϵ_{},J_{})+K_+(ϵ_+,J_+)g,$$
(55)
where $`K_\pm (ϵ_\pm ,J_\pm )`$ are suitable products of the currents, their derivatives and the corresponding parameters $`ϵ_\pm `$.
Whether this procedure ends or not is not so clear since any computation is rather intricate due to the transformation of the currents themselves. One can also note that the $`n>2`$ case will not yield symmetries of the ordinary $`SU(N)`$ non-abelian Wess-Zumino-Witten theory (apart from the $`SU(2)`$ case), since it does not preserve the unimodular condition of $`SU(N)`$, but they are bona fide transformations for fields with values in $`U(N)`$.
As a final remark, we mention that we have constructed the Wess-Zumino term associated with a particular cyclic cocycle of the noncommutative torus. But the procedure is general and given any $`2n`$-cyclic cocycle on an algebra $`𝒜`$ one can construct the associated Wess-Zumino term in a similar way. Furthermore, the ambiguity in the definition will still be measured by the coupling of a $`2n+1`$-cyclic cocycle with an element of $`K_1(C(S^1)𝒜)`$. This Wess-Zumino term may be added to the action of a principal chiral field, constructed with a Hochschild cocycle. In two dimensions, an analogue of the Polyakov-Wiegman identity still holds provided one uses a suitable scalar product in the LHS.
As a simple example, one can take the matrix algebra $`𝒜=M_n()`$ and the cyclic cocycle given by the trace. Then, for any unitary matrix $`g`$, it is easy to show that $`S_{WZ}[g]=k\mathrm{log}detg`$. In this very simple case, the ambiguity in defining the Wess-Zumino term is nothing but the ambiguity one encounters when defining the argument of a complex number of modulus one, which is arbitrary up to $`2i\pi `$. Furthermore, the Polyakov-Wiegman identity (we drop the kinetic term in this example) reduces to the statement that the argument of a product is the sum of the argument of its factors, up to $`2i\pi `$.
We end here our sketchy discussion of the classical aspects of the noncommutative Wess-Zumino-Witten model. We are aware that many interesting questions have been left aside. In our opinion the main question is to understand how far from a free field theory our model stands.
Finally, we mention that actions analogous to ours have been obtained in and .
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# A Polyvariant Binding-Time Analysis for Off-line Partial Deduction
## 1 Introduction
Partial evaluation and partial deduction are well-known techniques for specialising respectively functional and logic programs. While both depart from the same basic concept, there is quite a divergence between their application and overall approach. In functional programming, the most widespread approach is to use *off-line specialisers*. These are typically very simple and fast specialisers which take (almost) no control decisions concerning the degree of specialisation. In this context, the specialisation is performed as follows: First, a *binding-time analysis* (BTA) is performed on the program which annotates all its statements as either “reducible” or “non-reducible”. The annotated program is then passed to the off-line specialiser, which executes the statements marked reducible and produces residual code for the statements marked non-reducible. In logic programming, the *on-line* approach is almost the only one used. All work is done by a complex on-line specialiser which monitors the whole specialisation process and decides on the degree of specialisation while specialising the program. A few researchers have explored off-line specialisation, but lacking an appropriate notion of BTA, they worked with hand-annotated programs, something which is far from being practical. Until now, it was unclear how to perform BTA for logic programs.
The current paper remedies this situation. It develops a BTA for logic programs, not by translating the corresponding notions from functional programming to logic programming, but by departing from first principles. Given a logic program to be specialised, we develop a logic program which performs its on-line specialisation. The behaviour of this program is analysed and the results are used to take all decisions w.r.t. the degree of specialisation off-line. This turns the on-line specialiser into an off-line specialiser. A prototype has been built and the quality and speed of the off-line specialisation has been evaluated.
## 2 Background
### 2.1 Partial Deduction
In contrast to ordinary (full) evaluation, a partial evaluator receives a program $`P`$ along with only part of its input, called the static input. The remaining part of the input, called the dynamic input, will only be known at some later point in time. Given the static input $`S`$, the partial evaluator then produces a specialised version $`P_S`$ of $`P`$ which, when given the dynamic input $`D`$, produces the same output as the original program $`P`$. The goal is to exploit the static input in order to derive a more efficient program.
In the context of logic programming, full input to a program $`P`$ consists of a goal $`G`$ and evaluation corresponds to constructing a complete SLDNF-tree for $`P\{G\}`$. The static input is given in the form of a partially instantiated goal $`G^{}`$ (and the specialised program should be correct for all instances of $`G^{}`$).
A technique which produces specialised programs is known under the name of *partial deduction* . Its general idea is to construct a finite set of atoms $`𝒜`$ and a finite set of finite, but possibly incomplete SLDNF-trees (one for every<sup>1</sup><sup>1</sup>1Formally, an SLDNF-tree is obtained from an atom or goal by what is called an unfolding rule. atom in $`𝒜`$) which “cover” the possibly infinite SLDNF-tree for $`P\{G^{}\}`$. The derivation steps in these SLDNF-trees correspond to the computation steps which have been performed beforehand and the specialised program is then extracted from these trees by constructing one specialised clause per non-failing branch.
In partial deduction one usually distinguishes two levels of control: the global control, determining the set $`𝒜`$, thus deciding which atoms are to be partially deduced, and the local control, guiding construction of the finite SLDNF-trees for each individual atom in $`𝒜`$ and thus determining what the definitions for the partially deduced atoms look like.
### 2.2 Off-line vs. On-line Control
The (global and local) control problems of partial evaluation and deduction in general have been tackled from two different angles: the so-called on-line versus off-line approaches. The on-line approach performs all the control decisions during the actual specialisation phase. The off-line approach on the other hand performs a (binding-time) analysis phase prior to the actual specialisation phase. This analysis starts from a description of which parts of the inputs will be “*static*” (i.e. sufficiently known) and provides *binding-time annotations* which encode the control decisions to be made by the specialiser, so that the specialiser becomes much more simple and efficient.
Partial evaluation of functional programs has mainly stressed off-line approaches, while supercompilation of functional and partial deduction of logic programs have concentrated on on-line control.
On-line methods, usually obtain better specialisation, because no control decisions have to be taken beforehand, i.e. at a point where the full specialisation information is not yet available. The main reasons for using the off-line approach are to make specialisation itself more efficient and, due to a simpler specialiser algorithm, enable effective self-application (specialisation of the specialiser) .
Few authors discuss off-line specialisation in the context of logic programming , mainly because so far no automated binding-time analysers have been developed. This paper aims to remedy this problem.
## 3 Towards BTA for partial deduction
### 3.1 An on-line specialiser
The basic idea of BTA in functional programming is to model the flow of static input: the arguments of a function call flow to the function body, the result of a function flows back to the call expression. The expressions are annotated reducible when enough of their parameters are static, i.e. will be known at specialisation time, to allow the (partial) computation of the expression. Modelling the dataflow gives a system of inequalities over variables in a domain $`\{static,dynamic\}`$ whose least solution yields the best annotation.
This approach does not immediately translate to logic programs. Problems are that the dataflow in unification is bidirectional and that the degree of instantiation of a variable can change over its lifetime (see also ).
We follow a different approach and reconstruct binding-time analysis from first principles. We start with a Prolog program which performs the unfolding decisions of an on-line specialiser. However, whereas real on-line specialisers base their unfolding decisions on the history of the specialisation phase, ours bases its decisions solely on the actual arguments of the call (which can be more easily approximated off-line). This is in agreement with the off-line specialisers for functional languages which base their decision to evaluate or residualise an expression on the availability of the parameters of that expression. The next step will be to analyse the behaviour of this program (the binding-time analysis) and to use the results to make the unfolding decisions at compile time.
First we develop the on-line specialiser. Assuming that for each predicate $`p/m`$ a *test* predicate $`\mathrm{𝑢𝑛𝑓𝑜𝑙𝑑}\mathrm{\_}p/m`$ exists which decides whether to unfold a call or not, we obtain an on-line specialiser by replacing each call $`p(\overline{t})`$ by
$$(\mathrm{𝑢𝑛𝑓𝑜𝑙𝑑}\mathrm{\_}p(\overline{t})p(\overline{t});\mathrm{𝑚𝑒𝑚𝑜𝑖𝑠𝑒}\mathrm{\_}p(\overline{t}))$$
A call to $`\mathrm{𝑚𝑒𝑚𝑜𝑖𝑠𝑒}\mathrm{\_}p(\overline{t})`$ informs the specialiser that the call $`p(\overline{t})`$ has to be residualised. The specialiser has to check whether (a generalisation of) $`p(\overline{t})`$ has already been specialised —if not it has to initiate the specialisation of (a generalisation of) $`p(\overline{t})`$— and has to perform appropriate renaming of predicates to ensure that residual code calls the proper specialised version of the predicate it calls.
###### Example 1 (Funny append)
Consider the following on-line specialiser for a variant, $`\mathrm{𝑓𝑢𝑛𝑛𝑦𝑎𝑝𝑝}/\mathit{3}`$ of the $`\mathrm{𝑎𝑝𝑝𝑒𝑛𝑑}/\mathit{3}`$ predicate in which the first two arguments of the recursive call have been swapped:
funnyapp(\[\],X,X). funnyapp(\[X|U\],V,\[X|W\]) :- ( unfold\_funnyapp(V,U,W) -\> funnyapp(V,U,W) ; memoise\_funnyapp(V,U,W) ). unfold\_funnyapp(X,Y,Z) :- ground(X).
Specialising this program for a query funnyapp(\[a,b\],L,R) results in the specialised clause (the residual call is renamed as funnyapp\_1)
* funnyapp(\[a,b\],L,\[a|R1\]) :- funnyapp\_1(L,\[b\],R1).
Specialising the funnyapp program for the residual call funnyapp(L,\[b\],R1) gives (after renaming) the clauses
* funnyapp\_1(\[\],\[b\],\[b\]).
* funnyapp\_1(\[X|U\],\[b\],\[X,b|R\]) :- funnyapp\_2(U,\[\],R).
Once more specialising, now for the residual call funnyapp(U,\[\],R), gives
* funnyapp\_2(\[\],\[\],\[\]).
* funnyapp\_2(\[X|U\],\[\],\[X|U\]).
This completes the specialisation. Note that the sequence of residual calls is terminating in this example. In general, infinite sequences are possible. They can be avoided by generalising some arguments of the residual calls before specialising.
In the above example, instead of using ground(X) as condition of unfolding, one could also use the test nilterminated(X).This would allow to obtain the same level of specialisation for a query funnyappend(\[X,Y\],L,R). This test is another example of a so called *rigid* or *downward closed* property: if it holds for a certain term, it holds also for all its instances. Such properties are well suited for analysis by means of abstract interpretation.
### 3.2 From on-line to off-line
Turning the on-line specialiser into an off-line one requires to determine the $`\mathrm{𝑢𝑛𝑓𝑜𝑙𝑑}\mathrm{\_}p/n`$ predicates during a preceding analysis and to decide on whether to replace the $`(\mathrm{𝑢𝑛𝑓𝑜𝑙𝑑}\mathrm{\_}p(\overline{t})p(\overline{t});memoise(p(\overline{t})))`$ construct either by $`p(\overline{t})`$ or by $`memoise(p(\overline{t}))`$. The decision has to be based on a safe estimate of the calls $`\mathrm{𝑢𝑛𝑓𝑜𝑙𝑑}\mathrm{\_}p(\overline{t})`$ which will occur during the specialisation. Computing such safe approximations is exactly the purpose of *abstract interpretation* .
:- $`\{`$grnd(L1)$`\}`$ fap1(L1,L2,R) $`\{`$grnd(L1)$`\}`$. fap1(\[\],X,X). fap1(\[X|U\],V,\[X|W\]) :- $`\{`$grnd(X,U)$`\}`$ ( unf\_fap1(V,U,W) $`\{`$grnd(X,U,V)$`\}`$ -\> $`\{`$grnd(X,U,V)$`\}`$ fap2(V,U,W) $`\{`$grnd(X,U,V,W)$`\}`$ ; $`\{`$grnd(X,U)$`\}`$ memo\_fap3(V,U,W) $`\{`$grnd(X,U)$`\}`$ ) $`\{`$grnd(X,U)$`\}`$. unf\_fap1(X,Y,Z) :- $`\{`$grnd(Y)$`\}`$ ground(X) $`\{`$grnd(X,Y)$`\}`$. memo\_fap3(X,Y,Z).
By adding to the code of Example 1 the fact memoise\_funnyapp(X,Y,Z)., and an appropriate handling of the Prolog built-in ground/1, one can run a goal-dependent *polyvariant* groundness analysis (using e.g. PLAI coupled with the set-sharing domain) for a query where the first argument is ground and obtain the above annotated program. The annotated code for the version fap2 is omitted because it is irrelevant for us. Indeed, inspecting the annotations for unf\_fap1 we see that the analysis cannot infer the groundness of its first argument. So we decide off-line not to unfold, we cancel the test and the then branch and simplify the code into:
:- $`\{`$grnd(L1)$`\}`$ fap1(L1,L2,R) $`\{`$grnd(L1)$`\}`$. fap1(\[\],X,X). fap1(\[X|U\],V,\[X|W\]) :- $`\{`$grnd(X,U)$`\}`$ memo\_fap3(V,U,W) $`\{`$grnd(X,U)$`\}`$.
The residual call to funnyappend has a different call pattern than the original call: its second argument is now ground. Thus we perform a second analysis and obtain (the annotated code for fap4 is omitted):
:- $`\{`$grnd(L2)$`\}`$ fap3(L1,L2,R) $`\{`$grnd(L2)$`\}`$. fap3(\[\],X,X). fap3(\[X|U\],V,\[X|W\]) :- $`\{`$grnd(V)$`\}`$ ( unf\_fap2(V,U,W) $`\{`$grnd(V)$`\}`$ -\> $`\{`$grnd(V)$`\}`$ fap4(V,U,W) $`\{`$grnd(V)$`\}`$ ; $`\{`$grnd(V)$`\}`$ memo\_fap5(V,U,W) $`\{`$grnd(V)$`\}`$ ) $`\{`$grnd(V)$`\}`$. unf\_fap2(X,Y,Z) :- $`\{`$grnd(X)$`\}`$ ground(X) $`\{`$grnd(X)$`\}`$. memo\_fap5(X,Y,Z).
This time, the annotations for unf\_fap2 show that the groundness test will definitely succeed. So we decide off-line always to unfold and only keep the *then* branch. Moreover, the fap4 call has the same call pattern as the original call to funnyapp, so we also rename it as fap1. This yields the second code fragment:
:- $`\{`$grnd(L2)$`\}`$ fap3(L1,L2,R) $`\{`$grnd(L2)$`\}`$. fap3(\[\],X,X). fap3(\[X|U\],V,\[X|W\]) :- $`\{`$grnd(V)$`\}`$ fap1(V,U,W) $`\{`$grnd(V)$`\}`$
Applying the specialiser on these two code fragments for a query fap1(\[a,b\],L,R) gives the same specialised code as in Example 1. However, this time, no calls to unfold\_funnyapp have to be evaluated during specialisation.
### 3.3 Automation
To weave the step by step analysis sketched above in a single analysis, a special purpose tool has to be built. We implemented a system based on the abstract domain POS, also called PROP . It describes the state of the program variables by means of *positive* boolean formulas, i.e., formulas built from $`,`$ and $``$. Its most popular use is for groundness analysis. In that case, the formula $`X`$ expresses that the program variable $`X`$ is (definitely) bound to a ground term, $`XY`$ expresses that $`X`$ is bound to a ground term iff $`Y`$ is, so an eventual binding of $`X`$ to a ground term will be propagated to $`Y`$. This domain is extended with $`\mathrm{𝑓𝑎𝑙𝑠𝑒}`$ as bottom element and is ordered by boolean implication. Groundness analysis corresponds to checking the rigidity<sup>2</sup><sup>2</sup>2A term is rigid w.r.t. a norm if all its instances have the same size w.r.t. the norm. of program variables w.r.t. the *termsize norm*<sup>3</sup><sup>3</sup>3The termsize norm includes all subterms in the measure of the term. and abstracts a unification such as $`X=[Y|Z]`$ by the boolean formula $`XYZ`$. However POS can also be used with other semi-linear norms . In e.g. normalised programs, it only requires to redefine the abstraction of the unifications. For example, with the *listlength norm*<sup>4</sup><sup>4</sup>4The listlength norm includes only the tail of the list in the measure of the list (and measures other terms as 0); nil-terminated lists are rigid under this norm., unification of $`X=[Y|Z]`$ is abstracted as $`XZ`$, and a formula $`X`$ means that the program variable $`X`$ is bound to a term with a bounded listlength, i.e. either the term is a nil-terminated list, or has a main functor which is not a list constructor.
The analyser has to decide the outcome of the $`\mathrm{𝑢𝑛𝑓𝑜𝑙𝑑}\mathrm{\_}p`$ test and has to decide which branch to take for further analysis while doing the analysis. Also it has to launch the analysis of the generalisations of the memoised calls. The generalisation we employ is to replace an argument which is not rigid under the norm used in the analysis by the abstraction of a fresh variable. These requirements exclude the direct use of the abstract compilation technique in the way advocated by e.g. . One problem of the scheme of is that it handles constructs $`(ground(X)p(\overline{t});\mathrm{𝑚𝑒𝑚𝑜𝑖𝑠𝑒}(p(\overline{t})))`$ too inaccurately. The boolean formula is represented as a truth table, i.e. a set of tuples, and the analyser processes the truth table a tuple at a time. Therefore it cannot infer in a program point that $`X`$ is true, i.e. that $`X`$ is definitely ground, so it can never conclude that the else branch cannot be taken. The other problem is that the analyses launched for the *memoised* calls should not interfere (i.e. output should not flow back) with the analysis of the clauses containing the memoised calls. Note that defining $`\mathrm{𝑚𝑒𝑚𝑜𝑖𝑠𝑒}`$ as $`\mathrm{𝑚𝑒𝑚𝑜𝑖𝑠𝑒}\mathrm{\_}p(X_1,\mathrm{},X_n):\mathrm{𝑐𝑜𝑝𝑦}(X_1,Y_1),\mathrm{},\mathrm{𝑐𝑜𝑝𝑦}(X_n,Y_n)`$, and abstracting $`\mathrm{𝑐𝑜𝑝𝑦}(X,Y)`$ as $`XY`$ does not work: The abstract success state of executing $`p(Y_1,\mathrm{},Y_n)`$ will update the abstractions of $`X_1,\mathrm{},X_n`$.
Our prototype binding-time analyser currently consists of $`800`$ lines of Prolog code and uses XSB as a generic tool for semantic-based program analysis . The boolean formulas from the POS domain are represented by their truth tables. This representation enables abstract operations to have straightforward implementations based on the projection and equi-join operations of the relational algebra. The disadvantage is that the size of truth tables quickly increases with the number of variables in a clause. The use of dedicated data structures like BDDs to represent the boolean formulas as in often results in better performance but at the expense of substantial programming efforts.
The main part of the analyser can be seen as a source-to-source transformation (i.e. abstract compilation) that given the program $`P`$ to be analysed, produces an abstract program $`P^\alpha `$ with suitable annotations. The abstract program can be directly run under XSB (using tabling to ensure termination). The execution leaves the results of the analysis in the XSB *tables*. Each predicate $`p/n`$ of $`P`$ is abstracted to a predicate $`p^\alpha /2`$ whose arguments carry input and output *sets* of tuples. The core part of setting up the analysis is then to define the code for the abstract interpretation of each call (at a program point PP<sub>#</sub> of interest):
$$(\mathrm{𝑢𝑛𝑓𝑜𝑙𝑑}\mathrm{\_}p(\overline{X})p(\overline{X});\mathrm{𝑚𝑒𝑚𝑜𝑖𝑠𝑒}\mathrm{\_}p(\overline{X})).$$
(1)
is abstracted by the following code fragment:
* $`project(`$Args,TPP<sub>in</sub>,TC$`),`$
* ( $`\mathrm{𝑢𝑛𝑓𝑜𝑙𝑑}\mathrm{\_}p(`$TC$`)`$ -\>
* $`\mathrm{𝑢𝑛𝑓𝑜𝑙𝑑}(`$TC,PP$`{}_{\mathrm{\#}}{}^{})`$, $`p^\alpha (`$TC,TR$`)`$
* ; TR=TC, $`generalise(`$TC,TCG$`),`$ $`memo(`$TCG,PP$`{}_{\mathrm{\#}}{}^{}),`$ $`p^\alpha (`$TCG,\_$`)`$ ),
* $`equi\mathrm{\_}join(`$Args,TPP<sub>in</sub>,TR,TPP$`{}_{out}{}^{}),`$
Predicates $`unfold/2`$ and $`memo/2`$ which abstract the behaviour of each call in the form of (1) above are *tabled* predicates which have no effect on the computation, but only record information containing the results of the analysis. Their arguments are the current abstraction and the current program point. This information is then dumped from the XSB tables and is fed to the off-line specialiser. The variable TPP<sub>in</sub> holds the truth table which represents the abstraction of the program state in the point prior to the call. The call to $`project/3`$ projects the truth table on the positions Args of the variables $`\overline{X}`$ participating in the call. The result is TC (Tuples of the Call). The predicate $`unfold\mathrm{\_}p/1`$ (currently supplied by the user for each predicate $`p/n`$ to be analysed) inspects TC to decide whether there is sufficient information to unfold the call. If it succeeds the *then* branch is taken which analyses the effects of unfolding $`p/n`$. This is done by executing $`p^\alpha /2`$ with TC as abstraction of the call state. The analysis returns TR as abstraction of the program state reached after unfolding $`p/n`$. If the call to $`unfold\mathrm{\_}p/1`$ fails, the call is memoised, and the program state remains unchanged, so TR = TC. The generalisation of the memoised call also needs to be analysed; therefore the else branch first generalises the current state TC into TCG by erasing all dependencies for non-rigid arguments<sup>5</sup><sup>5</sup>5A position is rigid if it has an “s” in each tuple e.g. $`generalise(`$\[p(s,s,s), p(s,d,d)\],TCG$`)`$ yields TCG = \[p(s,s,s), p(s,s,d), p(s,d,s), p(s,d,d)\]. and then calls $`p^\alpha /2`$ with TCG as initial state, but takes care not to use the returned abstract state as the bindings resulting from specialising memoised calls do not flow back. These actions effectively realise the intended functionality of $`\mathrm{𝑚𝑒𝑚𝑜𝑖𝑠𝑒}\mathrm{\_}p/1`$. Finally, the new program state TR over the variables $`\overline{X}`$ has to be propagated to the other program variables described by the initial state TPP<sub>in</sub>. This is achieved simply by taking the equi-join over the Args of TPP<sub>in</sub> and TR. The new program state is described by TPP<sub>out</sub>.
One of our examples (see Section 4.2) uses two different norms in the $`\mathrm{𝑢𝑛𝑓𝑜𝑙𝑑}`$ tests: the term norm which tests for groundness, and the listlength norm which tests for the boundedness of lists (whether lists are nil-terminated). This does not pose a problem for our framework, we simply use a truth table which encodes two boolean formulas, one for the term norm and one for the listlength norm.
## 4 Some Experiments and Benchmarks
We first discuss the parser and liftsolve examples from .
### 4.1 The parser example
A small generic parser for languages defined by grammars of the form $`S::=aS|X`$ ($`X`$ is a placeholder for a terminal symbol as well as the first argument to nont/3; arguments 2 and 3 represent the string to be parsed as a difference list):
* nont(X,T,R) :- t(a,T,V),nont(X,V,R).
* nont(X,T,R) :- t(X,T,R).
* t(X,\[X|Es\],Es).
A termination analysis can easily determine that calls to t/3 always terminate and that calls to nont/3 terminate if their second argument is ground. One can therefore derive the following unfold predicates:
* unfold\_t(X,S1,S2).
* unfold\_nont(X,T,R) :- ground(T).
Performing our analysis for the entry point :- $`\{grnd(X)\}`$ nont(X,\_,\_) we obtain the following annotated program (dynamic arguments \[i.e. non-ground ones\] and non-reducible predicates \[i.e. memoised ones\] are underlined):
* nont(X,T,R) :- t(a,T,V), nont(X,V,R).
* nont(X,T,R) :- t(X,T,R).
* t(X,\[X|Es\],Es).
Feeding this information into the off-line system logen and specialising nont(c,T,R), we obtain:
* nont\__0(\[a$`|`$B\],C) :- nont\__0(B,C).
* nont\__0(\[c$`|`$D\],D).
Analysing the same specialiser for :- $`\{grnd(T)\}`$ nont(\_,T,\_) yields:
* nont(X,T,R) :- t(a,T,V), nont(X,V,R).
* nont(X,T,R) :- t(X,T,R).
* t(X,\[X|Es\],Es).
Feeding this information into logen and specialising nont(X,\[a,a,c\],R) yields:
* nont\__0(c,\[\]).
* nont\__0(a,\[c\]).
* nont\__0(a,\[a,c\]).
### 4.2 The liftsolve example
The following program is a meta-interpreter for the ground representation, in which the goals are “lifted” to the non-ground representation for resolution. To perform the lifting, an accumulating parameter is used to keep track of the variables that have already been encountered and generated. The predicate mng and l\_mng transform (a list of) ground terms (the first argument) into (a list of) non-ground terms (the second argument; the third and fourth arguments represent the incoming and outgoing accumulator respectively). The predicate solve uses these predicates to “lift” clauses of a program in ground representation (its first argument) and then use them for resolution with a non-ground goal (its second argument) to be solved.
solve(GrP,\[\]). solve(GrP,\[NgH|NgT\]) :- non\_ground\_member(term(clause,\[NgH$`|`$NgBdy\]),GrP), solve(GrP,NgBdy), solve(GrP,NgT). non\_ground\_member(NgX,\[GrH$`|`$\_GrT\]) :- make\_non\_ground(GrH,NgX). non\_ground\_member(NgX,\[\_GrH$`|`$GrT\]) :- non\_ground\_member(NgX,GrT). make\_non\_ground(G,NG) :- mng(G,NG,\[\],\_Sub). mng(var(N),X,\[\],\[sub(N,X)\]). mng(var(N),X,\[sub(N,X)$`|`$T\],\[sub(N,X)$`|`$T\]). mng(var(N),X,\[sub(M,Y)$`|`$T\],\[sub(M,Y)$`|`$T1\]) :- N `\==` M, mng(var(N),X,T,T1). mng(term(F,Args),term(F,IArgs),InS,OutS) :- lmng(Args,IArgs,InS,OutS). lmng(\[\],\[\],Sub,Sub). lmng(\[H|T\],\[IH$`|`$IT\],InS,OutS) :- mng(H,IH,InS,InS1), lmng(T,IT,InS1,OutS).
The following unfold predicates can be derived by a termination analysis:
* unfold\_lmng(Gs,NGs,InSub,OutSub) :- ground(Gs), bounded\_list(InSub).
* unfold\_mng(G,NG,InSub,OutSub) :- ground(G), bounded\_list(InSub).
* unfold\_make\_non\_ground(G,NG) :- ground(G).
* unfold\_non\_ground\_member(NgX,L) :- ground(L).
* unfold\_solve(GrP,Query) :- ground(GrP).
Analysing the specialiser for the entry point solve(ground,\_) we obtain:
* solve(GrP,\[\]).
* solve(GrP,\[NgH$`|`$NgT\]) :-
non\_ground\_member(term(clause,\[NgH$`|`$NgBdy\]), GrP),
solve(GrP,NgBdy), solve(GrP,NgT).
* non\_ground\_member(NgX,\[GrH$`|`$\_GrT\]) :- make\_non\_ground(GrH,NgX).
* non\_ground\_member(NgX,\[\_GrH$`|`$GrT\]) :- non\_ground\_member(NgX,GrT).
* make\_non\_ground(G,NG) :- mng(G,NG, \[\],\_Sub).
* mng(var(N),X,\[\],\[sub(N,X)\]).
* mng(var(N),X,\[sub(N,X)$`|`$T\], \[sub(N,X)$`|`$T\]).
* mng(var(N),X,\[sub(M,Y)$`|`$T\], \[sub(M,Y)$`|`$T1\]) :- N `\==` M, mng(var(N),X,T,T1).
* mng(term(F,Args),term(F,IArgs), InS,OutS) :- lmng(Args,IArgs, InS,OutS).
* lmng(\[\],\[\],Sub,Sub).
* lmng(\[H$`|`$T\],\[IH$`|`$IT\], InS,OutS) :- mng(H,IH,InS,InS1), lmng1(T,IT, InS1,OutS).
* lmng1(\[\],\[\],Sub,Sub).
* lmng1(\[H$`|`$T\],\[IH$`|`$IT\], InS,OutS) :- mng1(H,IH,InS,InS1), lmng1(T,IT, InS1,OutS).
* mng1(var(N),X,\[\],\[sub(N,X)\]).
* mng1(var(N),X, \[sub(N,X)$`|`$T\],\[sub(N,X)$`|`$T\]).
* mng1(var(N),X, \[sub(M,Y)$`|`$T\],\[sub(M,Y)$`|`$T1\]) :- N $`\backslash `$== M, mng1(var(N),X,T,T1).
* mng1(term(F,Args),term(F,IArgs), InS,OutS) :- lmng1(Args,IArgs, InS,OutS).
One can observe that the call lmng1(T,IT, InS1,OutS) has not been unfolded. Indeed, the third argument InS1 is considered to be dynamic (non-ground) and the call to unfold\_lmng will thus not always succeed. However, based on the termination analysis, it is actually sufficient for termination if the third arguments to mng and lmng are bounded lists (as the listlength norm can be used in the termination proof). If we use our prototype to also keep track of bounded lists we obtain the desired result: the call lmng1(T,IT,InS1,OutS) can be unfolded as the first argument is ground and third argument can be inferred to be a bounded list. By feeding the so obtained annotations into logen we obtain a specialiser which removes (most of) the meta-interpretation overhead. E.g. specialising
* solve(\[term(clause,\[term(q,\[var(1)\]), term(p,\[var(1)\])\]),
* term(clause,\[term(p,\[term(a,\[\])\])\])\],G)
yields the following residual program:
* solve\__0(\[\]).
* solve\__0(\[term(q,\[B\])|C\]) :- solve\__0(\[term(p,\[B\])\]),solve\__0(C).
* solve\__0(\[term(p,\[term(a,\[\])\])|D\]) :- solve\__0(\[\]),solve\__0(D).
### 4.3 Some Benchmarks
We now study the efficiency and quality of our approach on a set of benchmarks. Except for the parser benchmark all benchmarks come from the dppd benchmark library . We ran our prototype analyser, bta, that performs binding-time analysis and fed the result into the off-line compiler generator logen in order to derive a specialiser for the task at hand. The ecce on-line partial deduction system has been used for comparison (settings are the same as for ecce-x in , i.e. a mixtus like unfolding, a global control based upon characteristic trees but no use of conjunctive partial deduction). The interested reader can consult to see how ecce compares with other systems.
All experiments were conducted on a Sun Ultra-1 running SunOS 5.5.1. ecce and logen were run using Prolog by BIM 4.1.0. bta was run on XSB 1.7.2.
In Table 1 one can see a summary of the transformation times. The columns under bta contain: the time to abstract and compile the program + the time for execution of the abstracted program (both under XSB). The column under logen contains the time to generate the specialiser with logen using the so obtained annotations. Observe, that for any given initial annotation, this has only to be performed once: the so obtained specialiser can then be used over and over again for different specialisation tasks. E.g. the same specialiser was used for the liftsolve.app and liftsolve.app4 benchmark. The ‘\*’ for liftsolve.app indicates the time for the abstract compilation only producing code for the groundness analysis. The extra arguments and instructions for the bounded list analysis were added by hand (but will be generated automatically in the next version of the prototype). The column under PD gives the time for the off-line specialisation. The last column of the table contains the ratio of running ecce over running the specialisers generated by bta + logen. As can be seen, the specialisers produced by bta + logen run 28 – 880 times faster than ecce. We conjecture that for larger programs (e.g liftsolve with a very big object program) this difference can get even bigger. Also, for 3 benchmarks the combined time of running bta + logen and then the so obtained specialiser was less than running ecce, i.e. our off-line approach fares well even in “one-shot” situations. Of course, to arrive at a fully automatic (terminating) system one will still have to add the time for the termination analysis, needed to derive the “unfold” predicates.
Table 2 compares the efficiency of the specialised programs (for the run time queries see ; for the parser example we ran $`\mathrm{𝑛𝑜𝑛𝑡}(c,[a^{17},c,b],\left[b\right])`$ 100 times). As was to be expected, the programs generated by the on-line specialiser ecce outperform those generated by our off-line system. E.g. for the match.kmp benchmark ecce is able to derive a Knuth-Morris-Pratt style searcher, while off-line systems (so far) are unable to achieve such a feat. However, one can see that the specialised programs generated by bta + logen are still very satisfactory. The most satisfactory application is liftsolve.app (as well as liftsolve.app4), where the specialiser generated by bta + logen runs 167 (resp. 880) times faster than ecce while producing residual code of equal (resp. almost equal) efficiency. In fact, the specialiser compiled the append object program from the ground representation into the non-ground one in just 0.006 s (to be compared with e.g. the compilers generated by sage which run in the order of minutes). Furthermore, the time to produce the residual program and then running it is less than the time needed to run the original program for the given set of runtime queries. This nicely illustrates the potential of our approach for applications such as runtime code generation, where the specialisation time is (also) of prime importance.
## 5 Discussion
We have formulated a binding-time analysis for logic programs, and have reported on a prototype implementation and on an evaluation of its effectiveness. To develop the binding-time analysis, we have followed an original approach: Given a program $`P`$ to be analysed we transform it into an on-line specialiser program $`P^{}`$, in which the unfolding decision are explicitly coded as calls to predicates unfold\_p. The on-line specialiser is different from usual ones in the sense that it — like off-line specialisers — uses the availability of arguments to decide on the unfolding of calls. Next, we apply abstract interpretation —a binding-time analysis— to gather information about the run-time behaviour of $`P^{}`$. The information in the program points related to unfold\_p allows to decide whether the test will definitely succeed —in which case the unfolding branch is retained— or will possibly fail —in which case the branch yielding residual code is retained. The resulting program now behaves as an off-line specialiser as all unfolding decisions have been taken at analysis time.
An issue to be discussed in more detail is the termination of the specialisation. First, a specialiser has a global control component. It must ensure that only a finite number of atoms are specialised. In our prototype, we generalise the residual calls before generating a specialised version: arguments which are not rigid<sup>6</sup><sup>6</sup>6I.e., “static” from functional programming becomes “rigid w.r.t. a given norm.” w.r.t. the norm used in the unfolding condition are replaced by fresh variables. This works well in practice but is not a sufficient condition for termination. In principle one could define the memoise\_p predicates as:
* memoise\_p($`\overline{X}`$) :- copy\_term($`\overline{X}`$,$`\overline{Y}`$), generalise($`\overline{Y}`$,$`\overline{Z}`$), p($`\overline{Z}`$).
and then generalise such that quasi-termination of the program, where calls to p are tabulated, can be proven. In practice, the built-in copy\_term/2 and the built-ins needed to implement generalise/2 will make this a non-trivial task. Secondly, there is the local control component. It must ensure that the unfolding of a particular atom terminates. This is decided by the code of the transformed program. Defining the unfold\_p predicates by hand is error-prone and consequently not entirely reliable. In principle, one could replace the calls memoise\_p by true and apply off-the-shelf tools for proving termination of logic programs . Whether these will do well depends on how well they handle the $`\mathrm{𝑖𝑓}\mathrm{𝑡ℎ𝑒𝑛}\mathrm{𝑒𝑙𝑠𝑒}`$ construct used in deciding on the unfolding and the built-ins used in the rigidity test (e.g. the analysis has to infer that X is bounded and rigid w.r.t. the norm in the program point following a test ground(X)). It is likely that small extensions to these tools will suffice to apply them successfully in proving termination of the unfolding<sup>7</sup><sup>7</sup>7After a small extension by its author, the system of could handle small examples. However, so far we have not done exhaustive testing., at least when the unfolding conditions are based on rigidity tests with respect to the norms used by those termination analysis tools.
A more interesting approach for the local control problem is to automatically generate unfolding conditions by program analysis. Actually, one could apply a more general scheme for handling the unfolding than the one used so far. Having for each predicate p/n the original clauses with head p/n and transformed clauses with head pt/n, the transformed clauses could be derived from the original by replacing each call q/m by:
* ( terminates\_q($`\overline{t}`$) -\> q($`\overline{t}`$)
; ( unfold\_q($`\overline{t}`$) -\> qt($`\overline{t}`$) ; memoise\_q($`\overline{t}`$) ) )
In , Decorte and De Schreye describe how the constraint-based termination analysis of can be adapted to generate a finite set of “most general” termination conditions (e.g. for append/3 they would generate rigidity w.r.t. the listlength norm of the first argument and rigidity w.r.t. the listlength norm of the third argument as the two most general termination conditions; for our funnyapp/3 they would generate rigidity of the first and second argument w.r.t. the listlength norm as the most general termination condition.). These conditions can be used to define the terminates\_q predicates. If they succeed, the call q($`\overline{t}`$) can be executed with the original code and is guaranteed to terminate. Moreover, as they are based on rigidity, they are very well suited to be approximated by our binding-time analysis. Actually, in all our benchmarks programs, we were using termination conditions for controlling the unfolding, so in fact we could have further improved the speed of the specialiser by not checking the condition on each iteration but using the above scheme.
Generating unfold\_q definitions is a harder problem. It is related to the generation of “safe” (i.e. termination ensuring) delay declarations in languages such as MU-Prolog and Gödel. This is a subtle problem as discussed in . For example, the condition (nonvar(X); nonvar(Z)) is not safe for a call append(X,Y,Z); execution, and in our case unfolding, could go on infinitely for some non-linear calls (e.g. append(\[a|L\],Y,L)). Also the condition nonvar/1 is not rigid. (For funnyapp/3 we had rigid conditions, however this is rather the exception than the rule.) A safe unfolding condition for append(X,Y,Z) is linear(append(X,Y,Z)), (nonvar(X); nonvar(Z)). Linearity is well suited for analysis (e.g. ), but a test nonvar(X) is not. Moreover, unless X is ground, the test is typically not invariant over the different iterations of a recursive predicate. A solution could be to switch to a hybrid specialiser: deciding the linearity test at analysis-time and the simple nonvar tests at run-time. But as said above, perhaps due to lack of a good application (for languages with delay, speed is more important than safety), there seems to be no work on generating such conditions.
Another hybrid approach is taken in a recent work independent of ours . This work also starts from the termination condition. When it is violated, the size of the term w.r.t. the norm used in the termination condition and the maximal reduction of the size in a single iteration is used to compute the number of unfolding steps. The program is transformed and calls to be unfolded are given an extra argument initialised with the allowed number of unfolding steps. An on-line test checks the value of the counter and the call is residualised when the counter reaches zero.
### Acknowledgements
M. Bruynooghe and M. Leuschel are supported by the Fund for Scientific Research - Flanders Belgium (FWO). K. Sagonas is supported by the Research Council of the K.U. Leuven. Some of the present ideas originated from discussions and joint work with Jesper Jørgensen, and from the PhD. work of Dirk Dussart , to both of whom we are very grateful. We thank Bart Demoen, Stefaan Decorte, Bern Martens, Danny De Schreye and Sandro Etalle for interesting discussions, ideas and comments.
## Appendix 0.A Specialised Programs generated by bta + logen
### 0.A.1 Parser
Original program:
```
nont(X,T,R) :- t(a,T,V),nont(X,V,R).
nont(X,T,R) :- t(X,T,R).
t(X,[X|Es],Es).
```
Partial deduction query:
```
nont(c,X,Y).
```
Specialised program (where nont(c,X,Y) has been renamed to nont\__0(X,Y)):
```
nont__0([a|B],C) :- nont__0(B,C).
nont__0([c|D],D).
```
### 0.A.2 Liftsolve.app
Original program:
```
solve(GrRules,[]).
solve(GrRules,[NgH|NgT]) :-
non_ground_member(term(clause,[NgH|NgBody]),GrRules),
solve(GrRules,NgBody),
solve(GrRules,NgT).
non_ground_member(NgX,[GrH|_GrT]) :-
make_non_ground(GrH,NgX).
non_ground_member(NgX,[_GrH|GrT]) :-
non_ground_member(NgX,GrT).
make_non_ground(G,NG) :- mng(G,NG,[],Sub).
mng(var(N),X,[],[sub(N,X)]).
mng(var(N),X,[sub(N,X)|T],[sub(N,X)|T]).
mng(var(N),X,[sub(M,Y)|T],[sub(M,Y)|T1]) :-
N \== M, mng(var(N),X,T,T1).
mng(term(F,Args),term(F,IArgs),InSub,OutSub) :-
l_mng(Args,IArgs,InSub,OutSub).
l_mng([],[],Sub,Sub).
l_mng([H|T],[IH|IT],InSub,OutSub) :-
mng(H,IH,InSub,IntSub),
l_mng(T,IT,IntSub,OutSub).
```
Partial deduction query:
```
solve([term(clause,[term(app,[term(null,[]),var(l),var(l)])]),
term(clause,[term(app,[term(cons,[var(h),var(x)]),var(y),
term(cons,[var(h),var(z)])]),term(app,[var(x),var(y),var(z)])])],
[term(app,[X1,X2,X3])]).
```
Specialised program:
```
solve__0([]).
solve__0([term(app,[term(null,[]),B,B])|C]) :-
solve__0([]), solve__0(C).
solve__0([term(app,[term(cons,[D,E]),F,term(cons,[D,G])])|H]) :-
solve__0([term(app,[E,F,G])]), solve__0(H).
```
### 0.A.3 Liftsolve.app4
Same original program as liftsolve.app.
Partial deduction query:
```
solve([term(clause,[term(app,[term(null,[]),var(l),var(l)])]),
term(clause,[term(app,[term(cons,[var(h),var(x)]),var(y),
term(cons,[var(h),var(z)])]),
term(app2,[var(x),var(y),var(z)])]),
term(clause,[term(app2,[term(null,[]),var(l),var(l)])]),
term(clause,[term(app2,[term(cons,[var(h),var(x)]),var(y),
term(cons,[var(h),var(z)])]),
term(app3,[var(x),var(y),var(z)])]),
term(clause,[term(app3,[term(null,[]),var(l),var(l)])]),
term(clause,[term(app3,[term(cons,[var(h),var(x)]),var(y),
term(cons,[var(h),var(z)])]),
term(app4,[var(x),var(y),var(z)])]),
term(clause,[term(app4,[term(null,[]),var(l),var(l)])]),
term(clause,[term(app4,[term(cons,[var(h),var(x)]),var(y),
term(cons,[var(h),var(z)])]),
term(app,[var(x),var(y),var(z)])])],
[term(app,[_X,_Y,_Z])])
```
Specialised program:
```
solve__0([]).
solve__0([term(app,[term(null,[]),B,B])|C]) :-
solve__0([]),solve__0(C).
solve__0([term(app,[term(cons,[D,E]),F,term(cons,[D,G])])|H]) :-
solve__0([term(app2,[E,F,G])]),solve__0(H).
solve__0([term(app2,[term(null,[]),I,I])|J]) :-
solve__0([]),solve__0(J).
solve__0([term(app2,[term(cons,[K,L]),M,term(cons,[K,N])])|O]) :-
solve__0([term(app3,[L,M,N])]),solve__0(O).
solve__0([term(app3,[term(null,[]),P,P])|Q]) :-
solve__0([]),solve__0(Q).
solve__0([term(app3,[term(cons,[R,S]),T,term(cons,[R,U])])|V]) :-
solve__0([term(app4,[S,T,U])]),solve__0(V).
solve__0([term(app4,[term(null,[]),W,W])|X]) :-
solve__0([]),solve__0(X).
solve__0([term(app4,[term(cons,[Y,Z]),A_1,term(cons,[Y,B_1])])|C_1]) :-
solve__0([term(app,[Z,A_1,B_1])]),solve__0(C_1).
```
### 0.A.4 Depth
Original program:
```
depth( true, 0 ).
depth( (_g1,_gs), _depth ) :-
depth( _g1, _depth_g1 ),
depth( _gs, _depth_gs ),
max( _depth_g1, _depth_gs, _depth ).
depth( _goal, s(_depth) ) :-
prog_clause( _goal, _body ),
depth( _body, _depth ).
```
Partial deduction query:
```
depth(member(X,[a,b,c,m,d,e,m,f,g,m,i,j]),D).
```
Specialised program:
```
depth__0(true,0).
depth__0(member(B,C),s(D)) :-
depth__0(append(E,[B|F],C),D).
depth__0(append([],G,G),s(H)) :-
depth__0(true,H).
depth__0(append([I|J],K,[I|L]),s(M)) :-
depth__0(append(J,K,L),M).
```
### 0.A.5 Match.Kmp
Original program:
```
match(Pat,T) :- match1(Pat,T,Pat,T).
match1([],Ts,P,T).
match1([A|Ps],[B|Ts],P,[X|T]) :-
A\==B,match1(P,T,P,T).
match1([A|Ps],[A|Ts],P,T) :-
match1(Ps,Ts,P,T).
```
Partial deduction query:
```
match([a,a,b],R).
```
Specialised program:
```
match1__4(B,C).
match1__3([B|C],[D|E]) :- \==(b,B), match1__1(E,E).
match1__3([b|F],G) :- match1__4(F,G).
match1__2([B|C],[D|E]) :- \==(a,B), match1__1(E,E).
match1__2([a|F],G) :- match1__3(F,G).
match1__1([B|C],[D|E]) :- \==(a,B), match1__1(E,E).
match1__1([a|F],G) :- match1__2(F,G).
match__0(B) :- match1__1(B,B).
```
### 0.A.6 Regexp.r1
Original program:
```
generate(empty,T,T).
generate(char(X),[X|T],T).
generate(or(X,Y),H,T) :- generate(X,H,T).
generate(or(X,Y),H,T) :- generate(Y,H,T).
generate(cat(X,Y),H,T) :- generate(X,H,T1), generate(Y,T1,T).
generate(star(X),T,T).
generate(star(X),H,T) :- generate(X,H,T1), generate(star(X),T1,T).
```
Partial deduction query:
```
generate(cat(star(or(char(a),char(b))),
cat(char(a),cat(char(a),char(b)))),X1,[])
```
Specialised program:
```
generate__3([a|B],B).
generate__4([b|B],B).
generate__2(B,C) :- generate__3(B,C).
generate__2(D,E) :- generate__4(D,E).
generate__1(B,B).
generate__1(C,D) :- generate__2(C,E), generate__1(E,D).
generate__6(B,C) :- generate__3(B,D), generate__4(D,C).
generate__5(B,C) :- generate__3(B,D), generate__6(D,C).
generate__0(B,C) :- generate__1(B,D), generate__5(D,C).
```
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# Do shifting Bragg peaks of cuprate stripes reveal fractionally charged kinks ?
thanks: J. Zaanen acknowledges the hospitality of the Institute of Theoretical Physics at the University of California, Santa Barbara where part of this work was done. This research was supported in part by the National Science Foundation under Grant No. PHY94-07194.
## Abstract
The stripe phases found in correlated oxides can be viewed as an ordering of the solitons associated with doping the Mott-insulating state. Inspired by the recent observation that the stripes tilt away from the main axis of the crystal lattice in the regime $`x1/8`$, we propose that a new type of stripe phase is realized in the large doping regime. This new phase should be viewed as a doped version of the microscopically insulating $`x1/8`$ stripes. The topological excitations associated with the extra doping are fractionally charged kinks along the stripes whose motions make the stripe fluctuate. We argue that the directional degree of freedom of the kinks might order, causing the stripe phase to tilt. Quantitative predictions follow for the doping dependence of the tilt angle, which in turn can be used to determine the fundamental charge quantum associated with the stripe phase.
In the last few years, overwhelming experimental evidence has accumulated that the holes in the cuprate materials, such as the high $`T_c`$ superconductors, form various types of striped phases tranquadaover ; sciencerep ; scienceper . The building blocks of such phases consist of domain walls which can be viewed as string-like configurations of holes that separate hole-free antiferromagnetic domains. Important information about the microscopic origin of stripes is contained in the doping dependence of the inverse of the stripe spacing, shown in Fig. 1A. This spacing can be obtained directly from the the distance $`\delta `$ between the measured Bragg-peaks (Fig. 1B). It was experimentally found that $`\delta `$ is to a good approximation directly proportional to doping $`x`$ for a variety of stripe systems like the nickelates Nickelates and manganites Manganites . It is also found in cuprate stripe phases, but only in the doping regime where $`x1/8`$, while for $`x1/8`$, $`\delta `$ becomes roughly $`x`$ independent — see Fig. 1A Yamadaplot ; Uchida . In the regime $`x<0.05`$, the stripes become diagonal Wakimoto . Quite recently it was found Matsushita ; Lee ; Kimura that in a number of cuprates the spin peaks tilt away from the $`(1,0)`$ or $`(0,1)`$ axis in reciprocal space. This puzzling behavior corresponds with a tilt of the stripe phase away from the main crystal axis of the perovskite planes in real space and seems characteristic for the $`x1/8`$ regime. We want to suggest here that this tilt might reflect the presence of fractionally charged solitons living on the stripes. This notion leads to an explicit relation between the $`\delta `$ versus $`x`$ dependence and the $`Y`$-shift, which, if confirmed experimentally, allows one to determine the topological charge of stripes.
The $`x1/8`$ stripes might be viewed as the condensation of the topological defects associated with doping the half-filled Mott-insulator ZaGu ; SteveK ; Pryadko (see Fig. 2A). We here suggest that the $`x1/8`$ state should be viewed in turn as the condensation of the topological defects associated with doping the stripe state itself. This state is as indicated in Fig. 2B. The stripes for $`x1/8`$ are locally not different from the $`x1/8`$ stripes sketched in Fig. 2A, except that once in a while the stripe ‘steps sideways’: it forms kinks. These kinks carry half the fundamental charge quantum associated with the stripe-insulator of the $`x1/8`$ regime and as they move a stripe always in the same direction they cause the orientation of the stripe phase to deviate from the lattice axis, explaining the Y-shift. Although not based on phase separation between two different types of stripe fillings, our proposal shares the essential idea, that the crossover around $`x=1/8`$ and the charge density and orientation of stripes are related, with a suggestion of White and Scalapino WhiteScal .
Our central assumption is that the cuprate stripe phase in the $`x1/8`$ regime is insulating on the microscopic scale. Experimentally, the cuprate stripe phases are metallic-like and a popular viewpoint is that the electronic state on the stripes is metallic. However, the $`\delta =Sx`$ relation, where $`S`$ is the linear slope in the small $`x`$ limit, argues strongly against this internal metallicity. Its meaning, already suggested by earlier mean field calculations, is simply that every hole stabilizes a piece of charged stripe with a length which is an integer multiple of the lattice constant. This special stability when the electron charge and the lattice are commensurate, shows that the internal electronic state of a stripe is insulating on short length scales. At the same time, stripes which are internally insulating on microscopic scales are not necessarily inconsistent with metallic behavior on macroscopic scales. The stripe metal might correspond with a quantum-disordered stripe insulator, characterized by a growth of the quantum fluctuations when length- and time scales are increasing. Although we will not address this issue in any further detail, the remainder should be read as a suggestion for a critical experimental test of this hypothesis.
We follow the literature with regard to the nature of the reference insulator in the $`x1/8`$ regime WiNa ; ZaOl . The slope $`S`$ of the $`\delta `$ versus $`x`$ relationship for small $`x`$ implies that one hole stabilizes two charge stripe unit cells. Accordingly, the dynamics associated with the electrons along the stripe can be associated with that of a quarter filled one-dimensional fermion system ZXShen . It is well known that such a system can become Mott-insulating by breaking translational symmetry with real space charge periods $`2a`$ and $`4a`$, or wavevectors $`4k_F`$ and $`2k_F`$ respectively, where $`k_F=\pi /(4a)`$ ($`a`$ is the lattice constant). Schematically, the ordering pattern on the stripe is like $`\mathrm{}00\mathrm{}`$ and $`\mathrm{}0000\mathrm{}`$ for the $`4k_F`$ and $`2k_F`$ stripes, respectively. Here $`0`$ denotes the hole and $``$ the presence of a spin, which we expect to be disordered due to quantum fluctuations. At present it is not known which type of density wave order is realized on the stripe. All that really matters for the remainder of our discussion is that the $`4k_F`$ stripe should be considered as the condensation of the charge $`e`$ of one hole WiNa , while the $`2k_F`$ stripe is associated with the charge quantum of a pair of holes, $`2e`$ ZaOl .
Let us now turn to the $`x1/8`$ regime. The (near) independence of $`\delta (x)`$ as function of doping $`x`$ implies that a fraction of the holes $`x^{}=x\delta (x)/Sx1/8`$ cannot be ‘absorbed’ by the $`x1/8`$ insulator. These excess holes should dope the $`x1/8`$ state and it is expected that these holes dope the ‘soft’ insulator associated with the density wave on the charge stripes instead of the magnetic domains, the remnants of the ‘hard’ insulator of half-filling. Doped charge density wave states are well understood in the context of conventional 1D systems Haldane . A key concept was introduced by Schrieffer Schrieffer in the study of polyacetyleen: the elementary excitations in such a system are not electrons but in fact parts of an electron. The reason is that these excitations can be viewed as electrons bound to topological defects in the density wave order parameter. That the charge of the electron fractionalizes in doped $`4k_F`$ density wave systems is easily seen by considering the strong coupling limit ZaOsvS . The undoped state with one (static) hole added at the central site is indicated in Fig. 3A. After a couple of hops a configuration is reached where the bare hole has decayed into two propagating excitations carrying half the hole charge ($`e/2`$) which are at the same time domain walls (kinks) in the density wave (Fig. 3B). Such an $`e/2`$ domain wall corresponds with a region of enhanced charge density, and hence an increased Coulomb energy. In contrast to a one-dimensional crystal, a stripe has the additional freedom of moving sidewards, thereby increasing the distance between the two holes associated with the domain wall ZaOsvS , as sketched in Fig. 3C. Thus, the fractionally charged domain walls also cause the stripe to step sidewards, modulating the position of the stripes in space, see Fig. 3D. If all the kinks are in the same direction, the net result is that the stripe takes on average an orientation in space which is tilted away from the lattice axis, very much like slanted phases that are found in lattice string models eskes .
Let us now consider a dense system of such ‘slanted’ stripes. In the presence of any stripe-stripe interaction this will condense at zero temperature in a stripe phase where all the stripes are tilted in the same direction: our explanation for the Y-shift. It is also expected that the kinks themselves order in a regular pattern at zero temperature. The argument is the usual one: the kink gas on a single stripe becomes at long wavelength a Luttinger liquid showing algebraic long range order. In the presence of any interaction between the Luttinger liquids on different stripes this will change into true long range order at zero temperature LiqCr .
The above considerations do not depend on the type of stripe density wave order, except that the charge $`q`$ of the kinks is different. How should the charge of the kinks be counted for $`2k_F`$ parent stripes? The rule is that the charge of the soliton is half the charge quantum associated with the parent. Since the $`2k_F`$ stripes can be considered as condensations of pairs of electrons, the charge associated with a charged kink in the $`2k_F`$ stripe becomes one electron charge. In terms of the strong coupling cartoon of Fig. 3, one should now associate two sites of the real lattice with one site in the lattice of Fig. 3C, so that Fig. 3D can be mapped back to the $`2k_F`$ case by inserting the doubled lattice and the doubled charges.
These general ideas can be illustrated with explicit calculations. Since we are considering fully ordered insulating states, mean-field theory is able to provide meaningful qualitative outcomes. The reference $`2k_F`$ and $`4k_F`$ stripe insulators actually both exist in the Hartree-Fock solutions of the Hubbard model as low lying, but weakly metastable states: the true mean-field ground state corresponds with filled stripes (1 hole per stripe unit cell) ZaOl ; Seibold . In Fig. 2D we show a typical example of a straight bond centred $`2k_F`$ stripe, calculated with the mean-field approximation, which strongly resembles the stripe patterns found in the density matrix renormalization group calculations by White and Scalapino for the $`tJ`$ model WhiteScalDMRG .
By constraining the distance between the stripes to be fixed, increasing the hole density and using appropriate boundary conditions, doped stripes can also be investigated in mean-field theory. We find that such doped stripes, if they exist as locally stable mean-field solutions, prefer localized fractionally charged kinks. A typical example is shown in Fig. 2B where the density of excess holes is indicated for a $`4k_F`$ stripe state at a hole density of $`x=0.139`$. In line with the qualitative arguments given above, every hole doped into the $`4k_F`$ $`x=1/8`$ reference insulator gives rise to two kinks, each carrying a charge of $`e/2`$ and moving the stripe sideways by one lattice constant. Although these ‘slanted’ stripes are quite stable solutions if the reference insulator is of the $`4k_F`$ variety, we did not manage to find stable solutions for the $`2k_F`$ stripes doped with additional holes — within mean field theory, such kinks tend to disintegrate.
In order to get some feeling of the relative stability of the $`4k_F`$ slanted phase we compared it with a straight stripe where an additional hole is simply inserted in the straight density wave ($`\mathrm{}000\mathrm{}`$). On the Hartree-Fock level, a stripe with two kinks of charge $`e/2`$ always has significantly lower energy (of order 0.01 $`t`$). We also checked that the sidesteps can be further stabilized by including a nearest-neighbour Coulomb interaction. The energetics associated with the relative transversal orientation of the kinks is a more delicate matter which will strongly depend on the distance between kinks. We believe that collective fluctuation effects will promote long range slanted order over zig-zag type stripe patterns, in much the same way as quantum fluctuations make lattice strings directed eskes .
Of course, our calculations merely serve to illustrate the principle — neither the question which possibility is realized in nature nor the question whether the ‘topological’ insulator is the proper reference state can be settled by these calculations. However, simple predictions follow which at least in principle can be checked by experiment. Most importantly, the charged kinks offer an explanation for the Y-shift which has already been observed in experiments. However, the kink notion implies that the superlattice peaks of the slanted stripe phase are located along the straight lines crossing the superlattice peaks associated with the $`1/8`$ phase as indicated in Fig. 1C. The reason is that the average distance between the stripes does not change when the stripes are making sidesteps, thereby leaving the distance between the peaks in reciprocal space along the horizontal axis in Fig. 1C unchanged.
Lattice commensuration effects make a transition to kinked stripes around $`x=1/8`$ most likely. However, due to complications associated with quenched disorder we have no prediction for whether there should be a sharp transition or a rapid but smooth cross-over. However, the scenario that there are no free holes, but that instead all holes go onto stripes, leads to an immediate relation between the $`\delta (x)`$ dependence and the $`Y`$-shift. Let us define the magnitude of the $`Y`$-shift as $`\delta k`$ and the fractional charge (in units of $`e`$) carried by the kink as $`q`$. Then $`x^{}=x\delta (x)/Sx1/8`$ is the hole density associated with doping the stripes in the $`x1/8`$ regime. It can be shown that the Y-shift obeys $`\delta k=\frac{x^{}}{2q}\frac{2\pi }{a}`$ and that $`\delta =(xx^{})\frac{2\pi }{a}`$ where $`a`$ is the lattice constant. So we have $`\delta +2q\delta k=\frac{2\pi }{a}x`$, independent of whether there is a crossover or sharp transition, and where it takes place. In Fig. 1D we show our prediction for the Y-shift as function of $`x`$ for the $`2k_F`$ and $`4k_F`$ cases, both for the case of a sharp transition at $`\frac{1}{8}`$ and for the same smooth crossover as in Fig. 1A. We have also indicated in this figure two experimental results for the Y-shift, where the $`x=0.15`$ point has been reported by Lee et al. Lee for an oxygen doped sample of $`La_2CuO_4`$ and the $`x=0.13`$ point is measured in a $`La_{2x}Sr_xCuO_4`$ sample by Matsushita et al. Matsushita . There is also one experimental indication of a Y-shift at dopings slightly less than 1/8 by Kimura et al. Kimura , which suggests that there is a smooth crossover. Clearly an extensive experimental study is called for to establish whether the Y-shift is due to kinks and whether our relationship between $`\delta k`$ and $`\delta (x)`$ holds. The existing data seem to favour a kink charge of $`e`$ over charge $`e/2`$, suggesting that the elementary charge quantum associated with the stripe phase is $`2e`$. Given the proximity of a superconductor characterised by charge quantum $`2e`$, such a stripe charge-quantisation would obviously be quite interesting footnote2 .
We argued that the charged kinks themselves should order. It should therefore in principle be possible to observe the superlattice reflections of this kink lattice in diffraction experiments. As we argued, the kinks on single stripes will repel but the kink order from stripe to stripe is less easy to establish. On the one hand, the kinks carry charge and under the assumption that the screening length is of order or larger than the average kink separation the kinks would tend to maintain a maximum separation, thereby forming a Wigner crystal as indicated in Fig. 2C. However, one could imagine that elastic deformation energies are minimised when kinks line up as indicated in Fig. 2B. The precise realization of the kink superlattice is therefore a subtle, quantitative matter. Observation of this kink superlattice may be a formidable experimental challenge. However, given the conceptual implications implied by such an observation it should be given a high priority.
In conclusion, inspired by the tilting of the stripes seen at higher dopings we have proposed a scenario of stripe phases characterized by a proliferation of fractionally charged solitons. Remarkably, by studying merely the structural characteristics of this phase which is realized at higher dopings, the elementary charge quantum numbers of the stripes can be established. This is undoubtedly a crucial piece of information in the quest for the understanding of high $`T_c`$ superconductivity WiNa .
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# Weak 𝜔-Categories as 𝜔-Hypergraphs
## 1 Introduction
J. Baez and J. Dolan recently proposed an important and impressive definition of weak $`n`$-categories. They utilize nonstandard $`n`$-cells with not just one but many $`n1`$-cells as their domains for taming coherence conditions. Authors’ primary motivation was to understand their idea along the famous slogan “categories are graphs with monoid structures”. Thus they investigated a suitable notion of $`n`$\- or $`\omega `$-dimensional graph-like structures which should include the underlying structures of Baez-Dolan-style weak $`\omega `$-categories.
In the way of pursuing such structures, they found a general notion of $`\omega `$-dimensional structures whose $`n`$-cells have many $`n1`$-cells not only in their domains but also in their codomains. This notion contains various categorical algebras: $`\omega `$-categories, bicategoreis, double categories, etc. Meanwhile, authors noticed that it can be thought of as a form of $`\omega `$-dimensional hypergraphs. Hypergraphs have been explored in mathematics, database theory, concurrency theory and graph rewriting as a device to represent complex notions. But their higher-dimensional extensions are still not known corresponding to $`n`$\- or $`\omega `$-graphs for ordinary graphs. Therefore such structures are named $`\omega `$-hypergraphs<sup>1</sup><sup>1</sup>1The definition of $`\omega `$-hypergraphs in this paper is not the most general form, because each node is shared by at most two hyperedges..
Thus the purpose of this paper is two-folded: One is to provide a general environment for representing various concepts, especially developing various category theories. Another is to give a definition of weak $`\omega `$-categories which respects saturatedness in the meaning of M. Makkai.
## 2 Trees and forests
Our main idea is to represent an $`n`$-cell as a tree with links and polarity. This is refinement of usual simplice (Figure 1).
We start with the definition of trees and forests.
###### Definition 2.1 ($`n`$-trees and $`n`$-forests)
For any natural number $`n0`$, an $`n`$-tree $`T`$ is a triple $`r^T,S^T,\pi ^T`$ consisting of
* $`S_n^T=\{r^T\}`$, whose element $`r^T`$ is called the root of $`T`$;
* $`S^T=_{0in}S_i^T`$, where $`S_i^T`$ is a finite set whose elements are called $`i`$-nodes or simply nodes;
* $`\pi ^T=_{0in1}\pi _i^T`$, where $`\pi _i^T`$ is a function from $`S_i^T`$ to $`S_{i+1}^T`$.
Also, an $`n`$-forest $`F`$ is a pair $`S^F,\pi ^F`$, consisting of
* $`S^F=_{0in}S_i^F`$, where $`S_i^F`$ is a finite set of $`i`$-nodes;
* $`\pi ^F=_{0in1}\pi _i^F`$, where $`\pi _i^F`$ is a function from $`S_i^F`$ to $`S_{i+1}^F`$.
###### Definition 2.2 (isomorphism of trees and forests)
For any natural number $`n0`$, a homomorphism of $`n`$-trees $`\sigma :TT^{}`$ is a map from $`S^T`$ to $`S^T^{}`$ such that, for every $`xS_i^T`$, $`\sigma (x)S_i^T^{}`$ and $`\sigma \pi ^T=\pi ^T^{}\sigma `$. A homomorphism $`\sigma `$ is an isomorphism when it is a bijection. A homomorphism and an isomorphism of $`n`$-forests are also defined in the same way.
###### Definition 2.3 (subtrees and subforests)
For an $`n`$-tree $`T=r^T,S^T,\pi ^T`$ and a $`k`$-node $`s`$ ($`0kn`$), a subtree with the root $`s`$ is defined as $`T|^s=s,S^{T|^s},\pi ^{T|^s}`$ where
* $`S^{T|^s}=_{0ik}S_i^{T|^s}`$, where $`S_i^{T|^s}=\{tS_i^T|(\pi ^T)^{ki}(t)=s\}`$;
* $`\pi ^{T|^s}=_{0ik1}\pi _i^{T|^s}`$ where $`\pi _i^{T|^s}=\pi _i|_{S_i^{T|^s}}`$
And for a $`k`$-node $`s`$ ($`1kn`$), a subforest under $`s`$ is defined as $`T|_s=S^{T|_s},\pi ^{T|_s}`$ where
* $`S^{T|_s}=_{0ik1}S_i^{T|_s}`$, where $`S_i^{T|_s}=\{tS_i^T|(\pi ^T)^{ki}(t)=s\}`$;
* $`\pi ^{T|_s}=_{0ik2}\pi _i^{T|_s}`$ where $`\pi _i^{T|_s}=\pi _i|_{S_i^{T|_s}}`$
Also in the same way, for an $`n`$-forest $`F=S^F,\pi ^F`$ and a $`k`$-node $`s`$ ($`0kn`$), a subtree with the root $`s`$ is defined as $`F|^s=s,S^{F|^s},\pi ^{F|^s}`$, and for a $`k`$-node $`s`$ ($`1kn`$), a subforest under $`s`$ as $`F|_s=S^{F|_s},\pi ^{F|_s}`$.
## 3 Shells
Shells play the same role as shape diagrams in the ordinary category theory. We will mutually inductively define a shell for each cell as a tree with polarity and links and one for each frame as a forest with polarity and links (Figure 3).
### 3.1 the base case
###### Definition 3.1 ($`0`$-cell shells and $`0`$-frame shells)
A $`0`$-cell shell $`\theta `$ is a singleton set with polarity. More precisely, it is $`r^\theta ,S^\theta ,\mathrm{},ϵ^\theta ,\mathrm{},\mathrm{}`$ where $`S^\theta =S_0^\theta =\{r^\theta \}`$ and $`ϵ^\theta `$ is a function from $`S^\theta `$ to $`\{1,1\}`$. Similarly, A $`0`$-frame shell $`\xi `$ is a set with polarity, that is, $`S^\xi ,\mathrm{},ϵ^\xi ,\mathrm{},\mathrm{}`$ where $`S^\xi =S_0^\xi `$ and $`ϵ^\xi `$ is a function from $`S^\xi `$ to $`\{1,1\}`$.
For the meaning of these definitions, see the following sections.
### 3.2 the induction step
Suppose that for the dimensions less than $`n`$, all staff has already been defined.
###### Definition 3.2 ($`n`$-cell shells)
An $`n`$-cell shell $`\theta `$ is
$$r^\theta ,S^\theta ,\pi ^\theta ,ϵ^\theta ,\mathrm{{\rm Y}}^\theta ,\{\sigma _{s,s^{}}^\theta \}_{s,s^{}\mathrm{{\rm Y}}^\theta }$$
consisting of the following data:
* $`\underset{¯}{\theta }=r^\theta ,S^\theta ,\pi ^\theta `$ is an $`n`$-tree, called the base $`n`$-tree of $`\theta `$;
* polarity: $`ϵ^\theta `$ is a function from $`S^\theta `$ to $`\{1,1\}`$;
* links: $`\mathrm{{\rm Y}}^\theta =_{0in2}\mathrm{{\rm Y}}_i^\theta `$, where $`\mathrm{{\rm Y}}_i^\theta S_i^\theta \times S_i^\theta `$;
* linking isomorphisms: for each $`s,s^{}\mathrm{{\rm Y}}_i^\theta `$, $`\sigma _{s,s^{}}`$ is a $`i`$-tree isomorphism from $`\underset{¯}{\theta }|^s`$ to $`\underset{¯}{\theta }|^s^{}`$,
which satisfy the following condition:
* mutuality: $`\theta |_{r^\theta }=S^{\theta |_{r^\theta }},\pi ^{\theta |_{r^\theta }},ϵ^{\theta |_{r^\theta }},\mathrm{{\rm Y}}^\theta ,\{\sigma _{s,s^{}}^\theta \}_{s,s^{}\mathrm{{\rm Y}}^\theta }`$ is an $`n1`$-frame shell, where
+ $`S^{\theta |_{r^\theta }},\pi ^{\theta |_{r^\theta }}=\underset{¯}{\theta }|_{r^\theta }`$
+ $`ϵ^{\theta |_{r^\theta }}=ϵ^\theta |_{S^{\theta |_{r^\theta }}}`$
###### Definition 3.3 ($`n`$-frame shells)
An $`n`$-frame shell $`\xi `$ is
$$S^\xi ,\pi ^\xi ,ϵ^\xi ,\mathrm{{\rm Y}}^\xi ,\{\sigma _{s,s^{}}^\xi \}_{s,s^{}\mathrm{{\rm Y}}^\xi }$$
consisting of the following data:
* $`\underset{¯}{\xi }=S^\xi ,\pi ^\xi `$ is an $`n`$-forest, called the base $`n`$-forest of $`\xi `$;
* polarity: $`ϵ^\xi `$ is a function from $`S^\xi `$ to $`\{1,1\}`$;
* links: $`\mathrm{{\rm Y}}^\xi =_{0in1}\mathrm{{\rm Y}}_i^\xi `$, where $`\mathrm{{\rm Y}}_i^\xi S_i^\xi \times S_i^\xi `$;
* linking isomorphisms: for $`s,s^{}\mathrm{{\rm Y}}_i^\xi `$, $`\sigma _{s,s^{}}`$ is a $`i`$-tree isomorphism from $`\underset{¯}{\xi }|^s`$ to $`\underset{¯}{\xi }|^s^{}`$,
which satisfy the following conditions:
* mutuality: for any $`sS_n^\xi `$, $`\xi |^s=s,S^{\xi |^s},\pi ^{\xi |^s},ϵ^{\xi |^s},\mathrm{{\rm Y}}^{\xi |^s},\{\sigma _{s,s^{}}^\xi \}_{s,s^{}\mathrm{{\rm Y}}^{\xi |^s}}`$ is an $`n`$-cell shell, where
+ $`s,S^{\xi |^s},\pi ^{\xi |^s}=\underset{¯}{\xi }|^s`$
+ $`ϵ^{\xi |^s}=ϵ^\xi |_{S^{\xi |^s}}`$;
+ $`\mathrm{{\rm Y}}^{\xi |^s}=_{0in2}\mathrm{{\rm Y}}_i^{\xi |^s}`$, where $`\mathrm{{\rm Y}}_i^{\xi |^s}=\mathrm{{\rm Y}}_i^\xi |_{S_i^{\xi |^s}\times S_i^{\xi |^s}}`$
* bijectivity: if $`s,t,s,t^{}\mathrm{{\rm Y}}_{n1}^\xi `$, then $`t=t^{}`$; and if $`s,t,s^{},t\mathrm{{\rm Y}}_{n1}^\xi `$, then $`s=s^{}`$;
* involution: if $`s,t\mathrm{{\rm Y}}_{n1}^\xi `$, then $`t,s\mathrm{{\rm Y}}_{n1}^\xi `$ and $`\sigma _{s,t}^\xi \sigma _{t,s}^\xi =\mathrm{Id}`$ and $`\sigma _{t,s}^\xi \sigma _{s,t}^\xi =\mathrm{Id}`$;
* conjugation: if $`s,s^{}\mathrm{{\rm Y}}_{n1}^\xi `$ and $`tS^\xi `$ such that $`(\pi ^\xi )^i(t)=s`$ for some $`i0`$, then $`ϵ(t)ϵ(\sigma _{s,s^{}}^\xi (t))=1`$ (this implies anti-reflexivity: if $`s,s^{}\mathrm{{\rm Y}}_{n1}^\xi `$, then $`ss^{}`$);
* correspondence of links: if $`s,s^{}\mathrm{{\rm Y}}_{n1}^\xi `$ and $`t,t^{}\mathrm{{\rm Y}}_i^\xi `$ for some $`in3`$ and $`(\pi ^\xi )^{ni}(t)=(\pi ^\xi )^{ni}(t^{})=s`$, then $`\sigma _{s,s^{}}^\xi (t),\sigma _{s,s^{}}^\xi (t^{})\mathrm{{\rm Y}}_i^\xi `$ and $`\sigma _{t,t^{}}^\xi \sigma _{s,s^{}}^\xi =\sigma _{\sigma _{s,s^{}}^\xi (t),\sigma _{s,s^{}}^\xi (t^{})}^\xi \sigma _{t,t^{}}^\xi `$;
* commutativity of links: for $`k2`$ and $`s_1^{},s_2`$, $`s_2^{},s_3`$,…, $`s_{k1}^{},s_k`$, $`s_k^{},s_1`$ in $`\mathrm{{\rm Y}}^\xi `$ such that $`\pi ^\xi (s_i)=s_i^{}`$ or $`\pi ^\xi (s_i^{})=s_i`$, if $`s^{}=(\sigma _{s_k^{},s_1}^\xi \sigma _{s_{k1}^{},s_k}^\xi \mathrm{}\sigma _{s_2^{},s_3}^\xi \sigma _{s_1^{},s_2}^\xi )(s)`$ is defined, then $`s^{}=s`$;
(that is, if $`s_i^{}`$ is of the smallest level between $`s_1^{}`$,…,$`s_k^{}`$, then the composition of such isomorphisms as above $`\sigma _{s_q^{},s_i}^\xi \mathrm{}\sigma _{s_i^{},s_p}^\xi `$ is defined and is the identity homomorphism of the subtree at $`s_i^{}`$).
* closedness: for every $`sS_{n1}^\xi `$, there exists a(n unique) node $`s^{}S_{n1}^\xi `$ such that $`s,s^{}\mathrm{{\rm Y}}_{n1}^\xi `$.
Closedness means globularity of higher dimensional cells. Note that every $`tS^\xi `$ is in $`S^{\xi |^s}`$ for just one $`sS_n^\xi `$; and also every $`t,t^{}\mathrm{{\rm Y}}_i^\xi `$ for $`in2`$ is in $`\mathrm{{\rm Y}}_i^{\xi |^s}`$ for just one $`sS_n^\xi `$. The latter is due to the closedness of frame shells at lower levels.
###### Proposition 3.1
For a cell shell $`\theta `$, if $`s,s^{}\mathrm{{\rm Y}}^\theta `$, then $`\pi ^i(s)=\pi ^i(s^{})`$ where $`i=1`$ or $`2`$.
###### Remark 3.1
Thus the situation of the correspondence of links for $`k=2`$ occur only when $`l_1=l_2=1`$ and either the parents of $`s_1^{}`$ and $`s_2`$ or those of $`s_2`$ and $`s_1^{}`$ are the same. And for an $`n`$-cell shell $`\theta `$ in $`n`$-frame shell, an outer link $`s,s^{}`$ of which $`s`$ or $`s^{}`$ is not in $`S_\theta `$, must be an $`n1`$-link.
###### Definition 3.4 ($`_n`$, $`()^{}`$)
For two $`n`$-frame shells
$`\xi `$ $`=S^\xi ,\pi ^\xi ,ϵ^\xi ,\mathrm{{\rm Y}}^\xi ,\{\sigma _{s,s^{}}^\xi \}_{s,s^{}\mathrm{{\rm Y}}^\xi }\text{and}`$
$`\xi ^{}`$ $`=S^\xi ^{},\pi ^\xi ^{},ϵ^\xi ^{},\mathrm{{\rm Y}}^\xi ^{},\{\sigma _{s,s^{}}^\xi ^{}\}_{s,s^{}\mathrm{{\rm Y}}^\xi ^{}}\text{,}`$
an isomorphism $`f`$ from $`\xi `$ to $`\xi ^{}`$ is an $`n`$-forest isomorphism (with its inverse $`f^1`$) such that
* $`ϵ^\xi (s)=ϵ^\xi ^{}(f(s))`$ ($``$ $`ϵ^\xi ^{}(t)=ϵ^\xi (f^1(t))`$);
* if $`s,s^{}\mathrm{{\rm Y}}^\xi `$, then $`f(s),f(s^{})\mathrm{{\rm Y}}^\xi ^{}`$ ($``$ if $`t,t^{}\mathrm{{\rm Y}}^\xi ^{}`$, then $`f^1(t),f^1(t^{})\mathrm{{\rm Y}}^\xi `$);
* $`f\sigma _{s,s^{}}^\xi =\sigma _{f(s),f(s^{})}^\xi ^{}f`$ ($``$ $`f^1\sigma _{t,t^{}}^\xi ^{}=\sigma _{f(t),f(t^{})}^\xi f^1`$).
When an isomorphism $`f`$ from $`\xi `$ to $`\xi ^{}`$ exists, we say that $`\xi `$ is isomorphic to $`\xi ^{}`$, and write $`f:\xi _n\xi ^{}`$, $`\xi _n\xi ^{}`$, or simply $`\xi \xi ^{}`$. Obviously $`_n`$ is an equivalence relation.
For an $`n`$-frame shell $`\xi `$, $`(\xi )^{}`$ is defined as $`S^\xi ,\pi ^\xi ,ϵ^{(\xi )^{}},\mathrm{{\rm Y}}^\xi ,\{\sigma _{s,s^{}}^\xi \}_{s,s^{}\mathrm{{\rm Y}}^\xi },`$ where $`ϵ^{(\xi )^{}}(s)=ϵ^\xi (s)`$. It is easy to check well-definedness, that is, $`(\xi )^{}`$ is in fact an $`n`$-frame shell, and $`((\xi )^{})^{}=\xi `$.
An $`n`$-cell shell can be seen as an $`n`$-frame shell. Thus we can define isomorphsims between $`n`$-cell shells similarly.
## 4 Diagrams and $`\omega `$-hypergraphs
Cell diagrams and frames are mutually inductively defined.
###### Definition 4.1 ($`i`$-cell)
We prepare a set of $`i`$-cells for each $`i\{0\}`$:
* $`\mathrm{\Sigma }_i=\mathrm{\Sigma }_{i,1}\mathrm{\Sigma }_{i,1}`$,
* a bijection $`()^{}:\mathrm{\Sigma }_i\mathrm{\Sigma }_i`$ such that for $`c\mathrm{\Sigma }_{i,k}`$ with $`k\{1,1\}`$, $`c^{}\mathrm{\Sigma }_{i,k}`$ and $`(c^{})^{}=c`$.
Elements of $`\mathrm{\Sigma }_i`$ are called $`i`$-cells; those of $`\mathrm{\Sigma }_{i,1}`$ positive $`i`$-cells; and those of $`\mathrm{\Sigma }_{i,1}`$ negative $`i`$-cells. $`c^{}`$ is called the conjugate of $`c`$.
### 4.1 the base case
###### Definition 4.2 ($`0`$-hypergraph, $`0`$-cell diagram and $`0`$-frame)
For consistency, let $`\mathrm{𝐅𝐫𝐦}_1`$ be $`\{\mathrm{}\}`$, the only one $`1`$-frame isomorphism the empty function $`\mathrm{}:\mathrm{}\mathrm{}`$ and $`_0:\mathrm{\Sigma }_0\mathrm{𝐅𝐫𝐦}_1`$ the unique function. A $`0`$-hypergraph is $`\mathrm{\Sigma }_0,_0`$. A $`0`$-cell diagram $`\eta `$ is $`r^\eta ,S^\eta ,\mathrm{},ϵ^\eta ,\mathrm{},\mathrm{},\lambda ^\eta ,\{\rho _{r^\eta }^\eta \}_{r^\eta S^\eta }`$, where $`\underset{¯}{\eta }=r^\eta ,S^\eta ,\mathrm{},ϵ^\eta ,\mathrm{},\mathrm{}`$ is a $`0`$-cell shell, $`\lambda ^\eta `$ is a function from $`S^\eta `$ to $`\mathrm{\Sigma }_0`$ such that $`\lambda ^\eta (r^\eta )\mathrm{\Sigma }_{0,ϵ^\eta (r^\eta )}`$, and $`\rho _{r^\eta }^\eta `$ is the empty function. A $`0`$-frame diagram, or simply a $`0`$-frame, $`\zeta `$ is $`S^\zeta ,\mathrm{},ϵ^\zeta ,\mathrm{},\mathrm{},\lambda ^\zeta ,\{\rho _s^\zeta \}_{sS^\zeta }`$, where $`\underset{¯}{\zeta }=S^\zeta ,\mathrm{},ϵ^\zeta ,\mathrm{},\mathrm{}`$ is a $`0`$-frame shell, $`\lambda ^\zeta `$ is a function from $`S^\zeta `$ to $`\mathrm{\Sigma }_0`$ such that for any $`sS^\zeta `$, $`\lambda ^\zeta (s)\mathrm{\Sigma }_{0,ϵ^\zeta (s)}`$, and each $`\rho _s^\zeta `$ is the empty function. A $`0`$-frame isomorphism from $`\zeta `$ to $`\zeta ^{}`$ is a $`0`$-frame shell isomorphism $`f:\underset{¯}{\zeta }\underset{¯}{\zeta ^{}}`$ (in fact, a bijection from $`S^\zeta `$ to $`S^\zeta ^{}`$) satisfying $`\lambda ^\zeta =\lambda ^\zeta ^{}f`$. $`\mathrm{𝐅𝐫𝐦}_0`$ is the set of $`0`$-frames.
### 4.2 the induction step
Suppose that $`n1`$ and that for the dimensions less than $`n`$, all staff has already been defined.
###### Definition 4.3 (boundary of $`n`$-cells)
As a parameter of definitions, a function $`_n:\mathrm{\Sigma }_n\mathrm{𝐅𝐫𝐦}_{n1}`$ satisfying $`(_n(c))^{}=_n(c^{})`$ are given (for the $`n1`$ dimension, $`()^{}`$ for frames have been defined). $`_n(c)`$ is called the boundary of $`c`$.
###### Definition 4.4 ($`n`$-hypergraph)
An $`n`$-hypergraph $`G=\mathrm{\Sigma },`$ consists of
* $`\mathrm{\Sigma }=_{0in}\mathrm{\Sigma }_i`$, and
* $`=_{1in}_i`$.
###### Definition 4.5 ($`n`$-cell diagram)
An $`n`$-cell diagram $`\eta `$ is
$$r^\eta ,S^\eta ,\pi ^\eta ,ϵ^\eta ,\mathrm{{\rm Y}}^\eta ,\{\sigma _{s,s^{}}^\eta \}_{s,s^{}\mathrm{{\rm Y}}^\eta },\lambda ^\eta ,\{\rho _s^\eta \}_{sS^\eta }$$
where
* $`\underset{¯}{\eta }=r^\eta ,S^\eta ,\pi ^\eta ,ϵ^\eta ,\mathrm{{\rm Y}}^\eta ,\{\sigma _{s,s^{}}^\eta \}_{s,s^{}\mathrm{{\rm Y}}^\eta }`$ is an $`n`$-cell shell, called the base $`n`$-cell shell of $`\eta `$;
* assignment of cells: $`\lambda ^\eta =_{0in}\lambda _i^\eta `$, where $`\lambda _i^\eta `$ is a function from $`S_i^\eta `$ to $`\mathrm{\Sigma }_i`$ such that for any $`sS_i`$, $`\lambda _i(s)\mathrm{\Sigma }_{i,ϵ(s)}`$
* identification in boundaries: for $`sS_i^\eta `$, $`\rho _s^\eta `$ is an $`i1`$-frame isomorphism from $`\eta |_s`$ to $`_i(\lambda _i^\eta (s))`$
which satisfy the following conditions:
* mutuality: $`\eta |_{r^\eta }=S^{\eta |_{r^\eta }},\pi ^{\eta |_{r^\eta }},ϵ^{\eta |_{r^\eta }},\mathrm{{\rm Y}}^\eta ,\{\sigma _{s,s^{}}^\eta \}_{s,s^{}\mathrm{{\rm Y}}^\eta },\lambda ^{\eta |_{r^\eta }},\{\rho _s^\eta \}_{sS^{\eta |_{r^\eta }}}`$ is an $`n1`$-frame, where
+ $`S^{\eta |_{r^\eta }},\pi ^{\eta |_{r^\eta }},ϵ^{\eta |_{r^\eta }},\mathrm{{\rm Y}}^\eta ,\{\sigma _{s,s^{}}^\eta \}_{s,s^{}\mathrm{{\rm Y}}^\eta }=\underset{¯}{\eta }|_{r^\eta }`$
+ $`\lambda ^{\eta |_{r^\eta }}=\lambda ^\eta |_{S^{\eta |_{r^\eta }}}`$.
###### Definition 4.6 ($`n`$-frame)
$`n`$-frame diagram or $`n`$-frame $`\zeta `$ is
$$S^\zeta ,\pi ^\zeta ,ϵ^\zeta ,\mathrm{{\rm Y}}^\zeta ,\{\sigma _{s,s^{}}^\zeta \}_{s,s^{}\mathrm{{\rm Y}}^\zeta },\lambda ^\zeta ,\{\rho _s^\zeta \}_{sS^\zeta }$$
where
* $`\underset{¯}{\zeta }=S^\zeta ,\pi ^\zeta ,ϵ^\zeta ,\mathrm{{\rm Y}}^\zeta ,\{\sigma _{s,s^{}}^\zeta \}_{s,s^{}\mathrm{{\rm Y}}^\zeta }`$ is an $`n`$-frame shell, called the base $`n`$-frame shell of $`\zeta `$;
* assignment of cells: $`\lambda ^\zeta =_{0in}\lambda _i^\zeta `$, where $`\lambda _i^\zeta `$ is a function from $`S_i^\zeta `$ to $`\mathrm{\Sigma }_i`$ such that for any $`sS_i`$, $`\lambda _i(s)\mathrm{\Sigma }_{i,ϵ(s)}`$
* identification in boundaries: for $`sS_i^\zeta `$, $`\rho _s^\zeta `$ is an $`i1`$-frame isomorphism from $`\zeta |_s`$ to $`_i(\lambda _i^\zeta (s))`$
which satisfy the following conditions:
* mutuality: for every $`sS_n^\zeta `$, $`\zeta |^s=s,S^{\zeta |^s},\pi ^{\zeta |^s},ϵ^{\zeta |^s},\mathrm{{\rm Y}}^{\zeta |^s},\{\sigma _{s,s^{}}^\zeta \}_{s,s^{}\mathrm{{\rm Y}}^{\zeta |^s}},\lambda ^{\zeta |^s},\{\rho _t^\zeta \}_{tS^{\zeta |^s}}`$ is an $`n`$-cell diagram, where
+ $`s,S^{\zeta |^s},\pi ^{\zeta |^s},ϵ^{\zeta |^s},\mathrm{{\rm Y}}^{\zeta |^s},\{\sigma _{s,s^{}}^\zeta \}_{s,s^{}\mathrm{{\rm Y}}^{\zeta |^s}}=\underset{¯}{\zeta }|^s`$
+ $`\lambda ^{\zeta |^s}=\lambda ^\zeta |_{S^{\zeta |^s}}`$
* compatibility on links:
+ for $`s,s^{}\mathrm{{\rm Y}}_{n1}^\zeta `$, $`\lambda (s)=(\lambda (s^{}))^{}`$ and
+ for $`s,s^{}\mathrm{{\rm Y}}_{n1}^\zeta `$ and $`tS^{\zeta |_s}`$, $`\rho _s^\zeta (t)=\rho _s^{}^\zeta (\sigma _{s,s^{}}(t))`$
###### Proposition 4.1
For $`s,s^{}\mathrm{{\rm Y}}_{n1}^\zeta `$ and $`tS_k^\zeta `$ for some $`kn1`$ such that $`(\pi ^\zeta )^{nk1}(t)=s`$, $`(\lambda (t))^{}=\lambda (\sigma _{s,s^{}}^\zeta (t))`$.
Proof It is induced from the compatibility on links and the definition of $`()^{}`$ for $`n2`$-frames. $`\mathrm{}`$
###### Definition 4.7 ($`_n`$, $`\mathrm{𝐅𝐫𝐦}_n`$, $`()^{}`$)
For two $`n`$-frames
$`\zeta `$ $`=S^\zeta ,\pi ^\zeta ,ϵ^\zeta ,\mathrm{{\rm Y}}^\zeta ,\{\sigma _{s,s^{}}^\zeta \}_{s,s^{}\mathrm{{\rm Y}}^\zeta },\lambda ^\zeta ,\{\rho _s^\zeta \}_{sS^\zeta }\text{and}`$
$`\zeta ^{}`$ $`=S^\zeta ^{},\pi ^\zeta ^{},ϵ^\zeta ^{},\mathrm{{\rm Y}}^\zeta ^{},\{\sigma _{s,s^{}}^\zeta ^{}\}_{s,s^{}\mathrm{{\rm Y}}^\zeta ^{}},\lambda ^\zeta ^{},\{\rho _s^\zeta ^{}\}_{sS^\zeta ^{}}\text{,}`$
an isomorphism $`f`$ from $`\zeta `$ to $`\zeta ^{}`$ is an isomorphism of $`n`$-frame shells $`f:\underset{¯}{\zeta }\underset{¯}{\zeta ^{}}`$ (with its inverse $`f^1`$) such that
* for $`sS^\zeta `$, $`\lambda ^\zeta (s)=\lambda ^\zeta ^{}(f(s))`$ ($``$ $`\lambda ^\zeta ^{}(s^{})=\lambda ^\zeta (f^1(s^{}))`$),
* for $`sS^\zeta `$ and $`tS^{\zeta |_s}`$, $`\rho _s^\zeta (t)=\rho _{f(s)}^\zeta ^{}(f(t))`$ ($``$ $`\rho _s^{}^\zeta ^{}(t^{})=\rho _{f^1(s^{})}^\zeta (f^1(t^{}))`$).
When an isomorphism $`f`$ from $`\zeta `$ to $`\zeta ^{}`$ exists, we say that $`\zeta `$ is isomorphic to $`\zeta ^{}`$, and write $`f:\zeta _n\zeta ^{}`$, $`\zeta _n\zeta ^{}`$, or simply $`\zeta \zeta ^{}`$. Obviously $`_n`$ is an equivalence relation. The collection of all $`n`$-frames is denoted by $`\mathrm{𝐅𝐫𝐦}_n`$. For an $`n`$-frame $`\zeta `$, $`(\zeta )^{}`$ is defined as $`S^\zeta ,\pi ^\zeta ,ϵ^{(\zeta )^{}},\mathrm{{\rm Y}}^\zeta ,\{\sigma _{s,s^{}}^\zeta \}_{s,s^{}\mathrm{{\rm Y}}^\zeta },\lambda ^{(\zeta )^{}}`$ where $`ϵ^{(\zeta )^{}}(s)=ϵ^\zeta (s)`$ and $`\lambda ^{(\zeta )^{}}(s)=(\lambda ^\zeta (s))^{}`$. It is easy to check well-definedness, that is, $`(\zeta )^{}`$ is in fact an $`n`$-frame, and $`((\zeta )^{})^{}=\zeta `$.
An $`n`$-cell diagram can be seen as an $`n`$-frame. Thus we can define an isomorphism between $`n`$-cell diagrams similarly.
###### Remark 4.1
Indeed, conditions for $`\rho `$ in the definitions of cell diagrams and frames ensure the commutativity of links and other commutativity of their base shells (it is easy to check this). Therefore if we use shells only for diagrams, we need not introduce such commutativity. A main purpose to do it is to treat closure operations for shells. Due to commutativity, a closure becomes unique in a sense.
### 4.3 $`\omega `$-hypergraphs
###### Definition 4.8 ($`\omega `$-hypergraph)
An $`\omega `$-hypergraph $`G=\mathrm{\Sigma },`$ consists of
* $`\mathrm{\Sigma }=_{0i}\mathrm{\Sigma }_i`$, and
* $`=_{1i}_i`$.
###### Remark 4.2
Boundaries $`_i`$ depend on frames in the previous step of the inductive definition. Therefore as pointed out in , to formalize the definition of $`\omega `$-hypergraphs in a logical system, we need a sort of dependent choice axiom, $`\text{GDC}_\tau `$ in §4.4.3 or $`\text{DC}_1`$ in §8.2.3. The strength of this is in between the countable axiom of choice and the full axiom of choice .
## 5 Pasting diagrams and their closures
###### Definition 5.1 ($`n`$-pasting shells)
An $`n`$-pasting shell consists of the same data and conditions as an $`n`$-frame shell, but at the last induction step, the closedness condition is not required. That is, $`n1`$-nodes which do not appear in $`\mathrm{{\rm Y}}_{n1}^\xi `$ are allowed. We call them open nodes of the pasting shell. An $`n`$-pasting shell is positive or negative if for all $`sS_n^\xi `$, $`ϵ_n^\xi (s)=1`$ or $`1`$, respectively.
###### Definition 5.2 ($`n`$-pasting diagrams)
An $`n`$-pasting diagram $`\zeta `$ is defined in the same way as $`n`$-frame, but $`\underset{¯}{\zeta }`$ is an $`n`$-pasting shell instead of an $`n`$-frame shell. $`_n`$, $`\mathrm{𝐏𝐃}_n`$, $`()^{}`$ is also defined similarly. An $`n`$-pasting diagram is positive or negative if the base $`n`$-pasting shell is positive or negative, respectively.
###### Lemma 5.1
Consider an $`n`$-pasting shell $`\xi =S^\xi ,\pi ^\xi ,ϵ^\xi ,\mathrm{{\rm Y}}^\xi ,\{\sigma _{s,s^{}}^\xi \}_{s,s^{}\mathrm{{\rm Y}}^\xi }`$. Let a condition $`\mathrm{\Psi }(y_0,y_1,\mathrm{},y_m;x_0,x_1,\mathrm{},x_m)`$ ($`1m`$) be abbreviated that
* $`x_0,x_1,x_2,x_3,\mathrm{},x_{m1},x_m\mathrm{{\rm Y}}_{n2}^\xi `$,
* $`y_0,y_1,y_2,y_3,\mathrm{},y_{m1},y_m\mathrm{{\rm Y}}_{n1}^\xi `$,
* $`\pi ^\xi (x_i)=y_i`$ and
* $`\sigma _{x_{m1},x_m}^\xi \mathrm{}\sigma _{x_2,x_3}^\xi \sigma _{y_1,y_2}^\xi \sigma _{x_0,x_1}^\xi (x_0)=x_m`$.
Note that same nodes may be duplicated in parameters of $`\mathrm{\Psi }`$; in paticular, $`y_0`$ may be equal to $`y_m`$. Then
1. For every $`y_0,y_1,\mathrm{},y_m`$ and $`x_0,x_1,\mathrm{},x_m`$ satisfying $`\mathrm{\Psi }(y_0,y_1,\mathrm{},y_m;x_0,x_1,\mathrm{},x_m)`$, we have $`\mathrm{\Psi }(y_m,y_{m1},\mathrm{},y_0;x_m,x_{m1},\mathrm{},x_0)`$.
2. For every open node $`y_0`$ and its child $`x_0`$ , there uniquely exist $`y_0,y_1,\mathrm{},y_m`$ and $`x_0,x_1,\mathrm{},x_m`$ satisfying $`\mathrm{\Psi }(y_0,y_1,\mathrm{},y_m;x_0,x_1,\mathrm{},x_m)`$ and that $`y_m`$ is an open node ($`y_1,\mathrm{},y_{m1}`$ are not open by the second condition of $`\mathrm{\Psi }`$).
Proof (1) Trivial from the conditions of frame shells. (2) Starting from $`y_0`$ and $`x_0`$, we can uniquely fix a required sequence $`y_0,x_0,x_1,y_1,y_2,x_2,\mathrm{}`$ by the following process: For $`x_{2i}`$ ($`0i`$), $`x_{2i+1}`$ is uniquely determined by the bijectivity of links; then $`\pi ^\xi (x_{2i+1})=y_{2i+1}`$ and for $`y_{2i+1}`$, $`y_{2i+2}`$ is again uniquely determined by the bijectivity of links; therefore for $`x_{2i+1}`$, $`\sigma _{y_{2i+1},y_{2i+2}}(x_{2i+1})=x_{2i+2}`$ is also unique. Next, we show that this process necessarily gets to an open node. Every link $`x,x^{}\mathrm{{\rm Y}}_{n2}^\xi `$ appears at most once in the process because for a link to appear twice means the existence of a link $`y,y_0\mathrm{{\rm Y}}_{n1}^\xi `$ for some $`y`$, and this contradicts that $`y_0`$ is open. Since $`S_{n2}`$ is finite and so is $`\mathrm{{\rm Y}}_{n2}^\xi `$, the process starting from an open node $`y_0`$ reaches an open node $`y_m`$ for a finite $`m`$ and stops there. $`\mathrm{}`$
###### Proposition 5.2 (closer and closure of an $`n`$-pasting shell)
For any $`n`$-pasting shell $`\xi =S^\xi ,\pi ^\xi ,ϵ^\xi ,\mathrm{{\rm Y}}^\xi ,\{\sigma _{s,s^{}}^\xi \}_{s,s^{}\mathrm{{\rm Y}}^\xi }`$, we can construct a closer of $`\xi `$, an $`n`$-cell shell $`\widehat{\xi }=r^{\widehat{\xi }},S^{\widehat{\xi }},\pi ^{\widehat{\xi }},ϵ^{\widehat{\xi }},\mathrm{{\rm Y}}^{\widehat{\xi }},\{\sigma _{s,s^{}}^{\widehat{\xi }}\}_{s,s^{}\mathrm{{\rm Y}}^{\widehat{\xi }}}`$, and a closure of $`\xi `$, an $`n`$-frame shell $`\overline{\xi }=S^{\overline{\xi }},\pi ^{\overline{\xi }},ϵ^{\overline{\xi }},\mathrm{{\rm Y}}^{\overline{\xi }},\{\sigma _{s,s^{}}^{\overline{\xi }}\}_{s,s^{}\mathrm{{\rm Y}}^{\overline{\xi }}}`$ uniquely up to isomorphisms and polarity as follows: Let $`\{s_0,s_1,\mathrm{},s_k\}`$ be the set of open nodes of $`\xi `$. For each $`s_l`$, we prepare an $`n1`$-cell shell $`\tau _l=t_l,S^{\tau _l},\pi ^{\tau _l},ϵ^{\tau _l},\mathrm{{\rm Y}}^{\tau _l},\{\sigma _{t,t^{}}^{\tau _l}\}_{t,t^{}\mathrm{{\rm Y}}^{\tau _l}}`$, isomorphic to $`(\xi |^{s_l})^{}`$ via an isomorphsim $`f_l:\tau _l(\xi |^{s_l})^{}`$. Then the components of the closer $`\widehat{\xi }`$ are:
* $`S_n^{\widehat{\xi }}=\{r^{\widehat{\xi }}\}`$ where $`\{r^{\widehat{\xi }}\}`$ is a singleton set, and $`S_i^{\widehat{\xi }}=_{0lk}S_i^{\tau _l}`$ for $`0in1`$,
* $`\pi _{n1}^{\widehat{\xi }}=\pi _{r^{\widehat{\xi }}}`$ where $`\pi _{r^{\widehat{\xi }}}(t_l)=r^{\widehat{\xi }}`$ for $`0lk`$, and $`\pi _i^{\widehat{\xi }}=_{0lk}\pi _i^{\tau _l}`$ for $`0in2`$,
* $`ϵ^{\widehat{\xi }}=(_{0lk}ϵ^{\tau _l})ϵ_{r^{\widehat{\xi }}}`$ where $`ϵ_{r^{\widehat{\xi }}}(r^{\widehat{\xi }})=1`$ (the negative closer) or $`1`$ (the positive closer),
* $`\mathrm{{\rm Y}}^{\widehat{\xi }}=(_{0lk}\mathrm{{\rm Y}}^{\tau _l})\{f_l^1(x),f_l^{}^1(x^{})|\mathrm{\Phi }(s_l,s_l^{};x,x^{})\}`$, where $`\mathrm{\Phi }(s_l,s_l^{};x,x^{})`$ is abbreviated that there exist $`y_0=s_l,y_1,y_2,\mathrm{},y_m=s_l^{}`$ and $`x_0=x,x_1,x_2,\mathrm{},x_m=x^{}`$ satisfying $`\mathrm{\Psi }(y_0,y_1,\mathrm{},y_m;x_0,x_1,\mathrm{},x_m)`$,
* $`\{\sigma _{s,s^{}}^{\widehat{\xi }}\}_{s,s^{}\mathrm{{\rm Y}}^{\widehat{\xi }}}`$ is defined as
+ $`\sigma _{t,t^{}}^{\widehat{\xi }}=\sigma _{t,t^{}}^{\tau _l}`$ for $`t,t^{}\mathrm{{\rm Y}}^{\tau _l}`$,
+ $`\sigma _{f_l^1(x),f_l^{}^1(x^{})}^{\widehat{\xi }}=f_l^{}^1\sigma _{x_m,x}^\xi \mathrm{}\sigma _{x_2,x_3}^\xi \sigma _{y_1,y_2}^\xi \sigma _{x_0,x_1}^\xi f_l`$.
and the components of the closure $`\overline{\xi }`$ are:
* $`S^{\overline{\xi }}=S^\xi S^{\widehat{\xi }}`$,
* $`\pi ^{\overline{\xi }}=\pi ^\xi \pi ^{\widehat{\xi }}`$,
* $`ϵ^{\overline{\xi }}=ϵ^\xi ϵ^{\widehat{\xi }}`$,
* $`\mathrm{{\rm Y}}^{\overline{\xi }}=\mathrm{{\rm Y}}^\xi \{s_l,t_l,t_l,s_l|\mathrm{\hspace{0.17em}0}lk\}\mathrm{{\rm Y}}^{\widehat{\xi }}`$,
* $`\{\sigma _{s,s^{}}^{\overline{\xi }}\}_{s,s^{}\mathrm{{\rm Y}}^{\overline{\xi }}}`$ is defined as
+ $`\sigma _{s,s^{}}^{\overline{\xi }}=\sigma _{s,s^{}}^\xi `$ for $`s,s^{}\mathrm{{\rm Y}}^\xi `$,
+ $`\sigma _{t_l,s_l}^{\overline{\xi }}=f_l`$ and $`\sigma _{s_l,t_l}^{\overline{\xi }}=f_l^1`$,
+ $`\sigma _{t,t^{}}^{\overline{\xi }}=\sigma _{t,t^{}}^{\widehat{\xi }}`$ for $`t,t^{}\mathrm{{\rm Y}}^{\widehat{\xi }}`$.
Proof First, we will check conditions for $`f_l^1(x),f_l^{}^1(x^{})`$ and $`\sigma _{f_l^1(x),f_l^{}^1(x^{})}^{\overline{\xi }}`$. Other parts are rather easy:
* The mutuality condition is obvious. The bijectivity condition is shown by the Lemma 5.1 (2) and the involution condition by the Lemma 5.1 (1).
* The conjugation condition is derived from the following results: $`ϵ^\xi (x_{2i})ϵ^\xi (x_{2i+1})=1`$ by the conjugation of $`\xi `$, $`ϵ^\xi (x_{2i+1})ϵ^\xi (x_{2i+2})=1`$ by the definition of $`\sigma _{y_{2i+1},y_{2i+2}}^\xi `$, $`ϵ^{\tau _l}(f_l^1(x))ϵ^\xi (x)=1`$ and $`ϵ^{\tau _l^{}}(f_l^{}^1(x^{}))ϵ^\xi (x^{})=1`$.
* The correspondence of links condition for $`f_l^1(x),f_l^{}^1(x^{})`$ is shown by chaining the correspondence of links in $`\xi `$ and commutativity of $`f`$, $`f^{}`$.
* The closedness condition is straightforward from the Lemma 5.1 (2).
Next, we will check the commutativity of links condition for the closure. The following three cases are possible:
1. All links are in $`\mathrm{{\rm Y}}^\xi `$;
2. All links are in $`\mathrm{{\rm Y}}^{\widehat{\xi }}`$;
3. Links in $`\{s_l,t_l,t_l,s_l|\mathrm{\hspace{0.17em}0}lk\}`$ occur.
Commutativity for the case 1 is trivial from the definition of pasting diagrams. For the cases 2 and 3, we show commutativity of a path in Figure 5. The oval path (a) is commutative from the definition of pasting diagrams; the square (b) and (c) is from the construction of closers; and the square (d) is by pasting an alternation of the type (b) and (c) squares. Thus the outer path of arrows in this case is commutative by pasting (a)–(d). In the general case 3, the side trip (b)–(d) might occur several times. The largest roundabouts are paths running only through the closer. This implies commutativity for the case 2, that is, the commutativity of links condition for the closer. $`\mathrm{}`$
## 6 Directed $`\omega `$-hypergraphs
### 6.1 directed $`\omega `$-hypergraphs
###### Definition 6.1 (shape graph)
The (undirected) shape graph $`N_0^\xi ,N_1^\xi E_0^\xi ,E_1^\xi `$ of an $`n`$-frame $`\xi `$ is defined as follows:
* the body node set $`N_0^\xi `$ is $`S_n^\xi `$;
* the foot node set $`N_1^\xi `$ is $`S_{n1}^\xi `$;
* the leg edge set $`E_0^\xi `$ is $`\{\{s,t\}|sS_n^\xi ,tS_{n1}^\xi ,\pi ^\xi (t)=s\}`$;
* the link edge set $`E_1^\xi `$ is $`\{\{t,t^{}\}|t,t^{}\mathrm{{\rm Y}}_{n1}^\xi \}`$ (note that $`\pi ^\xi (t)\pi ^\xi (t^{})`$ from the definition of $`\mathrm{{\rm Y}}_{n1}^\xi `$).
If $`sN_0^\xi `$, $`tN_1^\xi `$ and $`\{s,t\}E_0^\xi `$ then we call $`t`$ a foot of $`s`$ and $`\{s,t\}`$ a leg of $`s`$. The shape graph $`N_0^\zeta ,N_1^\zeta E_0^\zeta ,E_1^\zeta `$ of an $`n`$-frame $`\zeta `$ is the shape graph of $`\underset{¯}{\zeta }`$.
The shape graph of $`n`$-cell shells or $`n`$-cell diagrams is defined as a special case of $`n`$-frame shells or $`n`$-frames.
###### Definition 6.2 ($`n`$-directed $`n`$-frame)
An $`n`$-frame $`\zeta `$ is $`n`$-directed if it satisfies the following conditions:
* headedness: for exactly one $`sS_n^\zeta `$ and every other $`s^{}S_n^\zeta `$, either $`\lambda _n^\zeta (s)`$ is positive and $`\lambda _n^\zeta (s^{})`$ is negative or $`\lambda _n^\zeta (s)`$ is negative and $`\lambda _n^\zeta (s^{})`$ is positive, where $`s`$ is called the positive or negative head of $`\zeta `$, respectively;
* connectedness: its shape graph is connected;
* acyclicity: the graph obtained from its shape graph by getting rid of a body node corresponding to the head, its legs and feet, and link edges connected to them, is acyclic (indeed, this graph is a tree).
An $`n`$-frame with the positive head is said to be positively $`n`$-directed, and that with the negative head be negatively $`n`$-directed. The same $`n`$-frame can be both positively and negatively $`n`$-directed. An $`n+1`$-cell whose boundary is such an $`n`$-frame is called a simple $`n+1`$-cell.
###### Definition 6.3 (directed $`n`$\- and $`\omega `$-hypergraph)
A directed $`\omega `$-hypergraph is an $`\omega `$-hypergraph which satisfies the following condition:
* directedness: for each $`i1`$, the boundary of any positive $`i`$-cell is a positively $`i1`$-directed $`i1`$-frame, and that of any negative $`i`$-cell is a negatively $`i1`$-directed $`i1`$-frame.
For each $`n0`$, a directed $`n`$-hypergraph is also defined as an $`n`$-hypergraph satifying the same condition.
### 6.2 directed shells
In the category theory besed on $`n`$\- or $`\omega `$-hypergraphs, directed $`n`$-cell shells and directed $`n`$-frame shells play the role of shape diagrams in the usual theory. They are defined by adding some conditions to the induction step of the definitions of $`n`$-cell shells and $`n`$-frame shells
###### Definition 6.4 (directed $`n`$-cell shell)
An additional condition is as follows:
* directedness: If $`ϵ(r)=1`$, then $`\theta |_r`$ is a positively directed $`n1`$-frame shell and if $`ϵ(r)=1`$, then it is a negatively directed one.
###### Definition 6.5 (directed $`n`$-frame shell)
Additional conditions are as follows:
* headedness: for exactly one $`sS_n^\xi `$ and every other $`s^{}S_n^\xi `$, either $`ϵ_n^\xi (s)=1`$ and $`ϵ_n^\xi (s^{})=1`$ or $`ϵ_n^\xi (s)=1`$ and $`ϵ_n^\xi (s^{})=1`$, where $`s`$ is called the positive or negative head of $`\xi `$, respectively;
* connectedness: its shape graph is connected;
* acyclicity: the graph obtained from its shape graph by getting rid of a body node corresponding to the head, its legs and feet, and link edges connected to them, is acyclic (indeed, this graph is a tree).
An $`n`$-frame shell with the positive head is said to be positively directed, and that with the negative head be negatively directed. The same $`n`$-frame shell can be both positively and negatively directed. If for an $`n+1`$-cell shell $`\theta `$ with root $`r`$, $`\theta |_r`$ is such an $`n`$-frame shell, then it is called a simple $`n+1`$-cell shell.
###### Proposition 6.1
An $`n`$\- or $`\omega `$-hypergraph $`\mathrm{\Sigma },`$ is a directed $`n`$\- or $`\omega `$-hypergraph iff for each positive $`i`$-cell $`c`$, $`\underset{¯}{_i(c)}`$ is a positivery directed $`i1`$-frame shell and for each negative $`i`$-cell, it is a negatively directed one.
Proof By induction on dimensions. $`\mathrm{}`$
###### Definition 6.6 (directed $`n`$-frame and directed $`n`$-cell diagram)
An $`n`$-frame $`\zeta `$ is a positively or negatively directed $`n`$-frame if $`\underset{¯}{\zeta }`$ is a positively or negatively directed $`n`$-frame shell, respectively. Also a positively or negatively directed $`n`$-cell diagram is defined in the same way.
###### Corollary 6.2
In any directed $`n`$\- or $`\omega `$-hypergraph, an $`n`$-frame is a positively or negatively directed $`n`$-frame iff it is an positively or negatively $`n`$-directed $`n`$-frame, respectively.
###### Definition 6.7 ($`\mathrm{𝐃𝐅𝐫𝐦}_k`$)
For a directed $`n`$\- or $`\omega `$-hypergraph, the category (groupoid) whose objects are all directed $`k`$-frames and whose arrows are all isomorphisms is denoted by $`\mathrm{𝐃𝐅𝐫𝐦}_k`$. The collection of all directed $`k`$-frames is also denoted by $`\mathrm{𝐃𝐅𝐫𝐦}_k`$.
In the rest of this paper, we will mainly use the usual diagramatic notations for shells and diagrams (Figure 7).
### 6.3 examples
###### Example 6.1 (hypergraph in rewriting)
An (directed) hypergraph used in hypergraph rewritinng is a directed $`1`$-hypergraph.
###### Example 6.2 ($`\omega `$-multigraph)
An $`\omega `$-multigraph is a directed $`\omega `$-hypergraph $`\mathrm{\Sigma },`$ such that $`\mathrm{\Sigma }_{0,1}`$ is a singlton set and that any $`c\mathrm{\Sigma }_{1,1}`$ is a simple $`1`$-cell.
###### Example 6.3 (doublegraph)
Doublegraphs are underlying graph-like structures for double categories. They are obtained by splitting $`1`$-cells into vertical cells and horizontal cells:
$$\mathrm{\Sigma }_0=\mathrm{\Sigma }_{0,1}\mathrm{\Sigma }_{0,1}$$
$$\mathrm{\Sigma }_1=\mathrm{\Sigma }_{1,1}\mathrm{\Sigma }_{1,1}$$
$$\mathrm{\Sigma }_{1,1}=\mathrm{\Sigma }_{1,1}^v\mathrm{\Sigma }_{1,1}^h$$
$$\mathrm{\Sigma }_{1,1}=\mathrm{\Sigma }_{1,1}^v\mathrm{\Sigma }_{1,1}^h$$
and if $`c\mathrm{\Sigma }_{1,k}`$, then $`c^{}\mathrm{\Sigma }_{1,k}`$, etc. $`2`$-cells are as follows:
This appoach can be easily extended to multiple categories.
###### Example 6.4 (fc-multigraph)
fc-multigraphs are underlying graph-like structures for fc-multicategories introduced by T. Leinster . Similar notions also appear in . They are a mixture of $`2`$-multigraphs and double graphs. 0-cells and 1-cells are the same as double graphs. $`2`$-cells are as follows
### 6.4 directed pasting shells and diagrams
###### Definition 6.8 (boundary graphs)
The shape graph of an $`n`$-pasting shell $`\xi `$ is defined in the same way of $`n`$-frame shells. The boundary graph of an $`n`$-pasting shell $`\xi `$ is defined as the shape graph of $`(\widehat{\xi }|_{r^{\widehat{\xi }}})^{}`$.
Note that the boundary graph of an $`n`$-cell shell $`\theta `$ as a special case of $`n`$-pasting shells, matches with the shape graph of $`\theta |_{r^\theta }`$.
###### Definition 6.9 (directed $`n`$-pasting shells)
A directed $`n`$-pasting shell $`\xi `$ is an $`n`$-pasting shell consisting of directed $`n`$-cell shells satisfying the following conditions:
* homegeneity: it is positive or negative as an $`n`$-pasting shell.
* connectedness: its shape graph is connected;
* acyclicity: its shape graph is acyclic (indeed, this graph is a tree).
###### Definition 6.10 (directed $`n`$-pasting diagrams)
A directed $`n`$-pasting diagram is an $`n`$-pasting diagram whose base $`n`$-pasting shell is a directed $`n`$-pasting shell. A positive or negative directed $`n`$-pasting diagram is trivially defined, respectively.
###### Proposition 6.3
For any directed $`n`$-frame $`\zeta `$ and its head node $`h`$, we can uniquely split it into an directed $`n`$-cell diagram $`\mathrm{cod}(\zeta )`$, called the codomain of $`\zeta `$, and an $`n`$-pasting diagram $`\mathrm{dom}(\zeta )`$, called the domain of $`\zeta `$, where
* $`\mathrm{cod}(\zeta )=\zeta |^h`$
* $`\mathrm{dom}(\zeta )`$ is obtained by deleting, from the data of $`\zeta `$, $`\zeta |^h`$ and indice $`s,s^{}`$ of which $`s`$ or $`s^{}`$ is in $`\mathrm{{\rm Y}}^{\zeta |^h}`$.
Of course, if $`\zeta `$ is positively directed, then $`\mathrm{cod}(\zeta )`$ is positive and $`\mathrm{dom}(\zeta )`$ is negative, and if negatively directed, then $`\mathrm{cod}(\zeta )`$ is negative and $`\mathrm{dom}(\zeta )`$ is positive.
###### Proposition 6.4 (closure of a directed $`n`$-pasting shell)
For any negative (resp. positive) directed $`n`$-pasting shell $`\xi `$, (1) its positive (resp. negative) closer $`\widehat{\xi }`$ is a directed $`n`$-cell shell and (2) its positive (resp. negative) closure $`\overline{\xi }`$ is a directed $`n`$-frame shell.
Proof (1) We have to check that $`\widehat{\xi }|_{r^{\widehat{\xi }}}`$ is a positively directed $`n1`$-frame shell. (i) Headedness: From homogeneity, connectedness and acyclicity of $`\xi `$, the set $`O_\xi `$ of open nodes of $`\xi `$ contains just one positive node. Therefore $`\widehat{\xi }|_{r^{\widehat{\xi }}}`$ contains just one negative node. (ii) Connectedness and acyclicity: By induction on the construction of directed $`n`$-pasting shells. (a) The boundary graph of a directed $`n`$-cell shell satisfies the connectedness and acyclicity conditions. (b) If you make a directed $`n`$-pasting shell $`\xi ^{}`$ by linking a directed $`n`$-pasting shell $`\xi `$ and a directed $`n`$-cell shell $`\theta `$ at one open $`n1`$-node (with satisfying the homogeneity, connectedness and acyclicity for $`\xi ^{}`$), then replacement of the boundary graph occurs (Figure 8) and the resulting boundary graph of $`\xi ^{}`$ also satisfies the connectedness and acyclicity. The acyclicity is obtained by reduction to absurdity. Suppose the existance of cycles and consider the graph obtained by deleting link edges of cycles from the boundary graph, and its polarity of foot nodes, it contradicts the directedness of the boundary $`n1`$-frame shell $`\theta |_{r^\theta }`$ of each $`n`$-cell shell $`\theta `$ in the directed $`n`$-pasting shell. (c) Any directed $`n`$-pasting shell can be constructed by finitely iterating this process, and then the boundary graph of the closer satisfies the connectedness and acyclicity. Thus $`\widehat{\xi }|_{r^{\widehat{\xi }}}`$ is a positively directed $`n1`$-frame shell. Since $`ϵ^{\widehat{\xi }}(r^{\widehat{\xi }})=1`$, $`\widehat{\xi }`$ is a positively directed $`n`$-cell shell. The negative case is in parallel.
(2) From the homogeneity, connectedness and acyclicity of $`\xi `$, the headedness, connectedness and acyclicity of $`\overline{\xi }`$ is obvious. $`\mathrm{}`$
## 7 Weak $`\omega `$-categories
### 7.1 $`\omega `$-identity, $`\omega `$-invertibility and $`\omega `$-universality
We will coinductively define three notions: $`\omega `$-identity, $`\omega `$-invertibility and $`\omega `$-universality. All $`n`$-dimensional notions depend on $`n+1`$\- or $`n+2`$-dimensional ones. The reader unfamiliar with coinductive definitions may think of only the case in which coinduction steps terminate.
One source of our idea is Michael Makkai’s work on anabicategories . At a glance, as Makkai pointed out in , saturated anabicategories could be regarded as 2-dimensional weak cateogires of Baez-Dolan. But we don’t think of them to be equivalent notions for some reasons:
1) In anabicategories, two composite arrows for the same composable sequence of arrows are equivalent. While in Baez-Dolan’s there are two opposite universal 2-cells between two composite arrows by virtue of balancedness, no explicit relation between them appears. In fact, we can prove that they are equivalences in a sense, because Baez and Dolan only think of finite dimensional cases. But we cannot prove it in that way for infinite dimensional cases. Therefore we need to characterize those opposites as a sort of equivalences.
2) Different from anabicategories, composites of empty sequence are introduced in Baez-Dolan’s and expected to play the role of identities. But as well as the above, we cannot prove the property of identity in infinite dimensional cases. Thus we also have to define identities explicitly.
3) In (not necessarily saturated) anabicategories, an object isomorphic to a composite might not be a composite. It suggests that if we introduce a sort of equivalences, it is natural to treat equivalences and composition separatedly and add a saturatedness condition.
###### Definition 7.1 ($`\omega `$-identical cells)
An $`n`$-cell $`c`$ is $`\omega `$-identical if it is simple and satisfies $`\mathrm{dom}(c)\mathrm{cod}(c)`$ and the following conditions:
$$\begin{array}{cccc}f\hfill & & \alpha \hfill & \\ \text{}\hfill & & \text{}\hfill & \alpha \text{}\omega \text{-universal}\hfill \end{array}$$
and
$$\begin{array}{cccc}g\hfill & & \beta \hfill & \\ \text{}\hfill & & \text{}\hfill & \beta \text{}\omega \text{-universal.}\hfill \end{array}$$
###### Definition 7.2 ($`\omega `$-invertibility, $`\omega `$-equivalence, $``$)
A pair of $`n`$-cells $`f`$ and $`g`$ is an $`\omega `$-invertible pair if both $`f`$ and $`g`$ are simple and satisfy $`\mathrm{dom}(f)\mathrm{cod}(g)`$ and $`\mathrm{dom}(g)\mathrm{cod}(f)`$ and the following condtions:
$$\begin{array}{cc}\alpha ,i\hfill & \\ \text{}\hfill & i\text{}\omega \text{-identical,}\hfill \end{array}$$
and
$$\begin{array}{cc}\beta ,j\hfill & \\ \text{}\hfill & j\text{}\omega \text{-identical.}\hfill \end{array}$$
$`f`$ and $`g`$ are called $`\omega `$-invertible. We say that two $`n1`$ cells $`\lambda (r^{\mathrm{dom}(f)})`$ and $`\lambda (r^{\mathrm{cod}(f)})`$ are $`\omega `$-equivalent and write it as $`\lambda (r^{\mathrm{dom}(f)})\lambda (r^{\mathrm{cod}(f)})`$.
###### Definition 7.3 ($`\omega `$-universal cells)
An $`n`$-cell $`u`$ is $`\omega `$-universal if for any $`n`$-cell $`f`$ of
there exist an $`n`$-cell $`g`$ and an $`\omega `$-universal $`n+1`$-cell $`\alpha `$ to be
and for two such pairs, $`g`$ and $`\alpha `$, $`h`$ and $`\beta `$
there exist an $`\omega `$-invertible pair of $`n+1`$-cells, $`\gamma `$ and $`\delta `$, and two $`\omega `$-universal $`n+2`$-cells, $`\mathrm{\Phi }`$ and $`\mathrm{\Psi }`$, such that
### 7.2 weak $`\omega `$-categories
###### Definition 7.4 (weak $`\omega `$-categories)
A directed $`\omega `$-hypergraph is a weak $`\omega `$-category if it satisfies the following conditions:
* existence of closers and occupants: For any $`n`$-pasting diagram $`P`$,
there exist an $`\omega `$-universal $`n+1`$-cell $`\alpha `$ and an $`n+1`$-cell $`h`$ such that $`\mathrm{dom}(\alpha )_nP`$ and $`\lambda (r^{\mathrm{cod}(\alpha )})=h`$:
Moreover, if $`P`$ is an empty pasting diagram, then for each $`n1`$-cell diagram $`x`$, there exist such $`\alpha `$ and $`h`$ satisfying $`\mathrm{dom}(h)_{n1}\mathrm{cod}(h)_{n1}x`$.
We call $`\alpha `$ an occupant for $`P`$ and $`h`$ a closer of $`P`$.
* weak uniqueness of closers and occupants: For two such pairs as above, $`\alpha `$ and $`h`$, $`\beta `$ and $`k`$
there exist an $`\omega `$-invertible pair $`\gamma `$ and $`\delta `$ and two $`\omega `$-universal cells $`\mathrm{\Phi }`$ and $`\mathrm{\Psi }`$ such that
* saturation of closers and occupants: For an $`\omega `$-universal cell $`\alpha `$ and an $`\omega `$-invertible pair $`\gamma `$ and $`\delta `$
there exists $`\omega `$-universal cells $`\beta `$, $`\mathrm{\Phi }`$, $`\mathrm{\Psi }`$ such that
* $`\omega `$-identical closers (1): For any $`\omega `$-universal $`n+1`$-cell $`\alpha `$ as follows, an $`n`$-cell $`i`$ is $`\omega `$-identitical:
* $`\omega `$-identical closers (2): For any $`\omega `$-identitical $`n`$-cell $`i`$, there is an $`\omega `$-universal $`n+1`$-cell $`\alpha `$ as above.
* $`\omega `$-universal closers: Any closer for a $`n`$-pasting diagram made of $`\omega `$-universal $`n`$-cells is $`\omega `$-universal.
A weak $`\omega `$-categories is defined when this coinductive definition makes sense.
###### Proposition 7.1
An $`\omega `$-identical cell is $`\omega `$-invertible.
###### Proposition 7.2
Every two $`\omega `$-identical cells are $`\omega `$-equivalent.
Again, the reader unfamiliar with coinduction may think of weak $`n`$-categories.
###### Definition 7.5 (weak $`n`$-categories)
A weak $`\omega `$-category is a weak $`n`$-category if for each $`k`$ higher than $`n`$, all simple $`k`$-cells are $`\omega `$-invertible.
From the axioms above, we can recognize that identity is independent from the definition of composition. In fact, to define our weak $`\omega `$-categories, we can exclude empty pasting diagram and related axioms and add an axiom for existance of $`\omega `$-identical $`n+1`$-cells for each $`n`$-cells. This is slightly simpler than those defined above. And we conjecture that our definition would be equivalent to that of J. Penon , and furthermore that if we abandon saturatedness and for each $`n`$-cell we choose just one $`\omega `$-identical $`n+1`$-cell whose domain and codomain are that $`n`$-cell, then the category of our small weak $`\omega `$-categories (with suitable functors) would be isomorphic to the category of Penon’s.
Acknowledgements: The first author thanks Prof. Sounders Mac Lane for his comments and powerful encouragement for his talks at CT97 and CT99, and Prof. Michael Makkai for thoughtful explanation of his ideas. We also thank A. Higuchi for helpful discussion. The first author’s work is partially supported by Kyoto Sangyo University.
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# Solutions to WDVV from generalized Drinfeld-Sokolov hierarchies
## 1 Introduction and discussion
Equations of Witten, Dijkgraaf, Verlinde and Verlinde \[WDVV\]
$$\begin{array}{c}\hfill c_{ijk}(t)\eta ^{kl}c_{lmn}(t)=c_{njk}(t)\eta ^{kl}c_{lmi}(t),\text{where}c_{ijk}(t)=_i_j_kF(t)\\ \hfill \eta _{mn}^1=_1_m_nF(t)\text{ – constant,}\underset{i=1}{\overset{n}{}}d_it_i_iF(t)=\left(3d\right)F(t)\end{array}$$
(1.1)
describe deformations of two dimensional topological conformal quantum field theory. They have appeared to be intimately related to integrable systems. This was recognized by B. Dubrovin \[Du1\], who showed that each solution to WDVV equations gives rise to dispersionless bi Hamiltonian integrable system, that is with Poisson structures of hydrodynamic type \[DN\] and Hamiltonians not depending on field derivatives. This result was extended by one loop correction \[DuZ\], where it was shown that so constructed bi Hamiltonian structure is always a $`𝒲`$ algebra, i.e. it contains a Virasoro subalgebra.
On the other hand, due to B. Dubrovin \[Du4\], it is possible, under certain assumption, to recover solution to WDVV equation from bi Hamiltonian structure of hydrodynamic type
$$\{u^i(x),u^j(y)\}=g^{ij}\left(u(x)\right)\delta ^{}\left(xy\right)+\mathrm{\Gamma }_k^{ij}\left(u(x)\right)u_k^{}(x)\delta \left(xy\right).$$
If matrix $`g`$ does not degenerate identically, it defines flat metrics \[DN\] on the target space.
Whitham averaging method \[DN\] provides Poisson structures of this type. This method allows to describe slow modulated $`m`$-phase solutions of non linear Hamiltonian system of equations. If one needs local Poisson structures of hydrodynamic type we should apply zero phase averaging being simply dispersionless limit.
The dispersionless limit of Drinfeld - Sokolov hierarchies \[DS\], associated to non twisted affine Lie algebra $`\widehat{𝔤}`$ was shown \[Kr\] to provide polynomial solutions to WDVV equations (1.1). These solutions, viewed as Frobenius structures on orbits of Coxeter groups $`𝒢`$ \[Du2\], correspond to Weyl group $`𝒢=𝒲\left(𝔤\right)`$ of underlying simple Lie algebra $`𝔤`$. They correspond \[Du5\] to the orbit of braid group on Stocks matrices $`S`$ passing through “Coxeter” point, i.e. $`S+S^{tr}`$ defines a Carter graph \[C\] of Coxeter conjugacy class in $`𝒲\left(𝔤\right)`$. Orbits given by primitive conjugacy classes were conjectured \[Du5\] to correspond to algebraic solutions of WDVV equations (1.1).
In this article we address the dispersionless limit of generalized Drinfeld - Sokolov hierarchies \[GHM, BGHM\], associated with $`(i)`$ non twisted affine Lie algebra $`\widehat{𝔤}`$, $`(ii)`$ its Heisenberg subalgebra $`^{[w]}`$, corresponding \[KP\] to regular primitive conjugacy class $`[w]`$ in $`𝒲\left(𝔤\right)`$, and $`(iii)`$ regular element $`\mathrm{\Lambda }`$ from it. The naive limit is shown to exist only for standard Drinfeld-Sokolov hierarchies, due to failure to parameterize the whole phase space by densities $`𝒩_i`$ of annihilators of the first Poisson structure. As a consequence, some fields generically evolve fast. We combine Hamiltonian reduction with Whitham averaging \[M\], restricting Poisson structures to the subvariety parameterized by annihilators’ densities and then taking the limit $`ϵ0`$. So constructed dispersionless limit is proved to give rise to algebraic solution to WDVV equations (1.1). Monodromy group of the Frobenius manifold coincides with $`𝒲\left(𝔤\right)`$.
The paper is organized as follows. In section 2 we gather, to make the paper self-content, some preliminary information about gradations of Kac-Moody algebra and about its Heisenberg subalgebras. Section 3 reviews basic facts about generalized Drinfeld - Sokolov hierarchies. The properties of regular primitive conjugacy classes, relevant for the paper, are recollected in section 4. Section 5 brings together facts about type I integrable hierarchies with primitive regular conjugacy class $`[w]`$ and regular element $`\mathrm{\Lambda }`$. Quite rich underlying finite dimensional geometry is analysed in section 6. Finite group $`R`$, acting non linearly on Miura variables, and a ring of $`R`$-invariant polynomials are studied. The dispersionless limit of the hierarchies in consideration is examined in section 7. The non degeneracy of obtained metrics is proved, their flat coordinates are found and it is shown that Frobenius structure can be always extracted from this averaged Poisson structure. Finally, section 8 treats the simplest integrable hierarchy with non Coxeter regular primitive conjugacy class, i.e. $`𝔤=D_4`$ and $`[w]=[D_4(a_1)]`$. Corresponding solution to WDVV equations is found and finite group $`R`$ is analysed.
## 2 Preliminaries
Let $`𝔤`$ be a simple Lie algebra of rank $`r`$, and $`\widehat{𝔤}=𝔤[z,z^1]d`$ – its affine Lie algebra. This algebra is well-known to be graded, with gradations being in one-to-one correspondence \[K\] with finite order inner automorphisms of $`𝔤`$.
###### Definition 2.1.
Let $`𝐬=(s_0,s_1\mathrm{},s_r)`$ – a sequence of non negative relatively prime integers. Let $`N_𝐬=_{i=1}^rk_is_i`$, where $`k_i`$ – are Kac labels (such that $`\alpha _{\text{max}}=_{i=1}^rk_i\alpha _i`$ and $`k_0=1`$). We then define the finite order automorphisms $`\sigma `$ of $`𝔤`$ in some Cartan-Weyl basis by
$$\sigma \left(H\right)=H,\sigma \left(E_\alpha \right)=e^{2\pi i\delta _𝐬\alpha }E_\alpha ,$$
(2.1)
with
$$\delta _𝐬=\frac{1}{N_𝐬}\underset{k=1}{\overset{r}{}}\frac{2}{\alpha _i^2}s_k\omega _k;\alpha _i^{}\omega _j=\delta _{ij}.$$
###### Proposition 2.1 (\[K\]).
Such automorphisms exhaust the finite order inner automorphisms of $`𝔤`$ up to conjugacy.
An inner automorphisms of $`𝔤`$ can be used to define a new $``$-gradation of $`\widehat{𝔤}`$.
###### Definition 2.2.
A gradation of type $`𝐬`$ is defined via derivation $`d_𝐬`$
$$d_𝐬=N_𝐬\left(z\frac{d}{dz}+\mathrm{ad}H_{\delta _𝐬}\right),$$
(2.2)
where $`H_{\delta _𝐬}𝔥`$ such that $`[H_{\delta _𝐬},E_{\alpha _k}]=\left(\alpha _k\delta _𝐬\right)E_{\alpha _k}=s_kE_{\alpha _k}/N_𝐬`$.
Let us denote $`\widehat{𝔤}_k\left(𝐬\right)`$ the eigenspace of derivation $`d_𝐬`$ with an eigenvalue $`k`$:
$$\widehat{𝔤}=\underset{k}{}\widehat{𝔤}_k\left(𝐬_w\right),[d_𝐬,\widehat{𝔤}_k\left(𝐬_w\right)]=k\widehat{𝔤}_k\left(𝐬_w\right).$$
Homogeneous gradation corresponds then to $`𝐬_h=(1,0,\mathrm{},0)`$, while principal to $`𝐬_p=(1,1,\mathrm{},1)`$. One can introduce a partial ordering on the set of gradation.
###### Definition 2.3.
We say that $`𝐬𝐬^{}`$ if $`s_i0`$ whenever $`s_i^{}0`$.
###### Proposition 2.2 (\[GHM\]).
If $`𝐬𝐬^{}`$ then the following is true
(i) $`\widehat{𝔤}_0(𝐬)\widehat{𝔤}_0(𝐬^{})`$,
(ii) $`\widehat{𝔤}_{>0}(𝐬)\widehat{𝔤}_0(𝐬^{})`$ and $`\widehat{𝔤}_{<0}(𝐬)\widehat{𝔤}_0(𝐬^{})`$,
(iii) $`\widehat{𝔤}_{>0}(𝐬^{})\widehat{𝔤}_{>0}(𝐬)`$ and $`\widehat{𝔤}_{<0}(𝐬^{})\widehat{𝔤}_{<0}(𝐬)`$.
###### Theorem 2.1 (\[KP\]).
Given a Kac-Moody algebra $`\widehat{𝔤}`$, its maximal nilpotent subalgebras are called Heisenberg subalgebras. Heisenberg subalgebras are in one-to-one correspondence with conjugacy classes of Weyl group $`𝒲\left(𝔤\right)`$ of $`𝔤`$.
Hence, they will be denoted by $`^{[w]}`$. Let us recall the essentials of their construction.
Let $`w`$ be a representative of a given conjugacy class $`[w]`$ acting naturally on $`𝔥`$ in some Cartan-Weyl basis marked by a prime to distinguish it from another basis (2.1) in use. Its action on Cartan subagebra $`𝔥`$ is known to be inner in $`𝔤`$. Then, let $`Y𝔤`$ be such an element that $`w=\mathrm{exp}\left(2\pi i\mathrm{ad}Y\right)_{|𝔥}`$. Let $`\widehat{w}`$ be an obvious extension of $`w`$ on the whole algebra $`𝔤`$:
$$\widehat{w}H_\alpha ^{}=H_{w(\alpha )}^{},\widehat{w}E_\alpha ^{}=\psi _\alpha E_{w(\alpha )}^{},$$
(2.3)
where $`\psi _\alpha `$ are structure constant compatible factors taking values in $`S^1`$ and may be chosen to be real. It can be shown \[KP\] that the order of this extension $`N_{\widehat{w}}`$ may only be either the same as or twice as the order of $`w`$. The following construction does not, however, depend on this ambiguity.
Let $`a𝔤`$ be a representative of some orbit of $`\widehat{w}`$. We write $`a=_\xi a_\xi `$ as a sum of $`\widehat{w}`$ eigenvectors $`a_\xi `$ with distinct eigenvalues $`\xi =\mathrm{exp}\left[2\pi i\frac{k}{N_{\widehat{w}}}\right]`$ and put $`\stackrel{~}{a}(k)\widehat{𝔤}`$ for any $`k`$ to be
$$\stackrel{~}{a}(k)=\mathrm{exp}\left[i\phi \frac{k}{N_{\widehat{w}}}\right]\mathrm{exp}\left[i\phi \mathrm{ad}Y\right]a_\xi +\delta _{k,0}\left(a|Y\right)c,$$
(2.4)
where $`c`$ is the central element of Kac-Moody algebra. It is obvious that $`\stackrel{~}{a}(k)`$ is well defined on $`S^1`$ and hence depends on $`\phi `$ through $`z=\mathrm{exp}(i\phi )`$ only. It verifies commutation relations of Kac-Moody algebra:
$$[\stackrel{~}{a}(k),\stackrel{~}{b}(l)]=\stackrel{~}{[a,b]}\left(k+l\right)+k\delta _{k,l}\left(a|b\right)c.$$
The image of $`𝔥`$ under this map is a subalgebra of $`\widehat{𝔤}`$ isomorphic to an infinite dimensional Heisenberg algebra:
$$[\stackrel{~}{h}_1(k),\stackrel{~}{h}_2(l)]=k\delta _{k,l}\left(h_1|h_2\right)c.$$
Factorizing out the center $`𝔷=c`$ of Kac-Moody algebra $`\widehat{𝔤}`$, the Heisenberg subalgebras become maximal commutative ones. A practical realization of this construction faces the difficulty to find $`Y`$ explicitly.
Associated to each Heisenberg subalgebra $`^{[w]}`$ there is a distinguished gradation $`𝐬_w`$ with the property that $`^{[w]}`$ admits $`𝐬_w`$ grading \[GHM\]:
$$^{[w]}=\underset{kE}{}_k^{[w]}\left(𝐬_w\right),$$
with
$$\begin{array}{cc}\hfill E& =\left\{k+N_{\widehat{w}}\right|kI(w)\},\hfill \\ \hfill I(w)& =\{k_n_+,n=1,\mathrm{},r|\mathrm{\hspace{0.17em}0}k_1k_2\mathrm{}k_r<N_{\widehat{w}},\hfill \\ & e^{2\pi ik_n/N_{\widehat{w}}}\text{ – eigenvalue of }w\}.\hfill \end{array}$$
(2.5)
It is clear, that $`𝐬_w`$ is the gradation corresponding to automorphism $`w`$ used to construct $`^{[w]}`$. Indeed, let us choose $`Y𝔤`$ so that spectrum of $`N_{\widehat{w}}\mathrm{ad}Y`$ on $`𝔥`$ coincides with $`I\left(w\right)`$. Then
$$d_𝐬=N_{\widehat{w}}\left(z\frac{d}{dz}+\mathrm{ad}Y\right),d_𝐬\left(\stackrel{~}{h}\left(k\right)\right)=k\stackrel{~}{h}\left(k\right).$$
###### Definition 2.4.
An element $`\mathrm{\Lambda }\widehat{𝔤}`$ is called *semisimple* if $`\widehat{𝔤}=\mathrm{Ker}\left(\mathrm{ad}\mathrm{\Lambda }\right)\mathrm{Im}\left(\mathrm{ad}\mathrm{\Lambda }\right)`$ and $`\mathrm{Ker}\left(\mathrm{ad}\mathrm{\Lambda }\right)\mathrm{Im}\left(\mathrm{ad}\mathrm{\Lambda }\right)=\mathrm{}`$.
###### Definition 2.5.
A semisimple element $`\mathrm{\Lambda }^{[w]}`$ is called *regular* if $`\mathrm{Ker}\left(\mathrm{ad}\mathrm{\Lambda }\right)=^{[w]}`$.
## 3 Generalized integrable hierarchies
Let $`[w]`$ be some conjugacy class in $`𝒲\left(𝔤\right)`$, $`𝐬_w`$ – corresponding gradation. Pick up any $`𝐬𝐬_w`$ and some constant semisimple element $`\mathrm{\Lambda }^{[w]}`$ of certain $`𝐬_w`$ grade $`iE_+`$. Let us define the matrix Lax operator \[GHM\]
$$=_x+\mathrm{\Lambda }+q(x),$$
(3.1)
where $`q(x)`$ takes value in
$$Q\left(i\right)=\widehat{𝔤}_0\left(𝐬\right)\widehat{𝔤}_{<i}\left(𝐬_w\right).$$
(3.2)
###### Proposition 3.1 (\[GHM, DS, W\]).
There exist a unique formal series
$$T(q)\mathrm{Im}\left(\mathrm{ad}\mathrm{\Lambda }\right)\widehat{𝔤}_{<0}\left(𝐬_w\right)$$
that
$$L=\mathrm{exp}\left[\mathrm{ad}T\right]=_x+\mathrm{\Lambda }+h\left(q\right),h(q)=\underset{k<i}{}h_k(q)$$
(3.3)
where $`h\left(q\right)_{<i}^{[w]}\left(𝐬_w\right)`$. Both $`P_k^{(𝐬_w)}T(q)`$ and $`P_k^{(𝐬_w)}h(q)`$ are polynomials in $`q(x)`$ and its $`x`$-derivatives for any $`k`$.
There is a gauge symmetry acting on (3.1) and mapping $`Q(i)`$ to $`Q(i)`$:
$$\mathrm{exp}\left[\mathrm{ad}n(x)\right],$$
(3.4)
where $`n(x)`$ takes values in
$$P=\widehat{𝔤}_0\left(𝐬\right)\widehat{𝔤}_{<0}\left(𝐬_w\right).$$
(3.5)
###### Proposition 3.2 (\[DS, FGMS\]).
If $`\mathrm{\Lambda }`$ is chosen so that
$$\mathrm{Ker}\left(\mathrm{ad}\mathrm{\Lambda }\right)P=\mathrm{}\mathrm{ad}\mathrm{\Lambda }:PQ,$$
(3.6)
then elements of $`Q^{\text{can}}\left(i\right)^{}`$, where
$$Q(i)=Q^{\text{can}}(i)[\mathrm{\Lambda },P],$$
(3.7)
are generators of the ring of gauge invariant polynomials in $`q(x)`$ and its derivatives; (3.7) holds as equality of vector spaces. Lax operator $`L`$ in (3.3) is gauge invariant.
###### Definition 3.1.
Denote by $``$ the phase space, spanned by the elements of $`Q^{can}\left(i\right)^{}`$ and by $`\left(\right)`$ the space of functionals of the type
$$\phi \left[q\right]=_{S^1}𝑑xf(x,q,q^{},\mathrm{}).$$
###### Proposition 3.3 (\[GHM\]).
Dimension of phase space $``$ is independent on auxiliary gradation $`𝐬`$ and equals to
$$dim=dimQ^{\text{can}}(i)^{}=\underset{k=0}{\overset{i1}{}}dim\widehat{𝔤}_k\left(𝐬_w\right)$$
(3.8)
###### Remark 3.1.
$`dimdim\widehat{𝔤}_0(𝐬_w)dim\widehat{𝔤}_0(𝐬_p)=r`$.
Condition (3.6) is automatically satisfied \[FGMS\] if $`\mathrm{\Lambda }`$ is chosen to be regular.
###### Theorem 3.1 (\[GHM\]).
Given $`b_{>0}^{[w]}\left(𝐬_w\right)`$, let $`𝒜\left(b\right)=\mathrm{exp}\left[\mathrm{ad}T\right]b`$. One defines two sets of time flows
$$\frac{}{t_b}=[P_0^{(𝐬)}𝒜\left(b\right),]\frac{}{t_b^{}}=[P_0^{(𝐬_w)}𝒜\left(b\right),]$$
(3.9)
being commutative within each set
$$b_1,b_2^{[w]}[\frac{}{t_{b_1}},\frac{}{t_{b_2}}]=0,[\frac{}{t_{b_1}^{}},\frac{}{t_{b_2}^{}}]=0.$$
Both time flows (3.9) preserve the phase space of gauge invariants $``$. They coincide there and retain their commutativity property.
###### Definition 3.2.
Introduce the pairing on the space of functions in $`C^{\mathrm{}}(S^1,\widehat{𝔤})`$
$$A,B=_{S^1}𝑑x\underset{k}{}\eta (P_k^{(𝐬)}A(x),P_k^{(𝐬)}B(x))_{\widehat{𝔤}_0(𝐬)}$$
(3.10)
for some gradation $`𝐬`$ and Killing form $`\eta `$ on $`\widehat{𝔤}_0(𝐬)`$.
If $`\eta _{\widehat{𝔤}_0(𝐬)}`$ is properly normalized, the pairing does not depend on the gradation $`𝐬`$ chosen. We will assume this normalization chosen further on.
###### Proposition 3.4 (\[DS, GHM\]).
Gauge invariant functionals $`H_b\left[q\right]=b,h(q)`$ are integrals of flows (3.9).
###### Remark 3.2.
$`H_b0`$ if $`\mathrm{deg}_{𝐬_w}\left(b\right)i`$.
###### Definition 3.3.
For any $`\phi `$ define gradient $`d_q\phi \widehat{𝔤}_0\left(𝐬\right)/\widehat{𝔤}_{<i}\left(𝐬_w\right)`$ by
$$\frac{d}{d\epsilon }\phi \left[q+\epsilon r\right]_{|_{\epsilon =0}}=r,d_q\phi rC^{\mathrm{}}(S^1,Q(i))$$
###### Theorem 3.2 (\[BGHM\]).
Let $`𝐬=𝐬_h`$, then
(i) there is a one parameter family of Hamiltonian structures on the gauge equivalence classes of the generalized Drinfeld-Sokolov hierarchy given by
$$\{\phi ,\psi \}_\lambda =\mathrm{\Lambda }+q,[d_q\phi ,d_q\psi ]_{R^{(𝐬)}}d_q\phi ,\left(d_q\psi \right)^{}$$
(3.11)
where $`R^{(𝐬)}=\left(P_0^{(𝐬)}P_{<0}^{(𝐬)}\right)/2\lambda /z`$. Expanding in powers of $`\lambda `$, $`\{,\}_\lambda =\{,\}_2+\lambda \{,\}_1`$, we obtain two coordinated Hamiltonian structures on $``$:
$`\{\phi ,\psi \}_1`$ $`=`$ $`d_q\phi ,z^1[d_q\psi ,],`$
$`\{\phi ,\psi \}_2`$ $`=`$ $`\mathrm{\Lambda }+q,[d_q\phi ,d_q\psi ]_Rd_q\phi ,\left(d_q\psi \right)^{},`$
where $`R=\left(P_0^{(𝐬)}P_{<0}^{(𝐬)}\right)/2`$. Under the time evolution in the coordinate $`t_b`$, the following recursion relation holds:
$$\frac{\phi }{t_b}=\{\phi ,H_{zb}\}_1=\{\phi ,H_b\}_2.$$
(3.12)
(ii) Hamiltonians $`H_b`$ with $`i<\mathrm{deg}_{𝐬_w}b<0`$ are Casimirs of (3.11). Hamiltonians $`H_b`$ with $`\mathrm{deg}_{𝐬_w}b=0`$ are Casimirs of the first bracket only.
Hierarchies (3.12) with $`\mathrm{\Lambda }`$ regular were dubbed in ref. \[GHM\] the hierarchies of type I.
###### Remark 3.3.
The hierarchies with grade one element $`\mathrm{\Lambda }`$ are distinguished by the fact that pencil of Poisson structures has no local annihilators. This is true also for the first structure as long as conjugacy class $`[w]`$ is chosen to be non degenerate, that is fixing no vector in $`𝔥`$ and consequently $`0I(w)`$.
###### Corollary 3.1.
Hamiltonians $`H_b`$, with $`b^{[w]}`$ of positive $`𝐬_w`$ degree belonging to $`I(w)`$, annihilate the first Poisson structure.
There are $`r`$ independent annihilators of $`\{,\}_1`$, that yet generate nontrivial flows with respect to $`\{,\}_2`$.
###### Proof.
Choose $`b`$ such that $`0\mathrm{deg}_{𝐬_w}\left(zb\right)<N_w`$. Then $`N_w\mathrm{deg}_{𝐬_w}\left(b\right)<0`$ and $`H_b`$ annihilates $`\{,\}_2`$ or vanishes identically.
To prove independence, it is enough to notice that equation (3.3) determining $`h(q)`$, projected on eigenspaces of $`d_{𝐬_w}`$ reads
$$h_k+[\mathrm{\Lambda },T_{k+1}]=q_k+\text{Polynomial}(h_{i<k},T_{jk},q_{i<k}).$$
Thus, $`h_k`$ with $`kI(w)`$ are linear in Drinfeld-Sokolov variable of scaling degree $`k+1`$ with successive terms involving variables of lower scaling weights and this proves independence.
Notice, that $`𝐬=𝐬_h`$ corresponds to the largest gauge group, while $`𝐬=𝐬_w`$ – to no gauge freedom at all. The former choice leads to generalized Drinfeld-Sokolov hierarchies, while the latter – to modified generalized Drinfeld-Sokolov hierarchies.
###### Theorem 3.3 (\[BGHM\]).
<sup>1</sup><sup>1</sup>1In \[BGHM\] it was proved for partially modified hierarchies, that is for any $`𝐬`$ such that $`𝐬_h𝐬𝐬_w`$.
Let $`_m`$ be the phase space of modified hierarchy, and $`\{,\}_m`$ – its second Poisson bracket:
$$\{\phi ,\psi \}_m=q_m+\mathrm{\Lambda },[d_{q_m}\phi ,d_{q_m}\psi ]_{R_m}d_{q_m}\phi ,\left(d_{q_m}\psi \right)^{},$$
(3.13)
The Miura mapping is a (non invertible) Hamiltonian mapping,
$$\mu :(_m,\{,\}_m)(,\{,\}_2),$$
such that it defines a reduction of the dynamical equations of the hierarchy to those of the modified hierarchy.
Drinfeld-Sokolov hierarchies \[DS\] are recovered picking up Coxeter conjugacy class $`[w_c]`$ with representative $`w_c=_{k=1}^rr_{\alpha _k}`$, associated principal gradation ($`𝐬_w=𝐬_p`$) and regular element $`\mathrm{\Lambda }=zE_{\alpha _{max}}+_{k=1}^rE_{\alpha _k}`$ of grade one.
## 4 Regular primitive conjugacy classes of $`𝒲\left(𝔤\right)`$
Conjugacy classes in Weyl group $`𝒲\left(𝔤\right)`$ of simple Lie algebra $`𝔤`$ were uniformly described by R. Carter \[C\] in terms of Carter graphs.
###### Definition 4.1.
Conjugacy class $`[w]𝒲\left(𝔤\right)`$ is called *non degenerate* if $`det(1w)0`$ and is called *primitive* if $`det(1w)=det(1w_c)=det(K)`$, $`K`$ – being the Cartan matrix.
Primitive conjugacy classes are distinguished as they have no representative in any proper Weyl subgroup $`W^{}𝒲\left(𝔤\right)`$. All primitive conjugacy classes were listed by R. Carter \[C\] together with their characteristic polynomials $`det(tw)`$.
###### Definition 4.2 (\[DF\]).
Conjugacy class $`[w]𝒲\left(𝔤\right)`$ is called *regular* if associated Heisenberg subalgebra $`^{[w]}`$ admits regular element $`\mathrm{\Lambda }`$.
Regular conjugacy classes were elegantly studied by T. Springer in ref. \[Sp\], though another definition of regularity was used.
###### Definition 4.3 (\[Sp\]).
If $`G`$ is a finite reflection group in a finite dimensional vector space $`V`$, then $`vV`$ is called *regular* if no nonidentity element of $`G`$ fixes $`v`$.
An element $`gG`$ is *regular* if it has a regular eigenvector.
Two definitions were shown \[DF\](see also sec. 9 of \[Sp\]) to be equivalent. The authors of the reference \[DF\] studied generalized hierarchies associated with classical Lie algebras $`𝔤`$ and regular conjugacy classes. They synthesized available information on regular primitive conjugacy classes.
###### Theorem 4.1 (\[Sp\]).
Let $`[w]𝒲\left(𝔤\right)`$ be a regular conjugacy class of order $`N`$, i.e. $`w[w],w^N=1`$. Then eigenvalues of $`w`$ are $`\mathrm{exp}\left(2\pi ik_n/N\right)`$, where $`k_nI(w_c)`$ are exponents. The elements of the root system $`R`$ of $`𝔤`$ are permuted by $`w[w]`$ in orbits of length $`N`$.
So $`\{1,N1\}I(w)`$, that implies $`𝐬_{w}^{}{}_{0}{}^{}=1`$.
###### Theorem 4.2 (\[Sp\]).
Let $`[w]𝒲\left(𝔤\right)`$ be regular primitive conjugacy class. Let $`\mathrm{\Lambda }=I_++zC_{(N1)}_1^{[w]}(𝐬_w)`$ be grade one regular element. Then
(i) $`I_{}𝔤`$ of $`𝐬_w`$ grade $`1`$, such that $`\rho =N_w\delta _{𝐬_w}`$, $`I_+`$ and $`I_{}`$ form $`sl_2`$ subalgebra
$$[\rho ,I_\pm ]=\pm I_\pm [I_+,I_{}]=2\rho ;$$
(4.1)
(ii) Eigenvalues of $`\mathrm{ad}\rho `$ on $`𝔤`$ are integers;
(iii)
$$dim𝔤(0)=dim𝔤(\pm 1)=\frac{2}{N}\underset{k=1}{\overset{r}{}}p_i=\frac{rh}{N},p_iI\left(w_c\right)$$
(4.2)
###### Corollary 4.1.
The extension $`\widehat{w}`$ of $`w`$ to the whole algebra $`𝔤`$ is of order $`N`$.
###### Proof.
The order of extension does not depend on the basis chosen. In the basis of (2.1) $`\widehat{w}=\mathrm{exp}\left[2\pi i\mathrm{ad}\rho /N\right]`$. Then $`\widehat{w}^N=1`$ as a straightforward consequence of statement $`(ii)`$ of the theorem. ∎
Dimensions of $`𝔤(k)`$ may be read off from the character of this (reducible) representation of $`sl_2`$ on $`𝔤`$:
$$\chi _{[w]}\left(q\right)=\mathrm{Tr}\left(q^{\mathrm{ad}\rho }\right)=r+\underset{\alpha R}{}q^{(\alpha ,\rho )}=\underset{k=(N1)}{\overset{N1}{}}dim𝔤(k)q^k.$$
(4.3)
To determine the $`𝐬`$ type of regular primitive conjugacy class we compare $`\chi _{\left[\widehat{w}\right]}`$ in the basis of (2.1) and in that of (2.3):
$$\mathrm{Tr}w=r+2\underset{\alpha R^+}{}\mathrm{cos}\left(2\pi \delta _{𝐬_w}\alpha \right)=\chi _{[w]}\left(e^{2\pi i/N_w}\right)$$
(4.4)
The knowledge of the characteristic polynomials of $`[w]`$ due to \[C\] allows to compute the order $`N_w`$ of $`[w]`$. Then one searches for $`𝐬_w`$ with $`N_w=N`$ and $`𝐬_{w}^{}{}_{0}{}^{}=1`$, satisfying (4.4) and (4.2). This and some more (see sec. 5) information about regular primitive conjugacy classes relevant for construction of generalized hierarchies is collected in appendix A.
## 5 Type-I, $`i=1`$ hierarchies
Here and further on we restrict ourselves to hierarchies with a regular primitive conjugacy class and a regular element $`I_++zC_{(N1)}=\mathrm{\Lambda }^{[w]}`$ of $`𝐬_w`$ grade one. Notice, that whatever auxiliary gradation $`𝐬`$ we choose, $`Q(i)𝔤`$ and we are working with finite dimensional Lie algebra. Adjoint action of $`\rho `$ induces gradation of $`𝔤`$:
$$𝔤=\underset{k=N+1}{\overset{N1}{}}𝔤_k.$$
Notice, that $`\widehat{𝔤}_0\left(𝐬_w\right)`$ coincides with $`𝔤_0`$.
Consider first homogeneous auxiliary gradation $`𝐬=𝐬_h`$. As follows from definition 3.3 and proposition 2.2 gradients have no negative $`𝐬_w`$ grade part: $`P_{<0}^{(𝐬_w)}d_q\phi =0mod\widehat{𝔤}_{<1}(𝐬_w)`$. So the bi Hamiltonian structure (3.11) simplifies to
$$\begin{array}{cc}\hfill \{\phi ,\psi \}_1& =d_q\phi ,[d_q\psi ,C_{(N1)}],\hfill \\ \hfill \{\phi ,\psi \}_2& =d_q\phi ,[d_q\psi ,_x+I_++q].\hfill \end{array}$$
(5.1)
As we have no $`R`$ matrix in commutator now, we can recognize in the second bracket a Kirillov-Poisson bracket corresponding to untwisted affinization $`\widehat{𝔤}_x`$ in $`x`$ of $`𝔤`$. Introducing $`J=I_++q`$ we obtain
$$\{\phi ,\psi \}_{KM}=d_J\phi ,[d_J\psi ,_x+J].$$
(5.2)
This fact enabled authors of ref. \[BFRFW\] to invent a practical algorithm to compute Hamiltonian structure of Drinfeld-Sokolov hierarchies and their construction can be repeated here. They showed that the second Poisson structure (5.1) may be obtained by Hamiltonian reduction of (5.2) (see also \[CFMP\]). Let us embed $`Q^{can}\widehat{𝔤}_x`$ and lift flows on $`Q^{can}`$ to flows on $`\widehat{𝔤}_x`$. One may choose functional $`\mathrm{\Psi }`$ on $`\widehat{𝔤}_x`$ coinciding with $`\psi `$ on $`Q^{can}`$ such that $`\mathrm{\Psi }`$ flows in $`\widehat{𝔤}_x`$ preserve $`Q^{can}`$. Then
$$\delta q=[d_J\mathrm{\Psi }_{|Q^{can}},_x+I_++q]=[d_q\psi (x)+r,_x+I_++q(x)]Q^{can},$$
(5.3)
where $`\eta (r,Q^{can})=0`$. Condition (5.3) determines $`r`$ uniquely as a function of $`q`$ and its $`x`$-derivatives.
The grading of Kac-Moody algebra induces a grading on the Poisson structure, namely – the following theorem holds true.
###### Theorem 5.1 (\[BGHM\]).
The dynamical equation (3.9) of the hierarchy generated by the Hamiltonian $`H_b`$ ( $`b_j^{[w]}(𝐬_w)`$) with respect to the second Hamiltonian structure, is invariant under the following rescaling
$$x\lambda x,t_b\lambda ^jt_b,q_k\lambda ^{(k+1)}q_k,$$
(5.4)
where $`q_k`$ is the component of $`q(x)`$ with $`𝐬_w`$ grade $`k`$.
This is a consequence of existence of Virasoro algebra due to Sugawara construction.
###### Theorem 5.2 (\[FGMS, BFRFW\]).
The second Poisson bracket of (5.1) is a $`𝒲`$ algebra associated to $`sl_2`$ subalgebra (4.1) with Virasoro generator
$$w_2(x)=\frac{1}{2}\eta (I_++q,I_++q)+\eta (\rho ,q^{}),$$
(5.5)
satisfying Virasoro algebra
$$\{w_2(x),w_2(y)\}=\eta (\rho ,\rho )\delta ^{\prime \prime \prime }\left(xy\right)+2w_2(x)\delta ^{}\left(xy\right)+w_2^{}(x)\delta \left(xy\right).$$
(5.6)
There exists a minimal weight Drinfeld-Sokolov gauge \[BFRFW, DF\]
$$\mathrm{ad}I_{}\left(Q^{can}\right)=0$$
(5.7)
of conformal primaries:
$$\{w_{k,i}(x),w_2(y)\}_2=kw_{k,i}(x)\delta ^{}\left(xy\right)+w_{k,i}^{}(x)\delta \left(xy\right),$$
(5.8)
where
$$Q^{can}q^{can}=w_2\frac{I_{}}{\eta (I_+,I_{})}+w_{k,i}F_{k1,i},[\rho ,F_{k,i}]=kF_{k,i}.$$
Let us denote by $`𝖯𝗋_w`$ a set of scaling weights of fields of the $`𝒲`$-algebra in question. $`𝖯𝗋_w`$ is invariant with respect to scale preserving changes of coordinates, so in particular it does not depend on the Drinfeld-Sokolov gauge slice chosen.
###### Proposition 5.1.
$`I(w)𝖯𝗋_w`$. Multiplicities of $`N_w1`$ in $`I(w)`$ and in $`𝖯𝗋_w`$ are equal.
###### Proof.
Let us take densities of Hamiltonians $`H_b`$, $`\mathrm{deg}_{𝐬_w}bI(w)`$. Their scaling degree are correspondingly $`\mathrm{deg}_{𝐬_w}b+1`$. These densities are gauge invariant, as follows from proposition 3.2. Due to corollary to the theorem 3.2 they are independent. Thus they may be taken as a part of coordinates on $``$. So $`I(w)𝖯𝗋_w`$. In particular this implies that the multiplicity of $`N_w1`$ in $`I(w)`$ is less or equal to the one in $`𝖯𝗋_w`$.
However, as was shown in ref. \[Sp\], for every $`C𝔤(N_w+1)`$ there exist such $`I𝔤(1)`$ that $`I+C`$ is regular in $`𝔤`$. So we have an opposite inequality of multiplicities. This completes the proof. ∎
###### Remark 5.1.
One can choose $`F_{N_w1,1}=C_{(N_w1)}`$. Then the first Poisson structure can be always read off the second by shifting $`w_{N_w,1}`$ by $`\lambda `$ and taking the linear term.
The character (4.3) enables us to determine $`\mathrm{ad}\rho `$ eigenspace decomposition of $`Q^{can}`$, and thus $`𝖯𝗋_w`$. Indeed,
$$\left(1q\right)\chi _{[w]}=\underset{k=N+1}{\overset{1}{}}dimQ_k^{can}q^k+𝒪\left(q\right),$$
(5.9)
as follows from injectivity (3.6) of mapping $`\mathrm{ad}I_+:Q_kQ_{k+1}`$ for negative $`k`$ and definition of $`Q^{can}`$ (3.7). Then $`𝖯𝗋_w`$ is the set of positive numbers $`k`$ such that $`dimQ_k^{can}0`$, every number occurring $`dimQ_k^{can}`$ times.
###### Proposition 5.2.
Let
$$𝒰\left(w\right)=\left\{1k_1\mathrm{}k_{2n}<N_w1|k_i\underset{w}{𝖯𝗋}\text{but }k_iI(w)\right\},$$
(5.10)
where $`2n=dimr`$. Then $`k_i+k_{2n+1i}=N_w1`$.
###### Proof.
Brute force checking of data presented in Appendix A. We lack more elegant proof so far. ∎
Now, consider modified hierarchy, i.e. $`𝐬=𝐬_w`$. Its second Poisson structure (3.13) simplifies as well
$$\{\phi ,\psi \}_m=q_m,[d_{q_m}\phi ,d_{q_m}\psi ]d_{q_m}\phi ,\left(d_{q_m}\psi \right)^{}.$$
The phase space (3.2) of modified hierarchies $`(𝔤,\mathrm{\Lambda },𝐬_w)`$ reads
$$=\left(Q^m(1)\right)^{}=\left(\widehat{𝔤}_0(𝐬_w)\right)^{}\widehat{𝔤}_0(𝐬_w).$$
(5.11)
Let $`X_i`$ be a basis of $`\widehat{𝔤}_0(𝐬_w)𝔤`$ with Grahm matrix $`𝒦`$ being restricted Killing form and let $`q_m=_i\nu ^iX_i`$, then
$$\{\nu ^i(x),\nu ^j(y)\}_m=\left(𝒦^1\right)^{ij}\delta ^{}\left(xy\right)+f_{}^{ij}{}_{k}{}^{}\nu ^k(x)\delta \left(xy\right),$$
(5.12)
where $`f_{}^{ij}{}_{k}{}^{}`$ are structure constants of $`\widehat{𝔤}_0(𝐬_w)^{}`$. Notice, that $`𝔥\widehat{𝔤}_0(𝐬_w)`$ due to proposition 2.2, and thus matrix $`f_{}^{ij}{}_{k}{}^{}\nu ^k`$ is of corank $`r`$.
Given $`[w]`$ one obtains $`\widehat{𝔤}_0(𝐬_w)`$ from extended Dynkin diagram removing nodes with $`𝐬_{w}^{}{}_{k}{}^{}0`$ and padding with $`u(1)`$ to maintain the rank as follows from (2.2).
We collect all pertinent information about generalized Drinfeld-Sokolov hierarchies with regular primitive conjugacy class $`[w]`$ and regular element $`\mathrm{\Lambda }^{[w]}`$ of $`𝐬_w`$ grade one in appendix A.
## 6 Finite dimensional geometry behind the hierarchy
Let us introduce the dispersion parameter $`ϵ`$ by rescaling all derivatives $`ϵ`$. Then, the leading term of $`ϵ0`$ expansion of a Poisson bracket on loop space will verify Jacobi identity, thus furnishing another Poisson structure:
$$\{w^i(x),w^j(y)\}^{(0)}=\frac{1}{ϵ}A^{ij}(w)\delta \left(xy\right).$$
(6.1)
We write upper indices to emphasize the covariant nature of objects. The leading bracket (6.1) obviously defines finite dimensional one and we thus end up with pencil of finite dimensional Poisson brackets corresponding to $`𝐬=𝐬_h`$ integrable hierarchy:
$$\overline{\{w^i,w^j\}}_\lambda =A^{ij}\left(w\right)\lambda B^{ij}\left(w\right).$$
(6.2)
Notice, that densities of annihilators of the first Poisson bracket (5.1) $`𝒩^a=𝒩^a\left(w\right)`$, where $`1ar`$ are Casimirs of the first finite bracket:
$$\overline{\{f,𝒩^a\}}_1=0f,$$
(6.3)
and are in involution with respect to the second
$$\overline{\{𝒩^a,𝒩^b\}}_2=0.$$
(6.4)
Eqs. (6.3) and (6.4) show that Poisson tensor $`A\lambda B`$ is of corank $`r`$ for all $`\lambda `$ and generic point of $``$.
Due to Corollary 3.1, one may choose $`w^i=(𝒩^a,u^A)`$ as coordinates on $``$, where $`u^A`$ are Drinfeld-Sokolov variables with scaling weights $`k_A+1`$, $`k_A𝒰\left(w\right)`$.
Assume that the coordinate $`𝒩^r`$ is chosen to be linear in Drinfeld-Sokolov variable $`w_{N_w,1}`$ (cf. Remark 5.1), and so the pencil (6.1) can be obtained by means of shifting of $`𝒩^r`$ by $`\lambda `$.
###### Lemma 6.1.
Matrix $`B^{AB}(u,𝒩)`$ is non degenerate with constant determinant.
###### Proof.
As follows from scaling weight grading of Poisson structures
$$\mathrm{deg}_{sc}\left(A^{AB}\lambda B^{AB}\right)=\mathrm{deg}u^A+\mathrm{deg}u^B1.$$
Since $`\mathrm{deg}_{sc}\lambda =N_w`$
$$B^{AB}=\{\begin{array}{cc}0\hfill & \text{ if }\mathrm{deg}_{sc}u^Au^B<N_w\hfill \\ \text{const. }\hfill & \text{ if }\mathrm{deg}_{sc}u^Au^B=N_w\hfill \\ B^{AB}\hfill & \text{ if }\mathrm{deg}_{sc}u^Au^B>N_w\hfill \end{array}$$
Recall that we have assumed ascending ordering of $`\mathrm{deg}_{sc}u^A`$, and so matrix $`B`$ is lower triangular with respect to anti diagonal $`\mathrm{deg}_{sc}u^Au^B=N_w`$. This proves that $`detB`$ equals to product of those constants. Since, by the choice of coordinates, we know that $`\left|\begin{array}{c}B^{AB}\end{array}\right|`$ is nondegenerate we conclude, that all those constants are nonzero. ∎
###### Remark 6.1.
Notice, that $`𝒩^1=w_2`$ annihilates the pencil of finite dimensional brackets: $`\overline{\{f,𝒩^1\}}_\lambda =0`$ for any $`f`$.
###### Proposition 6.1.
Let $`𝒢^{A,a}(u,𝒩)=\overline{\{u^A,𝒩^a\}}_2`$. Consider a variety $`_r`$ defined by equations $`𝒢^{A,2}(u,𝒩)=0`$ and suppose that matrix $`\left|\begin{array}{c}\frac{𝒢^{A,2}}{u^B}\end{array}\right|`$ is not degenerate on $`_r`$. Then
$$𝒢^{B,b}(u,𝒩)_{|_r}=0B,b.$$
###### Proof.
By the implicit function theorem one may resolve the system of polynomial equations $`𝒢^{A,2}(u,𝒩)=0`$ locally with respect to variables $`u`$. Consider the Hamiltonian flows in involution generated by $`𝒩^a`$:
$$\frac{\phi }{t_a}=\overline{\{\phi ,𝒩^a\}}_2.$$
Then one has
$$\frac{}{t_a}\frac{}{t_2}u^A=\frac{}{t_2}\frac{}{t_a}u^A\frac{𝒢^{A,2}}{u^B}𝒢^{B,a}=\frac{𝒢^{A,a}}{u^B}𝒢^{B,2}.$$
Restricting on $`_r`$ and using the non degeneracy of $`\frac{𝒢^{A,2}}{u^B}`$ the result follows. ∎
###### Remark 6.2.
So, the equation $`𝒢^{A,2}=0`$, subject to condition of Proposition 6.1 defines an algebraic variety $`_r`$ of stationary points of all Hamiltonians in involution with respect to Poisson bracket (6.2). $`_r`$ “corresponds” to the kernel of Poisson tensor of (6.2).
As an immediate consequence of eqs. (6.3) and (6.4) we have the following
###### Proposition 6.2.
Polynomials $`𝒢^{A,a}(u,𝒩)`$ do not depend on $`𝒩^r`$.
As follows from generalized Drinfeld-Sokolov construction, this finite dimensional bracket is but Kirillov-Kostant bracket
$$\overline{\{\phi ,\psi \}}_\lambda =\left(\mathrm{\Lambda }+q|[d_q\phi ,d_q\psi ]\right)$$
(6.5)
restricted on the space of $`\mathrm{ad}P`$ invariant functions $`\phi \left(\stackrel{~}{q}\right)=\phi \left(q\right)`$:
$$\mathrm{\Lambda }+\stackrel{~}{q}=\mathrm{exp}\left[\mathrm{ad}n\right]\left(\mathrm{\Lambda }+q\right)nP$$
(6.6)
###### Lemma 6.2.
Drinfeld-Sokolov coordinates $`w^i(q)`$, in dispersionless limit, generate the ring of $`\mathrm{ad}P`$ invariant polynomials in $`q`$.
###### Proof.
Consider a ring of polynomial functions in $`q`$. It is obviously generated by elements of $`Q^{}`$. Using the adjoint action of nilpotent group with Lie algebra being $`P`$ one may always reduce $`\stackrel{~}{q}`$ to canonical form $`qQ^{can}`$ (see eq. (3.7)). Indeed, projecting eq. (6.6) on $`\mathrm{ad}\rho `$ eigenspaces we get
$$\stackrel{~}{q}_i=q_i+[n_{i+1},I_+]+\text{polynomial}(n_i,\mathrm{},n_1;q_{i1},\mathrm{},q_1)$$
(6.7)
where $`[\rho ,q_i]+iq_i=0`$ and $`0i<N_w`$. This allows to resolve, starting from $`i=0`$ and proceeding inductively, for $`\stackrel{~}{q}`$ and $`n(q)`$ as polynomials in $`q`$. Recapitulating, $`(\stackrel{~}{q},n)`$ generate a ring of polynomials in $`q`$. The subring of gauge invariant polynomials is, then, obviously generated by $`\stackrel{~}{q}`$, i.e. by Drinfeld-Sokolov coordinates. ∎
Now consider the case of modified hierarchy, $`𝐬=𝐬_w`$. The leading term Poisson structure gives Kirillov-Kostant bracket on $`\widehat{𝔤}_0(𝐬_w)^{}`$:
$$\overline{\{\nu ^i,\nu ^j\}}_m=f_{}^{ij}{}_{k}{}^{}\nu ^k.$$
(6.8)
The Miura map (3.3) provides a map from modified hierarchy ($`𝐬=𝐬_w`$) into $`𝐬=𝐬_h`$ one. Thus we have polynomial expressions $`w^i\left(\nu \right)=w^i\left(q_m\right)`$ for dispersionless Drinfeld-Sokolov variables. Notice, that modified hierarchies have their “gauge algebra” $`P`$ empty, and yet Miura coordinates $`\nu `$ do not generate the ring of gauge invariant polynomials. The reason is that, imposing $`\stackrel{~}{q}Q_m`$, equations (6.7) do not have a unique solution. Following ideas of ref. \[BFRFW\] we have
###### Proposition 6.3.
There is a finite subgroup $`R\mathrm{exp}\left(\mathrm{ad}P\right)`$ acting on $`Q_m=𝔤_0`$.
$$\mathrm{ord}R\underset{k𝖯𝗋_w}{}\left(k+1\right).$$
(6.9)
###### Proof.
Let $`nP`$ corresponds to group element fixing $`Q_m`$, i.e.
$$I_++\stackrel{ˇ}{q}_m=\mathrm{exp}\left[\mathrm{ad}n\right]\left(I_++q_m\right)=\underset{k=0}{\overset{N_w+1}{}}\frac{1}{k!}\mathrm{ad}^kn\left(I_++q_m\right).$$
(6.10)
Projecting on $`\mathrm{ad}\rho `$ eigenspaces we gain the system of equations (6.7). Since $`q_m𝔤_0`$, we have
$$\stackrel{ˇ}{q}_m=q_m+[n_1,I_+].$$
(6.11)
The rest of equations (6.7) for $`i>0`$ determines $`n`$. Given positive $`i`$ there are $`dim𝔤_i`$ scalar equations. Due to injectivity of map $`\mathrm{ad}I_+:𝔤_k𝔤_{k+1}`$ for negative $`k`$ we may unambiguously solve for $`dimg_{i1}`$ unknowns contained in $`n_{i+1}`$ and remain with $`m_i=dim𝔤_idimg_{i1}`$ scalar equations for $`n_j`$, where $`j<i`$. Notice, that by definition, $`m_i0`$ iff $`i𝖯𝗋_w`$ with $`m_i`$ being the multiplicity of $`i`$ in $`𝖯𝗋_w`$. Starting with $`i=1`$ case
$$[n_2,I_+]+[n_1,q_m+\frac{1}{2}[n_1,I_+]]=0,$$
(6.12)
we solve for $`n_2=n_2(n_1,q_m)`$. Proceeding further with excluding $`n_k`$, $`k>1`$, we end up with $`m_i`$ equations of order $`i+1`$, for each distinct $`iPr_w`$, to determine $`dim`$ unknowns contained in $`n_1`$. Notice, that the number of equations $`\mathrm{ord}𝖯𝗋_w`$ equals to the number of unknowns $`dim𝔤_1=dim`$ and by Bezout theorem we obtain no more than $`_{k𝖯𝗋_w}\left(k+1\right)`$ solutions $`n_1=n_1\left(q_m\right)`$, which determine that number of nonlinear transformations $`q_m\stackrel{ˇ}{q}_m`$ of $`Q_m`$. If only those $`dim`$ equations for $`n_1`$ are independent we obtain an equality in eq. (6.9). ∎
Since $`n_1g_1`$ completely determines any transformation from $`R`$, as follows from the proof of Proposition 6.3, it is tempting to make an anzäts for the simplest transformations as $`n𝔤_1`$. It proves to be consistent only for simply laced Lie algebras $`𝔤`$. We may then rewrite eq. (6.10) as follows
$$\stackrel{ˇ}{q}_m=q_m+[n,I_+]+\underset{k=1}{\overset{N_w}{}}\frac{1}{k!}\mathrm{ad}^kn\left(\frac{1}{k+1}[n,I_+]+q_m\right).$$
To make the sum, contributing to negative $`\mathrm{ad}\rho `$ eigenspaces, vanish we need that
$$[n,[n,H]]=0,H\widehat{𝔤}_0\left(𝐬_w\right),$$
(6.13a)
$$[n,q_m+\frac{1}{2}[n,I_+]]=0.$$
(6.13b)
To solve equations (6.13) we make use of Weyl group $`𝒲\left(𝔤_0\right)`$ of semisimple Lie subalgebra $`𝔤_0𝔤`$.
###### Lemma 6.3.
Subspaces $`𝔤_k𝔤`$, $`k`$, are stable under natural action of $`𝒲\left(𝔤_0\right)𝒲\left(𝔤\right)`$. Each orbit $`𝒵𝔤_1`$ of $`𝒲\left(𝔤_0\right)`$ is a *commutative* subalgebra of $`𝔤`$, if $`𝔤`$ is simply laced.
###### Proof.
The action of Weyl group $`𝒲\left(𝔤_0\right)`$ are inner in $`𝔤_0`$ and in $`𝔤`$. For any root $`\alpha `$, such that $`E_\alpha 𝔤_0`$, one has a reflection $`r_\alpha `$ in the hyperplane perpendicular to the root acting on $`𝔤`$ as follows:
$$\widehat{r}_\alpha =\mathrm{exp}\left[\mathrm{ad}E_\alpha \right]\mathrm{exp}\left[\frac{2}{\alpha ^2}\mathrm{ad}E_\alpha \right]\mathrm{exp}\left[\mathrm{ad}E_\alpha \right].$$
(6.14)
These reflections act canonically on Cartan subalgebra $`𝔥𝔤_0𝔤`$. Since all elements of $`𝒲\left(𝔤_0\right)`$ can be expressed as a product of these, we conclude that $`𝒲\left(𝔤_0\right)`$ stabilizes $`𝔤_k`$. Take a root $`\beta `$, such that $`E_\beta 𝔤_1`$, and consider the $`𝒲\left(𝔤_0\right)`$ orbit $`𝒵`$ that passes through it. Fix $`w𝒲\left(𝔤_0\right)`$ and assume that $`w`$ does not fix $`\beta `$. It suffices to prove that $`\beta +w\left(\beta \right)`$ is never a root.
Since $`𝒲\left(𝔤_0\right)`$ stabilizes $`𝔤_1`$, $`\gamma \stackrel{\mathrm{def}}{=}w\left(\beta \right)\beta `$ is a linear combination of simple roots of $`𝔤_0`$. For simply laced Lie algebras $`𝔤`$, $`\gamma `$ itself is a root such that $`E_\gamma 𝔤_0`$. Assume that $`\beta +w\left(\beta \right)=2\beta +\gamma `$ is a root, then
$$\frac{2\left(\beta |\gamma \right)}{\left(\beta |\beta \right)}2.$$
(6.15)
This is impossible for any two roots $`\beta `$ and $`\gamma `$ of simply laced simple Lie algebra $`𝔤`$. ∎
###### Corollary 6.1.
For any orbit $`𝒵`$, and any element $`E_\beta 𝒵`$, $`\mathrm{ad}E_\beta `$ maps $`𝔤_0`$ on $`𝒵`$.
###### Proof.
Consider the set $`𝒮\stackrel{\mathrm{def}}{=}\left\{[X,Y]\right|X𝔤_0,Y𝒵\}`$. Since $`𝔥𝔤_0`$ we have $`𝒵𝒮`$. The set $`𝒮`$ is stable under the homomorphism $`𝒲\left(𝔤_0\right)`$ because $`𝔤_0`$ and $`𝒵`$ are stable. So $`𝒮=𝒵`$ and, hence, $`\mathrm{ad}E_\beta `$ maps $`𝔤_0`$ in $`𝒵`$. The surjectivity is obvious. ∎
###### Remark 6.3.
Having the simple roots chosen in $`𝔥^{}`$, it is plain to see, that different orbits $`𝒵`$ may be labeled by simple roots of $`𝐬_w`$ degree one, i.e. such $`\alpha `$ that $`E_\alpha 𝔤_1`$.
Choose an orbit $`𝒵𝔤_1`$, and let $`\{X_q\}`$ denote the set of root vectors that form the basis in $`𝒵`$. We then solve (6.13a) by letting $`n=x_qX_q𝒵`$. Due to Corollary 6.1 we have
###### Lemma 6.4.
Condition (6.13b) gives $`dim𝒵`$ homogeneous quadratic equations for parameters $`x_i`$, that posess non trivial solutions.
Let us speculate a bit on the structure of the group $`R`$. Fix $`q_m𝔤_0`$, some orbit $`𝒵`$, and let $`n\left(q_m\right)𝒵`$ be a solution of eq. (6.13b), that gives the map $`n:q_m\stackrel{ˇ}{q}_m`$. Apply to the result another transformation with $`\stackrel{ˇ}{n}𝒵`$. We get another map $`\stackrel{ˇ}{n}:\stackrel{ˇ}{q}_m\stackrel{ˇ}{\stackrel{ˇ}{q}}_m`$. Notice that their composition would be another map given by $`n^{}\left(q_m\right)=n\left(q_m\right)+\stackrel{ˇ}{n}\left(\stackrel{ˇ}{q}_m\left(q_m\right)\right)`$, $`n^{}:q_m\stackrel{ˇ}{\stackrel{ˇ}{q}}_m`$. Both $`n`$ and $`n^{}`$ will verify eqs. (6.13) and obviously $`nn^{}`$. It is clear, that further iterations will again bring another solution of (6.13b). We therefore arrive to
###### Proposition 6.4.
For each $`𝒲\left(𝔤_0\right)`$ orbit $`𝒵𝔤_1`$ there is subgroup $`R_𝒵`$ of $`R`$.
Notice, that if $`dim𝒵=1`$, $`n^{}`$ must vanish and our simplest transformation must be reflections, i.e. $`R_𝒵=_2`$.
Any transformation $`nP`$ from $`R`$ is uniquely determined by its $`𝔤_1`$ projection $`n_1`$. That $`n_1`$ may be uniquely split $`n_1=_kn_1^{(k)}`$ with $`n_1^{(k)}𝒵_k`$ in distinct $`𝒲\left(𝔤_0\right)`$ orbits. It suggests that
$$\mathrm{exp}\left[\mathrm{ad}n\left(q_m\right)\right]\left(I_++q_m\right)=\underset{i}{}\mathrm{exp}\left[\mathrm{ad}n_1^{k_i}\left(q_m\right)\right]\left(I_++q_m\right).$$
To ensure that $`R`$ is generated by the simplest transformations one must prove that each factor preserves $`𝔤_0`$. We, however, do not have the proof.
###### Remark 6.4.
In the case of standard Drinfeld-Sokolov hierarchy we have $`𝔤_0=𝔥`$. Then eq. (6.13a) implies that $`n=x_kE_{\alpha _k}𝔤_1`$ and eq. (6.13b) requires
$$x_k=\frac{2}{\left(\alpha _k|\alpha _k\right)}\underset{i=1}{\overset{r}{}}\left(\alpha _k|\alpha _i\right)\nu ^ir_k\left(\nu ^j\right)=\nu ^j\frac{2\delta _{jk}}{\left(\alpha _k|\alpha _k\right)}\underset{i=1}{\overset{r}{}}\left(\alpha _k|\alpha _i\right)\nu ^i.$$
(6.16)
The same transformation of $`\nu `$ results from the Weyl reflection, acting on $`𝔤_0=𝔥`$, corresponding to simple root $`\alpha _k`$. Hence, polynomials $`w^i\left(\nu \right)`$ result to be Coxeter polynomials \[BFRFW\].
###### Corollary 6.2.
Polynomials $`w^i\left(\nu \right)`$ of degrees $`k_i+1`$ with $`k_i𝖯𝗋_w`$ are invariant with respect to discrete group $`R`$ and generate the ring of $`R`$ invariant polynomials on $`Q_m`$.
###### Proof.
By Lemma 6.2 $`w^k(q)`$ generate the ring of gauge invariant polynomials in $`q`$. Restricting them on $`Q_m`$ and restricting adjoint group to $`R`$ the proof follows. ∎
###### Remark 6.5.
Polynomials $`w^i\left(\nu \right)`$ are not however $`𝒲\left(𝔤_0\right)`$ invariant. This means that they are not, for non trivial $`𝒲\left(𝔤_0\right)`$, a restriction to $`𝔤_0`$ of $`Ad`$ invariant polynomials on $`𝔤`$.
Notice, $`R`$ invariance of $`w^i\left(\nu \right)`$ implies that Kirillov-Kostant brackets is equivariant with respect to action of $`R`$, i.e.
$$\overline{\{\stackrel{ˇ}{\nu }^i,\stackrel{ˇ}{\nu }^j\}}_m=f_{}^{ij}{}_{k}{}^{}\stackrel{ˇ}{\nu }^k=\frac{\stackrel{ˇ}{\nu }^i}{\nu ^m}f_{}^{mn}{}_{l}{}^{}\nu ^l\frac{\stackrel{ˇ}{\nu }^j}{\nu ^n}.$$
The $`R`$ invariant variety $`_r`$ is defined in $`Q_m`$, much the same, as fixed point of all Hamiltonians in involution:
$$F^k\left(\nu \right)=\overline{\{\nu ^k,𝒩^2\}}_m=f_{}^{kl}{}_{n}{}^{}\nu ^n\frac{𝒩^2}{\nu ^l}.$$
(6.17)
Indeed, (6.17) implies $`𝒢^{A,2}=0`$ and so defines $`_r`$. It means, that one can restrict action of $`R`$ group on $`_r`$.
To construct “good” coordinates on $`_r`$ we consider Casimirs of Kirillov-Kostant bracket (6.8). Recall, that $`𝔤_0`$ is semisimple Lie algebra, i.e $`𝔤_0=_i𝔤_{}^{}{}_{}{}^{(i)}`$, where $`𝔤_{}^{}{}_{}{}^{(i)}`$ is either simple or abelian. Let $`J_{i,k}(\nu )`$ be algebraically independent $`\mathrm{ad}𝔤_0`$ invariant polynomials of degree $`k`$ corresponding to $`𝔤_{}^{}{}_{}{}^{(i)}`$, i.e. invariant also under $`\mathrm{ad}𝔤_{}^{}{}_{}{}^{(i)}`$. They obviously annihilate (6.8).
Introduce $`r`$ “abelian” coordinates $`\mu `$ on $`𝔥^{}=_i\left(𝔥^{(i)}\right)^{}`$ defined by equations
$$J_{i,k}\left(\mu \right)_{|𝔥}=J_{i,k}\left(\nu \right).$$
(6.18)
Coordinates $`\mu `$ have scaling weight one, as well as Miura coordinates.
###### Proposition 6.5.
Coordinates $`\mu `$ do not depend on the choice of $`\mathrm{ad}𝔤_0`$ invariant polynomials $`J_{i,k}`$.
###### Proof.
Indeed, let us choose another set of algebraically independent polynomials
$$f_j\left(\nu \right)=f_j\left(J\left(\nu \right)\right).$$
They are known to be of the same degrees. Using them we define new coordinates $`\mu ^{}`$:
$$f_j\left(\mu ^{}\right)|_𝔥^{}=f_j\left(\nu \right)f_j\left(J\left(\mu ^{}\right)\right)|_𝔥^{}=f_j\left(J\left(\nu \right)\right)=f_j\left(J\left(\mu \right)\right)|_𝔥^{}.$$
Starting from functions $`f`$ of the lowest degree and proceeding up we conclude that $`\mu ^{}`$ may be chosen to coincide with $`\mu `$. ∎
We may complete $`\mu `$ by those Miura coordinates $`\nu `$ that correspond to $`X_i𝔥`$ in $`q_m=\nu ^iX_i`$, to have coordinates on the whole phase space $``$. Let us denote them $`\eta ^i`$ so that $`\eta ^i=\mu ^i`$ for $`1ir`$.
###### Lemma 6.5.
Group $`R`$ admits restriction on the “abelian” coordinates $`\mu `$.
###### Proof.
Indeed, Casimirs $`J_{i,k}`$, being invariant with respect to adjoint action of $`𝔤_0`$, are not left invariant by action of group $`R`$. Let $`\stackrel{ˇ}{J}_{i,k}\left(\nu \right)=J_{i,k}\left(\stackrel{ˇ}{\nu }\left(\nu \right)\right)`$. Then
$$\begin{array}{cc}\hfill 0& =f_{}^{lm}{}_{n}{}^{}\stackrel{ˇ}{\nu }^n\frac{J_{i,k}\left(\stackrel{ˇ}{\nu }\right)}{\stackrel{ˇ}{\nu }^m}=\overline{\{\stackrel{ˇ}{\nu }^l,J_{i,k}\left(\stackrel{ˇ}{\nu }\right)\}}_m\hfill \\ \hfill 0& =f_{}^{lp}{}_{n}{}^{}\nu ^n\frac{\stackrel{ˇ}{\nu }^m}{\nu ^p}\frac{J_{i,k}\left(\stackrel{ˇ}{\nu }\right)}{\stackrel{ˇ}{\nu }^m}=f_{}^{lp}{}_{n}{}^{}\nu ^n\frac{\stackrel{ˇ}{J}_{i,k}\left(\nu \right)}{\nu ^m}.\hfill \end{array}$$
This implies that $`\stackrel{ˇ}{J}_{i,k}\left(\nu \right)`$ results $`\mathrm{ad}\widehat{𝔤}_0(𝐬_w)`$ invariant and hence depends on $`\mu `$ only.
In the next section we shall make use of the following
###### Lemma 6.6.
Define the following matrix
$$\left(𝕂^1\right)^{mn}=\frac{\eta ^m}{\nu ^i}\left(𝒦^1\right)^{ij}\frac{\eta ^n}{\nu ^j}.$$
(6.19)
Then its submatrix $`\left(𝕂^1\right)^{kl}`$ with $`1k,lr`$ is constant non degenerate matrix.
###### Proof.
Since $`J_{i,k}`$ are $`\mathrm{ad}𝔤_0`$ invariant polynomials and so is Killing metrics, we may make use of Chevalley theorem
$$\frac{J_1\left(\nu \right)}{\nu ^i}\left(𝒦^1\right)^{ij}\frac{J_2\left(\nu \right)}{\nu ^j}=F_{1,2}\left(J\left(\nu \right)\right).$$
(6.20)
stating that $`F`$ is a polynomial in $`\mathrm{ad}𝔤_0`$ invariant polynomials $`J_{i,k}\left(\nu \right)`$. Now using definition of coordinates $`\mu `$ (6.18) and changing variables to $`\eta `$ we obtain
$$F_{1,2}\left(J\left(\mu \right)\right)=\frac{J_1\left(\eta \right)}{\eta ^i}\left(𝕂^1\right)^{ij}\frac{J_2\left(\eta \right)}{\eta ^j}=\underset{m,n=1}{\overset{r}{}}\frac{J_1\left(\mu \right)}{\mu ^m}\left(𝕂^1\right)^{mn}\frac{J_2\left(\mu \right)}{\mu ^n}.$$
(6.21)
By another Chevalley theorem $`J\left(\mu \right)`$ are invariant with respect to Weyl group of $`𝔤_0`$. This means that $`r\times r`$ submatrix $`𝕂^1`$ is inverse of Killing metrics of rank $`r`$ semisimple Lie algebra $`𝔤_0`$ restricted to Cartan subalgebra. The latter coincides with Killing form of algebra $`𝔤`$. Thus it is non degenerate due to celebrated Cartan’s criterion of semisimplicity. ∎
###### Corollary 6.3.
$`\left(𝕂^1\right)^{ij}=\left(𝒦^1\right)^{ij}`$ for $`1i,jr`$.
###### Lemma 6.7.
Group $`R`$ acts linearly on coordinates $`\mu `$.
###### Proof.
Consider $`\stackrel{ˇ}{J}\left(\mu \right)=J\left(\stackrel{ˇ}{\mu }\left(\mu \right)\right)`$. Rewrite (6.20) as follows
$$\begin{array}{cc}\hfill F_{1,2}\left(\stackrel{ˇ}{J}\left(\mu \right)\right)& =F_{1,2}\left(J\left(\stackrel{ˇ}{\mu }\right)\right)=\underset{m,n=1}{\overset{r}{}}\frac{J_1\left(\stackrel{ˇ}{\mu }\right)}{\stackrel{ˇ}{\mu }^m}\left(𝕂^1\right)^{mn}\frac{J_2\left(\stackrel{ˇ}{\mu }\right)}{\stackrel{ˇ}{\mu }^n}\hfill \\ & =\underset{m,n,k,l=1}{\overset{r}{}}\frac{\stackrel{ˇ}{J}_1\left(\mu \right)}{\mu ^m}\frac{\mu ^n}{\stackrel{ˇ}{\mu }^l}\left(𝕂^1\right)^{kl}\frac{\mu ^n}{\stackrel{ˇ}{\mu }^l}\frac{\stackrel{ˇ}{J}_2\left(\mu \right)}{\mu ^n}.\hfill \end{array}$$
Since functions $`F_{1,2}\left(J\right)`$ are not altered we have the same pairing and due to Lemma 6.6 $`R`$ action preserves constant non degenerate submatrix of the matrix $`𝕂^1`$. It is only possible if $`R`$ acts linearly. ∎
###### Remark 6.6.
Since the action of $`R`$ on $`\mu `$ is defined via action on polynomials $`J\left(\mu \right)`$, it can be determined only up to $`\mathrm{ad}𝔤_0`$, i.e. only as
$$\stackrel{ˇ}{J}\left(\mu \right)=J\left(A\mu \right)$$
for some constant matrix $`A`$, determined up to transformations
$$Aw_1Aw_2,\text{where}w_1,w_2𝒲\left(𝔤_0\right).$$
###### Lemma 6.8.
Non trivial transformations from $`R_𝒵`$ correspond to the same matrix $`A_𝒵`$ modulo this equivalence.
###### Proof.
Fix an $`𝒲\left(𝔤_0\right)`$ orbit $`𝒵`$ in $`𝔤_1`$. It suffices to show that $`𝒲\left(𝔤_0\right)`$ permutes solutions of eq. (6.13b).
Notice, that $`𝒲\left(𝔤_0\right)`$ invariant polynomials $`J`$, corresponding to abelian constituents of $`𝔤_0`$ are linear and may be chosen to be corresponding Miura variables as follows from (6.16). Then, due to Lemma 6.7 $`R`$ groups acts on them linearly in terms of $`\mu `$. Solving eqs. (6.18) for $`\mu \left(\nu \right)`$ we get just $`\mathrm{ord}𝒲\left(𝔤_0\right)`$ solutions as follows from Bezout theorem. They correspond to different solutions of (6.13b). ∎
###### Corollary 6.4.
Eqs. (6.13b) have non more than $`\mathrm{ord}𝒲\left(𝔤_0\right)`$ non trivial solutions.
###### Corollary 6.5.
Matrix $`A_𝒵`$ may be chosen to be a reflection, i.e. $`A_𝒵^2=1`$.
###### Proof.
Since $`R_𝒵`$ is a subgroup, there exist two transformations product of which is identity. Since they correspond to the same matrix $`A_𝒵`$, the proof follows. ∎
###### Theorem 6.1.
Let $`𝒲\left(𝔤_0\right)`$ be Weyl group of $`𝔤_0`$. Assume $`𝔤`$ simply laced. Then $`R|__r𝒲\left(𝔤_0\right)𝒲\left(𝔤\right)`$.
###### Proof.
Weyl group $`𝒲\left(𝔤_0\right)`$ preserves metrics $`𝕂^1`$, and is generated by reflections (6.16) corresponding to simple roots of zero $`𝐬_w`$ grade. Group $`R`$, restricted to $`_r`$, also preserves the metrics and is generated by reflections associated with simple roots of $`𝐬_w`$ grade 1, due to Lemmas 6.3 and 6.4. We conclude that $`R|__r𝒲\left(𝔤_0\right)`$ is a finite group generated by $`r`$ transformations associated to simple roots, that preserve $`𝕂^1`$. Since it is an inverse of the Killing form of algebra $`𝔤`$ and $`𝒲\left(𝔤_0\right)`$ acts canonically (6.16), it follows that $`R|__r𝒲\left(𝔤_0\right)`$ is isomorphic to Weyl group $`𝒲\left(𝔤\right)`$. ∎
## 7 Dispersionless limit
As was proved by I. Krichever \[Kr\], the dispersionless limit, or zero phase Whitham averaging, of the standard Drinfeld-Sokolov hierarchies provides solutions to WDVV equations.
Introducing dispersion parameter $`ϵ`$ the Poisson structure of hierarchies reads
$$\begin{array}{cc}\hfill \{w^i(x),w^j(y)\}& =\underset{k0}{}ϵ^{k1}\{w^i(x),w^j(y)\}^{(k)}\hfill \\ & =\frac{1}{ϵ}A^{ij}(w)\delta (xy)+g^{ij}(w)\delta ^{}(xy)\hfill \\ & +\mathrm{\Gamma }_k^{ij}(w)\left(w^k(x)\right)^{}\delta (xy)+𝒪(ϵ).\hfill \end{array}$$
(7.1)
###### Remark 7.1.
For the standard Drinfeld-Sokolov hierarchies $`A0`$ because $`dim=r`$ and hence annihilators of the first Poisson structure may be chosen as coordinates on the whole $``$. Note that due to Corollary 3.1 we only have $`rdim`$ independent annihilators, thus appearance of $`A`$ term should be, generally, expected.
Let us pick up $`w^i=(𝒩^a,u^A)`$ as coordinates on $`\left(Q^{can}\right)^{}`$ as in section 6. We assume that they are obtained by ultralocal change of variables from Drinfeld-Sokolov variables, i.e. contain no derivative terms.
Due to eq. (6.4) hierarchy time flows of $`𝒩^a`$ have polynomial $`ϵ`$ expansion, but dynamics of $`u^A`$ coordinates does not enjoy this property
$`{\displaystyle \frac{𝒩^a}{t_b}}`$ $`=`$ $`\{𝒩^a,H_b\}_2=_xA^a(b)+𝒪\left(ϵ\right),`$
$`{\displaystyle \frac{u^A}{t_b}}`$ $`=`$ $`\{u^A,H_b\}_2={\displaystyle \frac{1}{ϵ}}𝒢_\alpha ^{A,b}(u,𝒩)+_xA^A\left(b\right)+𝒪\left(ϵ\right).`$ (7.2)
Hence brackets $`\{𝒩^a,𝒩^b\}`$ admit dispersionless limit, while others do not. $`G^{A,b}`$’s come from $`A`$ term in (7.1) and thus are responsible for fast dynamics of $`u`$ coordinates. If $`u`$ coordinates evolved so as to vanish $`𝒢^{A,a}`$ identically, we would obtain a well defined dispersionless limit of the hierarchy. As we have seen in section 6 it happens on the algebraic subvariety $`_r`$.
We thus supplement Dubrovin-Novikov prescriptions \[DN\] for the restriction of the Poisson structure (7.1) on the slow - modulated zero phase solutions by requirement of additional restriction on $`_r`$. Dirac bracket provides restriction of the Hamiltonian structure on $``$ to $`_r`$. We review briefly, for reader’s convenience, the construction of Dirac bracket, referring to \[MR\] for details.
Given constraint equation $`𝒢^{A,2}=0`$ defining $`_r`$ and local coordinates $`𝒩^a`$ there, we introduce new ones
$$\stackrel{˘}{𝒩}^a\left(x\right)=𝒩^a\left(x\right)+\underset{A}{}𝑑y\tau _A^a(x,y)𝒢^{A,2}\left(w(y)\right).$$
such that $`\stackrel{˘}{𝒩}^a|__r=𝒩^a|__r`$ and with $`\tau `$ subject to condition
$$\{𝒩^a\left(x\right),𝒢^{A,2}\left(w(y)\right)\}|__r=0.$$
Looking for solution of $`\tau `$ as formal $`ϵ`$ series $`\tau =_{m0}ϵ^m\tau ^{(m)}`$, the equation above for $`\tau `$ amounts to the following
$$\{𝒩^a(x),𝒢^{A,2}(y)\}_{|__r}^{(k)}+\underset{m=0}{\overset{k}{}}𝑑z\tau _{}^{(m)}{}_{B}{}^{a}(x,z)\{𝒢^{B,2}(z),𝒢^{A,2}(y)\}_{|__r}^{(km)}=0.$$
(7.3)
Due to Proposition 6.1 we obtain that $`\tau ^{(0)}=0`$ provided that the matrix
$$\{𝒢^{B,2}(z),𝒢^{A,2}(y)\}^{(0)}|__r$$
is nondegenerate.
###### Definition 7.1.
Dirac bracket on $`_r`$ is defined as
$$\begin{array}{c}\{𝒩^a(x),𝒩^b(y)\}_D=\{\stackrel{˘}{𝒩}^a(x),\stackrel{˘}{𝒩}^b(y)\}|__r=\hfill \\ \hfill \{𝒩^a(x),𝒩^b(y)\}|__r𝑑z_1𝑑z_2\tau _A^a(x,z_1)\{𝒢^{A,2}\left(x\right),𝒢^{B,2}\left(z_2\right)\}|__r\tau _B^b(y,z_2)\end{array}$$
(7.4)
Dirac bracket verifies \[MR\] Jacobi identity. As an immediate consequence of (7.4) we have the following
###### Lemma 7.1.
If $`\tau ^{(0)}=0`$ then
$$\{𝒩^a(x),𝒩^b(y)\}_D^{(k)}=\{𝒩^a(x),𝒩^b(y)\}^{(k)}|__rk=0,1.$$
(7.5)
Because of (6.4) the $`ϵ0`$ expansion of Dirac bracket (7.4) starts with $`k=1`$ term, and thus the bracket admits dispersionless limit.
###### Corollary 7.1.
Dispersionless limit of bi Hamiltonian structure is bi Hamiltonian.
We thus arrive to the following theorem
###### Theorem 7.1.
Consider Hamiltonian dynamical system admitting constant solutions. Let some $`𝒩^a`$ be the densities of the local commuting integrals of the system, considered as the parameters of the full family of the constant solutions. Let $`u`$ denote the rest of dynamical variables. Assume that the matrix $`\overline{\{𝒢^{A,2},𝒢^{B,2}\}}_2|__r`$ does not degenerate identically. Then the dispersionless limit of the Hamiltonian structure restricted to $`_r`$, given by Dubrovin-Novikov formula
$$\{𝒩^a(x),𝒩^b(y)\}^{}=g^{ab}\left(𝒩(x)\right)\delta ^{}(xy)+\mathrm{\Gamma }_c^{ab}\left(𝒩(x)\right)\left(𝒩^c(x)\right)^{}\delta (xy),$$
(7.6)
satisfies the Jacobi identity and does not depend on the choice of $`𝒢`$.
This result is a particular case of theorem due to A. Maltsev \[M\] who proved, using Dirac reduction procedure, that Dubrovin-Novikov averaging procedure yields, under certain assumptions, a Poisson structure on the space of $`m`$-phased solutions of dynamical equations of original Hamiltonian system.
###### Lemma 7.2.
Assumptions of Theorem 7.1 verify for the hierarchies $`(𝔤,[w],\mathrm{\Lambda })`$ with regular primitive conjugacy class $`[w]`$ and grade one regular element $`\mathrm{\Lambda }`$.
###### Proof.
Due to (5.6), $`H_\mathrm{\Lambda }=𝑑xw_2\left(x\right)`$ is the momentum for the hierarchies in question, so their dynamical equations admit constant solutions, because of theorem 3.1.
Now we address the non degeneracy statement.
$$𝒜^{AB}=\overline{\{𝒢^{A,2},𝒢^{B,2}\}}_2=\frac{𝒢^{A,2}}{w^i}\overline{\{w^i,w^j\}}_2\frac{𝒢^{B,2}}{w^j}.$$
Restricting on $`_r`$ and using (6.4) and $`\overline{\{u^A,𝒩^a\}}_2|__r=0`$ we conclude
$$𝒜^{AB}|__r=\frac{𝒢^{A,2}}{u^C}(u\left(𝒩\right),𝒩)\overline{\{u^C,u^D\}}_2|__r\frac{𝒢^{B,2}}{u^D}(u\left(𝒩\right),𝒩).$$
Thus, due to assumption of Proposition 6.1, it suffices to prove the non degeneracy of $`𝔸^{CD}=\overline{\{u^C,u^D\}}_2|__r`$. The latter is obvious because
$$𝔸^{CD}=\stackrel{~}{𝔸}^{CD}+𝒩^rB^{CD}|__r,$$
where $`\stackrel{~}{𝔸}`$ does not depend on $`𝒩^r`$. Restriction on $`_r`$ does not alter the linear dependence of $`𝔸`$ on $`𝒩^r`$ because of Proposition 6.2. Due to Lemma 6.1 $`detB0`$ and is constant. It, thus, remains constant after restriction to $`_r`$. ∎
###### Proposition 7.1 (\[DN, Du4\]).
Given Poisson structure of hydrodynamic type (7.6), $`g^{ab}`$ is the flat covariant metrics as long as $`detg0`$ and $`\mathrm{\Gamma }`$ is its connection, related to Levi-Civita connection by the following relation
$$\mathrm{\Gamma }_c^{ab}=g^{ad}\mathrm{\Gamma }_{dc}^b.$$
Non degeneracy of matrix $`g^{ab}\left(𝒩\right)`$ obtained by Dirac reduction on $`_r`$ of Poisson structure of generalized integrable hierarchies in consideration is not an obvious fact and needs to be proved.
Miura map (3.3) provides us with polynomial expressions of Drinfeld-Sokolov coordinates in terms of Miura ones and its derivatives:
$$w_ϵ^i\left(\nu \right)=w^i\left(\nu \right)+ϵ\underset{j=1}{\overset{dim}{}}\left(\nu ^j\right)^{}w_j^{(1)}\left(\nu \right)+𝒪\left(ϵ^2\right).$$
This change of coordinates provides a map of Miura bracket (5.12) to the second Poisson structure (5.1). Under this map we obtain for the first two matrix in (7.1) the following expressions
$`A^{ij}\left(w\right)`$ $`=`$ $`{\displaystyle \frac{w^i}{\nu ^k}}f_{}^{kl}{}_{m}{}^{}\nu ^m{\displaystyle \frac{w^j}{\nu ^l}},`$
$`g^{ij}\left(w\right)`$ $`=`$ $`{\displaystyle \frac{w^i}{\nu ^k}}\left(𝒦^1\right)^{kl}{\displaystyle \frac{w^j}{\nu ^l}}+\left(w_j^{(1)}{\displaystyle \frac{w^i}{\nu ^k}}w_i^{(1)}{\displaystyle \frac{w^j}{\nu ^l}}\right)f_{}^{kl}{}_{m}{}^{}\nu ^m.`$ (7.7)
As an immediate consequence of definition of $`_r`$ we have
###### Proposition 7.2.
On $`_r`$ (7.7) simplifies to the following expression:
$$g^{ab}\left(𝒩\right)=\underset{k,l=1}{\overset{dim}{}}\frac{𝒩^a}{\nu ^k}\left(𝒦^1\right)^{kl}\frac{𝒩^b}{\nu ^l}.$$
(7.8)
###### Proposition 7.3.
$`\{,\}_1^{}`$ can be read off $`\{,\}_2^{}`$ by appropriate shift of $`𝒩_r`$.
###### Proof.
Due to Proposition 6.2 and Lemma 7.1 dispersionless limit of the second bracket remains linear in $`𝒩^r`$. The result follows. ∎
Note, that $`𝒩_1=w_2`$ satisfies Fuchs algebra
$$\{𝒩^1,𝒩^1\}_2^{}=2𝒩^1\delta ^{}+\left(𝒩^1\right)^{}\delta .$$
Thus we have arrived to
###### Proposition 7.4.
Zero phase Whitham averaging maps graded bi Hamiltonian structure (5.1) into graded bi Hamiltonian structure of hydrodynamic type.
Due to Lemma 7.1 $`𝒩^a`$ will remain annihilators of the first bracket:
$$\{𝒩^a(x),𝒩^b(y)\}_1^{}=\eta ^{ab}\delta ^{}\left(xy\right).$$
(7.9)
###### Remark 7.2.
Due to scaling weight grading, and chosen field ordering $`\eta ^{ab}`$ is anti-diagonal matrix.
Indeed, $`\eta ^{ab}`$ vanishes if $`\mathrm{deg}_{sc}\left(𝒩^a𝒩^b\right)N_w+2`$ and $`\mathrm{deg}_{sc}\left(𝒩^a𝒩^{r+1a}\right)=N_w+2`$ due to $`\mathrm{deg}_{sc}𝒩^a=k_a+1`$ and ascending ordering of $`k_aI(w)`$. Notice, that this is consistent with previously assigned $`𝒩^1`$ being $`w_2`$ and $`𝒩^r`$ linear in $`w_{N_w,1}`$, because of
$$\{𝒩^k,𝒩^1\}_2^{}=\mathrm{deg}_{sc}\left(𝒩^k\right)𝒩^k\delta ^{}+\left(𝒩^k\right)^{}\delta .$$
(7.10)
###### Proposition 7.5.
Dispersionless limit of Dirac restriction of Poisson structure of modified hierarchy reads
$$\{\mu ^i,\mu ^j\}^{}=\left(𝕂^1\right)^{ij}\delta ^{}\left(xy\right).$$
(7.11)
###### Proof.
$`\mu `$ are Casimirs of finite dimensional Kirillov-Kostant bracket (6.8) and thus due to Lemma 7.1 and Lemma 6.6 the result follows. ∎
###### Theorem 7.2.
Metrics $`g^{ab}\left(𝒩\right)`$ is not identically degenerate. Coordinates $`\mu `$ are flat coordinates for this metrics.
###### Proof.
Let us choose coordinates $`\eta `$ as follows. Take the first $`r`$ coordinates to be $`\mu `$ and choose the rest to be canonical for finite dimensional bracket. The advantage of this choice is the simplicity of constraint equations (6.17)
$$\frac{𝒩^2}{\eta ^k}|__r=0\stackrel{\text{ by Prop. }\text{6.1}}{}\frac{𝒩^a}{\eta ^k}|__r=0r+1kdim.$$
(7.12)
Due to Lemma 7.1 and Proposition 7.2 we have
$$\begin{array}{cc}\hfill g^{ab}& =\underset{k,l=1}{\overset{dim}{}}\frac{𝒩^a}{\nu ^k}\left(𝒦^1\right)^{kl}\frac{𝒩^b}{\nu ^l}|__r=\underset{k,l=1}{\overset{dim}{}}\frac{𝒩^a}{\eta ^k}\left(𝕂^1\right)^{kl}\frac{𝒩^b}{\eta ^l}|__r\hfill \\ & =\underset{k,l=1}{\overset{r}{}}\frac{𝒩^a}{\mu ^k}\left(𝕂^1\right)^{kl}\frac{𝒩^b}{\mu ^l}|__r.\hfill \end{array}$$
Both set of coordinates $`𝒩`$ and $`\mu `$ are local coordinates on $`_r`$, and so $`J(\mu )=det\left|\begin{array}{c}\frac{𝒩^a}{\mu ^i}\end{array}\right|`$ does not degenerate at generic point of $`_r`$. This proves non degeneracy.
The following identity along with Lemma 6.6 prove that $`\mu `$ are indeed flat coordinates of metrics $`g`$:
$$\frac{𝒩^a|__r}{\mu ^k}=\frac{𝒩^a(\mu ,\eta \left(\mu \right))}{\mu ^k}=\frac{𝒩^a}{\mu ^k}|__r+\underset{i=r+1}{\overset{dim}{}}\frac{𝒩^a}{\eta ^i}|__r\frac{\eta ^i}{\mu ^k}|__r=\frac{𝒩^a}{\mu ^k}|__r.$$
As a byproduct we conclude that brackets (7.6) and (7.11) define the same geometry on $`_r`$:
$$\{𝒩^a\left(x\right),𝒩^b\left(y\right)\}_2^{}=\frac{𝒩^a}{\mu ^m}\left(x\right)\{\mu ^m\left(x\right),\mu ^n\left(y\right)\}_m^{}\frac{𝒩^b}{\mu ^n}\left(y\right).\text{on }_r$$
(7.13)
This conclusion may be drawn also noting that both coordinates set are densities of mutually commuting integrals of corresponding exact brackets, the property that survives averaging. ∎
###### Corollary 7.2.
$`𝒩^a|__r\left(\mu \right)`$ are invariant with respect to linear action of $`R`$ group on $`_r`$.
###### Proof.
Indeed, due to Lemma 6.5 $`R`$ admits a restriction of $`_r`$ and preserves the latter. Hence, $`𝒩^a`$ were $`R`$ invariant in $``$ and so they remain restricted on $`_r`$. ∎
Following K. Saito \[Sa\] and using Theorem 7.2 we obtain the following
###### Proposition 7.6.
Metrics $`\eta ^{ab}`$ is non degenerate.
###### Proof.
Consider the following polynomial in $`\lambda `$
$$\begin{array}{cc}\hfill P\left(\lambda \right)& =det\left|g^{ab}\left(𝒩\right)\lambda \eta ^{ab}\right|=det\left|g^{ab}(𝒩^1,\mathrm{},𝒩^{r1},N^r\lambda )\right|\hfill \\ & =det\left|\begin{array}{c}\eta ^{ab}\end{array}\right|\left(𝒩^r\lambda \right)^n+\underset{n=0}{\overset{r1}{}}c_n(𝒩^1,\mathrm{},𝒩^{r1})\left(𝒩^r\lambda \right)^n.\hfill \end{array}$$
(7.14)
When all $`𝒩`$ but $`𝒩^r`$ vanish it simplifies to $`P(\lambda )=det\left|\begin{array}{c}\eta ^{ab}\end{array}\right|\left(𝒩^r\lambda \right)^n`$. At this point $`J\left(\mu \right)0`$ and thus $`g^{ab}`$ is non degenerate. Indeed, let $`\mu `$ be eigenvector of some representative, which always exists, of $`[w]`$ in $`R`$ with eigenvalue $`\xi =\mathrm{exp}\left[2i\pi /N\right]`$. Then, due to homogeneity of $`𝒩^a(\mu )`$ and from their $`R`$ invariance, we conclude that only variables of degree $`N`$ may differ from zero at this point of $`_r`$. If there are more than one such variable, then eigenvalue $`\xi `$ is degenerate and we can always choose $`\mu `$ so that only $`𝒩^r`$ does not vanish. Due to regularity of $`[w]`$, vector $`\mu `$ is not left fixed by any transformation from $`R`$, and thus the results follows. ∎
###### Proposition 7.7.
So obtained pencil of Hamiltonian structures provides us with quasihomogeneous \[Du4\] flat pencil of metrics.
###### Proof.
Let $`g`$ be the metrics of the second Poisson structure and let $`\eta `$ – of the first. Introduce function $`\tau =𝒩^1/N_w`$, and introduce the following vector fields
$$E^a=g^{ab}_b\tau e^a=\eta ^{ab}_b\tau .$$
(7.15)
Notice, that with this choice $`e^a=\delta _{a,r}`$ and $`𝔏_eg=\eta `$ and $`𝔏_e\eta =0`$ as follows from considerations above, and where we have assumed $`\eta `$ to be chosen anti diagonal with all nonzero entries being $`N_w`$. Then $`E^a=\mathrm{deg}_{sc}\left(𝒩^a\right)/N_w𝒩^a`$, as follows from (7.10). One sees immediately that $`[e,E]=e`$. The second Poisson structure is scaling weight graded and $`N_wE`$ is scaling weight Euler vector field. Thus $`g`$ must be an eigenvector of $`𝔏_E`$ : $`𝔏_Eg=(d1)g`$. In ref. \[Du4\] such flat pencils were called quasihomogeneous of degree $`d`$. ∎
###### Theorem 7.3.
Given $`\{,\}_\lambda ^{}`$ one may associate to it a solution to WDVV.
###### Proof.
Since obtained pencil of Hamiltonian structures satisfies Jacobi identity we have a flat pencil of metrics. It is quasihomogeneous as was shown in proposition 7.7. Thus, following ref. \[Du4\], it is enough to show that the degree of quasihomogeneity $`d1`$.
As was said in the proof of proposition 7.7 $`E^a=\frac{\mathrm{deg}_{sc}𝒩^a}{N_w}N^a`$. Due to quasihomogeneity of flat pencil we have $`𝔏_E\eta =(d2)\eta `$. But
$$\begin{array}{cc}\hfill 𝔏_E\eta ^{ab}& =E^c_c\eta ^{ab}\eta ^{ac}_cE^b\eta ^{cb}_cE^a\hfill \\ & =\left(\frac{\mathrm{deg}_{sc}𝒩^b}{N_w}+\frac{\mathrm{deg}_{sc}𝒩^a}{N_w}\right)\eta ^{ab}\hfill \\ & =\frac{N_w+2}{N_w}\eta ^{ab}.\hfill \end{array}$$
In the last line we have used the fact that $`\eta ^{ab}`$ vanishes unless $`\mathrm{deg}_{sc}\left(𝒩^a𝒩^b\right)=N_w+2`$, as follows from scaling weight grading of Poisson structures. From this we obtain
$$d=1\frac{2}{N_w}.$$
Since $`N_w`$ is finite we obtain that $`d<1`$. Note that for the Coxeter conjugacy class $`N_w=h`$ \- Coxeter number, and we recover the formula of B. Dubrovin, obtained while constructing polynomial solutions to WDVV equations on the orbits of Coxeter groups \[Du2\].
Practically, we can find Frobenius potential $`F(𝒩)`$ from the following relations
$$\begin{array}{cc}\hfill g^{ab}\left(𝒩\right)& =\left(d1d_ad_b\right)\eta ^{ac}\eta ^{bd}_c_dF\left(𝒩\right),\hfill \\ \hfill \mathrm{\Gamma }_c^{ab}& =\left(\frac{3d}{2}d_a\right)\eta ^{ad}\eta ^{bf}_d_f_cF,\hfill \end{array}$$
(7.16)
where $`d_a\delta _a^b=_aE^b`$. ∎
Thus, $`𝒩^a`$ are Saito coordinates \[Sa\] on Frobenius manifold \[Du3\] being flat coordinates for the metric $`\eta `$. However, flat coordinates $`\mu `$ of the intersection metrics $`g`$ are also very important. They clarify the geometric origin of the Frobenius structure. In general it is a challenging task, given a flat metric, to find its flat coordinates. But in the case in question we were lucky to use the theory of integrable systems.
## 8 Example: $`[w]=D_4(a_1)`$
To illustrate the developed technique we consider the example served as the motivation of the present work. Let $`𝔤=D_4`$ – the simplest classical Lie algebra where non Coxeter primitive conjugacy class occurs (see appendix A). Luckily it enjoys regularity property.
Take $`[w]=D_4(a_1)`$. We have readily that $`I(w)=(1,1,3,3)`$ and the set of conformal weights $`𝖯𝗋_w=(1,1,1,2,3,3)`$. Miura coordinates form a Kac Moody algebra of $`\widehat{𝔤}_0(𝐬_w)=u(1)^3su(2)`$, thus we shall have three exact Casimirs and one will be computed in dispersion parameter expansion.
Positive roots of $`D_4`$ read
$$\begin{array}{cc}\hfill _+& =\{\alpha _1,\alpha _2,\alpha _3,\alpha _4,\alpha _1+\alpha _2,\alpha _2+\alpha _3,\alpha _2+\alpha _4,\alpha _1+\alpha _2+\alpha _3,\hfill \\ & \alpha _1+\alpha _2+\alpha _4,\alpha _2+\alpha _3+\alpha _3,\alpha _1+\alpha _2+\alpha _3+\alpha _4,\alpha _1+2\alpha _2+\alpha _3+\alpha _4\}.\hfill \end{array}$$
Let us denote Lie algebra elements $`E_\alpha `$, $`\alpha _+`$ by its decomposition on simple roots. So if $`\alpha =\alpha _{max}`$ we write $`X_{12234}`$, and for $`E_{\alpha _1+\alpha _2}`$ write $`X_{12}`$. Similarly for negative roots, substituting $`X`$ with $`Y`$. We fix Cartan-Weyl basis (2.1). So
$$\rho =\underset{i,j=1}{\overset{r}{}}K_{ij}^1s_jH_{\alpha _i}=2H_{\alpha _1}+3H_{\alpha _2}+2H_{\alpha _3}+2H_{\alpha _4}.$$
The Heisenberg subalgebra $`^{[w]}`$ is spanned by $`z^k\mathrm{\Lambda }_{i,1}`$, $`z^k\mathrm{\Lambda }_{i,2}`$ for $`i\{1,3\}`$ and $`k`$.
$$\begin{array}{cc}\hfill \mathrm{\Lambda }_{1,1}& =X_1+X_3+zY_{12234}+X_{12}+X_{23}+X_{24},\hfill \\ \hfill \mathrm{\Lambda }_{2,1}& =X_1X_3+X_4zY_{1234}X_{12}+X_{23},\hfill \\ \hfill \mathrm{\Lambda }_{3,1}& =X_{1234}\frac{z}{2}\left(Y_1Y_3+2Y_4Y_{12}+Y_{23}\right),\hfill \\ \hfill \mathrm{\Lambda }_{3,2}& =X_{12234}\frac{z}{2}\left(Y_1+Y_3+Y_{12}+Y_{23}+2Y_{24}\right).\hfill \end{array}$$
(8.1)
These basis was just guessed, verifying linear independence and regularity. It turned out easier than proceed as in (2.4).
Let us choose $`\mathrm{\Lambda }=\mathrm{\Lambda }_{1,1}`$ for our integrable hierarchy. The $`sl_2`$ subalgebra constituents read $`\rho ,I_+,I_{}`$, where $`I_{}=3Y_1+3Y_3+Y_{12}+Y_{23}+4Y_{24}`$ and $`I_+=P_0^{𝐬_h}\mathrm{\Lambda }`$.
### The hierarchy
Let us fix the minimal weight gauge (5.7) as follows
$$\begin{array}{cc}\hfill q^{can}& =\left(w_2\frac{1}{10}u_2\frac{1}{2}v_2\right)\frac{1}{12}I_{}+u_2\frac{1}{20}\left(3Y_13Y_3+12Y_4\right)+\hfill \\ & +v_2\frac{1}{20}\left(3Y_1+3Y_3+4Y_{24}\right)+w_3\frac{1}{6}\left(Y_{123}Y_{234}2Y_{124}\right)+\hfill \\ & +Y_{1234}u_4+Y_{12234}w_4.\hfill \end{array}$$
Hamiltonians, annihilators of the first Poisson structure, read
$$\begin{array}{cc}\hfill H_{\mathrm{\Lambda }_{1,1}}& =𝑑xw_2,H_{\mathrm{\Lambda }_{1,2}}=𝑑xu_2,\hfill \\ \hfill H_{\mathrm{\Lambda }_{3,1}}& =dx(u_4+\frac{37}{3}u_2^2\frac{1}{6}w_2u_2+\frac{1}{64}w_2^2+\frac{1}{6}w_2v_2\hfill \\ & \frac{7}{6}u_2v_2+\frac{7}{12}v_2^2),\hfill \\ \hfill H_{\mathrm{\Lambda }_{3,2}}& =𝑑x\left(w_4\frac{7}{3}u_2^2+\frac{7}{24}w_2u_2+\frac{7}{3}u_2v_2\right).\hfill \end{array}$$
(8.2)
We shall denote the densities of this annihilators as $`𝒩_1`$, $`𝒩_2`$, $`𝒩_3`$ and $`𝒩_4`$ respectively.
As was explained in section 5 fixing the gauge somehow, we are able to compute both Poisson structures exactly, but the output is enormous to be presented here. On the other hand we need then to pass to $`𝒩,w_3,v_2`$ coordinates and eliminate auxiliary coordinates $`w_3`$ and $`v_2`$. To do so we need to know $`𝒲`$ algebra at least up to $`𝒪\left(ϵ\right)`$ after rescaling. But we choose to we omit these intermediate steps due to space restrictions and present the answer.
The simplest Hamiltonian generating $`G`$ terms is $`H_{\mathrm{\Lambda }_{1,2}}`$. Its flows with respect to the second Hamiltonian structure read
$$\begin{array}{cc}\hfill \frac{𝒩_1}{t_{\mathrm{\Lambda }_{1,2}}}& =𝒩_2^{},\frac{𝒩_2}{t_{\mathrm{\Lambda }_{1,2}}}=\frac{1}{3}\left(v_2𝒩_2+5𝒩_1\right)^{},\hfill \\ \hfill \frac{𝒩_3}{t_{\mathrm{\Lambda }_{1,2}}}& =\frac{1}{3}\left(𝒩_1𝒩_2+3𝒩_4\frac{1}{4}𝒩_2^2+\frac{1}{4}v_2𝒩_2+ϵ\frac{2}{5}w_3^{}\right)^{},\hfill \\ \hfill \frac{𝒩_4}{t_{\mathrm{\Lambda }_{1,2}}}& =\frac{1}{3}\left(𝒩_3+\frac{1}{4}𝒩_2^2+\frac{1}{12}𝒩_1\left(𝒩_1𝒩_2+v_2\right)\right)^{},\hfill \\ \hfill \frac{v_2}{t_{\mathrm{\Lambda }_{1,2}}}& =\frac{12}{ϵ}w_3+\frac{1}{3}\left(5𝒩_1+22𝒩_2+2v_2\right)^{},\hfill \\ \hfill \frac{w_3}{t_{\mathrm{\Lambda }_{1,2}}}& =\frac{1}{ϵ}\left(4𝒩_3+\frac{1}{12}v_2^2+\frac{1}{3}𝒩_1v_2\frac{1}{3}𝒩_1𝒩_2\frac{1}{6}𝒩_2v_2\frac{1}{6}𝒩_2^2+\frac{1}{4}𝒩_1^2\right)+\hfill \\ & +\frac{3}{20}ϵ\left(𝒩_2v_2\right)^{\prime \prime }.\hfill \end{array}$$
(8.3)
Following the recipe, we take $`1/ϵ`$ terms as constraints. Using them we obtain equation for the phase space subvariety $`_r`$ of slow motion
$$w_3=0,v_2=𝒩_22𝒩_1\pm \mathrm{\Delta },\mathrm{\Delta }=\sqrt{𝒩_1^2+3𝒩_2^2+48𝒩_3}.$$
(8.4)
This finally leads to the following restricted bi Hamiltonian structure
$$\begin{array}{cc}\hfill \{𝒩_i(x),𝒩_j(y)\}_1^{}& =4\delta _{i+j,5}\delta ^{}\left(xy\right).\hfill \\ \hfill \{𝒩_i(x),𝒩_j(y)\}_2^{}& =\left[\gamma ^{ij}\left(𝒩(x)\right)+\gamma ^{ji}\left(𝒩(y)\right)\right]\delta ^{}\left(xy\right),\hfill \end{array}$$
(8.5)
where matrix $`\gamma (𝒩)`$ reads
$$\gamma ^{i,1}=\{𝒩_1,𝒩_2,3𝒩_3,3𝒩_4\},\gamma ^{1,i}=\{𝒩_1,𝒩_2,𝒩_3,𝒩_4\},$$
$$\gamma ^{2,2}=\frac{1}{3}\left(𝒩_1+2\mathrm{\Delta }\right),\gamma ^{2,3}=𝒩_4+\frac{1}{6}𝒩_1𝒩_2+\frac{1}{12}𝒩_2\mathrm{\Delta },\gamma ^{3,2}=3\gamma ^{2,3},$$
$$\gamma ^{2,4}=\frac{1}{3}\left(𝒩_3\frac{1}{12}𝒩_1^2+\frac{1}{4}𝒩_2^2+\frac{1}{12}𝒩_1\mathrm{\Delta }\right),\gamma ^{4,2}=3\gamma ^{2,4},$$
$$\gamma ^{3,3}=\frac{1}{2}\left(𝒩_1𝒩_3+\frac{3}{32}𝒩_1𝒩_2^2+\frac{7}{288}𝒩_1^3\right)+\frac{1}{288}\left(𝒩_1^2+12𝒩_2^2+48𝒩_3\right)\mathrm{\Delta },$$
$$\gamma ^{4,3}=\frac{1}{2}\left(𝒩_2𝒩_3+\frac{7}{96}𝒩_2𝒩_1^2+\frac{1}{32}𝒩_2^3+\frac{1}{48}𝒩_1𝒩_2\mathrm{\Delta }\right),\gamma ^{3,4}=\gamma ^{4,3},$$
$$\gamma ^{4,4}=\frac{1}{6}\left(𝒩_1𝒩_3+\frac{19}{288}𝒩_1^3+\frac{7}{32}𝒩_1𝒩_2^2+\frac{1}{144}\left(4𝒩_1^2+3𝒩_2^2+48𝒩_3\right)\mathrm{\Delta }\right).$$
The metric $`g=\gamma +\gamma ^{tr}`$ is invertible,can be checked to be flat and forms, obviously, together with $`\eta `$ a flat pencil.
### Modified hierarchy
We consider the modified hierarchy to exemplify discrete group $`R`$. Choosing coordinates as in (5.12), but indexing them with subscript to facilitate reading of following formulae, we have
$$𝒦^1=\left(\begin{array}{ccc}K_{4\times 4}^1& & \\ \\ & 0& 1\\ & 1& 0\end{array}\right),f\nu =\left(\begin{array}{cccccc}0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& \nu _5& \nu _6\\ 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0\\ 0& \nu _5& 0& 0& 0& \omega \\ 0& \nu _6& 0& 0& \omega & 0\end{array}\right),$$
where $`\omega =\left(\nu _12\nu _2+\nu _3+\nu _4\right)`$. Notice, that quadratic Casimir of $`su(2)`$ constituent of $`\widehat{𝔤}_0(𝐬_w)`$ read
$$J_2=\left(\nu _12\nu _2+\nu _3+\nu _4\right)^2+4\nu _5\nu _6.$$
(8.6)
According to (6.18) we introduce “abelian” coordinates $`\mu `$:
$$\begin{array}{cc}\hfill \mu _i=\nu _i,& i=1,3,4,\hfill \\ \hfill \left(\mu _12\mu _2+\mu _3+\mu _4\right)^2& =\left(\nu _12\nu _2+\nu _3+\nu _4\right)^2+4\nu _5\nu _6.\hfill \end{array}$$
(8.7)
It is invariant with respect to $`su(2)`$ Weyl group, acting on $`\mu `$ according to (6.16):
$$\mu _i\mu _i,i=1,3,4,\mu _2\mu _1+\mu _3+\mu _4\mu _2.$$
(8.8)
Then group $`R`$ are generated by the following elementary reflections corresponding to three simple roots of $`𝐬_w`$ degree one: $`\alpha _1`$, $`\alpha _3`$ and $`\alpha _4`$. Each corresponds gauge transformation (6.6) with $`n=x_1E_{\alpha _k}+x_2E_{\alpha _k+\alpha _2}`$, with $`k=1,3,4`$. For each $`x`$ we find quadratic equation (6.13b) and thus we obtain, for instance, the following transformation for $`\alpha _1`$:
$$\begin{array}{cc}\hfill R_1^\pm :\nu _1& \frac{\nu _1+\nu _3+\nu _4\pm \sqrt{J_2}}{2},\nu _3\nu _3,\nu _4\nu _4,\hfill \\ \hfill R_1^\pm :\nu _2& \frac{1}{2\left(\omega +\nu _5\nu _6\right)}(2\nu _2\nu _52(2\nu _22(\nu _3+\nu _4)+\nu _5)\nu _6+\hfill \\ & +(\nu _22(\nu _3+\nu _4)3\nu _6)\omega +\omega ^2\pm \sqrt{J_2}(\nu _1+\nu _2\nu _3\nu _4+\nu _6)),\hfill \\ \hfill R_1^\pm :\nu _5& \frac{\left(2\nu _1\nu _2+\nu _5\right)\left(\omega +2\nu _5\sqrt{J_2}\right)}{2\left(\omega +\nu _5\nu _6\right)},\hfill \\ \hfill R_1^\pm :\nu _6& \frac{\left(\nu _1+\nu _2\nu _3\nu _4+\nu _6\right)\left(\omega 2\nu _6\pm \sqrt{J_2}\right)}{2\left(\omega +\nu _5\nu _6\right)}.\hfill \end{array}$$
As illustration to Lemma 6.7 note that
$$J_2\left(R_1^\pm \left(\nu \right)\right)=\left(\frac{3\mu _1\mu _3\mu _4\pm \sqrt{J_2}}{2}\right)^2$$
One can check that $`R^\pm `$ are not reflections, because, for example,
$$\left(R_1^+\right)^2=\{\begin{array}{cc}\mathrm{Id}_{\mathrm{Q}_\mathrm{m}^{}}\hfill & \text{ if }\mathrm{\hspace{0.17em}3}\nu _1\nu _3\nu _4\pm \sqrt{J_2}0,\hfill \\ R_1^{}\hfill & \text{ otherwise.}\hfill \end{array}$$
We can get rid of sign $`\pm `$, by making use of Weyl group $`𝒲\left(𝔤_0\right)`$ of $`sl(2)`$ action (8.8). We define action of $`R_1`$ as $`R_1^+`$ for $`\nu _2,\nu _5,\nu _6`$ variables, choosing the following solution of eq. (8.7)
$$\mu _2=\frac{1}{2}\left\{\sqrt{J_2}+\mu _1+\mu _3+\mu _4\right\}.$$
This choice yields for $`\mu `$ coordinates
$$R_1:\mu _2\mu _2,\mu _3\mu _3,\mu _4\mu _4,\mu _1\mu _1+\mu _2.$$
Then $`R_1^{}=R_2R_1R_2`$, where $`R_2`$ acts on $`\mu `$ variables by (8.8) and acts trivially on $`\nu _{2,5,6}`$. It can be explicitly checked that so defined operation $`R_1`$ is a reflection $`R_1^2=Id`$. It means that we have well defined reflections on Riemann surface (8.7) over $`Q_m^{}`$.
The same way we define reflections $`R_3`$ and $`R_4`$. Recall that shortcut $`\omega `$ stands for $`\omega =\mu _1+\mu _3+\mu _42\nu _2`$
$$\begin{array}{cc}\hfill R_3:\mu _1& \mu _1,\mu _2\mu _2\mu _4\mu _4\mu _3\mu _3+\mu _2,\hfill \\ \hfill R_3:\nu _2& \frac{1}{2\left(\omega \nu _5+\nu _6\right)}[2\nu _5\nu _6+\nu _6\omega +\omega ^22\mu _3(2\nu _6+\omega )+\hfill \\ & +\nu _2(2\nu _5+4\nu _6+3\omega )+\sqrt{J_2}(\mu _1\nu _2\mu _3+\mu _4+\nu _6)],\hfill \\ \hfill R_3:\nu _5& \frac{\left(2\mu _3\nu _2\nu _5\right)\left(\omega +2\nu _5+\sqrt{J_2}\right)}{2\left(\omega \nu _5+\nu _6\right)},\hfill \\ \hfill R_3:\nu _6& \frac{\left(\mu _1\nu _2\mu _3+\mu _4+\nu _6\right)\left(\omega +2\nu _6+\sqrt{J_2}\right)}{2\left(\omega \nu _5+\nu _6\right)}.\hfill \end{array}$$
$$\begin{array}{cc}\hfill R_4:\mu _1& \mu _1,\mu _2\mu _2,\mu _4\mu _4,\mu _3\mu _3+\mu _2,\hfill \\ \hfill R_4:\nu _2& \frac{\mu _1+2\nu _2+\mu _33\mu _4+\sqrt{J_2}}{2},\nu _6\nu _6,\hfill \\ \hfill R_4:\nu _5& \frac{\left(\nu _2+2\mu _4\right)\left(\sqrt{J_2}\omega \right)}{2\nu _6}.\hfill \end{array}$$
Notice, that $`\mu _k=\nu _k`$ for $`k=1,3,4`$ and we obtain that $`R`$ acts linearly on these $`\mu `$. It is easy to check for these linear transformations, and it certainly needs symbolic computation program to verify that
$$R_k^2=1,\left(R_1R_3\right)^2=\left(R_1R_4\right)^2=\left(R_3R_4\right)^2=1,$$
$$\left(R_2R_1\right)^3=\left(R_2R_3\right)^3=\left(R_2R_4\right)^3=1.$$
That is they generate Weyl group of $`D_4`$ Lie algebra.
### Solution to WDVV
By theorem 7.3 we can extract a solution to WDVV from bi Hamiltonian structure (8.5). Since scaling degree of fields $`𝒩_1,\mathrm{},𝒩_4`$ are $`2,2,4,4`$ respectively, we can find Euler vector field
$$E=(\frac{1}{2}𝒩_1,\frac{1}{2}𝒩_2,𝒩_3,𝒩_3),$$
and the grade $`d=1/2`$. Following \[Du4\] Frobenius potential $`F(𝒩)`$ can be extracted from (7.16). It should be noted that $`\mathrm{\Gamma }_k^{ij}=_k\gamma ^{j,i}`$, making two relations equivalent. We thus find the claimed free energy
$$\begin{array}{cc}\hfill \frac{1}{4}F\left(𝒩\right)& =𝒩_2𝒩_3𝒩_4+\frac{1}{2}𝒩_1𝒩_4^2+\frac{\mathrm{\Delta }^5}{2^53^45}+\frac{1}{6}𝒩_1𝒩_3^2\frac{1}{108}𝒩_3𝒩_1^3+\hfill \\ & +\frac{1}{12}𝒩_1𝒩_2^2𝒩_3+\frac{19}{2^83^45}𝒩_1^5+\frac{7}{2^73^3}𝒩_1^3𝒩_2^2+\frac{1}{32^8}𝒩_1𝒩_2^4.\hfill \end{array}$$
(8.9)
It may be explicitly checked to verify WDVV equations (1.1).
## Acknowledgement
I would like to thank B. Dubrovin for posing me this problem and for a lot of valuable advice. I am indebted to A. Maltsev for explaining me his article and for enlightening discussions. I would like also to thank G. Falqui for stimulating talks.
I would like to thank International School for Advanced Studies, Italy, where this work has been done, for hospitality and creative atmosphere.
## Appendix A Regular primitive conjugacy classes and their properties
Here we collect information about regular primitive conjugacy classes of Weyl group $`𝒲\left(𝔤\right)`$ for simple Lie algebras $`𝔤`$. Classes are labeled by the type of Coxeter diagram. Recall, that it coincides with the Dynkin diagram for the Coxeter conjugacy class.
### $`𝐀_𝐧`$
$`[w]`$ $`=`$ $`A_nN_w=n+1𝐬_w=𝐬_p`$
$`I(w)`$ $`=`$ $`(1,2,3,\mathrm{},n1,n)=\underset{w}{𝖯𝗋}\widehat{𝔤}_0(𝐬_w)=u(1)^n`$
### $`𝐁_𝐧,𝐂_𝐧`$
$`[w]`$ $`=`$ $`B_n,C_nN_w=2n𝐬_w=𝐬_p`$
$`I(w)`$ $`=`$ $`(1,3,5\mathrm{},2n1)=\underset{w}{𝖯𝗋}\widehat{𝔤}_0(𝐬_w)=u(1)^n`$
### $`𝐃_𝐧`$
We pick four last roots to form $`D_4`$ subalgebra.
$`[w]`$ $`=`$ $`D_nN_w=2n2𝐬_w=𝐬_p`$
$`I(w)`$ $`=`$ $`(1,3,\mathrm{},2n1;n1)=\underset{w}{𝖯𝗋}\widehat{𝔤}_0(𝐬_w)=u(1)^n`$
$`[w]`$ $`=`$ $`D_{2n}(a_{n1}),D_{2n}(b_{n1})N_w=2n`$
$`𝐬_w`$ $`=`$ $`(1,\underset{2n2\text{ times}}{\underset{}{1,0,1,0,\mathrm{},1,0}},1,1)`$
$`I(w)`$ $`=`$ $`(\underset{2n\text{ numbers}}{\underset{}{1,1,3,3,\mathrm{},2n3,2n3,2n1,2n1}})`$
$`\underset{w}{𝖯𝗋}`$ $`=`$ $`(\underset{n1\text{ groups in 4 elements}}{\underset{}{\stackrel{}{1,1,1,2},\mathrm{},\stackrel{}{2n3,2n3,2n3,2n2}}},2n1,2n1)`$
$`\widehat{𝔤}_0(𝐬_w)`$ $`=`$ $`u(1)^{n+1}su(2)^{n1}`$
### $`𝐆_\mathrm{𝟐}`$
$`[w]`$ $`=`$ $`G_2N_w=6𝐬_w=𝐬_p`$
$`I(w)`$ $`=`$ $`(1,5)=\underset{w}{𝖯𝗋}\widehat{𝔤}_0(𝐬_w)=u(1)^2`$
### $`𝐅_\mathrm{𝟒}`$
$`\alpha _1^2=\alpha _2^2=2`$,$`\alpha _3^2=\alpha _4^2=4`$ with double bond between the second and the third roots.
$`[w]`$ $`=`$ $`F_4N_w=12𝐬_w=𝐬_p`$
$`I(w)`$ $`=`$ $`(1,5,7,11)=\underset{w}{𝖯𝗋}\widehat{𝔤}_0(𝐬_w)=u(1)^4`$
$`[w]`$ $`=`$ $`F_4(a_1)N_w=6`$
$`𝐬_w`$ $`=`$ $`(1,1,0,1,0)`$
$`I(w)`$ $`=`$ $`(1,1,5,5)`$
$`\underset{w}{𝖯𝗋}`$ $`=`$ $`(1,1,1,2,3,4,5,5)`$
$`\widehat{𝔤}_0(𝐬_w)`$ $`=`$ $`u(1)^2su(2)^2`$
### $`𝐄_\mathrm{𝟔}`$
Roots $`\alpha _6,\alpha _2,\alpha _3,\alpha _4`$ form $`D_4`$ subalgebra.
$`[w]`$ $`=`$ $`E_6N_w=12𝐬_w=𝐬_p`$
$`I(w)`$ $`=`$ $`(1,4,5,7,8,11)=\underset{w}{𝖯𝗋}\widehat{𝔤}_0(𝐬_w)=u(1)^6`$
$`[w]`$ $`=`$ $`E_6(a_1)N_w=9`$
$`𝐬_w`$ $`=`$ $`(1,1,1,0,1,1,1)`$
$`I(w)`$ $`=`$ $`(1,2,4,5,7,8)`$
$`\underset{w}{𝖯𝗋}`$ $`=`$ $`(1,2,3,4,5,5,7,8)`$
$`\widehat{𝔤}_0(𝐬_w)`$ $`=`$ $`u(1)^5su(2)`$
$`[w]`$ $`=`$ $`E_6(a_2)N_w=6`$
$`𝐬_w`$ $`=`$ $`(1,1,0,1,0,1,0)`$
$`I(w)`$ $`=`$ $`(1,1,2,4,5,5)`$
$`\underset{w}{𝖯𝗋}`$ $`=`$ $`(1,1,1,2,2,2,3,3,4,4,5,5)`$
$`\widehat{𝔤}_0(𝐬_w)`$ $`=`$ $`u(1)^3su(2)^3`$
### $`𝐄_\mathrm{𝟕}`$
$`[w]`$ $`=`$ $`E_7N_w=18𝐬_w=𝐬_p`$
$`I(w)`$ $`=`$ $`(1,5,7,9,11,13,17)=\underset{w}{𝖯𝗋}\widehat{𝔤}_0(𝐬_w)=u(1)^7`$
$`[w]`$ $`=`$ $`E_7(a_1)N_w=14`$
$`𝐬_w`$ $`=`$ $`(1,1,1,0,1,1,1,1)`$
$`I(w)`$ $`=`$ $`(1,3,5,7,9,11,13)`$
$`\underset{w}{𝖯𝗋}`$ $`=`$ $`(1,3,5,5,7,8,9,11,13)`$
$`\widehat{𝔤}_0(𝐬_w)`$ $`=`$ $`u(1)^6su(2)`$
$`[w]`$ $`=`$ $`E_7(a_4)N_w=6`$
$`𝐬_w`$ $`=`$ $`(1,0,0,1,0,0,1,0)`$
$`I(w)`$ $`=`$ $`(1,1,1,3,5,5,5)`$
$`\underset{w}{𝖯𝗋}`$ $`=`$ $`(1,1,1,1,1,1,2,2,2,2,3,3,3,3,3,4,4,4,5,5,5)`$
$`\widehat{𝔤}_0(𝐬_w)`$ $`=`$ $`u(1)^2su(2)su(3)^2`$
### $`𝐄_\mathrm{𝟖}`$
$`[w]`$ $`=`$ $`E_8N_w=30𝐬_w=𝐬_p`$
$`I(w)`$ $`=`$ $`(1,7,11,13,17,19,23,29)=\underset{w}{𝖯𝗋}\widehat{𝔤}_0(𝐬_w)=u(1)^8`$
$`[w]`$ $`=`$ $`E_8(a_1)N_w=24`$
$`𝐬_w`$ $`=`$ $`(1,1,1,0,1,1,1,1,1)`$
$`I(w)`$ $`=`$ $`(1,5,7,11,13,17,19,23)`$
$`\underset{w}{𝖯𝗋}`$ $`=`$ $`(1,5,7,9,11,13,14,17,19,23)`$
$`\widehat{𝔤}_0(𝐬_w)`$ $`=`$ $`u(1)^7su(2)`$
$`[w]`$ $`=`$ $`E_8(a_2)N_w=20`$
$`𝐬_w`$ $`=`$ $`(1,1,1,0,1,0,1,1,1)`$
$`I(w)`$ $`=`$ $`(1,3,7,9,11,13,17,19)`$
$`\underset{w}{𝖯𝗋}`$ $`=`$ $`(1,3,5,7,8,9,11,11,13,14,17,19)`$
$`\widehat{𝔤}_0(𝐬_w)`$ $`=`$ $`u(1)^6su(2)^2`$
$`[w]`$ $`=`$ $`E_8(a_3),E_8(b_3)N_w=12`$
$`𝐬_w`$ $`=`$ $`(1,1,0,1,0,0,1,0,0)`$
$`I(w)`$ $`=`$ $`(1,1,5,5,7,7,11,11)`$
$`\underset{w}{𝖯𝗋}`$ $`=`$ $`(1,1,1,2,3,4,5,5,5,5,6,6,7,7,7,8,9,10,11,11)`$
$`\widehat{𝔤}_0(𝐬_w)`$ $`=`$ $`u(1)^3su(2)^3su(3)`$
$`[w]`$ $`=`$ $`E_8(a_5),E_8(b_5)N_w=15`$
$`𝐬_w`$ $`=`$ $`(1,1,0,1,0,1,0,1,0)`$
$`I(w)`$ $`=`$ $`(1,2,4,7,8,11,13,14)`$
$`\underset{w}{𝖯𝗋}`$ $`=`$ $`(1,2,3,4,5,5,7,7,7,8,9,9,11,11,13,14)`$
$`\widehat{𝔤}_0(𝐬_w)`$ $`=`$ $`u(1)^4su(2)^4`$
$`[w]`$ $`=`$ $`E_8(a_6)N_w=10`$
$`𝐬_w`$ $`=`$ $`(1,0,0,1,0,0,1,0,0)`$
$`I(w)`$ $`=`$ $`(1,1,3,3,7,7,9,9)`$
$`\underset{w}{𝖯𝗋}`$ $`=`$ $`(1,1,1,2,3,3,3,3,3,4,4,4,5,5,5,6,6,6,7,7,7,8,9,9)`$
$`\widehat{𝔤}_0(𝐬_w)`$ $`=`$ $`u(1)^2su(2)^2su(3)^2`$
$`[w]`$ $`=`$ $`E_8(a_8)N_w=6`$
$`𝐬_w`$ $`=`$ $`(1,0,0,0,1,0,0,0,0)`$
$`I(w)`$ $`=`$ $`(1,1,1,1,5,5,5,5)`$
$`\underset{w}{𝖯𝗋}`$ $`=`$ $`(\underset{\text{10 times}}{\underset{}{1,\mathrm{},1}},\underset{\text{10 times}}{\underset{}{2,\mathrm{},2}},\underset{\text{10 times}}{\underset{}{3,\mathrm{},3}},\underset{\text{6 times}}{\underset{}{4,\mathrm{},4}},\underset{\text{4 times}}{\underset{}{5,\mathrm{},5}})`$
$`\widehat{𝔤}_0(𝐬_w)`$ $`=`$ $`u(1)su(4)su(5)`$
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# Quantum Computation and the localization of Modular Functors 11footnote 1Based on lectures prepared for the joint Microsoft/University of Washington celebration of mathematics April 2000 and the AMS meeting on mathematics in the new millennium UCLA, August 2000.
## 1 The Picture Principle
Reality has the habit of intruding on the prodigies of purest thought and encumbering them with unpleasant embellishments. So it is astonishing when the chthonian hammer of the engineer resonates precisely to the gossamer fluttering of theory. Such a moment may soon be at hand in the practice and theory of quantum computation. The most compelling theoretical question, $`\mathrm{`}\mathrm{`}`$localization,” is yielding an answer which points the way to a solution of Quantum Computing’s (QC) most daunting engineering problem: reaching the accuracy threshold for fault tolerant computation.
After Shor’s discovery \[S1\] of a polynomial time factoring algorithm in the quantum model QC, skeptics properly questioned whether a unitary evolution could ever be induced to process information fault tolerantly. The most obvious tricks, such as making a backup copy, useful in a dissipative system (e.g. pencil and paper) are unavailable in quantum mechanics. To overcome these difficulties, a remarkable theoretical framework based on $`\mathrm{`}\mathrm{`}`$stabilizer codes,” $`\mathrm{`}\mathrm{`}`$transversal gates,” $`\mathrm{`}\mathrm{`}`$cat-state-ancilli, ” and nested concatenations of these was erected \[S2\], \[S3\], \[A,B-O\], \[K1\], and \[KLZ\]. While the result is a consistent recipe for fault-tolerant quantum computation, the accuracy threshold which would allow this combinatorial behemoth to overcome its own overhead has been estimated as about $`10^6`$, one i.i.d. error per one million physical gate operations and requiring gates accurate also to one part in a million. This places a formidable task before the engineer and physicist. But within the year the beginnings of a new idea on fault tolerance had been generated by Kitaev \[K2\].
While the term is not yet present in that paper the idea is to construct (first mathematically) a $`\mathrm{`}\mathrm{`}`$quantum medium” and to store quantum states as topological structures within the medium and (eventually) manipulate these states, that is, apply gates to them, by topological transformations of the medium. For our purposes, we define a quantum medium as a collection of many finite level systems coupled together by a Hamiltonian $`H`$ obeying a strong locality condition: The individual systems are located in a $`2`$dimensional lattice or a more irregular cellulation of a surface $`\mathrm{\Sigma }`$. We postulate a constant $`d>0`$ so that $`H=\mathrm{\Sigma }\overline{H}_k`$ and each $`\overline{H}_k=H_k`$id, where the identity is on all tensor factors(= subsystem) not located within some ball $`B_{\mathrm{}}`$ of diameter $`d`$ in the lattice. For example, the Heisenberg magnet with $`H=J\underset{a,b=\text{edge}}{\mathrm{\Sigma }}\stackrel{}{\sigma }_a\stackrel{}{\sigma }_b`$ is a quantum medium of diameter $`=1`$. (But engineer be warned; localizing $`H^{\mathrm{}}`$ within balls of diameter $`=d`$ implies $`n`$ary interaction for $`nd^2`$. Controlling effective $`n`$ary terms for $`n2`$ will be tricky in the extreme and probably will require enforcing symmetries to cancel lower order terms.) Kitaev’s $`\mathrm{`}\mathrm{`}`$toric code” \[K2\] in which quantum states are stored as first homology of a torus, can be counted as having $`d=2`$; they require $`4`$ary interactions.
We study here a partial generalization of the toric code which also stores quantum information in a degenerate ground state $`V(\mathrm{\Sigma })`$ of a quantum medium. The medium is on a disk with point-like defects which we treat as punctures. The dimension of $`V(\mathrm{\Sigma })`$, $`\mathrm{\Sigma }`$ the punctured disk, grows exponentially with the number of punctures. Transformations of $`\mathrm{\Sigma }`$, that is braidings (up to isotopy) of the punctures in space-time, $`\mathrm{\Sigma }\times R`$, operate unitarily on $`V(\mathrm{\Sigma })`$. Other work (\[K2\], \[P\], and \[K,B\]) also explores the realization of elements of computation by braiding anyonic $`\mathrm{`}\mathrm{`}`$quasi-particles” or $`\mathrm{`}\mathrm{`}`$defects” of a quantum medium.
The vision is that stability of computation, at least sufficient to reach the $`10^6`$ threshold for $`\mathrm{`}\mathrm{`}`$software” error correction, is to be realized by the discreteness of algebraic topology: two $`Z_2`$homology cycles are never $`\mathrm{`}\mathrm{`}`$close,” two words in the braid group are equal or distinct. More exactly, it is geometry not topology which will confer stability. Working in a lattice model one may calculate \[K2\] that the perturbation Hamiltonian $`P`$ must be raised to the length scale $`L`$ before nonzero terms, $`<\zeta |P^L|\eta >,\zeta ,\eta `$ ground state $`(H)`$, are encountered and so the splitting of the ground state is estimated to be proportional to $`e^{\mathrm{\Omega }(L)}`$. The length scale in the previous two examples are: $`L=`$ (length of shortest essential cycle); and in the anyonic context, the closest that two defects are allowed to come to each other during braiding. The $`\mathrm{`}\mathrm{`}`$engineering goal” is to construct a physical quantum medium on a material disk whose ground state admits many localized excitations ($`\mathrm{`}\mathrm{`}`$anyons”) whose braidings effect computationally universal unitary transformations of the ground state. It is further hoped that actual $`\mathrm{`}\mathrm{`}`$errors,” the result of unwanted noisy excitations, are to be removed automatically by some relaxation process in which the system is coupled to a cold bath by another much weaker Hamiltonian $`H^{}`$. The mathematicians first cut at the engineering goal is to produce a mathematical quantum medium with these properties and this is accomplished by the theorem below. This $`\mathrm{`}\mathrm{`}`$first cut” is not yet interesting to experimentalists since the Hamiltonian contains summands which have as many as $`30`$ nontrivial indices, but it represents an exact existence theorem. The question for physicist is whether this phase can also be represented perturbatively with a simple Hamiltonian, perhaps a RVB model \[A\], \[N,S\]. This would be a major step toward physical realization.
###### Theorem 1.1.
Consider a rectangle $`R`$ of Euclidian square lattice consisting of $`15`$ boxes by $`30n`$ boxes. Associate a $`2`$level spin system $`^2`$ with each of the $`e:=960n+36`$ box edges in $`R`$. The disjoint union of these spin systems has Hilbert space $`(^2)^e=:X`$. There is a time dependent local Hamiltonian $`H_t=\left(\underset{𝑘}{\mathrm{\Sigma }}\overline{H}_{k,t}\right)`$ with fewer than $`2000n`$ terms and each $`H_k`$ having $`30`$ or fewer indices, supported in at most a $`5\times 3`$ rectangle of boxes - $`\mathrm{`}\mathrm{`}`$diameter $`=5`$.” For $`t=0`$, the ground states of $`H_0`$ form a sub-Hilbert space $`WX`$, and geometrically determines $`3n`$ exceptional points or $`\mathrm{`}\mathrm{`}`$defects” spaced out along the midline of $`R`$. Within $`W`$ there is a $`\mathrm{`}\mathrm{`}`$computational” sub-Hilbert space $`V(^2)^n`$, $`VW.`$ $`W`$ may be identified with the $`SU(2)`$Witten-Chern-Simons modular functor at level $`ł=r2=3`$ of the $`3n`$punctured disk with the fundamental representation of $`SU(2)`$ labeling each of the $`3n+1`$ boundary components. The Braid group $`B(3n)`$ of the defects acts unitarily on $`W`$ according to the Jones’ representation at level $`=5`$. Any quantum algorithm can be efficiently simulated on $`V`$ by restricting the action of $`B(3n)`$ to a $`\mathrm{`}\mathrm{`}`$computational subspace.”
The representation is implemented adiabatically by gradually deforming $`H_t`$ to $`H_{t+1}`$ and then to $`H_{t+2}`$ and so on. The passage from $`H_t`$ to $`H_{t+1}`$ involves turning off an exceptional term $`\overline{H}_{k,t}`$ which defines a defect site and turning on a new term $`\overline{H}_{k,tH}`$, which determines an alternative, adjacent, site for the defect at time $`t+1`$. Each braid generator can be implemented in $`4(r+1)`$ times steps. We believe, based on a conjectural energy gap, that the geometry confers stability to this implementation which increases exponentially, error = $`e^{\mathrm{\Omega }(L)}`$, under refinement of the lattice on $`R`$ by a factor of $`L`$, while the number of time step needed for a computation increases only linearly in $`L`$.
###### Comments 1.2.
* The second paragraph of the theorem should be read as a defensible physical proposition, whereas the first paragraph is mathematics.
* Our Hamiltonian may be too complicated to prove the persistence of an energy gap above the ground state in the thermodynamic limit. But based on an analogy with a simpler system the gap is conjectured and will be discussed at the end of the proof.
* The passage from the Jones’ representation to computation on $`V`$ is the subject of \[FLW1\] and \[FLW2\] where it is proved that universality holds for $`r=5`$ and $`r7`$. Functorially $`V`$ is a tensor summand of a subspace of $`W`$ but by fixing a reference vector in the complementary tensor factor we regard $`V`$ simply as a subspace of $`W`$.
* The idea of anyonic computation is taken from \[K2\] and in a more speculative form \[Fr\]. The new ingredient is the implementation of a computationally complete modular functor by a local Hamiltonian. Witten’s approach \[Wi\] to CSr was Lagrangian and so nonlocal; it yields an identically zero Hamiltonian under Legendre transform, \[FKW\] and \[A\]. This lecture, in contrast, supplies a Hamiltonian interpretation for CS5 (We may replace $`5`$ by any $`r7`$ in the statement at the expense of scaling the constants in the theorem by $`\frac{r}{5}`$ or $`\frac{r^2}{25}`$ according to whether they scale as lengths or areas.).
* We know of two works in progress with a similar objective. Kitaev and Bravyi \[K,B\] study a local model for the weaker functor CS4 on high genus surfaces, and Kitaev and Kupperberg \[K,K\] have an approach to construct local Hamiltonians generally for modular functors on surfaces of any genus which (unlike CS5) are quantum doubles \[D\]. Their approach has the advantage that the local contributions to the Hamiltonian can be arranged to commute so that an energy gap will be rigorously established. In contrast, an interesting feature of the present paper is that topological a combinatorial means yield an exact determination of a ground state defined by non-commuting terms. This is not usually possible. Finally, we will see that our local construction for $`H`$ extends to the higher genus surfaces if CSr is replaced by any modular functors of the form $`VV^{}`$. The simple topological reason for this may illuminate the analysis of \[K,K\].
* Shortly, we will give the reader a completely pictorial understanding of CSr on planar surfaces.
So far, we have only discussed the $`\mathrm{`}\mathrm{`}`$engineering”: the quest to specify $`H`$ (which will be described in the proof). Let us take a brief digression from that sulfurous underworld of grinding gears to the Elysian fields of abstract thought. The Witten-Chern-Simons theory descends from the signature (= Pontryagin form) in dimension 4 and every step of the desent to lower dimension leads to deeper abstraction until mathematical wit is well nigh exhausted as the point (dimension $`=0`$) is reached. To tell this story in its barest outline, we restrict to $`G=SU(2)`$, and borrow from Atiyah \[A\], Freed \[F\], and Walker \[W\]. The signature of a closed $`4`$manifold is an integer as is the Pontryagin class of an $`SU(2)`$ bundle over a closed $`4`$manifold. An $`SU(2)`$bundle over a closed $`3`$manifold is topologically trivial but if endowed with a connection acquires a secondary $`\mathrm{`}\mathrm{`}`$ Chern-Simons” class in the circle $`=R/𝒵`$. Quantizing \[Wi\] at level $`ł`$, leads to the topological Jones-Witten-Chern-Simons invariant $``$ which is morally an average of the classical Chern-Simons invariant over all connections. The invariant for a closed surface $`\mathrm{\Sigma }`$ (with some additional structure) is a finite dimensional vector space $`V`$; and each $`3`$manifold bounding $`\mathrm{\Sigma }`$ determines a vector $`vV`$. Before dividing by gauge symmetry, the vector space $`\overline{V}`$ is the infinite dimensional space of sections of the associated complex line bundle to a natural $`S^1`$bundle over the space of $`SU(2)`$ connections $`A`$ on $`SU(2)`$ bundles over $`\mathrm{\Sigma }`$. A $`3`$manifold $`Y`$ with connection, $`\overline{A}`$, on a bundle extending the bundle over the boundary, $`(Y,\overline{A})=(\mathrm{\Sigma },A)`$, determines a map $`f\{(Y^{},\overline{A^{}})|(Y^{},\overline{A^{}})=(\mathrm{\Sigma },A)\}S^1`$ by integrating the Chern-Simons form over $`YY^{}`$. The consistent choices for such functionals constitute the total space of this $`\mathrm{`}\mathrm{`}`$natural” $`S^1`$bundle. In general, a map $`f`$ is $`\mathrm{`}\mathrm{`}`$consistent” if it obeys the additivity properties of the Chern-Simons integral: $`f(Y^{})f(Y^{\prime \prime })=\text{C.S.}(Y^{}Y^{\prime \prime })`$. Symplectic reduction followed by quantization as explained in \[A\] produces a finite dimensional $`V`$ from $`\overline{V}`$ with $`v(Y)V`$ depending only on the topology of $`Y`$. The definition of the Witten-Chern-Simons invariant for a surface with boundary is a collection of vector spaces indexed by certain labelings. For a $`1`$manifold the invariant seems to be a certain type of $`\mathrm{`}\mathrm{`}2`$category” while the correct definition for a point is but dimly perceived and the object of current research. Several authors assert that it is unnecessary to finish the progression, that we can be content with a theory whose smallest building blocks are $`\mathrm{`}\mathrm{`}`$pairs of pants” (three-punctured- spheres). The invariant for these while technically a vector in a $`2`$vector space is easily understood in terms of sets of vector spaces parameterized by $`\mathrm{`}\mathrm{`}`$labelings” of the boundary circles so no unusual categorical abstractions need be mastered. The reason for this assertion is that using a handle body decomposition all closed $`3`$manifold invariants can be calculated from gluing along surfaces with smooth boundary; gluings along faces with corners on the boundary, which one would encounter computing from a cellulation, can be avoided. But the Freed-Walker program rejects this advice on two grounds. First localizing $`V(\mathrm{\Sigma })`$ not merely to $`\mathrm{`}\mathrm{`}`$pants,” but to cells (i.e. neighborhoods of points) may give more natural consistency conditions, to replace the $`14`$ consistency equations of \[W\]; which in turn could eventually lead to classification of modular functors and a conceptual understanding. Second, to paraphrase Edmund Hillary, we should localize to points $`\mathrm{`}\mathrm{`}`$because they are there.”
The hyperbole of the first paragraph can now be made sound. CS5 is a universal model for quantum computation and for the physicist/engineer to implement it, a local Hamiltonian $`H`$ must be described. For the pure mathematician to be satisfied with his understanding of CS5 it must be localized to points. The two objectives are certainly similar in spirit and possibly identical. To clarify the connection, we introduce an intermediate concept, undoubtedly plebeian, but dear to a topologist. We would like when possible to describe a vector in a modular functor as a linear combinations of $`\mathrm{`}\mathrm{`}`$admissible” pictures up to $`\mathrm{`}\mathrm{`}`$equivalence.” This, after all, is exactly how we understand homology: $`vH_1(\mathrm{\Sigma },Z_2)`$ is an equivalence class of admissible pictures. To be admissible the picture must be a closed $`1`$manifold, the equivalence relation is bordism. Both $`\mathrm{`}\mathrm{`}1`$manifold-ness” and $`\mathrm{`}\mathrm{`}`$bordism” can be defined by local conditions which are the combinatorial analogs of $`\mathrm{`}\mathrm{`}`$closed” and $`\mathrm{`}\mathrm{`}`$co-closed” familiar from de Rham’s theory of differential forms. In Kitaev’s toric code these condition are imposed by vertex and face operators $`A_v`$ and $`B_f`$ respectively. There is a subtle shift here from the usual way of thinking of homology as equivalence classes of cycles to the $`\mathrm{`}\mathrm{`}`$harmonic” representative which is merely the equally weighted average of all cycles in the homology class. In this way quotients and equivalence classes are never encountered and homology is located within cycles, within chains, just as a C.S.S. code space is located within the fixed space of stabilizers built from products of $`\sigma _z`$’s and further within the fixed space of stabilizers, $`\mathrm{\Pi }\sigma _x`$’s.
To generalize from homology, we should think of a picture as (linear combinations of) anything we can draw on a surface $`\mathrm{\Sigma }`$. If helpful, we allow various colors and/or notational labels, framing fields, etc$`\mathrm{}`$, and even additional dimensions bundled over $`\mathrm{\Sigma }`$. But in the present case no such embellishments are required. What is important that if we move the surface by a diffeomorphism, the picture should also move and move canonically. Thus if $`\mathrm{\Sigma }`$ is a torus it would not suite our purposes to draw the picture of $`vV(\mathrm{\Sigma })`$ in a solid torus $`T`$, $`T=\mathrm{\Sigma }`$: a meridial Dehn twist on $`\mathrm{\Sigma }`$ extends over $`T`$, twisting the picture, but a longitudinal Dehn twist does not have any obvious way to act on a picture drawn in $`T`$. (To anticipate, a modular functor will have an $`S`$matrix which can transform a picture in one (call it the $`\mathrm{`}\mathrm{`}`$inside”) solid torus to a picture in the dual ($`\mathrm{`}\mathrm{`}`$outside”) solid torus where longitudinal a Dehn twist does act. But resorting to the $`S`$matrix does not solve our problem since its input and output pictures are on a scale of the injectivity radius of the surfaces and hence nonlocal.) We demand that $`\mathrm{`}\mathrm{`}`$admissibility” and $`\mathrm{`}\mathrm{`}`$equivalence” of pictures be locally determined, i.e. decided on the basis of restriction to small patches on $`\mathrm{\Sigma }`$. To make the connection with lattice models, we consider $`\mathrm{\Sigma }`$ discretized as a cell complex; the conditions must span only clumps of cells of constant combinatorial diameter. As in the example of harmonic $`1`$cycles, $`\mathrm{`}\mathrm{`}`$equivalence” is a slight misnomer: what we impose instead are invariance condition on the (linear combinations of) admissible pictures representing any fixed $`vV`$ which ensure that the stabilized vectors are in fact equally weighted superpositions of all admissible pictures representing $`v`$.
Now consider the question, perhaps the first question a geometric topologist should ask about a modular functor $`V(\mathrm{\Sigma })`$; Can you draw a (local) picture of it on $`\mathrm{\Sigma }`$ so that the mapping class group of $`\mathrm{\Sigma }`$ acts on $`V(\mathrm{\Sigma })`$ by the obvious induced action on pictures?
We should not expect it to be easy to discover the local rules for the pictures associated to a given modular functor $`V`$ and in fact they may not exist in much generality. Recall that a three manifold $`Y`$ bounding $`\mathrm{\Sigma }`$, $`Y=\mathrm{\Sigma }`$ determines a vector $`v(Y)V(\mathrm{\Sigma })`$ so we might think of our proposed picture $`P\left(v(Y)\right)`$ drawn on $`\mathrm{\Sigma }`$ as some ghostly recollection of $`Y`$. The present understanding of modular functors is closely related to surgery formulas on links, but to think in this way we must choose a $`\mathrm{`}\mathrm{`}`$base point” $`3`$manifold $`Y_0`$ with $`Y_0=\mathrm{\Sigma }`$ to hold the links. This choice seems to create an asymmetry which should not be present in $`P\left(v(Y)\right)`$. Thus for a pictorial representation of $`V`$ which is derived from surgery, we expect only part of the mapping group $``$ that part extending over $`Y_0`$ will act locally. To localize V, this problem must be overcome.
Let us propose a meta theorem or $`\mathrm{`}\mathrm{`}`$principle” that solving the $`\mathrm{`}\mathrm{`}`$picture problem,” which we call $`\mathrm{`}\mathrm{`}`$combinatorial localization,” should imply both the Freed-Walker program, which we call $`\mathrm{`}\mathrm{`}`$algebraic localization” and the design problem for the Hamiltonian $`H`$ which we call $`\mathrm{`}\mathrm{`}`$physical localization.”
Figure 1
The solid arrow is asserted with some confidence at least as a mathematical statement; the dotted arrow is speculative. While the solid arrow seems unlikely to have a literal converse: ground states of even simple Hamiltonians in dimension $`2`$ are too complicated to draw pictures of; conceivably the dotted arrow might be an equivalence constituting a culmination of the Freed-Walker program.
## 2 Combinatorial localization of CS5 on marked disks, and the proof of the theorem.
We show how to represent CS5 (and by extension all CSr) on a disk with marked points by local pictures. Since the representation of quantum computing within CS5 \[FLW\] only used the braid group acting on a disk with marked points, this partial solution to the combinatorial localization problem will suffice to prove the theorem (once we have explained the solid arrow in figure 1).
For any $`r2`$, CS$`r`$ has a combinatorial localization on any cellulated disk with marked labeled points, (labels $`ϵ\{0,1,\mathrm{},r2\}`$ lie on the marked points and disk boundary) provided the cellulation has bounded combinatorics and the marked points stay sufficiently far from each other and the boundary. For a concrete statement, let us take the cellulated disk to be a rectangle $`R`$ with a square Euclidean cellulation. We suppose that all marked points are at least $`r`$ lattice spacings from the boundary and $`9r`$ from each other. The marked points and $`R`$ are all assigned the label $`1`$ (the irreducible $`2`$ dimensional representation of $`sl(2,)_q`$). In this circumstances it is easy to build a trivalent $`\mathrm{`}\mathrm{`}r`$collared rooted tree” $`T_r`$ for the disk with marked points as shown in figure 2.
Figure 2
All straight segments of the tree are to be more than $`3r`$ lattice bonds in length; the root is on $`R`$ and the leaves are the marked points. The $`r`$collard condition is that an $`\frac{r}{2}`$ relative regular neighborhood $`N(T)`$ of lattice cells - the region within the dashed line - should be imbedded in $`R`$.
The box counts in the statement of the theorem are designed to permit a (discontinuous) family of $`T_r`$’s to be found for at all times during braiding. We say that the boxing of $`R`$ is roomy relative to the location of the marked points if it has this property. The key Lemma 2.1 will show that for roomy boxing that two discrete pictures, which we regard as smoothly equivalent are in fact combinatorially equivalent. More precisely, the infinity of smooth averaging operators acting on the space of combinatorial pictures has exactly the same joint fixed set as a finite subset of combinatorial operators.
Let us begin with a geometric interpretation of CS5$`(\mathrm{\Sigma })=:V(\mathrm{\Sigma })`$. For a closed surfaces $`\mathrm{\Sigma }`$ it is implicit in \[K,L\]. Let $`\mathrm{\Sigma }`$ bound a handle body $`H`$. A general $`3`$manifold $`Y`$ with boundary $`\mathrm{\Sigma }`$ can now be represented as a $`\mathrm{`}\mathrm{`}`$blackboard framed” surgery diagram in $`H`$. The special cabling morphism $`w`$ of the Temperley-Lieb category (See chapter 12 \[K,L\] or \[R,T\]) when composed into the surgery diagram yields a linear combinations of $`1`$ manifolds, each labeled by $`\mathrm{`}\mathrm{`}1`$”. We may write $`H`$ as a planar surface cross interval, $`H\mathrm{\Sigma }\mathrm{\_}\times I`$, so that $`\mathrm{\Sigma }=\mathrm{\Sigma }\mathrm{\_}_{}\mathrm{\Sigma }\mathrm{\_}`$, where $`\mathrm{\Sigma }\mathrm{\_}`$ denotes $`\mathrm{\Sigma }\mathrm{\_}`$ with its orientation reversed. Now projecting these $`1`$manifolds to $`\mathrm{\Sigma }\mathrm{\_}`$, we see a linear combination of immersed $`1`$labeled $`1`$manifolds with overcrossings indicated at double points. This pictures determines the vector $`v(Y)`$. The Kauffman relations at a root of unity, in our case $`e^{2\pi i/5}`$, allow extensive simplification of these pictures via the recoupling formalism. In fact each $`vV`$ can be encoded as a labeling of a fixed (framed, imbedded, and vertex planar) trivalent graph, which is a spine for a $`\mathrm{\Sigma }\mathrm{\_}`$.
It is an important observation of Walker’s (personal communication) and Gelca’s \[G\] that this description can be extended to labeled surfaces with boundary. (Verification follows directly from the gluing axiom.) In the case of a disk with $`n`$ marked points $`(D,n)`$ \- treating marked points as crushed boundary components - the modular functor with $`n+1`$ labels $`\stackrel{}{\mathrm{}}\text{,}V_{\stackrel{}{\mathrm{}}}(D,n)`$ has as its basis $`q`$admissible labelings with boundary condition on a fixed trivalent tree imbedded in $`D`$, rooted on $`D`$, with leaves on the marked points. The boundary condition is that the label on the root is the label given on $`D^2`$ and each leaf has the label associated to its marked point. As in \[FLW\], we only need consider the case where all labels $`=1`$.
The (framed) braid group acts on the labeled tree $`T`$ via its imbedding in the disk. To see the induced action on $`V(D,n)`$ (we drop labeling subscripts), perturb the imbedding of $`T`$(rel its endpoints) by pushing it downward into a three ball $`D\times [0,1]`$, where we think of $`D`$ identified with $`D\times 0`$. Now implement any desired braid $`b`$ as a diffeormorphism of $`D\times [0,ϵ]`$ where $`ϵ>0`$ is small with respect to the previous push. Viewed from above, $`b(T)`$ has overcrossings but the recoupling $`(6j)`$ rules (and isotopies) allow $`b(T)`$ to be described in the original basis of $`q`$admissible labelings on $`T`$ (with root and leaves still carrying the label 1). For example the simplest Kauffman relations, on strands of $`b(T)`$ labeled by $`\mathrm{`}\mathrm{`}1`$” read:
$$\text{}=e^{\pi i/10})(e^{\pi i/10}\text{}\text{ and }=e^{\pi i/5}+e^{\pi i/5}=:d.$$
A detailed example: the effect of a single braid generator, is given immediately following the statement of Lemma 2.1 to elucidate the recoupling of braids.
There is a topological observation inherent in inducing the braid action on $`V(D,n)`$. By capping off, any diffeomorphism of a planar surface extends to the two sphere and can be extended further to a diffeomorphism of the $`3`$ball $`B^3`$. The action on $`V`$ comes from projecting this topological extension acting on labeled trivalent trees back into the original planar surface (after crushing the inner boundary components to points). In fact, it is the correspondence between $`3`$manifolds and diagrams which proves that we have correctly specified the action on the functor, for we have $`v(\overline{f}Y)=f_{}v(Y)`$ where $`\overline{f}|_{Y=\mathrm{\Sigma }}=f`$. Generally, when a surfaces $`\mathrm{\Sigma }`$ has genus $`>0`$ there will be no way of including it in the boundary of a $`3`$manifold $`M`$ so that all diffeomorphisms of $`\mathrm{\Sigma }`$ extend over $`Y`$. However it is a triviality that any diffeomorphism of $`\mathrm{\Sigma }`$ extends over $`\mathrm{\Sigma }\times I`$ by product with id<sub>I</sub>. Now let this extension act on the appropriate equivalence classes of framed $`q`$admissibly labeled trivalent graphs imbedded in $`\mathrm{\Sigma }\times I`$ projected back into $`\mathrm{\Sigma }`$ to define the action on any $`SU(N)`$level $`=r`$ modular functor $`V`$. Thus the $`\mathrm{`}\mathrm{`}`$doubled” functor $`V(\mathrm{\Sigma })V^{}(\mathrm{\Sigma })=V(\mathrm{\Sigma }\overline{\mathrm{\Sigma }})=V\left((\mathrm{\Sigma }\times I)\right)`$ has a combinatorial localization, i.e. is describable by local pictures. This may have some relation to unpublished work of Kitaev and Kupperberg (private communication) on local descriptions for Drinfeld doubles.
We set aside for later study the problem of devising combinatorial local rules for the necessary elementary equivalences of such trees $`T`$: $`6j`$moves, ribbon equivalence, vertex half-twist equivalence, and regular homotopy.
One would hope to define a quantum medium for CS5 of individual systems with levels to record labels $`0,1,2,3`$ (and possible additional levels to store other information) and terms $`H_k`$ with at most $`6`$ indices (as in a $`6j`$symbol) corresponding to these elementary equivalences. While this count seems correct in the smooth setting, there the crude Hilbert space is infinite dimensional which may create new difficulties. We have not been able to find a discrete setting in which all the equivalences are expressed efficiently. For the purpose of this lecture, we stay with discrete models for quantum media built from $`2`$level systems, but to do this we accept terms $`H_k`$ with up to $`30`$ indices.
The fundamental $`2`$dimensional representation of $`SU(2)`$ generates $`SU(2)`$’s complex representation ring and as a result recoupling theory achieves a very simple result: an element $`vV(D,n)`$ is a linear combination of imbedded $`1`$manifolds each labeled by $`\mathrm{`}\mathrm{`}1`$”, i.e. the standard $`2`$dimensional representation and given the boundary condition: each $`1`$manifold of the linear combination meets each marked point (and $`D`$) once. Thus $`\mathrm{`}\mathrm{`}`$manifoldness” and the $`\mathrm{`}\mathrm{`}`$boundary condition” define admissibility for our picture. this makes good sense combinatorically in the lattice of $`R`$, as well as, smoothly. We point out that our notion of $`1`$manifold is strict: at each vertex $`0`$ or $`2`$ edges (not $`4`$) should be occupied.
It is time to define the local equivalence moves between pictures. We are working within the Temperley-Lieb category modulo the relation that the $`(r1)^{\text{th}}=4^{\text{th}}`$ Jones-Wenzl projector is trivial. This is our most interesting relation. As a smooth equivalence relation this has only one form but combinatorially, we need to impose two versions of it according to how the output endpoints are grouped. We denote these by and . The second picture stands for : in conventional projector notation (\[K,L\]).
A second relation says that removing a circle which bounds a disk free from punctures multiples the diagram by the scalar $`\frac{1}{d}`$, $`d=e^{\pi i/5}+e^{\pi i/5}`$. A third relation replaces the undercrossing that arise through braiding with legitimate morphisms if the category. In terms of smooth pictures, the relation replaces the $`\mathrm{`}\mathrm{`}`$virtual” uncrossing in the middle diagram with a two term sum:
Figure 3
The middle picture is $`\mathrm{`}\mathrm{`}`$ virtual”; it is not actually an admissible picture to be assigned a weight. This relation requires a little care and lattice space to discretize since we do not want to permit the intermediate picture:
Figure 4
which would represent the wrong boundary data at the indicated defect. Recall that each defect is labeled by $`1`$ representing the $`2`$dimensional irreducible representation of $`sl(2,)_q`$ which is recorded by a single line leaving the defect.
Finally, a fourth class of equivalence permits isotopy. Again the reader should note that enough neighboring sites should be observed by the appropriate $`H_k`$ to preserve imbeddedness. For example, cases 1 and 2 are allowable, case 3 is not.
Figure 5
There will be isotopy relations for arc endpoints as well. For example, in cases $`1`$ and $`2`$ of Figure 3 imagine the open circle filled to become an end point and the shorter of the two line segments meeting it deleted. Morally, we should define operators $`\overline{H}_k`$ which enforce the average of the initial $`I`$ and final $`F`$ configuration of cases $`1`$ and $`2`$. However there is a detail, to get the overall phase correct, and not settle for merely a projective representation, we must fix a base point direction: say the positive ray emanating from each endpoint at $`45`$ degrees and find positive semidefinite $`\overline{H}_k`$’s which assign zero norm to $`\frac{1}{\sqrt{2}}(I_1e^{\pi i/10}F_1)`$ and $`\frac{1}{\sqrt{2}}[I_2F_2]`$ in cases $`1`$ and $`2`$ respectively. These operators correspond to asserting equivalences: $`I_1e^{\pi i/10}F_1`$ and $`I_2F_2`$. The general rule is that a state obtained by clockwise (counterclockwise) isotopy through the base point direction must be adjusted by the phase $`+()e^{i\pi /2r}`$ before being averaged. Similarly there is an isotopy relation for the arc end point on the boundary circle of the disk $`D`$. Here some point on the boundary is chosen an phase is adjusted by $`(+)e^{i\pi /2r}`$ as this point is crossed clockwise (counterclockwise).
Let us return to the raltions $`=0=`$ . Combinatorically the first may be written out with the left hand side a $`3\times 3`$ lattice square foliated by parallel straight lines (of label $`=1`$). Wenzl’s \[We\], recursion formula, yields an identity equating $`4`$ parallel lines with a linear combination of $`13`$ $`\mathrm{`}\mathrm{`}`$smaller” terms each containing $`\mathrm{`}\mathrm{`}`$turn arounds.” The form of the relation is shown below:
Figure 6.0
In its other incarnation, the 4th Jones-Wenzl’s projector relation $`=0`$ looks like this:
Figure 6.1
The coefficients $`a_i`$ are rational functors of $`d`$ which can be computed from the Wenzl’s recursion relation for projectors (see pg. 18 \[K, L\] or \[We\]). Figure 6.0 is merely the lattice counterpart of the more familiar smooth relation, Figure 6.0, which may be applied within any diagram (at $`r=5`$) whenever four $`1`$ labeled lines are found running parallel. Obviously Figure 6.1 also has a smooth counterpart.
Figure 6.0
The admissibility conditions and the above four classes of $`\mathrm{`}\mathrm{`}`$equivalences” must be rewritten as operators $`A_i`$ and $`B_j`$ respectively; collectively denoted $`\overline{H}_k`$. Let $`G`$ denote the ground state of the soon-to-be-defined Hamiltonian $`H=\underset{𝑘}{\mathrm{\Sigma }}\overline{H}_k`$. Let $`V`$ denote the CSr modular functor of the disk with $`3n`$ marked points and all labels $`=1`$. Via recoupling, we may describe $`V`$ in the fashion of homology. Set $`P=`$ \[admissible pictures\] and write: $`V=V_s=P/_s`$, where $`_s`$ is the smooth-category equivalence relation corresponding to our four combinatorical equivalences: $`_c`$. Lemma 2.1 will prove that under the $`\mathrm{`}\mathrm{`}`$roomy hypothesis” $`_s`$ and $`_c`$ induce identical equivalence classes of admissible pictures (which of course are combinatorial objects). So we may also write $`V=V_c=P/_c`$. Our goal is to tailor $`H`$ so that the ground states $`gG`$ correspond bijectively to linear functionals $`\varphi :V`$ under the map $`\varphi \underset{\underset{\text{admissible pictures}}{p}}{\mathrm{\Sigma }}\varphi (p)(p)`$. This will identify $`G`$ with $`V^{}`$, but since $`V`$ has a canonical nonsingular Hermitian inner product $`(`$\[Wi\] and \[K,L\]$`)`$ this also gives an isomorphism $`GV`$.
The inner product $`<p_1,p_2>`$ is defined on pictures by imbedding the disk $`D`$ into the $`(x,y)`$plane, deforming $`p_1`$ upward rel endpoints and $`p_2`$ downward rel endpoints. The union of the deformed pictures $`\stackrel{~}{p}_1\stackrel{~}{p}_2`$ is a (vertically framed) link in $`R^3`$ and it Kauffman bracket is $`<p_1,p_2>`$. Note that the vertical framing is singular where $`p_1`$ and $`p_2`$ share a common lattice bonds meeting $`p_1`$ and $`p_2`$; here the convention is to bend such bonds of $`p_2`$ slightly clockwise at the endpoints internal to $`D`$ and counterclockwise at an endpoint on $`D`$.
The definition of the $`A_i`$ operators is quite obvious. Consider, a vertex $`v`$ in the interior of $`R`$. A Hermitian $`A_v`$ with $`4`$ indices whose ground state is spanned by classical states of valence $`0`$ or $`2`$ at $`v`$ is said to enforce $`\mathrm{`}\mathrm{`}1`$manifoldness” at $`v`$. Clearly the ground state of $`A_v`$ has dimension $`7`$. To enforce, instead, a $`\mathrm{`}\mathrm{`}`$defect” or marked point labeled by the fundamental representation, $`\mathrm{`}\mathrm{`}`$1” of $`SU(2)`$, we would use instead a Hermitian operator $`A_v^{}`$ with ground state spanned by the four classical states of valence $`=1`$ at $`v`$.
Turning now to $`\mathrm{`}\mathrm{`}`$relations” $`B_j`$ consider a box $`b`$ of $`R`$ centered in a $`3\times 3`$ square of boxes:
Figure 7
there are $`12`$ nonboundary edges $`\{e\}`$ (shown in bold). If $`\{c\text{’s}\}`$ are the nonempty (classical) manifold configuration of these edges, i.e. valence $`\{0,2\}`$ at each of the four internal vertices, and iff $`c_0`$ and $`c_1=c_0\text{xor}b\{c\text{’s}\}`$, set $`d=\frac{1}{\sqrt{2}}\left(c_0c_1\right)`$ and let $`\{d\}`$ be the set of such vectors. Let $`B_b=\underset{dϵ\{d\}}{\mathrm{\Sigma }}|d><d|`$ be the Hermitian operator with $`12`$ indices on $`(^2)^{\{e\}}`$ whose ground state is orthogonal to span $`\{d\}`$. $`B_b`$ is the operator which $`\mathrm{`}\mathrm{`}`$allows isotopy across $`b`$.”
To remove circles which bound disks we need, in the presence of isotopy, only introduce operators which deletes a box. This operator may be written as $`|\theta ><\theta |`$ where $`\theta `$ is a unit vector proportional to $`|\text{box}>+\left(e^{\pi i/5}+e^{\pi i/5}\right)|\varphi >`$.
We postpone the definition of the operator corresponding to figures 3 and 4 since this must involve the dynamics $`\mathrm{`}\mathrm{`}t`$” of $`H_t`$. Some trick is needed to avoided adding new levels to our system to encode $`\mathrm{`}\mathrm{`}`$crossings.”
The projector corresponding to , Figure 6.0, requires a $`24`$index operator acting on a $`3\times 3`$ grid of edges or $`\mathrm{`}\mathrm{`}`$box” $`B`$ whose $`1`$ dimensional excited state is spanned by the vector obtained by putting all fourteen term in Figure 6.0 on the left hand side of the equation. Similarly the projector corresponding to , Figure 6.1 is a 30 index operator acting on the bonds of a region the shape of l.h.s. in Figure 6.1. This $`\mathrm{`}\mathrm{`}`$nobby box” $`B^{}`$ is a $`2\times 5`$ rectangle union an additional small box in the middle of one of the long sides.
Now we turn to the dynamics. Almost all conditions $`H_k`$ that combine to yield $`H`$ are permanent, only the end point operators $`A_v^{}`$ should change as we execute braiding. Because of the technical problem illustrated in figure 4; any lattice resolution into a superposition of two $`1`$manifolds as in figure 3 may cause collision with other strands. One way to deal with this problem is to locate the marked points on a second lattice $`𝐋^{}`$ consisting of the mid points of the edges in the original Lattice $`𝐋`$ of boxes in $`R`$. This means that we have to add additional $`2`$index $`A`$ operators holding equal the two classical states on both halves of the original edges, i.e. ground state $`(A)=(00>,|11>)`$, and that the end point operators $`A_w^{}`$ actually occur (with $`2`$dimensional ground states) on the finer lattice $`𝐋^{}`$, $`wϵ𝐋^{}`$. The dynamics consists of moving an endpoint diagonally on $`𝐋^{}`$, i.e. translating one unit horizontally or vertically in the structure of $`𝐋`$. In Figure 8 the endpoint $`w`$ is moved horizontally to $`w^{}`$ by replacing: $`\{A_w^{},A_w^{}\}`$ with $`\{A_w,A_w^{}^{}\}`$. If $`w`$ and $`w^{}`$ are immediately adjacent in $`𝐋`$ the operator swap will cause the end point to travel around a corner.
Figure 8
This operator swap can be performed gradually by slowly turning the appropriate terms on or off. If the adiabatic theory is applicable, and following the proof of Lemma 2.1 we discuss the heuristics for an energy gap (in the theromdynamic limit) for the family $`H_t`$, $`\psi _t`$ will be carried to a unique ground state $`\psi _t`$ of $`H_{t+1}`$. This ground as a functional on pictures is identical to $`\psi _t`$ provided pictures are identified according to the obvious isotopy rules (and phase rules at endpoints). If the lattice is refined by a linear factor $`L`$, tunneling to an undesired orthogonal ground state $`\psi _{t+1}^{}`$ should, by arguments analogous to those for the stability of homology classes \[K2\], have amplitudes scaling like $`ϵ(L)=e^{\mathrm{\Omega }(L)}`$. The mathematical description for adiabatic evolution of the system is via the natural connection $`A`$ on the tautological bundle over the complex Grassmannian $`X`$: The time evolution of $`G:=\{`$ground states $`(H_t)\}`$ defines a path in $`X`$ and $`A`$transport covers this motion with a unitary (i.e.isometric) identification $`G_0G_t`$, for all $`t0`$. After a braiding $`b`$ is completed at time $`t=T`$, the self-identification $`G_0G_T`$ is the representation of the braid b.
A ground state $`gGP`$ defines a functional $`g^{}`$ on $`P`$ via orthogonal projection. The $`\{B_j\}`$ have been chosen to correspond to $`_c`$ precisely so that a unique extension $`\varphi `$ exists:
and $`\varphi `$ satisfies $`g=\mathrm{\Sigma }\varphi (p)(p)`$. Conversely given a functional $`\varphi `$ on $`V_c`$ the $`g`$ associated to $`\varphi `$ by the formula above lies in the null space of each $`B_j`$, so in fact $`G=V_c^{}`$.
The important remaining point is to see that after braiding, when the marked points have been returned to there original sites set-wise, that the induced transformation on the ground state is precisely, up to error $`ϵ(L)`$, the unitary CS5 representation originally introduced by Jones \[J\] and studied in \[FLW\]. But this follows from the recoupling theory as presented in \[K,L\] provided we show that the combinatorial relations that we have imposed through Hermitian operators $`\{B_j\}`$ in fact are sufficient to span all the relations implied by the infinitely many smooth relations between pictures, that is $`V_s=V_c`$. For this the following lemma suffices.
###### Lemma 2.1.
Let $`\rho =\underset{𝑖}{\mathrm{\Sigma }}a_ip_i`$ be a linear relation between admissible combinatorial pictures in $`(R,\{3n\})`$ which holds under $`_s`$, the smooth recoupling theory associated to CS$`r`$. Provided that the configuration $`\{3n\}R`$ is roomy in the rectangle $`R`$, the same relation already holds under $`_c`$.
Before proving the lemma let us carry out a simple calculation to get a feel for how the action of braiding is computed via recoupling theory. If the reader wishes to try more complicated examples, the formulas on pages 93-100 of \[K,L\] are helpful. Here we compute the effect of a braid generator on a vector $`\psi _{}V:=`$CS5 ($`3`$punctured disk) where each boundary component has label$`=1`$ ( the $`2`$dimensional representation of $`sl(2,)_q)`$ and to account for phase each boundary has a marked base point.
$$\psi _{}^{}\text{ is the diagram: }\text{},$$
which as a labeled tree is:
$$\psi _{}^{}=\text{}$$
Let $`b`$ be the counterclockwise braided of the right most pair of punctures. Then $`b\psi _{}^{}`$ is represented by:
$$\text{}=\underset{\mathrm{`}\mathrm{`}\text{virtual picture”}}{\text{}}=$$
$$()A\text{}+A^1\text{},\text{ where }A=e^{2\pi i/10}.$$
Now $`\stackrel{2}{\text{}}`$ is our notation for the Jones-Wenzl idempotentent $`\stackrel{2}{\text{}}=\frac{1}{\sqrt{d^21}}\left(\text{}\frac{1}{d}\right)`$ where $`d=A^2A^2`$ and, as usual, the open ends in the above diagrams can be interpreted as permitting arbitrary (but constant) extension to the outside. Note: the orthogonality relations
* $`<\frac{1}{d},\frac{1}{d}>=\frac{1}{d^2}`$ $`=1`$
* $`<\frac{1}{\sqrt{d^21}}\left(\text{}\frac{1}{d}\right),\frac{1}{\sqrt{d^21}}\left(\text{}\frac{1}{d}\right)>=\frac{1}{d^21}\left(\text{}\frac{2}{d}\text{}+\frac{1}{d^2}\text{}\right)=\frac{1}{d^21}\left(d^2\frac{2}{d}d+\frac{1}{d^2}d^2\right)=1`$, and
* $`<\frac{1}{\sqrt{d^21}}(\text{}\frac{1}{d}),\frac{1}{d}>=\frac{1}{d\sqrt{d^21}}(\text{}\frac{1}{d}\text{})=\frac{1}{d\sqrt{d^21}}(d\frac{1}{d}d^2)=0`$
Normalizing, $`\psi _{}=\frac{1}{d}\psi _t^{}`$ is a unit vector, $`<\psi _{},\psi _{}>=1`$.
From the definition of $`\stackrel{2}{\text{}}`$ we have: $`=\sqrt{d^21}`$ $`\stackrel{2}{\text{}}\frac{1}{d}`$. So we use this to expand the two parallel lines in the second term of $`()`$ to get:
$`b\psi _{}=`$ $`A\psi _{}+{\displaystyle \frac{A^1}{d}}\left(\text{}\right)`$
$`=`$ $`A\psi _{}+{\displaystyle \frac{A^1}{d}}(\sqrt{d^21}\text{}+{\displaystyle \frac{1}{d}}\text{})()`$
$`=`$ $`A\psi _{}+{\displaystyle \frac{\sqrt{d^21}}{d}}A^1\psi _2+{\displaystyle \frac{A^1}{d}}\psi _{}`$
$`\text{where }\psi _2:=`$ $`\text{},`$
$`=`$ $`\left(A+{\displaystyle \frac{A^1}{d}}\right)\psi _{}+{\displaystyle \frac{\sqrt{d^21}}{d}}A^1\psi _2`$
As a check on unitarity note that under the sequelinear pairing,
$`<b\psi _{},b\psi _{}>=`$ $`\left(A+{\displaystyle \frac{A^1}{d}}\right)\left(A^1+{\displaystyle \frac{A}{d}}\right)+\left({\displaystyle \frac{\sqrt{d^21}}{d}}A^1{\displaystyle \frac{\sqrt{d^21}}{d}}A\right)`$
$`=`$ $`1+{\displaystyle \frac{A^2+A^2}{d}}+{\displaystyle \frac{1}{d^2}}+{\displaystyle \frac{d^21}{d^2}}=1`$
Proof of Lemma 2.1. The argument is based on the Birkhoff curve shortening principle where by a family of imbedded arcs and circles can be $`\mathrm{`}\mathrm{`}`$pulled tight” to a shorter geodesic position without crossings developing. We work combinatorically. By the $`\mathrm{`}\mathrm{`}`$roomy hypothesis” there is an $`r`$collared tree $`T:=T_rR`$. Assign a positive weight $`w(\beta )`$ to each bond $`\beta `$ (or $`1`$cell) of the cellutation of $`R`$ so that $`w`$ grows rapidly with distance to $`T`$: as a good first approximation, we may take $`w(\beta )=10^{\mathrm{\#}(\beta )}`$ where $`\mathrm{\#}(\beta )=`$minimum number of bonds joining $`\beta `$ to $`T`$. Now for any (classical) picture $`p_i`$ define its length $`L(p_i)=\underset{\beta ϵp_i}{\mathrm{\Sigma }}w(\beta )`$. Permitting combinatorial isotopy (rel the marked points) and the removal of small circles, but not the undercrossing, , or relations, we may pull $`p_i`$ tight by local moves to equivalent pictures (up to a scalar) which steadily reduce $`L(p_i)`$ until a local minimum is reached. Call this step $`\mathrm{`}\mathrm{`}`$pull tight”. Because of our weight function $`w`$, the new $`p_i`$ will try to lie mainly in a small neighborhood of $`T`$, and only occupancy of the bonds close to $`T`$ will force parts of the picture to lie farther away. Also the picture, seeking to occupy the bonds near $`T`$ efficiently will have its strands running parallel to $`T`$ in $`(r2)\times (r2)`$ lattice blocks $`\beta `$ a thwart the middle $`r`$bonds of each of the distinguished length $`3r`$ segments of $`T`$. Also at the trivalent vertices of $`T`$ near which sufficiently many strands pass, we would like to see copies of l.h.s. Figure 6.1. This will be true up to a small isotopy (across a few boxes) and can be made true on the nose by modifying the weight function $`w`$ by adding a small term proportional to the distance from each trivalent vertex out to a distance $`r`$ from that vertex. Now apply at some site in a $`B`$ or at some site $`B^{}`$ if the opportunity presents. This breaks $`p_i`$ into $`\underset{𝑗}{\mathrm{\Sigma }}b_{ij}q_{ij}`$ and for all $`j`$, $`L(q_{ij})<L(p_i)`$. Pull tight again to remove the slack created by the $`\mathrm{`}\mathrm{`}`$turn arounds” in Figure 6.0 or 6.1. Alternate pulling tight with applications of or until no further reductions in length can be made in this way. Call this cycle $`\mathrm{`}\mathrm{`}`$pull and cut”. With a slight abuse of notation let $`q_{ij}`$ denote one of the terminal classical states of this process. Now allow a single $`\mathrm{`}\mathrm{`}`$under crossing” move (Figure 3) to further reduce $`L(q_{ij})`$ if such a move is available. Now alternated the $`\mathrm{`}\mathrm{`}`$pull and cut” cycle with single under crossing moves until no daughter picture (still denote $`q_{ij}`$) can have its length reduced by further iteration of this process.
Manifestly, all the daughter pictures $`q_{ij}`$ now lie in $`N(T)`$ and pass through all boxes $`B`$ parallel to $`T`$ and with $`3`$ or few strands and pass through each $`B^{}`$ in a standard way according to some admissible triple as explained below. Note that $`(3,3,2)`$ is not admissible. At this point it is simple to formally reorganize the term of this sum $`\underset{ij}{\mathrm{\Sigma }}b_{ij},q_{ij}`$ as $`\mathrm{\Sigma }c_{\mathrm{}}T_{\mathrm{}}`$ where $`T_{\mathrm{}}`$ is an admissible labeling of $`T`$. As explained in \[K,L\], the leaves and root $`T`$ of $`t`$ are always labeled by $`1`$ (this is our choice) and the admissibility condition says that other edges (i.e. components of the intrinsic $`1`$skeleton of $`T`$) are labeled by $`a,b,c,d,\mathrm{}`$ taken from $`\{0,1,2,\mathrm{},r2\}`$ so that at each trivalent vertex of $`T`$ the following relations hold on the triple of incident labels $`a,b,`$ and $`c`$:
$`a`$ $`b+c`$
$`b`$ $`c+a`$
$`c`$ $`a+b`$
$`a+b+c`$ $`0\text{(mod 2), and}`$
$`a+b+c`$ $`<2(r1)=8.`$
An admissible labeling $`T_{\mathrm{}}`$ is interpreted as a linear combination of pictures by replacing each edge with the Jones-Wenzl projector corresponding to its label. The set of admissible labeled trees $`\{T_{\mathrm{}}\}`$ is an orthogonal basis for the modular functor $`V_s(R,\{3n\})`$, defined topologically using the smooth equivalence relation. (The subscript $`s`$ is to emphasis that the smooth relations are used in this definition; of course $`V_s=V`$.)
Because of the assumed $`p=\underset{i,j}{\mathrm{\Sigma }}b_{ij},q_{ij}=0V_s(R,\{3n\})`$, $`c_{\mathrm{}}=0`$ for all admissible $`\mathrm{}`$. But each $`q_{ij}`$ is an imbedded arc pairing $`x`$ of the $`\{`$leaves $``$ root$`\}`$ in $`N(T)`$ satisfying the additional admissibility restriction at each trivalent vertex of $`T`$. Such pairings are an alternative (though not orthogonal) basis for the modular functor $`V(R,\{3n\})`$ so collected in this basis we have for each pairing type $`x`$, we have $`\mathrm{\Sigma }b_{ij}^x,q_{ij}^x=0`$ where
$`b_{ij}^x=b_{ij}`$ $`\text{if }q_{ij}\text{ has type }=x`$
$`=0`$ $`\text{if }q_{ij}\text{ has type }x.`$
But all $`q_{ij}`$ of a fixed type are clearly combinatorally equivalent $`(_c)`$. Thus we have found a combinatorial path through applications of $`(_c)`$ from $`p`$ to the empty picture, or more precisely to a sum of zero times various pictures. $`\mathrm{}`$
Unlike \[K2\] the individual summands of $`H`$ do not commute. The ground states of $`H`$ has been computed topologically, however the spectrum spec$`(H)`$ is less accessible. The most important question is the existence of an energy gap above the ground state which is constant under lattice refinement, $`L\mathrm{}`$, i.e. in the thermodynamic limit. The following heuristics motivate the conjectured energy gap.
In finite classical systems such as random walk on a graph diffusion time is well known to scale inversely with the spectral gap of the Laplacian. Similarly, in some simple quantum mechanical systems where exact calculation is possible, the energy gap scales inversely to the diffusion time between classical states. In \[K2\] where direct calculation yields an energy gap above the ground state, the classical states are cycles and the $`\mathrm{`}\mathrm{`}`$diffusion” is through elementary bordisms. Since we have set up our ground state to be analogous to homology: $`GP/_c`$ with pictures playing the role of cycles and our $`\{B_j\}`$ playing the role of bordisms we expect similar diffusion properties and hence an energy gap. In lemma 2.1 the proof shows that equivalent pictures $`p_1`$ and $`p_2`$ are connected by a $`\mathrm{`}\mathrm{`}`$path” $`\gamma `$ of deformations ($`\mathrm{`}\mathrm{`}`$down” from $`p_1`$ to a neighborhood of $`T_r`$ and then back up to $`p_2`$). Rapid diffusion corresponds to observing that there are a plethora of such paths and in fact the procedure for finding $`\gamma `$ is highly under determined. More difficult would be a rigorous implication between diffusion and spec$`(H)`$. Extending the analogy with \[K2\], in both cases when the lattice is refined by a factor of $`L`$, a sequence of $`O(L)`$ local operators is required to transform between a pair of orthogonal ground states. So given the existence of an energy gap, the Hamiltonian $`H`$ will be stable to order $`O(L)`$ in perturbation theory; formally corresponds to tunneling amplitudes between orthogonal ground states which scale like $`e^{\mathrm{\Omega }(L)}`$.
There are several important open questions. The first is a rigorous treatment of the energy gap, but this is probably too difficult in the present model. Another is how to deal with errors in the form of actual rather than $`\mathrm{`}\mathrm{`}`$virtual” excitation which have already been discussed in the context of tunneling. Can a coupling to a could bath repair such errors or are more active measures required? For example, can broken endpoint pairs of a $`1`$manifold find each other and cancel through some imposed attraction (as suggested by Dan Gottesman in conversations) or merely through random walk? Nearby error pairs may be more serious in CS5 than in the toric codes since isotopy class not just homology needs to be preserved; the wrong reconnection pairing would result in an unrecoverable error. To make this unlikely, should additional terms be included into our Hamiltonian $`H`$ which could force distinct strands to be widely separated? This would put more weight on the simpler pictures, which are the ones that the quantum medium can most easily correct if damaged.
Kitaev’s very general notion of quantum media with its several antecedents in the study of quantum statistical mechanics looks likely to become a central object of study shared between theoretical physics, solid state physics, and topology. The main disappointment of the present investigation is the complexity of the local Hamiltonian $`H`$ used to construct stable universal topological quantum computation. One sees no easy road to radically simplifying it and still obtaining an exact description of CS5. However another path may be open. In our discussions, Kitaev has suggested (also see page 46 \[P\]) that simpler lattice Hamiltonians may renormalize in the scaling limit to topological modular functors. Perhaps the most interesting topological theories, such as CS5, because of their simplicity will have large $`\mathrm{`}\mathrm{`}`$basins of attraction” under renormalization and that identifiable universality classes of quantum media may not only exist mathematically but may even lie within the reach of engineers.
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# Very long storage times and evaporative cooling of cesium atoms in a quasi-electrostatic dipole trap
## Abstract
We have trapped cesium atoms over many minutes in the focus of a CO<sub>2</sub>-laser beam employing an extremely simple laser system. Collisional properties of the unpolarized atoms in their electronic ground state are investigated. Inelastic binary collisions changing the hyperfine state lead to trap loss which is quantitatively analyzed. Elastic collisions result in evaporative cooling of the trapped gas from 25 $`\mu `$K to 10 $`\mu `$K over a time scale of about 150 s.
The focus of a CO<sub>2</sub>-laser beam constitutes an almost perfect realization of a conservative trapping potential for neutral atoms . Atoms are confined in all three spatial dimensions by the optical dipole force pointing towards the maximum of the intensity . The CO<sub>2</sub>-laser wavelength of $`10.6\mu `$m is far below any optical transitions from the ground state which has important consequences. As one consequence, the optical potential becomes quasi-electrostatic, i.e. the static polarizability of the particle determines the depth of the trap (quasi-electrostatic trap, QUEST). Therefore, different atomic species and even molecules can be confined in the same trap. All substates of the electronic ground state experience the same trapping potential in contrast to magnetic traps. Atoms can thus be trapped in their absolute energetic ground state which excludes loss through inelastic binary collisions. Another consequence of the large laser detuning from resonance is the negligibly small photon scattering rate so that heating by the photon momentum recoil does not occur.
In this Rapid Communication we show that storage times of many minutes can be achieved in a focused-beam dipole trap with a low-cost, easy-to-use CO<sub>2</sub> laser usually employed for cutting and engraving of materials. The laser posseses neither frequency stabilization nor longitudinal mode selection. Despite the simplicity of the laser system, laser-noise induced heating rates, as first identified by Savard et al. , are found to be below 100 nK/s. Comparable storage times in a QUEST have recently been realized by O’Hara et al., who utilized an ultrastable, custom-made CO<sub>2</sub> laser . The simplicity of our trap setup in combination with the long storage times are ideal prerequisites for experiments on interesting collisional properties of the trapped gas. As one application of the trap, we have studied hyperfine-changing collisions of unpolarized cesium atoms. As another important application, evaporative cooling of the trapped gas is demonstrated, which has so far only once been observed in an optical dipole trap .
The trapping potential of a CO<sub>2</sub>-laser beam with a spatial intensity distribution $`I(𝐫)`$ is given by $`U(𝐫)=\alpha _{\mathrm{stat}}I(𝐫)/2\epsilon _0c`$ where $`\alpha _{\mathrm{stat}}`$ denotes the static polarizability of the atoms . For a focused beam of power $`P`$ and waist $`w`$ one gets a trap potential with a depth $`U_0=\alpha _{\mathrm{stat}}I_0/\epsilon _0c`$ with $`I_0=2P/\pi w^2`$. The CO<sub>2</sub> laser (Synrad 48-2WS) provides 25 W of power in a nearly TEM<sub>00</sub> transversal mode characterized by $`M^2=1.2`$. The laser beam is first expanded to a waist of 11 mm by a telescope, and then focused into the vacuum chamber by a lens of 254 mm focal length. The focus has a waist of 110 $`\mu `$m and a Rayleigh range $`z_\mathrm{R}`$ of 2.4 mm, yielding a trap depth 120 $`\mu `$K (in units of the Boltzmann constant $`k_\mathrm{B}`$). Gravity lowers the potential height along the vertical direction to 92 $`\mu `$K. The axial and radial oscillation frequencies in the harmonic approximation are given by $`\omega _\mathrm{z}=(2U_0/mz_\mathrm{R}^2)^{1/2}`$ and $`\omega _\mathrm{r}=(4U_0/mw^2)^{1/2}`$ with $`m`$ denoting the cesium mass. For our experimental values one gets an axial oscillation frequency $`\omega _\mathrm{z}/2\pi =8.1`$ Hz and a radial frequency $`\omega _\mathrm{r}/2\pi =254`$ Hz.
Atoms are transferred into the dipole trap from a magneto-optical trap (MOT) containing about $`10^6`$ atoms. The MOT is loaded from an atomic beam which is Zeeman-slowed in the fringe fields of the MOT magnetic quadrupole field. The main vacuum chamber at a background pressure of about $`2\times 10^{11}`$ mbar is connected to the oven chamber by a tube which is divided into differentially-pumped sections to assure a sufficiently large pressure gradient. To optimize transfer from the MOT into the dipole trap, the atoms are further cooled and compressed by decreasing the detuning of the MOT trapping laser from initially $`2\mathrm{\Gamma }`$ (natural linewidth $`\mathrm{\Gamma }/2\pi =5.3`$ MHz) to $`20\mathrm{\Gamma }`$ with respect to the 6<sup>2</sup>S<sub>1/2</sub>($`F=4`$) - 6<sup>2</sup>P<sub>3/2</sub>($`F=5`$) transition of the cesium D2 line. After 40 ms of compression, the distribution of atoms in the MOT has a rms radius of 120 $`\mu `$m corresponding to an mean density of $`10^{10}`$ atoms/cm<sup>3</sup>. The temperature is 25 $`\mu `$K as measured by ballistic expansion of the cloud after release from the MOT. After the laser beams and the magnetic field of the MOT were turned off, the atoms are trapped in the focus of the CO<sub>2</sub> laser beam which was present during the whole loading phase. Atoms are prepared in either the $`F=3`$ or the $`F=4`$ cesium hyperfine ground state by shuttering the MOT trapping laser 1 ms after or before the MOT repumping laser has been shuttered, respectively.
The number and spatial distribution of atoms trapped in the optical dipole trap are measured by taking an absorption image of the trapped atoms. A weak, resonant probe beam of 1 $`\mu `$W/cm<sup>2</sup> intensity is pulsed for 100 $`\mu `$s and absorption of the atomic cloud is imaged onto a CCD camera. Fig. 1(a) shows a typical absorption image of atoms trapped in the QUEST. The image is taken 5 s after transfer from the MOT. The transmitted intensity $`I_t`$ of the probe laser through the atomic sample is described by $`I_t(x,y)=I_0\mathrm{exp}\left(A\eta (x,y)\right)`$ where $`I_0`$ denotes the laser intensity and $`A`$ the absorption cross section for the resonant transition. The column density $`\eta `$ is given by the integral of the density distribution $`n(x,y,z)`$ along the direction of the laser beam.
By fitting a thermal equilibrium distribution to the data, we derive the mean density $`\overline{n}`$, the total number of trapped atoms $`N`$ and the temperature $`T`$. The maximum absorption of the trapped cloud is typically 17 % in the center of the distribution yielding a mean density of $`4\times 10^9`$ atoms/cm<sup>3</sup> in the dipole trap. The atoms have axially expanded into an rms extension of 750 $`\mu `$m while the radial rms extension is 30 $`\mu `$m. The temperature is the same as in the MOT before transfer indicating that the atoms are cooled into the dipole trap by the MOT.
Typically $`10^5`$ atoms are transferred into the dipole trap. Assuming sufficient ergodicity, the number of atoms transferred from the MOT into the dipole trap can be determined from the phase-space integral $`f_{\mathrm{MOT}}(𝐱,𝐩)\theta (U_0ϵ)𝑑𝐱𝑑𝐩`$ which describes the projection of the phase-space distribution $`f_{\mathrm{MOT}}(𝐱,𝐩)`$ in the MOT onto the trapping region of the dipole trap. The Heavyside step function $`\theta (U_0ϵ)`$ equals 1 when the total energy $`ϵ=U(𝐫)+p^2/2m`$ of atoms in the dipole trap is smaller than the trap depth, and 0 elsewhere. Taking our experimental parameters, one expects $`2\times 10^5`$ atoms to be transferred into the dipole trap in reasonable agreement with the actual value.
To measure the radial oscillation frequency of atoms in the dipole trap, we take advantage of the fast switching capability of the CO<sub>2</sub> laser. The laser can be turned off within about 200 $`\mu `$s by simply switching the RF power supply driving the gas discharge. By turning the laser off for a short time interval (around 1 ms), the trapped ensemble moves ballistically until the laser is turned on again (release-recapture). Part of the atoms will have escaped from the trap, while the recaptured atoms constitute a non-equilibrium distribution which oscillates at twice the oscillation frequency. When the first release-recapture cycle is followed by a second one, the number of finally recaptured atoms depends on the phase of the oscillation and thus on the delay time between the two cycles.
In Fig. 1(b), the number of recaptured atoms is plotted as a function of delay time between the two release-recapture cycles. One observes some cycles of coherent oscillation of the ensemble until the oscillation dephases mainly due to the anharmonicity of the potential. The oscillation period of 1.8 ms gives a radial oscillation frequency of 270 Hz in good agreement with the value expected for the trap parameters. The solid line in Fig. 1(b) shows the result of a Monte-Carlo simulation of the classical atom trajectories for the sequence of two release-recapture cycles with variable time delay between them.
The lifetime of atoms in the dipole trap is determined by recapturing the atoms back into the MOT after a variable storage time. For each loading and trapping cycle, three camera pictures are taken. The first picture gives the flourescence of atoms in the MOT shortly before transfer into the dipole trap. The second one shows the fluorescence distribution after recapture into the MOT. The last picture is taken after the atoms have been released from the recapture-MOT to provide the background which is then subtracted from the first two images. The number of atoms is determined from the integral over the fluorescence image. The particle number determined by this method is estimated to be accurate within a factor of two. By normalizing the number of recaptured atoms to the number of atoms initially in the MOT, fluctuations due to variations in the initial number of atoms can be cancelled out. Variations of the particle number in the MOT, however, are found to be below 10%, so that normalization was not used in the data presented here.
The decay of the number of stored atoms is shown in Fig. 2 for the two hyperfine ground states. Even after 10 min storage time, a few hundred atoms can be detected when prepared in the $`F=3`$ state. For atoms in the energetically higher $`F=4`$ hyperfine ground state, inelastic binary collisions between trapped atoms lead to additional trap loss . The energy of $`h\times 9.2`$ GHz released in the collisions is much larger than the trap depth, therefore both collision partners are ejected from the trap.
The decay curve for atoms in the $`F=4`$ state is fitted by the solution to the differential equation $`dN/dt=\gamma N\beta \overline{n}N`$ where $`\gamma `$ is the rate for trap loss through collisions with background gas and $`\beta `$ is the rate coefficient for inelastic binary collisions. The fit yields a decay constant $`\gamma ^1=165(25)`$ s. This value is consistent with the expected loss rates for collisions with background gas atoms at a pressure of $`10^{11}`$ mbar. No indication for laser-noise induced trap loss is found. From the fit one finds a decay coefficient $`\beta =2(1)\times 10^{11}`$ cm<sup>3</sup>/s with the error being mainly due to the uncertainty in the absloute particle number. A previous measurement of this quantity in an opto-electric trap gave a similar result.
Atoms in the energetic ground state $`F=3`$ can not undergo inelastic collisions. One would therefore expect a purely exponential decay. However, the decay curve in Fig. 2 shows a faster loss of particles at higher particle numbers. It was checked that more than 95% of the particles are initially prepared in the absolute ground state so that inelastic collisions between trapped particles can be excluded. At particle numbers below $`10^4`$, the decay curve approaches a pure exponential with a decay constant of $`\gamma ^1=140(20)`$ s in agreement with the value found for atoms in the $`F=4`$ state.
The faster initial trap loss can be attributed to evaporation of high-energetic atoms from the trap leading to cooling . At the initial temperature, the trap depth $`U_0`$ is only about $`5k_\mathrm{B}T`$. Therefore there is a certain probability that atoms leave the trap after an elastic collision. The rate for elastic collisions is given by $`\overline{n}\sigma \overline{v}`$ with the cross section $`\sigma `$ and the mean relative velocity $`\overline{v}=4(k_\mathrm{B}T/\pi m)^{1/2}`$. For evaporative cooling to be effective, the ratio between the elastic collision rate (providing thermalization and evaporation) and the rate for inelastic collisions (causing losses and heating) has to be large.
Up to now, only one experiment is reported on evaporative cooling in an optical dipole trap . In order to achieve sufficiently high densities, Adams et al. stored sodium atoms in a crossed-beam dipole trap which provides tight confinement in all three dimensions. In the simple focused-beam geometry used in our experiment it is rather surprising to find evaporative cooling since the achievable densities are rather low. The reason that evaporation actually takes place lies in the long storage times of our trap on the one hand, and the anomalously large elastic cross section of cesium at low temperatures on the other hand. Although the density of cesium atoms in the dipole trap is only of the order of $`10^9`$ atoms/cm<sup>3</sup>, one expects a thermalization time of only a few seconds due to the existence of a zero-energy resonance . This time scale is indeed more than an order of magnitude smaller than the storage time.
In Fig. 3, the temperature of the trapped atoms in the $`F=3`$ state is plotted versus storage time. Temperatures derived from the density distribution shown as the dots are in good agreement with measurements of ballistic expansion after release from the trap, which are shown as the three additional points in Fig. 3. One clearly observes cooling of the gas caused by evaporation of atoms from the trap at constant trap depth (plain evaporation). The final temperature of about 10$`\mu `$K corresponds to roughly 1/10 of the potential depth. Although the temperature is reduced by roughly a factor of 2 after 150 s, the phase-space density remains almost constant since the particle number diminishes by a factor of 10 at the same time. The temperature evolution for atoms in the $`F=4`$ state shows a similar behaviour, but with a slower decrease of the temperature. This is to be expected due to the faster initial density decrease through inelastic collisions and the according decrease in the rate of elastic collisions.
We apply a model developed by Luiten et al. to simulate the temperature and particle number evolution during evaporation. Given the shape of the potential and the corresponding density of states, the model provides two coupled differential equations for the evolution of temperature and particle number. For the true potential function of a focused Gaussian beam, the density of states diverges as the energy approaches the escape energy of the trap. We therefore approximate the trap potential by a three-dimensional Gaussian $`U(𝐫)U_0\mathrm{exp}(2(x/w)^22(y/w)^2)(z/z_\mathrm{R})^2)`$ for which the density of states remains finite.
The potential is fully determined by the trap parameters, therefore the only adjustable parameter in the model is the the rate $`\gamma `$ for inelastic collisions with background gas and the cross section $`\sigma `$ for elastic collisions. In the temperature range considered here, mainly s-wave collisions contribute to the cross section. The large scattering length of cesium atoms in the $`F=3`$ ground state results in a temperature-dependent effective cross section $`\sigma (T)=\pi ^2\mathrm{}^2/mk_\mathrm{B}T`$ (unitarity limit. Note, that the cross section in the unitarity limit is completely determined by the temperature. Since the model explicitly assumes a constant cross section for elastic collisions , we have fixed the effective cross section to the value $`\sigma =(930a_0)^2`$ ($`a_0`$ = Bohr radius) corresponding to a temperature of 15 $`\mu `$K.
The result of the simulation for $`\gamma ^1=130`$ s is shown by the solid line in Fig. 3. The model prediction for the time scale of temperature decrease is in reasonable agreement with the experimental data. As can be seen from Fig. 3, however, the model slightly overestimates the rate of evaporation. This discrepancy can be explained by the influence of gravity on the one hand, and the temperature dependence of the collision cross section on the other hand. Both effects are not contained in the model, adn are difficult to be included. Due to gravity, evaporation predominantly takes place in only one spatial dimension which slows down the cooling process . The temperature-dependent cross section gives rise to an increased thermalization time so that the rate of evaporative cooling is decreased. Nevertheless, the final temperature of the evaporation process is correctly reproduced by the model. The model also shows that the initial faster loss of particles in the $`F=3`$-state (see curve in Fig. 2) mainly stems from the evaporation process.
Employing a low-cost, off-the-shelf CO<sub>2</sub> laser without any longitudinal mode selection, our major concern was possible heating by laser-noise induced fluctuations of the trap potential. From the observation of evaporative cooling with a temperature decrease of about $`10\mu `$K over a time scale of 100 s, one can infer that heating rates by laser noise are much lower than 100 nK/s. From the storage times of several minutes one can infer a similar upper bound for the laser-noise induced heating rates.
Due to its simple ingredients, our setup represents a minimalistic version of an optical dipole trap providing very long storage times. The universality of the QUEST opens way to experiments aiming at fundamental questions. Atomic spin states can be prepared in the QUEST with very long coherence times which is of importance for applications in quantum computing and for the search of a permanent electric dipole moment of atoms. Another intriguing prospect is the formation and the storage of cold molecules. In the same dipole trap, we have recently also trapped lithium atoms achieving similar storage times as for cesium . Currently, we are studying the properties of a simultaneously stored mixture of both species in view of sympathetic cooling and photoassociation of cold heteronuclear dimers.
We gratefully acknowledge contributions of M. Nill in the early stage of the experiment and stimulating discussions with A. Mosk. We are indebted to D. Schwalm for encouragement and support.
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# Computational probes of Collective excitations in low-dimensional magnetism
## 1 Introduction
Notwithstanding the fact that experimental and theoretical studies of condensed matter systems are fundamentally complementary to each other, they share important features, which we wish to illuminate here in the context of quantum many-body dynamics. The object of study is a sample for experimental probes and a model for computational probes. Whether the model is regarded as a mathematical idealization of a real chunk of matter or whether the sample is viewed as a physical realization of a system defined with mathematical rigor may be a matter of perspective, but the influence exerted by theory and experiment is mutual. Any observation of note calls for an explanation, while any prediction of substance brings into motion attempts at verification.
When performing experimental or computational probes in quantum many-body dynamics, the goal of researchers is very much the same, namely trying to make sense of what, in general, is a jumble of fluctuations to which the fundamental degrees of freedom contribute in various combinations and configurations. Experimental probes begin with a measurement and computational probes with a calculation. The results of either probe then require an interpretation in terms and concepts that are common to both approaches.
Qualitatively, the specification of a quantum many-body system involves information on composition, interaction, and environment. Theoretically this information is encoded in the Hamiltonian and in the density operator, whereas experimentally part of this information is manifest or hidden in the sample and other parts are controlled by the instrumental setup.
Any system thus specified is subject to fluctuations, of which we distinguish three kinds. Depending on the circumstances, each kind may govern the dynamical properties of the system. Thermal fluctuations are likely to be dominant in any statistical ensemble at elevated temperatures. Parametric fluctuations are manifestations of quenched random inhomogeneities in composition (e.g. dilution) or interaction (e.g. random bonds or fields), leaving distinct marks in dynamical properties, especially at low temperatures.
In the absence of thermal and parametric fluctuations, quantum fluctuations remain present. They are a direct consequence of the (autonomous) quantum time evolution and the time-delayed projections that are part of any dynamical probe, be it experimental or computational. No zero point motion exists in classical Hamiltonian systems. Here the ground state corresponds to a stable fixed point in phase space, and all dynamical variables are constant.
Consider a (non-random) quantum many-body system in a pure quantum state. The quantum fluctuations then depend on three quantities: (i) the Hamiltonian $`H`$, (ii) the state, which we take here to be the ground state $`|G`$, and (iii) some dynamical variable, here expressed by operator $`A`$. The Hamiltonian governs the dynamics deterministically for as long as the system can be regarded in isolation. It makes no difference whether this time evolution is viewed in the Schrödinger or Heisenberg representation:
$`i\mathrm{}{\displaystyle \frac{d}{dt}}|\psi =H|\psi `$ $``$ $`|\psi (t)=e^{iHt/\mathrm{}}|\psi ,|\psi A|G`$ (1)
$`i\mathrm{}{\displaystyle \frac{dA}{dt}}=[A,H]`$ $``$ $`A(t)=e^{iHt/\mathrm{}}Ae^{iHt/\mathrm{}}.`$ (2)
The statistical element is introduced when we subject the time evolved state $`|\psi (t)`$ or, equivalently, the time evolved dynamical variable $`A(t)`$, to an observation. In practice, this step involves the observation or evaluation of a dynamic correlation function, which can then be viewed either as the projection of the time evolved state $`|\psi (t)`$ or as the product of the dynamical variable $`A(t)`$ at two different times evaluated in the stationary state $`|G`$:
$$A(t)A\stackrel{\mathrm{equil}.}{=}AA(t)=\{\begin{array}{cc}e^{iE_Gt/\mathrm{}}\psi (0)|\psi (t)\hfill & \\ G|Ae^{iHt/\mathrm{}}Ae^{iHt/\mathrm{}}|G.\hfill & \end{array}$$
(3)
The sum total of quantum fluctuations in a typical many-body system contains many different dynamical modes no matter in which state the system happens to be. Only once we pick a set of dynamical variables (one variable at a time) do we begin to gain insight into the role of specific modes taken from a multitude of quantum fluctuations. Criteria for choosing dynamical variables include experimental accessibility, computational amenability, or any specific purpose like the study of ordering tendencies or phase transitions.
On the microscopic level, dynamical probes are bound to be somewhat heavy-handed. No observation without perturbation! However, in all situations considered here, we shall assume that the interaction between the probe and the sample takes place in the regime of linear response. This assumption is quite realistic for neutron scattering under most circumstances. Consider the weakly perturbed Hamiltonian $`H(t)=H_0b(t)B`$, where the external field $`b(t)`$ couples to the dynamical variable $`B`$ of the system $`H_0`$. The linear response of any other dynamical variable $`A`$ to that perturbation,
$$A(t)A_0=_{\mathrm{}}^+\mathrm{}𝑑t^{}\stackrel{~}{\chi }_{AB}(tt^{})b(t^{}),$$
(4)
is then determined by Kubo’s formula for the response function,
$$\stackrel{~}{\chi }_{AB}(tt^{})=\frac{i}{\mathrm{}}\mathrm{\Theta }(tt^{})[A(t),B(t^{})]_0.$$
(5)
Its Fourier transform, the generalized susceptibility $`\chi _{AB}(\omega +iϵ)`$, has (for $`ϵ0`$) a symmetric real part and an antisymmetric imaginary part: $`\chi _{AB}(\omega )=\chi _{AB}^{}(\omega )+i\chi _{AB}^{\prime \prime }(\omega )`$. The latter is related to the Fourier transform of the correlation function $`A(t)A`$, named the structure function
$$S_{AA}(\omega )_{\mathrm{}}^+\mathrm{}𝑑te^{i\omega t}A(t)A=\frac{2\mathrm{}\chi _{AA}^{\prime \prime }(\omega )}{1e^{\beta \mathrm{}\omega }},$$
(6)
via the fluctuation-dissipation theorem. The spectral representation of this function, in particular the expression resulting in the limit $`T0`$, where thermal fluctuations are absent,
$$S_{AA}(\omega )\stackrel{T=0}{=}2\pi \underset{\lambda }{}|G|A|\lambda |^2\delta \left(\omega (E_\lambda E_G)/\mathrm{}\right),$$
(7)
demonstrates that observing quantum fluctuations in the linear response regime means observing collective modes and their transition rates from the ground state. In the following, we look at quantum fluctuations and the associated collective modes from this particular angle for a number of situations as investigated by a variety of methods.<sup>1</sup><sup>1</sup>1From here on we set $`\mathrm{}=1`$ to simplify the notation.
## 2 Magnons Excited from Valence-Bond-Solid State
Consider the one-dimensional (1D) $`s=1`$ Heisenberg model with bilinear and biquadratic exchange
$$H_\theta =J\underset{n=1}{\overset{N}{}}\left[\mathrm{cos}\theta (𝐒_n𝐒_{n+1})+\mathrm{sin}\theta \left(𝐒_n𝐒_{n+1}\right)^2\right].$$
(8)
At $`T=0`$ it has a phase with ferromagnetic long-range order (LRO) ($`\pi <\theta <3\pi /4`$ or $`\pi /2<\theta \pi `$), a phase with dimer LRO ($`3\pi /4<\theta <\pi /4`$), the Haldane phase with hidden topological LRO ($`\pi /4<\theta <\pi /4`$), and an obscure phase ($`\pi /4<\theta <\pi /2`$) that was named trimerized. These ordering tendencies help us choose suitable dynamical variables when we explore the spectrum of collective excitations. All dynamical variables used for this purpose will have the form of fluctuation operators,
$$F_Q^A\frac{1}{\sqrt{N}}\underset{n=1}{\overset{N}{}}e^{iQn}A_n,$$
(9)
where the operator $`A_n`$ acts locally at or near lattice site $`n`$. The associated structure function (7) is called the dynamic structure factor:
$$S_{AA}(Q,\omega )_{\mathrm{}}^+\mathrm{}𝑑te^{i\omega t}F_Q^A(t)F_Q^A^{}\stackrel{T=0}{=}\underset{\lambda }{}W_\lambda ^A\delta (\omega \omega _\lambda ),$$
(10)
where the sum runs over the dynamically relevant excitations $`|\lambda `$ with energy $`\omega _\lambda =E_\lambda E_G`$ and spectral weight $`W_\lambda ^A=2\pi |G|F_Q^A|\lambda |^2`$. To calculate $`S_{AA}(Q,\omega )`$ for $`H_\theta `$ we can employ special methods at the parameter values $`\theta =\pm \pi /4`$ where it is completely integrable or general methods otherwise. One suitable general method in the present context is the recursion method in combination with a finite-size continued-fraction analysis.
Here we introduce four different fluctuation operators $`F_Q^A`$. The local operators $`A_n`$ from which they are constructed are listed in Table 1.
$``$The spin fluctuations, probed by $`F_Q^S`$, represent Néel order parameter fluctuations for $`Q=\pi `$. They are expected to be strongest at the critical point $`\theta =\pi /4`$, where the $`Q=\pi `$ excitations are gapless.
$``$The dimer fluctuations, probed by $`F_Q^D`$, are also expected to be strongest (for $`Q=\pi `$) at $`\theta =\pi /4`$, which marks the onset of dimer LRO.
$``$The trimer fluctuations, probed by $`F_Q^T`$, are constructed from projection operators $`P_n^T`$ onto local trimer states $`|[1,2,3]`$. The state $`|[1,2,3]`$ happens to be the ground state of $`H_\theta `$ with $`N=3`$ for $`\mathrm{arctan}\frac{1}{3}\theta \pi /2`$, which was interpreted as suggesting that a trimerized phase might exist for $`N\mathrm{}`$.
$``$The center fluctuations, probed by $`F_Q^Z`$, are constructed from a modified spin operator and tune into existing period-three $`(+,0,)`$ or $`(,0,+)`$ patterns of local spin states. Finite-$`N`$ data indicate that for $`\pi /4\theta \pi /2`$, the fluctuations probed by $`F_{2\pi /3}^Z`$ are the strongest of all the ones listed.
The order parameters associated with the four fluctuation operators defined by Eq. (9) and the entries of Table 1 can be written in the form
$$P_A=\frac{1}{N}\underset{n=1}{\overset{N}{}}e^{iQ_An}A_n,$$
(11)
where $`Q_S=Q_D=\pi `$, and $`Q_T=Q_Z=2\pi /3`$. Each order parameter $`P_A`$ has a set of eigenvectors $`|\mathrm{\Phi }_k^A,k=1,2,\mathrm{},K`$ which represent the associated LRO in its purest form. The degree of degeneracy is $`K=2`$ for the Néel and dimer states, $`K=3`$ for the trimer states, and $`K=6`$ for center states.
The physical vacuum chosen here is very unlike any of the states $`|\mathrm{\Phi }_k^A`$. Within the Haldane phase, at the parameter value $`\theta _{VBS}=\mathrm{arctan}\frac{1}{3}`$, the ground state of $`H_\theta `$ is a realization of the valence-bond solid (VBS) wave function, which is non-degenerate and in which the Néel, dimer, trimer, and center ordering tendencies are all imperceptibly weak. In the VBS state, the spin 1 at each lattice site is expressed as a spin-1/2 pair in a triplet state. The singlet-pair forming valence bond involves one fictitious spin 1/2 from each of two neighboring lattice sites. The VBS state, which has total spin $`S_T=0`$, can then be regarded as a chain of valence bonds linking successive spin-1/2 pairs in this manner.
The topological LRO present in the Haldane phase and known to be strongest in the VBS state provides an environment, where a specific kind of elementary excitations can propagate freely and where the associated stationary states form a branch with well defined dispersion. We now probe these elementary excitations and composites thereof from different angles by the four fluctuation operators $`F_Q^A`$ defined in Table 1.
Only two of the $`F_Q^A`$ are fully rotationally invariant, $`[F_Q^D,S_T^i]=[F_Q^T,S_T^i]=0`$ for $`i=x,y,z`$. This produces the selection rules $`\mathrm{\Delta }S_T=0`$ in the dynamic dimer and trimer structure factors. The corresponding selection rules for the dynamic spin and center structure factors are $`\mathrm{\Delta }S_T=0,1`$ and $`\mathrm{\Delta }S_T=0,1,2`$, respectively, with the further restriction that transitions between singlets ($`S_T=0`$ states) are prohibited. At the VBS point, $`S_{DD}(Q,\omega )`$ and $`S_{TT}(Q,\omega )`$ thus couple exclusively to the $`S_T=0`$ excitation spectrum, and $`S_{SS}(Q,\omega )`$ exclusively to the $`S_T=1`$ excitation spectrum, whereas $`S_{ZZ}(Q,\omega )`$ couples to the $`S_T=1`$ and $`S_T=2`$ spectra.
In Fig. 1 we display $`\omega _\lambda ^A`$ versus $`Q`$ of the dynamically relevant spin, center, dimer, and trimer excitation spectra as obtained from the finite-size continued-fraction analysis.
The filtered access to the spectrum afforded by the four fluctuation operators gives us valuable clues about the nature of the elementary excitations that thrive in the VBS environment. The low-frequency region at $`Q\pi /2`$ in the spin and center spectra is dominated by a branch of states with $`S_T=1`$. They carry more than 95% of the spectral weight in $`S_{SS}(Q,\omega )`$ and $`S_{ZZ}(Q,\omega )`$. These triplets, which have been named magnons, remain invisible in the dimer and trimer spectra. Only composites of the magnon states may be observable in $`S_{DD}(Q,\omega )`$ and $`S_{TT}(Q,\omega )`$.
A dispersion of the general form $`\omega _M(Q)=J(a+b\mathrm{cos}Q)`$ for the 1-magnon branch can be inferred from the single-mode approximation of $`S_{SS}(Q,\omega )`$. Under the assumption that in the VBS environment the magnons are weakly interacting point particles, we can expect the existence of three kinds of 2-magnon scattering states formed by pairs of 1-magnon triplets: states with $`S_T=1`$, which contribute to $`S_{SS}(Q,\omega )`$ and $`S_{ZZ}(Q,\omega )`$, states with $`S_T=0`$, which contribute to $`S_{DD}(Q,\omega )`$ and $`S_{TT}(Q,\omega )`$, and states with $`S_T=2`$, which are observable in $`S_{ZZ}(Q,\omega )`$ only. Free 2-magnon states form a two-parameter continuum $`\omega _{2M}(k,Q)\omega _M(Q/2k)+\omega _M(Q/2+k)`$ in $`(Q,\omega )`$-space. The resulting continuum boundaries are $`\omega _\pm (Q)=2J[a\pm b\mathrm{cos}(Q/2)]`$.
The predicted coalescence of the 2-magnon continuum into one spectral line at $`Q=\pi `$ follows from the symmetry property $`\omega _M(Q)+\omega _M(\pi Q)=\mathrm{𝑐𝑜𝑛𝑠𝑡}`$ of the magnon dispersion. Interestingly, the finite-$`N`$ dimer spectrum does indeed collapse into a single spectral line at the $`N`$-independent excitation energy $`\omega _D=\sqrt{10}J`$, which carries all the spectral weight in $`S_{DD}(Q,\omega )`$.
We now use the exact 2-magnon excitation energy, $`\omega _\pm (\pi )=\omega _D`$, and the extrapolated value, $`\omega _M(\pi )=0.66433(2)J`$, of the 1-magnon excitation gap to fit the parameters $`a,b`$. The resulting values, $`a1.581,b0.917`$, used in $`\omega _M(Q)`$ and $`\omega _\pm (Q)`$ yield the dashed and solid lines in Fig. 1.
Both the 1-magnon dispersion and the 2-magnon spectral threshold are in very good agreement with all finite-$`N`$ data shown. Hence the magnon interaction is very weak at the bottom of the 2-magnon region in all three $`S_T`$ subspaces. The finite-$`N`$ data spilling out on top of the shaded areas in Fig. 1 suggest that at higher energies, the magnon interaction is repulsive, more strongly so in the $`S_T=0,2`$ subspaces than in the $`S_T=1`$ subspace.
This raises the possibility that bound 2-magnon states split off the top of the 2-magnon continuum of 2-magnon scattering states in the $`S_T=0,2`$ subspaces. The comparison of panels (a) and (b) at frequencies $`3\omega /J5`$ indeed suggests that the dynamically relevant finite-$`N`$ excitations are arranged in contrasting patterns. In panel (a) we have an arrangement of points which is typical of a two-parameter continuum. As $`N`$ increases, more points are added and spread roughly evenly along the frequency axis. In panel (b), by contrast, the data points are arranged in branches with an almost $`N`$-independent separation, which is characteristic for branches of bound states. The spin-1 compounds $`Ni(C_2H_8N_2)_2NO_2ClO_4`$ (NENP) and $`Ni(C_3H_{10}N_2)_2N_3(ClO_4)`$ (NINAZ), while not physical realizations of $`H_\theta `$ directly, nevertheless realize situations where the spectrum of collective excitations as probed by inelastic neutron scattering shares major features with the magnons excited from the VBS state.
## 3 Magnons Excited from Ferromagnetic State
We now turn to a more quantitative discussion of the difference between scattering states and bound states in the context of completely integrable situations, namely for the 1D $`s=\frac{1}{2}`$ Heisenberg ferromagnet:
$$H_F=J\underset{n=1}{\overset{N}{}}𝐒_n𝐒_{n+1}.$$
(12)
The ground state is $`(N+1)`$-fold degenerate. We select one of the ground-state eigenvectors, $`|F|\mathrm{}`$, as the physical vacuum for an exploration of collective excitations. As in the VBS case discussed in Sec. 2, the spin fluctuation operator probes transitions between the vacuum state and a branch of 1-magnon states. However, here the magnon dispersion is gapless: $`EE_F=J(1\mathrm{cos}k)`$.
Unlike in the VBS case, the 1-magnon states of $`H_F`$ are located in a separate invariant subspace. Only the 1-magnon states can contribute to the spin fluctuations. In the expression
$$|kS_k^{}|F,S_k^{}\frac{1}{\sqrt{N}}\underset{n=1}{\overset{N}{}}e^{ikn}S_n^{}$$
(13)
for the 1-magnon eigenvectors, the spin fluctuation operator plays the role of a magnon creation operator. That is not to say $`S_k^{}`$ is a true magnon creation operator. Multi-magnon superpositions, i.e. the states $`S_{k_1}^{}\mathrm{}S_{k_r}^{}|F`$, are a redundant set of non-orthogonal and non-stationary states, which are subject to two kinds of interactions:
$``$The kinematical interaction is caused by the restriction on the number of reversed spins at one lattice site.
$``$The dynamical interaction is caused by the off-diagonal part of $`H_F`$ in the basis of multi-magnon superpositions.
This distinction is quite natural in the framework of the Bethe ansatz as we shall see. The Bethe ansatz is an exact method for the calculation of eigenvectors and eigenvalues of completely integrable quantum many-body systems. The Bethe wave function of any eigenstate in the subspace with $`rN/2S_T^z`$ reversed spins relative to the magnon vacuum,
$$|\psi =\underset{1n_1<\mathrm{}<n_rN}{}a(n_1,\mathrm{},n_r)S_{n_1}^{}\mathrm{}S_{n_r}^{}|F.$$
(14)
has coefficients of the form
$`a(n_1,\mathrm{},n_r)={\displaystyle \underset{𝒫S_r}{}}\mathrm{exp}`$ $`\left(i{\displaystyle \underset{j=1}{\overset{r}{}}}k_{𝒫j}n_j+{\displaystyle \frac{i}{2}}{\displaystyle \underset{i<j}{\overset{r}{}}}\theta _{𝒫i𝒫j}\right)`$ (15)
determined by $`r`$ magnon momenta $`k_i`$ and one phase angle $`\theta _{ij}=\theta _{ji}`$ for each magnon pair. The sum $`𝒫S_r`$ is over the permutations of the labels $`\{1,2,\mathrm{},r\}`$. The consistency requirements for the coefficients $`a(n_1,\mathrm{},n_r)`$ inferred from the eigenvalue equation $`H|\psi =E|\psi `$ and the requirements imposed on the same coefficients by translational invariance can be cast in a set of equations for the momenta $`k_i`$ and phase angles $`\theta _{ij}`$:
$$2\mathrm{cot}\frac{\theta _{ij}}{2}=\mathrm{cot}\frac{k_i}{2}\mathrm{cot}\frac{k_j}{2},Nk_i=2\pi \lambda _i+\underset{ji}{\overset{r}{}}\theta _{ij}.$$
(16)
Every solution of these equations is specified by a set of $`r`$ Bethe quantum numbers $`\lambda _i\{1,2,\mathrm{},N1\}`$. Given a solution, the energy and wave number of the state it describes are
$$EE_F=J\underset{j=1}{\overset{r}{}}(1\mathrm{cos}k_j),k=\frac{2\pi }{N}\underset{i=1}{\overset{r}{}}\lambda _i.$$
(17)
In the subspace with $`r=1`$, we thus recover all $`N`$ 1-magnon states, one for each of the allowed values of $`\lambda _1`$. There exist $`N(N+1)/2`$ distinct 2-magnon superpositions of the 1-magnon states thus identified. However, this set of states must be accommodated in the $`r=2`$ subspace, whose dimensionality is only $`N(N1)/2`$. Inspection shows that of the $`N(N+1)/2`$ pairs of Bethe quantum numbers in the allowed range $`0\lambda _1\lambda _2N1`$, $`N`$ pairs do indeed not produce a solution of (16). The missing solutions are a consequence of the kinematical interaction between magnons.
The consequences of the dynamical 2-magnon interaction are illustrated in Fig. 2, which shows the complete $`r=2`$ spectrum $`(EE_F)/J`$ versus $`k`$ for $`k0`$ and $`N=32`$. Also shown in Fig 2 are the (fictitious) 2-magnon superpositions, where the $`k_1,k_2`$ in Eq. (17) are replaced by all combination of 1-magnon wave numbers. There are three classes of states.
The class $`C_1`$ contains $`N`$ states for which one of the two Bethe quantum numbers is zero. This means that one of the magnons has zero wave number. Its effect is a slight rotation of the magnon vacuum, in which the other magnon is as free to propagate as in the original vacuum. There is no dynamical interaction. All states in this class are, effectively, 1-magnon states.
The class $`C_2`$ of states has nonzero Bethe quantum numbers $`\lambda _2\lambda _12`$. There are $`N(N5)/2+3`$ such pairs. All of them yield a solution with real $`k_1,k_2`$, which makes them 2-magnon scattering states. A measure of the dynamical magnon interaction in Fig. 2 is the vertical displacement of any true scattering state $`()`$ from the nearest free-magnon pair (+). As $`N`$ increases, the energy correction diminishes for all class $`C_2`$ states and vanishes in the limit $`N\mathrm{}`$. The 2-magnon scattering states and the free 2-magnon states then form two-parameter continua with identical boundaries $`EE_F=2J[1\pm \mathrm{cos}(k/2)]`$.
The class $`C_3`$ of states has nonzero Bethe quantum numbers which either are equal $`\lambda _2=\lambda _1`$, or differ by unity $`\lambda _2=\lambda _1+1`$. There exist $`2N3`$ such pairs, but only $`N3`$ pairs yield solutions of (16). For the class $`C_3`$ states, the effects of the dynamical magnon interaction are much more prominent, and the interaction energy does not disappear when $`N\mathrm{}`$. In Fig. 2, these states form a branch of 2-magnon bound states with dispersion $`EE_F=\frac{1}{2}J(1\mathrm{cos}k)`$ below the continuum of 2-magnon scattering states.
The bound state character of the class $`C_3`$ states manifests itself in the enhanced probability that the two flipped spins are on neighboring sites of the lattice. This property of the wave function is best captured in the weight distribution $`|a(n_1,n_2)|`$ of basis vectors with flipped spins at sites $`n_1`$ and $`n_2`$. In Fig. 3 we have plotted $`|a(n_1,n_2)|`$ versus $`n_2n_1`$ for a sequence of class $`C_3`$ states between $`k=0`$ and $`k=\pi `$. The distribution is peaked at $`n_2n_1=1`$. Its width is controlled by the imaginary parts of $`k_1,k_2`$ in (15). The smallest width is observed in the bound state at $`k=\pi `$. In this case, all coefficients with $`n_2n_1+1`$ are zero, which implies that the two down spins are tightly bound together and have the largest binding energy.
The width of the distribution $`|a(n_1,n_2)|`$ increases as $`k`$ decreases, and the binding of the two down spins loosens. For finite $`N`$, the Bethe ansatz solutions switch from complex to real when the distribution has acquired a certain width. In scattering states the distribution $`|a(n_1,n_2)|`$ is always broad and tends to oscillate wildly. Some scattering states have a smooth distribution with a maximum for $`n_2=n_1+N/2`$, when the two down spins are farthest apart. The formation of bound states and scattering states of elementary excitations exist in many different contexts. But only in rare cases such as this one can they be investigated on the level of detail presented here.
## 4 Spinons Excited from Spin-Fluid State
Turning our attention to the Heisenberg antiferromagnet, we consider the Hamiltonian $`H_AH_F`$ with $`H_F`$ defined in Eq. (12). All the eigenvectors remain the same, but the energy eigenvalues have the opposite sign. The magnon vacuum $`|F`$ now is at the top of the excitation spectrum. The ground state $`|A`$ of $`H_A`$ is located in the invariant subspace with $`S_T^z=0`$. This subspace also contains the two Néel states $`|𝒩_1|\mathrm{},|𝒩_2|\mathrm{}`$, which are not eigenstates of $`H_A`$. In the framework of the Bethe ansatz, $`|A`$ can be obtained from $`|F`$ by exciting $`r=N/2`$ magnons with momenta $`k_i`$ and (negative) energies $`J(1\mathrm{cos}k_i)`$. The Bethe quantum numbers for this state are $`\{\lambda _i\}_A=\{1,3,5,\mathrm{},N1\}`$.
For reasons of computational convenience we rewrite the Bethe ansatz equations (16) in terms of the variables $`z_i\mathrm{cot}(k_i/2)`$:
$$N\mathrm{arctan}z_i=\pi I_i+\underset{ji}{}\mathrm{arctan}\left(\frac{z_iz_j}{2}\right),i=1,\mathrm{},r.$$
(18)
The associated Bethe quantum numbers $`N/2<I_iN/2`$ are integers for odd $`r`$ and half integers for even $`r`$. The relation between the sets $`\{\lambda _i\}`$ and $`\{I_i\}`$ depends on the configuration of the solution $`\{z_i\}`$ in the complex plane. Given the solution $`\{z_1,\mathrm{},z_r\}`$ of Eqs. (18) for a state specified by $`\{I_1,\mathrm{},I_r\}`$, its energy and wave number are
$$\frac{EE_F}{J}=\underset{i=1}{\overset{r}{}}\frac{2}{1+z_i^2},k=\pi r\frac{2\pi }{N}\underset{i=1}{\overset{r}{}}I_i,$$
(19)
with $`E_F=JN/4`$. For states with real $`\{z_i\}`$, Eqs. (18) can be converted into a convergent iterative process:
$$z_i^{(n+1)}=\mathrm{tan}\left(\frac{\pi }{N}I_i+\frac{1}{N}\underset{ji}{\overset{r}{}}\mathrm{arctan}\left[\frac{z_i^{(n)}z_j^{(n)}}{2}\right]\right)$$
(20)
with starting values $`z_i^{(1)}=\pi I_i/N`$. For the ground state $`|A`$ we have
$$\{I_i\}_A=\{\frac{N}{4}+\frac{1}{2},\frac{N}{4}+\frac{3}{2},\mathrm{},\frac{N}{4}\frac{1}{2}\}.$$
(21)
High-precision solutions $`\{z_i\}`$ can be obtained with little computational effort. The ground-state energy per site for $`N=4096`$, for example, reproduces the exact result $`(E_AE_F)/JN=\mathrm{ln}2`$ of the infinite chain to within 1 part in a million. The distribution of magnon momenta in the ground state $`|A`$ is broad and peaked at $`k_i=\pi `$:
$$\rho _0(k_i)=\left[8\mathrm{sin}^2\frac{k_i}{2}\mathrm{cosh}\left(\frac{\pi }{2}\mathrm{cot}\frac{k_i}{2}\right)\right]^1.$$
(22)
This state, which has obviously a very complicated structure when described in terms of magnons, will now be configured as the physical vacuum of $`H_A`$ for a different kind of elementary particle called the spinon.
A useful way to characterize the new physical vacuum is through the perfectly regular array (21) of Bethe quantum numbers as illustrated in the first row of Fig. 4.
The spectrum of $`H_A`$ can then be generated systematically in terms of the fundamental excitations as characterized by elementary modifications of this vacuum array.
In the subspace with $`S_T^z=1`$, a two-parameter set of states is obtained by removing one magnon from the state $`|A`$. In doing so we eliminate one of the $`N/2`$ Bethe quantum numbers from the set in the first row of Fig. 4 and rearrange the remaining $`I_i`$ in all configurations over the expanded range $`|I_i|\frac{1}{4}N`$. Changing $`S_T^z`$ by one means that the $`I_i`$ switch from half-integers to integers or vice versa. The number of distinct configurations with $`I_{i+1}I_i1`$ is $`N(N+2)/8`$. A generic configuration consists of three clusters with two gaps between them as shown in the second row of Fig. 4. The position of the gaps between the $`I_i`$-clusters determine the momenta $`\overline{k}_1,\overline{k}_2`$ of the two spinons, which, in turn, add up to the wave number of the two-spinon state relative to the wave number of the vacuum: $`Qkk_A=\overline{k}_1+\overline{k}_2`$.
A plot of the 2-spinon energies $`EE_A`$ versus wave number $`kk_A`$ for $`N=16`$ as inferred via (19) from the solutions $`\{z_i\}`$ is shown in Fig. 5 ($``$). The dots, which represent the corresponding data for $`N=256`$, produce a sort of density plot for the 2-spinon continuum which emerges in the limit $`N\mathrm{}`$. The exact lower and upper boundaries of the 2-spinon continuum are
$$\omega _L(Q)=\frac{\pi }{2}J|\mathrm{sin}Q|,\omega _U(Q)=\pi J|\mathrm{sin}\frac{Q}{2}|.$$
(23)
Like magnons, spinons carry a spin in addition to energy and momentum. Unlike magnons, which have spin 1, spinons are spin-$`1/2`$ particles. For even $`N`$, where all eigenstates have integer-valued $`S_T^z`$, spinons occur only in pairs. The spins $`s_1,s_2=\pm 1/2`$ of the two spinons in a 2-spinon eigenstate of $`H_A`$ can be combined in four different ways to form a triplet state or a singlet state.
We have already analyzed the spectrum of the 2-spinon triplet states with $`s_1=s_2=+1/2`$ $`(S_T=1,S_T^z=+1)`$. The 2-spinon triplets with $`s_1=s_2=1/2`$ $`(S_T=1,S_T^z=1)`$ are obtained from these states by a spin-flip transformation. One set of 2-spinon states with $`s_1=s_2`$ are the triplets with $`S_T^z=0`$. They are obtained via Bethe ansatz by a simple modification of the $`\{I_i\}`$ configurations from the kind shown in the second row of Fig. 4 and yields an equal number of (real) solutions. Symmetry requires that all three triplet components $`(S_T^z=0,\pm 1)`$ are degenerate.
The 2-spinon singlet states $`(S_T^z=S_T=0)`$ are characterized by one pair of complex conjugate solutions $`z_1=z_2^{}`$ in addition to the real solutions $`z_3,\mathrm{},z_{N/2}`$. Finding the $`I_i`$-configurations for a particular set of eigenstates with complex solutions and then solving the associated Bethe ansatz equations is, in general, delicate task. The 2-spinon singlets for $`N=16`$ and $`0<Q\pi `$ are shown as open circles in Fig. 5. As $`N`$ grows large, the effect of the complex solutions $`z_1=z_2^{}`$ on the energy relative to that of the real solutions $`z_3,\mathrm{},z_{N/2}`$ diminishes. In the limit $`N\mathrm{}`$, the 2-spinon singlets also form a continuum with boundaries (23). Yet the effect of the complex solutions will remain strong for other quantities including selection rules and transition rates.
The 2-spinon triplets play an important role in the low-temperature spin dynamics of quasi-1D antiferromagnetic compounds such as $`KCuF_3`$, $`Cu(C_6D_5COO)_23D_2O`$, $`Cs_2CuCl_4`$, and $`Cu(C_4H_4N_2(NO_3)_2)`$. They are the elementary excitations of $`H_A`$ which can be directly probed via inelastic neutron scattering. The 2-spinon singlets, in contrast, cannot be excited directly from $`|A`$ by neutrons because of selection rules. The singlet excitations are important nevertheless, but in a different context. Some quasi-1D antiferromagnetic compounds like $`CuGeO_3`$ are susceptible to a spin-Peierls transition, which involves a lattice distortion accompanied by an exchange dimerization. The operator which probes the dimer fluctuations in the ground state of $`H_A`$ couples primarily to the 2-spinon singlets and not at all to the 2-spinon triplets.
## 5 Spinons Excited from Antiferromagnetic State
Although the computational application of the Bethe ansatz yields the exact wave functions of the 2-spinon states for very large systems with little effort, the very structure of the Bethe wave function makes it hard to use this knowledge for the calculation of the transition rates pertaining to any fluctuation operator of interest. Nevertheless, there exist ways to extract useful lineshape information from the Bethe wave function for specific situations. Here we discuss an alternative approach, which exploits the higher symmetry of $`H_A`$ in the limit $`N\mathrm{}`$, described by the quantum group $`U_q(sl_2)`$. For example, the asymptotic degeneracy of the 2-spinon triplets and singlets noted previously is attributable to this higher symmetry.
The algebraic analysis of a completely integrable spin chain employed for the purpose of calculating dynamic structure factors for specific fluctuation operators requires the execution of the following program:
$``$ Span the infinite-dimensional physically relevant Hilbert space in the form of a separable Fock space of multiple spinon excitations.
$``$Generate the eigenvectors of the Hamiltonian in this Fock space by products of spinon creation operators (so-called vertex operators) from the ground state (physical vacuum).
$``$Determine the spectral properties (energy, momentum) of the multiple-spinon states accessible from the ground state by the selected fluctuation operator.
$``$Express the fluctuation operator of interest in terms of vertex operators.
$``$Evaluate matrix elements of products of vertex operators as are needed for the selected fluctuation operator in the spinon eigenbasis.
In the following, we outline how this program was carried out for the calculation of the exact 2-spinon part of one dynamic structure factor for the 1D $`s=\frac{1}{2}`$ $`XXZ`$ model,
$$H_{XXZ}=J\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}(S_n^xS_{n+1}^x+S_n^yS_{n+1}^y+\mathrm{\Delta }S_n^zS_{n+1}^z).$$
(24)
At $`T=0`$ $`H_{XXZ}`$ has a ferromagnetic phase for $`\mathrm{\Delta }1`$, a critical phase (spin-fluid) for $`1\mathrm{\Delta }<1`$, and an antiferromagnetic phase for $`\mathrm{\Delta }<1`$. The algebraic analysis operates in the massive phase stabilized by Néel LRO at $`\mathrm{\Delta }<1`$, but the isotropic limit $`\mathrm{\Delta }1^{}`$, which is equivalent to $`H_A`$, can be performed meaningfully.
Each spinon in the $`m`$-spinon eigenstate $`|\xi _m,ϵ_m;\mathrm{};\xi _1,ϵ_1_j`$ is characterized by a (complex) spectral parameter $`\xi _l`$ of unit length and a spin orientation $`ϵ_l=\pm 1`$. The spectral properties then follow from the rules governing the application of the translation operator and the Hamiltonian:
$`T|\xi _m,ϵ_m;\mathrm{};\xi _1,ϵ_1_j`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{m}{}}}{\displaystyle \frac{1}{\tau (\xi _i)}}|\xi _m,ϵ_m;\mathrm{};\xi _1,ϵ_1_{1j},`$ (25)
$`H|\xi _m,ϵ_m;\mathrm{};\xi _1,ϵ_1_j`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{m}{}}}e(\xi _i)|\xi _m,ϵ_m;\mathrm{};\xi _1,ϵ_1_j,`$ (26)
$`\tau (\xi )`$ $`=`$ $`e^{ip(\xi )}=\xi ^1{\displaystyle \frac{\theta _{q^4}(q\xi ^2)}{\theta _{q^4}(q\xi ^2)}},e(\xi )=J{\displaystyle \frac{1q^2}{4q}}\xi {\displaystyle \frac{d}{d\xi }}\mathrm{log}\tau (\xi ),`$ (27)
$$\theta _x(y)(x;x)(y;x)(xy^1;x),(y;x)\underset{n=0}{\overset{\mathrm{}}{}}(1yx^n).$$
(28)
Here $`q`$ is the deformation parameter of the quantum group $`U_q(sl_2)`$. Its value is determined by the exchange anisotropy, $`\mathrm{\Delta }=(q+q^1)/2`$. The twofold degenerate vacuum is represented by the vectors $`|0_0,|0_1`$, which transform into each other via translation: $`T|0_j=|0_{1j}`$. In the Ising limit $`\mathrm{\Delta }\mathrm{}`$, they become the pure Néel states $`|\mathrm{}\mathrm{},|\mathrm{}\mathrm{}`$.
By eliminating $`\xi `$ from the relations (25) and (26) we obtain the spinon energy-momentum relation
$$e_1(p)=I\sqrt{1k^2\mathrm{cos}^2p},I\frac{JK}{\pi }\mathrm{sinh}\frac{\pi K^{}}{K},$$
(29)
which is independent of the spin orientation. The elliptic integrals $`KK(k),K^{}K(k^{})`$ and their moduli $`k,k^{}\sqrt{1k^2}`$ depend on the anisotropy via $`q=\mathrm{exp}(\pi K^{}/K)`$.
The of 2-spinon spectrum with energies $`E(\xi _1,\xi _2)=e(\xi _1)+e(\xi _2)`$ and momenta $`P(\xi _1,\xi _2)=p(\xi _1)+p(\xi _2)`$ is then a two-parameter set with fourfold spin degeneracy $`ϵ_1,ϵ_2=\pm 1`$. In the isotropic limit, they are the 2-spinon triplets and singlets discussed in Sec. 4. In a finite system, the singlet-triplet degeneracy is removed (see Fig. 5). Exchange anisotropy splits up the triplet levels as well.
In the $`(Q,\omega )`$ plane, each set of 2-spinon excitations forms a continuum (see Fig. 6) of two sheets $`𝒞_\pm `$ with boundaries
$$\omega _0(Q)=\frac{2I}{1+\kappa }\mathrm{sin}Q,\omega _\pm (Q)=\frac{2I}{1+\kappa }\sqrt{1+\kappa ^2\pm 2\kappa \mathrm{cos}Q},$$
(30)
where $`\kappa \mathrm{cos}Q_\kappa =(1k^{})/(1+k^{})`$ is a convenient anisotropy parameter. The evaluation of the dynamic structure factor of any fluctuation operator $`F_Q^A`$, Eq. (9), requires that we know exact matrix elements $`{}_{j}{}^{}0|A_l|\xi _m,ϵ_m,\mathrm{},\xi _1,ϵ_1_{j^{}}^{}`$ for the associated local operators. In principle, they can be calculated exactly if we are able to express $`A_l`$ in terms of vertex operators. In practice, the exact results available thus far are limited to the dynamic spin structure factors
$$S_{\mu \mu }(Q,\omega )=\underset{n=\mathrm{}}{\overset{+\mathrm{}}{}}_{\mathrm{}}^{\mathrm{}}𝑑te^{i(\omega t+Qn)}S_n^\mu (t)S_0^\mu ,\mu =x,z.$$
(31)
In the $`m`$-spinon expansion $`S_{zz}(Q,\omega )=S_{zz}^{(0)}(Q,\omega )+S_{zz}^{(2)}(Q,\omega )+\mathrm{}`$ we know the matrix elements that determine the leading term,
$$2{}_{j}{}^{}0|S_n^z|0_{j}^{}=\frac{(q^2;q^2)^2}{(q^2;q^2)^2}(1)^{n+j}=\frac{2\pi }{K}\sqrt{k^{}}=2\overline{m}_z,$$
(32)
yielding $`S_{zz}^{(0)}(Q,\omega )=4\pi ^2\delta (\omega )\delta (Q\pi )\overline{m}_z^2,`$ where $`\overline{m}_z`$ is the staggered magnetization. No exact transition rates have been evaluated for the 2-spinon part $`S_{zz}^{(2)}(Q,\omega )`$ or any higher order term in the $`m`$-spinon expansion.
The corresponding $`m`$-spinon expansion of $`S_{xx}(Q,\omega )`$ starts with $`m=2`$, because the operator $`S_n^x`$ does not connect the vacuum sector with itself. All non-vanishing matrix elements which are needed for $`S_{xx}^{(2)}(Q,\omega )`$ can be expressed, as it turns out, by a single function $`X^j(\xi _2,\xi _1){}_{j}{}^{}0|\sigma _0^+|\xi _2,;\xi _1,_{j}^{}`$, which was determined by Jimbo and Miwa:
$$X^j(\xi _2,\xi _1)=\varrho ^2\frac{(q^4;q^4)^2}{(q^2;q^2)^3}\frac{(q\xi _1\xi _2)^{1j}\xi _1\gamma (\xi _1^2/\xi _2^2)\theta _{q^8}(\xi _1^2\xi _2^2q^{4j})}{\theta _{q^4}(\xi _1^2q^3)\theta _{q^4}(\xi _2^2q^3)},$$
(33)
where $`(x;y;z)_{n,m=0}^{\mathrm{}}\left(1xy^nz^m\right)`$ and
$`\gamma (\xi )`$ $``$ $`{\displaystyle \frac{(q^4\xi ;q^4;q^4)(\xi ^1;q^4;q^4)}{(q^6\xi ;q^4;q^4)(q^2\xi ^1;q^4;q^4)}},\varrho (q^2;q^2)^2{\displaystyle \frac{(q^4;q^4;q^4)}{(q^6;q^4;q^4)}}.`$ (34)
These ingredients yield the following expression to be evaluated:
$`S_{xx}^{(2)}(Q,\omega )`$ $`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle \frac{d\xi _1}{2i\xi _1}\frac{d\xi _2}{2i\xi _2}\delta [\omega E(\xi _1,\xi _2)]}`$ (35)
$`\times \{\delta [Q+P(\xi _1,\xi _2)]|X^0(\xi _2,\xi _1)+X^1(\xi _2,\xi _1)|^2`$
$`+\delta [Q\pi +P(\xi _1,\xi _2)]|X^0(\xi _2,\xi _1)X^1(\xi _2,\xi _1)|^2\}.`$
A compact rendition of the exact result reads
$$S_{xx}^{(2)}(Q,\omega )=\frac{\omega _0}{8I\omega }\left[1+\sqrt{\frac{\omega ^2\kappa ^2\omega _0^2}{\omega ^2\omega _0^2}}\right]\underset{c=\pm }{}\frac{\vartheta _A^2(\beta _{}^c)}{\vartheta _d^2(\beta _{}^c)}\frac{|\mathrm{tan}(Q/2)|^c}{W_c},$$
(36)
where
$$W_\pm =\sqrt{\frac{\omega _0^4}{\omega ^4}\kappa ^2\left(\frac{T}{\omega ^2}\pm \mathrm{cos}Q\right)^2},T=\sqrt{\omega ^2\kappa ^2\omega _0^2}\sqrt{\omega ^2\omega _0^2},$$
(37)
$$\beta _{}^c(Q,\omega )=\frac{1+\kappa }{2}F[\mathrm{arcsin}\left(\frac{2I\omega W_c}{\kappa (1+\kappa )\omega _0^2}\right),\kappa ],$$
(38)
$$\vartheta _A^2(\beta )=\mathrm{exp}\left(\underset{l=1}{\overset{\mathrm{}}{}}\frac{e^{\gamma l}}{l}\frac{\mathrm{cosh}(2\gamma l)\mathrm{cos}(t\gamma l)1}{\mathrm{sinh}(2l\gamma )\mathrm{cosh}(\gamma l)}\right),$$
(39)
$`\gamma =\pi K^{}/K`$, $`t2\beta /K^{}`$, and $`\vartheta _d(x)`$ is a Neville theta function. This function is plotted in Fig. 7 for $`(Q,\omega )𝒞_+`$. $`S_{xx}^{(2)}(Q,\omega )`$ has a square-root divergence at the portion $`\omega _0(Q)`$ of the spectral threshold. Along the portion $`\omega _+(Q)`$ of the lower boundary and along the entire upper boundary of $`𝒞_+`$, $`S_{xx}^{(2)}(Q,\omega )`$ has a square-root cusp. The line shapes for $`(Q,\omega )𝒞_{}`$ are inferred from the symmetry property $`S_{xx}^{(2)}(\pi Q,\omega )=S_{xx}^{(2)}(Q,\omega )`$.
The isotropic limit $`\mathrm{\Delta }1^{}`$ is delicate because of its singular nature. As the LRO in the spinon vacuum vanishes, the size of the Brillouin zone changes from $`(\pi /2,+\pi /2)`$ to $`(\pi ,+\pi )`$. A practical consequence of this phase transition is that we switch our perspective from considering both sheets $`𝒞_\pm `$ of 2-spinon excitations over the range $`(\pi /2,+\pi /2)`$ to considering only the sheet $`𝒞_+`$, now with boundaries (23), over the extended range $`(\pi ,+\pi )`$. With these subtleties taken into account, Eq. (36) reduces to
$$S_{xx}^{(2)}(Q,\omega )=\frac{1}{2}[\omega _U^2(Q)\omega ^2]^{1/2}e^{I(t)},$$
(40)
$$I(t)=_0^{\mathrm{}}𝑑x\frac{\mathrm{cosh}2x\mathrm{cos}xt1}{x\mathrm{sinh}2x\mathrm{cosh}x}e^x,\frac{\pi t}{4}=\mathrm{cosh}^1\sqrt{\frac{\omega _U^2(Q)\omega _L^2(Q)}{\omega ^2\omega _L^2(Q)}}.$$
(41)
This function, which is plotted in Fig. 8, has a square-root cusp at $`\omega _U(Q)`$ and a square-root divergence with logarithmic corrections at $`\omega _L(Q)`$.
## 6 Solitons Excited from Néel State
Near the Ising limit ($`\mathrm{\Delta }\mathrm{}`$), the exact result (36) can be expanded in powers of the anisotropy parameter $`\kappa `$. To leading order, we obtain
$$S_{xx}^{(2)}(Q,\omega )\frac{1}{2\mathrm{cos}^2Q}\sqrt{\mathrm{cos}^2Q\left(\frac{\mathrm{\Omega }}{\kappa }\right)^2\mathrm{\Theta }(\kappa |\mathrm{cos}Q||\mathrm{\Omega }|)}$$
(42)
with $`\mathrm{\Omega }\omega /2I1`$, which is identical (in that order) to the result obtained by Ishimura and Shiba from a first-order perturbation calculation about the Ising limit. In their calculation, the two vacuum states are approximated by the pure Néel states (see Fig. 9). The spinons are kink solitons (domain walls), which produce two adjacent up or down spins in an otherwise unperturbed Néel configuration. The only states which contribute to $`S_{xx}(Q,\omega )`$ in leading order of the perturbation calculation are states which contain two solitons with equal spin orientation. The states with solitons of opposite spin orientation dominate the spectrum of $`S_{zz}(Q,\omega )`$ except for the central peak at $`\omega =0,Q=\pi `$, which involves a transition between the two vacuum states.
The 2-soliton spectrum with boundaries $`|\mathrm{\Omega }|\kappa \mathrm{cos}Q`$, is akin to the generic 2-spinon spectrum shown in Fig. 6, but it does not capture any of the features that characterize the regime $`Q_\kappa Q\pi Q_\kappa `$ around the Brillouin zone boundary. The limitations of the perturbation results are much less severe near the center of the Brillouin zone. This illustrated in Fig. 10, which compares the line shapes of Eq. (36) and (42) at $`Q=\pi `$. Here a fairly strong departure from the Ising limit is required before a significant difference between the two results is produced.
Two intensively studied physical realizations of $`H_{XXZ}`$ at $`\mathrm{\Delta }<1`$ are $`\mathrm{CsCoCl}_3`$ and $`\mathrm{CsCoBr}_3`$. Spectroscopic data which probe both the quantum and thermal fluctuations of those materials are available from several neutron scattering experiments.
## Acknowledgments
Financial support from the URI Research Office (for G.M.) and from the DFG Schwerpunkt Kollektive Quantenzustände in elektronischen 1D Übergangsmetallverbindungen (for M.K.) is gratefully acknowledged.
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# On mesogranulation, network formation and supergranulation
## 1 Introduction
In the traditional view of solar convection seen at the sun’s surface, three scales play the main roles: granulation (1 Mm) which shows up as an intensity pattern most probably first seen by Herschel in 1801 BLD84, supergranulation ($`1530`$ Mm) which appears in (but not only) dopplergrams of the full disk of the sun as a pattern of horizontal velocities and which was first noticed by HartHart56b and confirmed by Leighton et al. LNS62, and mesogranulation ($`310`$ Mm) observed by November et al. NTGS81 on Doppler measurements of vertical velocities.
While the dynamics of granulation is rather well understood, its scale being controlled by the balance of radiative diffusion of heat and convection, the origin of the two other scales remains largely mysterious. The ionization of helium was often invoked to explain these two scales since the first and second ionizations of this atom occur at depths similar to the mesogranular and supergranular scale respectively. However, this geometrical explanation is based on a “laminar view” of solar convection which is, on the contrary, very strongly turbulent.
The aim of the present paper is to investigate the dynamics of scales larger than that of the granulation, using the horizontal flows given by granular motions determined by the new algorithms described in Roudier et al. (1999). Briefly, these algorithms, Local Correlation Tracking on binarized images (LCT<sub>bin</sub>) or Coherent Structure Tracking (CST), allow us to increase noticeably the spatial and temporal resolutions of the surface velocity fields; typically, we can bring the spatial grid size down to 0″7 and the time step down to 5 mn.
Thus, we first concentrate on mesogranulation (Sect. 2) and show the major role played at this scale by strong positive divergences (SPDs) and by averaging procedures. It turns out, indeed, that what has been described in previous work as mesogranulation results from a combination of a physical phenomenon (SPDs among which are found exploding granules) and a data processing effect applied to a turbulent flow (averaging). Since each author had his own technique for averaging data, results have been rather confusing and no clear-cut description of mesogranulation has emerged. We show here that when averaging is properly controlled, no quasi-steady flow can be detected in the mesoscale range.
We then proceed (Sect. 3) with the investigation of the transport properties of the mesoscale flows and show that the supergranulation scale appears when the positions of concentrations of corks are compared to the positions of network bright points. We conclude the paper with a discussion of a model which seems to explain many of the observations and the interactions of the three scales.
## 2 Mesogranulation
Since its discovery by November et al. NTGS81, mesogranulation has been sought using many different techniques (Doppler measurements, intensity variations, horizontal velocity fields and their divergence) with the idea that one could exhibit a quasi-steady cellular motion as clearly as granulation. However, the reports of observations aimed at pointing out this new feature of solar convection have never been clear-cut and always difficult to compare to each other.
In Table LABEL:meso\_table we summarize all the previous studies on mesogranulation. A common result of these studies is that they always find some features in the mesogranulation range of scales, i.e. within length scales between 4″ and 12″ (or 3 Mm and 10 Mm) and on time scales between 30 mn and 6 h. The features are patterns either in intensity, or in radial velocity, or in horizontal divergence, for example. The picture left by mesogranulation observations is therefore fuzzy: neither its characteristic size or time is well established and vary from one author to another.
### 2.1 Spectra and correlations
In order to better understand the situation, it is useful to consider the physics which leads to the above mentioned observations. This is obviously turbulent convection: motions seen at the sun’s surface result from the superposition of a large number of scales constituting the turbulent spectrum. The relevant questions are therefore: why should a scale like mesogranulation single out among other scales? And if it does, how would we recognize and characterize it?
Concerning granulation, the answers are known: it is the balance between radiative diffusion and advection which determines the size of granules and it emerges from other scales as the one with the highest contrast in intensity and with the largest fluctuations of velocity. The spectrum of turbulent kinetic energy shows a maximum around this scale.
At mesoscale, we do not know which mechanism could distinguish any particular scale; we may, however, try to recognize a scale which stands out and therefore search for a peak or a break of the slope in the kinetic energy spectrum. Such a feature in the spectrum is indeed the signature that some physical mechanism is injecting energy at a specific scale.
In previous work EMRN93, spectra of kinetic energy have not shown any characteristic feature at scales larger than granulation. In fact, large-scale features appear in turbulent flows with time averages. Indeed, let us suppose that the kinetic energy spectrum varies as $`E(k)k^\alpha `$ in the mesoscale range; using dimensional arguments, it turns out that the typical lifetime of structures of wavenumber $`k`$ is $`\tau _kk^{(\alpha 3)/2}`$ which means that the lifetime of turbulent structures increases with their size since $`\alpha <3`$<sup>1</sup><sup>1</sup>1The case $`\alpha =3`$ corresponds to a two-dimensional turbulence; the turnover time of eddies is then fixed by the background vorticity. In three-dimensional turbulence, scales outside the dissipative range are such that $`\alpha 5/3`$.. In other words, long time-averages show large-scale features and the longer the average the larger the scale. This point is clearly illustrated in Fig. 1 where we have computed kinetic energy spectra of the horizontal flow derived from granule tracking in Pic du Midi data RRMV99. These spectra $`E_\tau (k)`$ are defined as follows: The mean kinetic energy (defined by a time window of length $`\tau `$) at one point of the field reads:
$$\frac{1}{2}v^2_\tau =_0^{\mathrm{}}E_\tau (k)𝑑k$$
where $`E_\tau (k)`$ is related to the Fourier transform of the velocity components $`\stackrel{~}{v}_x`$, $`\stackrel{~}{v}_y`$ by
$$E_\tau (k)=_0^{2\pi }\left(|\stackrel{~}{v}_x|^2+|\stackrel{~}{v}_y|^2\right)k𝑑\theta ,$$
where $`k_x=k\mathrm{cos}\theta `$ and $`k_y=k\mathrm{sin}\theta `$ are the components of the wave vector. In this figure we clearly see the build-up of large-scales and the disappearance of small scales when the time averaging window is made longer. The average of the spectra of the velocity fields determined by LCT<sub>bin</sub> and a 15mn-window clearly shows a peak at a scale of $``$5″ (3500 km); this peak is the signature of the most energetic horizontal flows.
To further emphasize the turbulent nature of mesoscale features, we compute the correlation between successive time-averaged velocity fields as a function of time; we compute
$$C_x=\frac{\overline{\left(v_x_\tau \overline{v_x_\tau }\right)(t)\left(v_x_\tau \overline{v_x_\tau }\right)(t+n\tau )}}{\sqrt{\overline{\left(v_x_\tau \overline{v_x_\tau }\right)^2(t)}\overline{\left(v_x_\tau \overline{v_x_\tau }\right)^2(t+n\tau )}}}$$
as a function of $`n`$ ($`t`$ is arbitrary); in this expression overbars indicate spatial averages and $`_\tau `$ stands for a time average of length $`\tau `$. Results are plotted in Fig. 2. They show that the autocorrelation of the mean velocity fields is halved after one time step. This again emphasizes the role played by the time window which selects a spatial structure whose lifetime is precisely of the order of the time window, thus showing that no quasi-steady flow exists on a time scale longer than 15 mn.
### 2.2 Pitfalls of data processing
Turbulence, however, is not the only difficulty of this problem; data processing may also interfere and contribute to blur the results.
As a first instance, let us compute the kinetic energy spectra of horizontal flows, but using a larger spatial window (than in Fig. 1) for the determination of the velocity field. The result is shown in Fig. 4 for a 2″-resolving window. Very clearly the peak is now shifted to 12″ (8.7 Mm). This emphasizes the high sensitivity of the velocity field patterns (scales) to the choice of the sampling window.
Finally, we would like to mention another effect of data processing which can unduly extend the lifetime of mesoscale features. This is the use of sliding time windows. Indeed, as illustrated in Fig. 3 such windows can transform two independent time-evolving structures into a single moving structure. This is what happened when Muller et al. (1992) described a three-hour-living mesogranule which was used to show the supergranular flow. In fact, as shown by Fig. 5, an independent time sampling shows that no coherent structure lasts such a long time but new structures emerge after $``$30 mn.
We therefore interpret the results of previous work on mesogranulation as a consequence of uncontrolled averaging procedures. This explains the large variability of results which have been published on mesogranulation: they are all dependent on the way authors have combined their averages.
### 2.3 Strong positive divergences (SPDs)
The foregoing results, however, leave unanswered the question of the origin of kinetic energy in the mesoscale range; we may indeed wonder which flow patterns are contributing to the spectral peak at $`5`$″ in Fig. 1. A plot of the velocity field along with its divergence (Fig. 7) shows that the strong horizontal flows are generally associated with strong positive divergences (SPDs) among which the exploding granules are the most energetic. This is well illustrated by the two time sequences in Fig. 7 and Fig. 7. From Fig. 7, it is quite clear that the divergence field is highly variable, showing patterns which may extend up to 10″ (7.3 Mm). A histogram of the lifetime of these patterns (Fig. 8) gives a mean lifetime of 15 mn which is short.
Ending this section, we are therefore lead to the conclusion that no specific scale exists in the mesogranulation range except the scale of horizontal flows featured by SPDs. Hence we confirm the conclusion of Straus and Bonaccini SB97 that no gap in the kinematic energy spectrum separates granulation from mesogranulation which therefore must be considered just as the large-scale extension of granulation.
## 3 Supergranulation
### 3.1 Mean flows and SPDs
We may now wish to know whether supergranulation plays some part in the motion of granules. A first look at the spectrum of the 3-hour-average velocity field (Fig. 1) shows some energy at $`30`$″ (22 Mm). However, this spectral peak is only suggestive since the field of view is just 47.6″. Turning to real space (as opposed to spectral space), a plot of the velocity field (Fig. 9) shows that indeed a large-scale velocity field may exist. The rms velocity of the 3-hour-average field is $`230`$ m/s. On the unfiltered field we clearly see the imprint of SPDs with their typical size of 5″; this indicates that the origin of the mean flow may be found in the cumulative effects of SPDs. Furthermore, if this mean field is identified with supergranulation, we have at hand its origin: correlated SPDs.
### 3.2 Network formation from measured velocity fields
One way to confirm the view that a three-hour time average of horizontal flows shows the supergranulation scale, is to focus on its transport properties. This technique was already used in the past by Simon and Weiss SW89. It consists of integrating the trajectories of floating corks, initially uniformly distributed, and characterizing their spatial distribution after some time. Simon and Weiss used the SOUP data sequence, which lasted 28 mn, to determine a kinematic model of the horizontal velocity field. Then, they integrated the cork trajectories during a time very much longer (up to 16 h) than the data sequence. Our method is much closer to the data: with the techniques described in Roudier et al. (1999), we determine the velocity field evolution of the Pic du Midi data set with a high spatial and temporal resolution (0″7 and 5 mn). We then use this field to integrate the cork trajectories during the three hours of recorded data.
As shown in Fig. 10, corks are expelled from certain regions, whose size varies between 5 Mm and 15 Mm and again SPDs seem to play a major role. We cannot say that the cells thus formed are “supergranules” since they would continue to evolve if our data set were longer. If the concentration of corks is plotted, however, the supergranulation scale already appears. The clearest evidence of this is given by the remarkable coincidence between the regions with a high density of corks and the positions of network bright points (Fig. 11). This result should be compared with the one of Brandt et al. BRST94 who found a similarity between cork distribution and the Ca-network.
Referring back to the mean divergence field (Fig. 9), we also see that the regions with negative divergence are the ones where corks and network bright points tend to concentrate.
Finally, we also used corks to estimate a horizontal diffusivity $`D=<r^2>/4t`$ following Berger et al. BLST98. We found values between 50 km<sup>2</sup>/s and 100 km<sup>2</sup>/s but, as this diffusivity did not reach an asymptotic value at the end of our time-sequence, we think that the aforementioned values only show the amplitude of the transport but nothing about its physical origin (advection, diffusion, abnormal or turbulent diffusion).
### 3.3 Network formation from simulated velocity fields
In order to show the leading role of SPDs, we extracted from the data their positions in space and time as well as their mean radii. We used them to determine the amplitudes ($`V_n`$) of the model flow used by Simon and Weiss (1989) and Simon et al. STW91, namely
$$v=\underset{n}{}V_n\frac{|rr_n|}{R}e^{|rr_n|^2/R^2}e_n$$
where $`r_n`$ designates the position of SPDs, $`R`$ their range of influence, which lies between 0.5″ and 3″ with a mean value of 2.3″, and $`e_n`$ the direction of their flow.
Using such a model flow, we have repeated the calculation for the advection of corks and found a distribution very similar to the one obtained when using the full data (see Fig. 10).
It therefore turns out that SPDs give the main contribution to the horizontal transport of passive scalars.
## 4 Discussion
### 4.1 Mesoscale flows
The analysis of flow motions at mesoscale that we presented in Sect. 3 has shown that the scale of flow patterns is controlled by the size of the averaging time window. We interpret this result as evidence of the turbulent nature of the kinetic energy spectrum in this range of scales. We therefore refute the idea that some quasi-steady flow exists at mesoscale and show that previous identifications of mesogranules advected by a supergranular flow, as shown by Muller et al. (1992) for instance, are an artefact of the averaging procedure. We identify the only mechanism able to structure the horizontal flow field as being SPDs among which we find exploding granules. However, the scale controlled by SPDs is rather centered around 5″ (3.5 Mm) which is small compared to the scale generally attributed to mesogranulation (5-10 Mm); we conjecture that these latter scales are in fact built up by the nonlinear interactions of the former ones and form the (turbulent) continuation of the granular motion spectrum. Using two-dimensional simulations Ploner Ploner98 also concludes that mesoscale flows can result from nonlinear interactions of granular flows.
We have shown subsequently that the use of Pic du Midi velocity fields in computing trajectories of floating corks, shows that at the end of the time integration, which we take equal to the length of the data set, corks concentrate in regions which coincide with those occupied by network bright points. It thus turns out that in a rather short time the supergranulation scale emerges from the flow delineated by granule motions.
These results show that the traditional view which assumes supergranulation as a quasi-steady horizontal flow driven by the ionization of helium is certainly too simple. Obviously the build-up of the supergranulation scale as far as magnetic flux tubes are concerned, results from a process of turbulent transport in which SPDs play an important part.
### 4.2 A model for supergranulation
In order to explain the above mentioned observations and some others, we now present a model which seems to square with all the known constraints of supergranulation.
Let us consider the surface convection of the sun as a set of granules. Assume that each granule interacts nonlinearly, mainly with its nearest neighbours. The set of granules behave like a set of nonlinearly coupled oscillators. A general behaviour of such a system is that its energy can be focused on one or very few oscillators whose motion is then very much enhanced DP93. Such oscillators collecting the kinetic energy of the others would appear, in the solar context, as SPDs or even exploders. This may also be seen as a manifestation of intermittency of the turbulent solar flow. Now it has been noticed that such exploding granules appear generally in neighbouring places (see the mean divergences) which means that some temporal and spatial correlations exist between them, namely that a large-scale flow organizes their appearance (or vice-versa). This large-scale flow is obviously supergranulation; its origin should be found in a large-scale instability of the flow expressed by the set of granules.
Recent theoretical work GVF94; SSSF89 on turbulent flows has indeed shown that a large-scale perturbation of some given steady spatially periodic flow can be unstable in some bandwidth of wavenumbers. Typically two kinds of situations may occur: the original small-scale flow is invariant with respect to parity or not. If it is not, large-scale instabilities occur through an AKA<sup>2</sup><sup>2</sup>2Anisotropic Kinematic Alpha effect: it is the Reynolds stress dependence with respect to the mean velocity field; it occurs when the turbulence has helicity, i.e. lacks parity invariance. effect which is the equivalent of the $`\alpha `$-effect of turbulent MHD flows; if it is parity invariant, which is the case of solar granulation (flows are hardly helical), then instability appears through a negative eddy viscosity; a range of large-scale modes is then destabilized and some large-scale coherent flow starts. However, this new flow is rapidly hindered by the motions it induces: as the Reynolds stress distribution is modified, usually the turbulent viscosity comes back to positive values; in our case we are considering turbulent convection and any flow increasing the convective heat flux will be slowed by the overcooled material. A similar scenario was also envisaged by Krishan Krish91 using arguments based on inverse cascade of two-dimensional turbulence; indeed, large-scale instabilities are a way of realizing an inverse cascade. However, the mechanism invoked by Krishan is rather close to the AKA-effect which requires helical motions at supergranular scales; such motions do not seem to be observed. On the other hand, higher in the atmosphere, the stratification stabilizes the fluid and may produce two-dimensional motions which can ease energy transfer to large scales.
We therefore understand why large-scale (with respect to granulation scale) flows like supergranulation will appear. We also understand why supergranulation does not realize a perfect pavement of the whole sun as does granulation. The development of the large-scale instability is a stochastic process which results from a given arrangement of the granular flow. It has not the constraint, as any convective flow, to transport heat.
We also understand why the main downflows of a supergranule are not at its boundary as revealed by the MDI instrument on board of SOHO (see http://sohowww.nascom.nasa.gov/gallery/MDI/ mdi009.html,mdi009.ps or Zahn 1999). They are in fact associated with the most energetic granules i.e. the exploders, and therefore occur mainly in the bulk of the supergranule. The downflows on the boundaries should not be very intense since the supergranular flow must remain weak compared to the granular flows.
Hence, this model explains an interesting number of observations. Essentially, it states that supergranular flows are surface flows which do not result from a large-scale thermal instability.
## 5 Conclusion
To conclude this paper we would like to emphasize the point that the picture of surface solar convection which has been popular in the literature and where supergranulation advects mesogranulation which in turn advects granulation is too simple and misleading. It is obviously too simple because it tentatively describes turbulent convection with three scales instead of a continuous spectrum of scales. It is misleading as it reduces the nonlinear interaction between scales to a simple advection as for instance the kinematic models of Simon et al. 1991.
We have shown here that no quasi-steady flow could be identified at the mesogranulation scale and that after a three hour averaging the mean flow shows a component at the supergranulation scale while it keeps a small-scale (5″) component. This latter component seems to be the source of the former large-scale flow.
Our results therefore suggest a scenario where the large-scale supergranular flow is generated directly by the granular flow through a large-scale instability which fixes the scale, in space and time, of supergranulation. We thus conjecture that nonlinear interaction between flows at the granulation scale, in other words Reynolds stresses, are sufficient to drive flows at the supergranulation scale and that the energy released by the recombination of ionized helium plays no part. This scenario needs now to be tested for its various implications, theoretical as well as observational.
###### Acknowledgements.
We would like to thank very much Peter Brandt for very fruitful and helpful comments on the manuscript. We also appreciated the discussions and comments of our colleagues Keith Aldridge, Gérard Coupinot, Richard Muller and Sylvie Roques. We are very grateful to Richard Muller for providing us with the three-hour sequence of Pic du Midi. Special thanks are due to the Pic du Midi Observatory staff for their technical assistance.
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# A generalization of Hagopian’s theorem and exponents
## 1. Introduction
A space $`X`$ is a *continuum* if it is compact, connected and metrizable, and $`X`$ is *homogeneous* if given any $`x,yX`$ there is a homeomorphism $`h:XX`$ with $`h\left(x\right)=y`$. We show that if $`X`$ is a homogeneous continuum which (for some countable $`\kappa 1`$) admits a fiber bundle projection $`p:X𝐓^\kappa =\left(𝐑/𝐙\right)^\kappa `$ with totally disconnected fibers, then $`X`$ also admits a compatible abelian topological group structure. This generalizes a weakened version of the following theorem of Hagopian: If every subcontinuum of the homogeneous continuum $`X`$ is an arc, then $`X`$ is homeomorphic to a solenoid (which includes the circle as a possibility) \[H\]. For $`\kappa =1`$, the theorem currently under consideration follows from Hagopian’s theorem and applies to all compact one–dimensional minimal sets of flows (see \[AM\]). In higher dimensions, spaces admitting such a fiber bundle structure arise naturally in two settings: as the minimal sets of foliations and as the limit sets of discrete dynamical systems (see \[W\]).
A good example of an $`S^1`$ fiber bundle which is not homogeneous is the minimal set of a Denjoy flow on a torus: the path components are not “evenly spaced,” and so there is no Effros homeomorphism for sufficiently small numbers. If one follows path components far enough along, eventually they will spread apart. Our result shows that this situation persists in higher dimensions.
We then generalize the exponent group introduced in \[Cl\] to higher dimensions. In what follows we use the terminology and results of \[S\], and $`\pi ^\kappa :𝐑^\kappa 𝐓^\kappa `$ denotes the standard fibration with unique path lifting $`t_i_{i=1}^\kappa t_i\text{mod }1_{i=1}^\kappa `$ and $`d`$ is a metric for $`X.`$
## 2. Homogeneity
###### Theorem 2.1.
If $`X`$ is a homogeneous continuum and if $`X`$ admits a fiber bundle projection $`p:X𝐓^\kappa `$ with totally disconnected fibers, then $`X`$ admits an abelian topological group structure.
Proof: Suppose $`p:X𝐓^\kappa `$ is a fiber bundle projection with totally disconnected fiber $`F`$ for some $`\kappa >1`$. Then $`p`$ is a fibration with unique path lifting (see \[S\] 2.2.5, 2.7.14) and so the natural group action
$$\alpha :𝐑^\kappa \times 𝐓^\kappa 𝐓^\kappa $$
given by
$$\alpha (s,x)=\pi ^\kappa \left(s\right)+x$$
lifts uniquely to an action $`\stackrel{~}{\alpha }`$ on $`X`$
$$\begin{array}{ccc}𝐑^\kappa \times X\hfill & \stackrel{\stackrel{~}{\alpha }}{}\hfill & X\hfill \\ id\times p\hfill & & p\hfill \\ 𝐑^\kappa \times 𝐓^\kappa \hfill & \stackrel{\alpha }{}\hfill & 𝐓^\kappa \hfill \end{array}.$$
For, given any $`(s,x)𝐑^\kappa \times X`$, we have the interval $`I=[ts,\left(1t\right)s]`$ ($`t[0,1]`$) and the following commutative diagram
$$\begin{array}{ccccc}\left\{\mathrm{𝟎}\right\}\times X\hfill & & \stackrel{\mathrm{`}\mathrm{`}id\mathrm{"}}{}\hfill & \stackrel{}{}\hfill & X\hfill \\ \hfill & & & & p\hfill \\ I\times X\hfill & \stackrel{id\times p}{}\hfill & I\times 𝐓^\kappa \hfill & \stackrel{\alpha }{}\hfill & 𝐓^\kappa \hfill \end{array},$$
and so the fibration property allows us to define $`\stackrel{~}{\alpha }I\times X`$ as needed. To see that $`\stackrel{~}{\alpha }`$ is continuous, we can replace $`I`$ in the diagram with increasing open connected subsets of $`𝐑^\kappa `$ whose union is all of $`𝐑^\kappa `$.
Then for a given point $`eX`$,
$$\pi ^\kappa =p\stackrel{~}{\alpha }𝐑^\kappa \times \left\{e\right\}$$
is a fibration and $`p`$ has unique path lifting, and so it follows from lemma 2.1 of \[C\] that $`\stackrel{~}{\alpha }𝐑^\kappa \times \left\{e\right\}`$ is a fibration. Hence, the $`\stackrel{~}{\alpha }`$ trajectory of $`e`$ is the path component of $`e`$ in $`X`$, \[S\] 2.3.1. We now proceed to show that the path components of $`X`$ are dense in $`X`$.
The proof is similar to the proof of a related fact for matchbox manifolds given in \[AHO\]. Each point of $`X`$ is contained in an open set homeomorphic to $`V\times F`$ for some open subset $`V`$ of $`𝐓^\kappa `$. Since $`F`$ is compact and totally disconnected, there is a basis $``$ for $`X`$ consisting of sets homeomorphic to $`U\times Z`$ with $`U`$ open and connected in $`𝐓^\kappa `$ and $`Z`$ a closed and open subset of $`F.`$ Let $`BU\times Z`$ be any element of $``$. We wish to show that $`PC\left(B\right)`$ (the union of all the path components of $`X`$ which meet $`B`$) is both open and closed in $`X`$. For $`t𝐑^\kappa \left\{\mathrm{𝟎}\right\}`$, $`\stackrel{~}{\alpha }`$ induces the non-singular flow
$$\stackrel{~}{\alpha }\left[t\right]:𝐑\times XX;\text{ }\stackrel{~}{\alpha }\left[t\right](r,x)=\stackrel{~}{\alpha }(tr,x)\text{ (and similarly for }\alpha \text{).}$$
For $`xPC\left(B\right)`$ we may then choose $`s𝐑^\kappa `$ so that $`\stackrel{~}{\alpha }(s,x)=yB`$ and so that the corresponding linear flow $`\alpha \left[s\right]`$ (and hence $`\stackrel{~}{\alpha }\left[s\right]`$) is aperiodic, which is possible since $`\kappa >1`$ and since the collection of $`t𝐑^\kappa \left\{\mathrm{𝟎}\right\}`$ for which $`\alpha \left[t\right]`$ is aperiodic is dense. By a theorem of Bebutov (see \[NS\] V.2.15) we may construct a flow box joining $`y`$ and $`x`$ which contains both points in its interior. Hence, all points in a neighborhood of $`x`$ are also in $`PC\left(B\right)`$, demonstrating that it is open.
Suppose then that $`x\overline{PC\left(B\right)}`$ with $`\left(x_n\right)_nx`$ and $`\left\{x_n\right\}PC\left(B\right).`$ The path component of $`x_n`$ meets $`B`$ in a point $`b_n(u_n,z_n)`$ and since $`U`$ is path connected, we may assume that $`u_n=u`$ for some $`uU`$ and for all $`n`$. The compactness of $`Z`$ allows us to find a subsequence of $`\left(z_n\right)`$ converging to some $`zZ`$ with $`b(u,z)B`$. Then there is some $`\epsilon >0`$ so that $`B_d(b,\epsilon )B`$. Applying the Effros theorem \[E\],\[U\] to $`H\left(X\right)`$, the group of homeomorphisms of $`X`$ in the $`sup`$ metric, there is a $`\delta >0`$ so that for any $`v,wX`$ with $`d(v,w)<\delta `$ there is an $`h^{[v,w]}H\left(X\right)`$ with $`h^{[v,w]}\left(v\right)=w`$ and $`d(p,h^{[v,w]}\left(p\right))<\epsilon /2`$ for all $`pX`$ (denoted: $`\delta \stackrel{\text{ Eff}}{}\epsilon /2`$). We now choose $`n`$ so that $`d(b_n,b)<\epsilon /2`$ and $`d(x_n,x)<\delta `$ and a corresponding homeomorphism $`h^{[x_n,x]}`$. Then $`h^{[x_n,x]}\left(b_n\right)B_d(b,\epsilon )B`$ is a point in the path component of $`x`$, demonstrating that $`xPC\left(B\right),`$ and so $`PC\left(B\right)=X`$. Since $`B`$ was any such basis element, each path component of $`X`$ is dense.
With $`\stackrel{~}{\alpha }_x:𝐑^\kappa X`$ denoting the $`\stackrel{~}{\alpha }`$orbit of $`xX`$ and giving
$$A\stackrel{\text{def}}{=}\{\stackrel{~}{\alpha }_xxX\}$$
the topology it inherits from the collection of all maps $`𝐑^\kappa X`$ in the $`sup`$ metric, we proceed to show that $`h:XA;`$ $`h\left(x\right)\stackrel{\text{def}}{=}\stackrel{~}{\alpha }_x`$ is a homeomorphism. Since $`h`$ is clearly a bijection and $`X`$ is compact, it suffices to show that $`h`$ is continuous. Given $`\left(x_n\right)_nx`$ in $`X`$ and $`\epsilon >0`$ we need to find an $`N`$ so that $`\underset{s𝐑^\kappa }{sup}d(\stackrel{~}{\alpha }_{x_n}\left(s\right),\stackrel{~}{\alpha }_x\left(s\right))<\epsilon `$ for all $`nN`$.
Let $`d_\kappa `$ be a translation invariant metric for $`𝐓^\kappa `$. First we find a Lebesgue number $`\lambda >0`$ for a covering $`𝒪`$ of $`𝐓^\kappa `$ by open sets $`V`$ satisfying
$$p^1\left(V\right)V\times F\text{ }$$
and $`p\left(y\right)=v`$ for $`y(v,f).`$ Since $`X`$ and $`𝐓^\kappa `$ are compact, we may find a connected neighborhood $`U`$ of $`\mathrm{𝟎}𝐑^\kappa `$ so that
$$\underset{y𝐓^\kappa }{sup}\left\{\text{diam}\alpha \left(U\times \left\{y\right\}\right)\right\}<\lambda \text{ and }\underset{yX}{sup}\left\{\text{diam}\stackrel{~}{\alpha }\left(U\times \left\{y\right\}\right)\right\}<\frac{\epsilon }{3}$$
and so that $`\alpha \left(U\times \left\{e\right\}\right)`$ is the $`\eta `$neighborhood of $`e`$ in $`𝐓^\kappa `$ for some $`\eta >0`$. The translation invariance of $`d_\kappa `$ yields that $`\alpha \left(U\times \left\{y\right\}\right)`$ is the $`\eta `$neighborhood of $`y𝐓^\kappa `$. Since $`p`$ is uniformly continuous, there is a $`\tau >0`$ so that
$$d(y,y^{})<\tau d_\kappa (p\left(y\right),p\left(y^{}\right))<\eta .$$
Next we find
$$0<\delta \stackrel{\text{Eff}}{}\mathrm{min}\{\frac{1}{2}\epsilon ,\tau \}\text{ and }N\text{ so that }\left\{x_n\right\}_{nN}B_d(x,\delta ).$$
Given $`nN`$ there is an $`h^{[x_n,x]}H\left(X\right)`$ within $`\mathrm{min}\{\frac{1}{2}\epsilon ,\tau \}`$ of $`id_X`$. And so
$$p\left(x\right)=p\left(h^{[x_n,x]}\left(x_n\right)\right)\alpha \left(U\times \left\{p\left(x_n\right)\right\}\right)$$
since $`d(x_n,x)<\tau .`$ For $`s𝐑^\kappa `$, the translation invariance of $`d_\kappa `$ yields that $`d_\kappa (\alpha (s,p\left(x\right)),\alpha (s,p\left(x_n\right)))<\eta `$, while the choice of $`h^{[x_n,x]}`$ and the equality $`\alpha (s,p\left(x_n\right))=p\stackrel{~}{\alpha }(s,x_n)`$ yields that
$$d_\kappa (\alpha (s,p\left(x_n\right)),ph^{[x_n,x]}\stackrel{~}{\alpha }(s,\left(x_n\right)))<\eta .$$
Combining this with the equality $`\alpha (s,p\left(x\right))=p\stackrel{~}{\alpha }(s,x),`$ we obtain
$$()\{p\stackrel{~}{\alpha }(s,x),ph^{[x_n,x]}\stackrel{~}{\alpha }(s,\left(x_n\right))\}\alpha \left(U\times \left\{\alpha (s,p\left(x_n\right))\right\}\right)\text{ for all }s𝐑^\kappa .$$
For $`s𝐑^\kappa `$ let
$$U\left(s\right)\stackrel{\text{def}}{=}p^1\left(\alpha \left(U\times \left\{\alpha (s,p\left(x_n\right))\right\}\right)\right).$$
By its size, we know that $`\alpha \left(U\times \left\{\alpha (s,p\left(x_n\right))\right\}\right)`$ fits inside some $`V𝒪`$, and so
$$U\left(s\right)\alpha \left(U\times \left\{\alpha (s,p\left(x_n\right))\right\}\right)\times F$$
and for $`y(v,f)`$, $`p\left(y\right)=v`$. Then by $`()`$ we have
$$\{\stackrel{~}{\alpha }(s,x),h^{[x_n,x]}\stackrel{~}{\alpha }(s,\left(x_n\right))\}U\left(s\right).$$
Let
$$W\stackrel{\text{def}}{=}\{s𝐑^\kappa \stackrel{~}{\alpha }(s,\left(x\right))\text{ and }h^{[x_n,x]}\stackrel{~}{\alpha }(s,\left(x_n\right))\text{ are in the same component of }U\left(s\right)\}.$$
By construction, both $`W`$ and $`𝐑^\kappa W`$ are open in $`𝐑^\kappa `$ since the components of $`U\left(s\right)`$ are homeomorphic to $`U`$. Since $`h^{[x_n,x]}\left(x_n\right)=x`$, $`\mathrm{𝟎}W`$ and since $`𝐑^\kappa `$ is connected, $`W=𝐑^\kappa `$. Thus, with $`C\left(s\right)`$ denoting the component of $`U\left(s\right)`$ containing both $`\stackrel{~}{\alpha }(s,\left(x\right))`$ and $`h^{[x_n,x]}\stackrel{~}{\alpha }(s,\left(x_n\right))`$, and with
$$f_s\stackrel{\text{def}}{=}C\left(s\right)p^1\left(\alpha (s,p\left(x_n\right))\right)$$
we have that
$$\{\stackrel{~}{\alpha }(s,x),h^{[x_n,x]}\stackrel{~}{\alpha }(s,\left(x_n\right))\}\stackrel{~}{\alpha }\left(U\times \left\{f_s\right\}\right),$$
implying that $`d(\stackrel{~}{\alpha }(s,x),h^{[x_n,x]}\stackrel{~}{\alpha }(s,\left(x_n\right)))<\epsilon /3`$ by our initial choice of $`U`$. The choice of $`h^{[x_n,x]}`$ yields that $`d(\stackrel{~}{\alpha }(s,x),\stackrel{~}{\alpha }(s,\left(x_n\right)))<5\epsilon /6`$ . Since $`s`$ was any point of $`𝐑^\kappa `$ , we may finally conclude that
$$\underset{s𝐑^\kappa }{sup}d(\stackrel{~}{\alpha }_{x_n}\left(s\right),\stackrel{~}{\alpha }_x\left(s\right))<\epsilon \text{ for all }nN.$$
And so $`h`$ is a homeomorphism. It follows easily that $`\stackrel{~}{\alpha }`$ is uniformly Lyapunov stable in the strongest sense: for any given $`\epsilon >0`$ there is a $`\delta >0`$ so that
$$d(x,y)<\delta d(\stackrel{~}{\alpha }(s,x),\stackrel{~}{\alpha }(s,y))<\epsilon \text{ for all }s𝐑^\kappa .$$
Fixing a point $`eX`$, the density of the orbit of $`e`$ means that any $`a,bX`$ may be represented in the form $`a=lim_i\left\{\stackrel{~}{\alpha }(t_i^a,e)\right\}`$ and $`b=lim_i\left\{\stackrel{~}{\alpha }(t_i^b,e)\right\}`$. The operations
$$a=\underset{i}{lim}\left\{\stackrel{~}{\alpha }(t_i^a,e)\right\}\text{ and }a+b=\underset{i}{lim}\left\{\stackrel{~}{\alpha }(t_i^a+t_i^b,e)\right\}$$
then give $`X`$ a well-defined, abelian topological group structure compatible with the original topology of $`X`$. The proof is essentially identical with that found in \[NS\] V, 8.16 and is therefore omitted. $`\mathrm{}`$
It follows that such an $`X`$ is homeomorphic to either $`𝐓^\kappa `$ or the inverse limit of an inverse sequence of $`𝐓^\kappa `$ with epimorphic bonding maps (see \[C\]). And it follows directly that any such $`X`$ is bihomogeneous: the homeomorphism
$$x\left(a+b\right)x$$
switches $`a`$ and $`b`$.
## 3. Exponents
We now define the exponent group and explore its topological significance. In what follows, $`Hom(𝐑^\kappa ,𝐑^\kappa )`$ denotes the group of continuous homomorphisms $`𝐑^\kappa 𝐑^\kappa `$ with point-wise addition and $`[X;Y]`$ denotes the group of homotopy classes of maps $`XY`$, and $`\left[f\right]`$ denotes the homotopy class of $`f:XY`$.
###### Definition 3.1.
*Given a map* $`f:WX`$*, we define* $`\left\{w_i\right\}_{i=1}^{\mathrm{}}W`$ *to be an* $`f`$–sequence *if* $`\left\{f\left(w_i\right)\right\}_{i=1}^{\mathrm{}}`$ *converges in* $`X`$.
###### Definition 3.2.
*Given a map* $`f`$ *of* $`𝐑^\kappa `$ *into a metric space* $`X`$*, the* exponent group of$`f`$ , *denoted* $`_f,`$ *is*
$$\{AHom(𝐑^\kappa ,𝐑^\kappa )\left\{\pi ^\kappa \left(A\left(t_i\right)\right)\right\}_{i=1}^{\mathrm{}}\text{ }\text{converges in}\text{ }𝐓^\kappa \text{for all}\text{ }f\text{ –}\text{sequences}\text{ }\left\{t_i\right\}_{i=1}^{\mathrm{}}\}.$$
###### Lemma 3.3.
$`_f`$ is a subgroup of $`Hom(𝐑^\kappa ,𝐑^\kappa )`$.
Proof: The zero homomorphism is trivially in $`_f`$. If $`A`$ and $`B`$ are in $`_f`$ and if $`\left\{t_i\right\}_{i=1}^{\mathrm{}}`$ is an $`f`$sequence, then $`\left\{\pi ^\kappa \left(A\left(t_i\right)\right)\right\}_{i=1}^{\mathrm{}}`$ and $`\left\{\pi ^\kappa \left(B\left(t_i\right)\right)\right\}_{i=1}^{\mathrm{}}`$ converge in $`𝐓^\kappa `$. And since $``$ is continuous on $`Hom(𝐑^\kappa ,𝐑^\kappa )\times Hom(𝐑^\kappa ,𝐑^\kappa )`$, this implies that $`\left\{\pi ^\kappa \left(\left(AB\right)\left(t_i\right)\right)\right\}_{i=1}^{\mathrm{}}`$ converges in $`𝐓^\kappa `$. $`\mathrm{}`$
###### Theorem 3.4.
Given a map $`f`$ of $`𝐑^\kappa `$ into a metric space $`X`$, with
$$\iota :_f[\overline{f\left(𝐑^\kappa \right)};𝐓^\kappa ]$$
given by
$$A\stackrel{\iota }{}\left[f_A\right]\text{, where }f_A:\overline{f\left(𝐑^\kappa \right)}𝐓^\kappa \text{ sends }\underset{i}{lim}\left\{f\left(t_i\right)\right\}\text{ to }\underset{i}{lim}\left\{\pi ^\kappa \left(A\left(t_i\right)\right)\right\}$$
we have:
1. $`\iota `$ is a homomorphism (we give $`[\overline{f\left(𝐑^\kappa \right)};𝐓^\kappa ]`$ the group operation induced by point-wise addition of maps).
2. If $`\iota \left(A\right)=\iota \left(B\right)`$ and if $`𝒞`$ is any contractible topological subspace of $`𝐑^\kappa `$ satisfying the condition that the map $`\left(AB\right)_𝒞:𝒞\left(AB\right)\left(𝒞\right)`$; $`t\left(AB\right)\left(t\right)`$ has a continuous inverse, then for any $`f`$ –sequence $`\left\{t_i\right\}_{i=1}^{\mathrm{}}𝒞`$, the sequence $`\left\{\left(AB\right)\left(t_i\right)\right\}_{i=1}^{\mathrm{}}`$ converges in $`𝐑^\kappa `$.
3. If $`\overline{f\left(𝐑^\kappa \right)}`$ is compact, then $`\iota `$ is an embedding and $`_f`$ is countable.
4. If $`\iota \left(A\right)=\iota \left(B\right)`$ and if $`\kappa =n<\mathrm{}`$ and if there is an unbounded $`f`$–sequence $`\left\{t_i\right\}_{i=1}^{\mathrm{}}`$, then $`AB`$ is not invertible.
Proof: $`(1)`$ Let $`A_f`$ be given. If $`\underset{i}{lim}\left\{f\left(s_i\right)\right\}=\underset{i}{lim}\left\{f\left(t_i\right)\right\}=x`$ are two representations of a point in $`\overline{f\left(𝐑^\kappa \right)}`$, then $`\underset{i}{lim}\left\{\pi ^\kappa \left(A\left(s_i\right)\right)\right\}=\sigma `$ and $`\underset{i}{lim}\left\{\pi ^\kappa \left(A\left(t_i\right)\right)\right\}=\tau `$ both exist by the definition of $`_f`$. Then for $`i\{1,2,\mathrm{}\}`$ if we define $`u_{2i1}=s_i`$ and $`u_{2i}=t_i`$, we have that $`\underset{i}{lim}\left\{f\left(u_i\right)\right\}=x`$, implying that $`\underset{i}{lim}\left\{\pi ^\kappa \left(A\left(u_i\right)\right)\right\}`$ exists, which is only possible if $`\sigma =\tau `$. Thus, $`f_A`$ is a well–defined function. To see that $`f_A`$ is continuous, consider a convergent sequence $`\left\{\underset{i}{lim}\left\{f\left(t_i^j\right)\right\}\right\}_{j=1}^{\mathrm{}}=\left\{x^j\right\}_{j=1}^{\mathrm{}}x`$. Then with $`d`$ denoting a metric for $`X`$ and with $`d_\kappa `$ denoting a metric for $`𝐓^\kappa `$, for each $`j\{1,2,\mathrm{}\}`$ choose $`i_j`$ so that:
$$d(f\left(t_{i_j}^j\right),x^j)<\frac{1}{j}\text{ and }d_\kappa (f_A\left(x^j\right),\pi ^\kappa \left(A\left(t_{i_j}^j\right)\right))<\frac{1}{j}.$$
Then $`\underset{j}{lim}\left\{f\left(t_{i_j}^j\right)\right\}_{j=1}^{\mathrm{}}=x`$ and
$$f_A\left(x\right)=\underset{j}{lim}\left\{\pi ^\kappa \left(A\left(t_{i_j}^j\right)\right)\right\}_{j=1}^{\mathrm{}}=\underset{j}{lim}\left\{f_A\left(x^j\right)\right\}_{j=1}^{\mathrm{}},$$
demonstrating that $`f_A`$ is continuous. And if $`A,B_f`$ we have for $`x=\underset{i}{lim}\left\{f\left(t_i\right)\right\}`$
$$f_{A+B}\left(x\right)=\underset{i}{lim}\left\{\pi ^\kappa \left(\left(A+B\right)\left(t_i\right)\right)\right\}=\underset{i}{lim}\left\{\pi ^\kappa \left(A\left(t_i\right)\right)\right\}+\underset{i}{lim}\left\{\pi ^\kappa \left(B\left(t_i\right)\right)\right\}=f_A\left(x\right)+f_B\left(x\right)\text{,}$$
demonstrating that $`\iota `$ is a homomorphism.
$`(2)`$ Since $`\left[f_A\right]=\left[f_B\right]`$, we have that $`\left[f_{\left(AB\right)}\right]=\left[\text{constant map}\right]`$, and so there is a map $`g:(\overline{f\left(𝐑^\kappa \right)},f\left(\mathrm{𝟎}\right))(𝐑^\kappa ,\mathrm{𝟎})`$ lifting $`f_{\left(AB\right)}`$ making the following diagram commute
$$\begin{array}{ccccc}& & 𝐑^\kappa & & \\ & {}_{}{}^{g}& & ^{\pi ^\kappa }& \\ \overline{f\left(𝐑^\kappa \right)}& & \stackrel{f_{\left(AB\right)}}{}& & 𝐓^\kappa \end{array}.$$
Then for any $`t\left(AB\right)\left(𝒞\right)`$ we have
$$f_{\left(AB\right)}f\left(AB\right)_𝒞^1\left(t\right)=\pi ^\kappa \left(\left(AB\right)\left(AB\right)_𝒞^1\left(t\right)\right)=\pi ^\kappa \left(t\right),$$
and so we are led to the following commutative diagram:
$$\begin{array}{ccccccc}\left(AB\right)\left(𝒞\right)& \stackrel{\left(AB\right)_𝒞^1}{}& 𝒞& \stackrel{f}{}& f\left(𝒞\right)& \stackrel{g}{}& 𝐑^\kappa \\ \pi ^\kappa & & & & & f_{\left(AB\right)}& \pi ^\kappa \\ 𝐓^\kappa & & \stackrel{id}{}& & & & 𝐓^\kappa \end{array}\text{.}$$
And since, $`gf\left(AB\right)_𝒞^1`$ and $`id_{𝐑^\kappa }`$ both map $`\mathrm{𝟎}`$ to $`\mathbf{\hspace{0.17em}0}`$, we have that $`gf\left(AB\right)_𝒞^1=id_{𝐑^\kappa }`$ since both provide a lift of $`id_{𝐓^1}\pi ^\kappa _{\left(AB\right)\left(𝒞\right)}`$ and such a lift is uniquely determined. And so if $`\left\{t_i\right\}_{i=1}^{\mathrm{}}𝒞`$ is an $`f`$–sequence with $`\underset{i}{lim}\left\{f\left(t_i\right)\right\}=x\overline{f\left(𝐑^\kappa \right)}`$, we must have
$`g\left(x\right)`$ $`=`$ $`g\left(\underset{i}{lim}\left\{f\left(AB\right)_𝒞^1\left(\left(AB\right)_𝒞\left(t_i\right)\right)\right\}\right)`$
$`=`$ $`\underset{i}{lim}\left\{gf\left(AB\right)_𝒞^1\left(\left(AB\right)_𝒞\left(t_i\right)\right)\right\}=\underset{i}{lim}\left\{\left(AB\right)_𝒞\left(t_i\right)\right\}\text{ in }𝐑^\kappa `$
by the continuity of $`g`$.
$`(3)`$ Suppose then that $`\overline{f\left(𝐑^\kappa \right)}`$ is compact and that $`\iota \left(A\right)=\left[f_A\right]=\left[\text{ constant map}\right]`$ and that $`v𝐑^\kappa \mathrm{ker}A`$. Then with $`𝒞`$ denoting the vector subspace $`𝐑v𝐑^\kappa `$, we have that $`A_𝒞=\left(A0\right)_𝒞`$ is an isomorphism onto $`𝐑A\left(v\right)`$. And since $`\overline{f\left(𝐑^\kappa \right)}`$ is compact, there is a subsequence $`\left\{f\left(n_iv\right)\right\}_{i=1}^{\mathrm{}}`$ of the sequence $`\left\{f\left(nv\right)\right\}_{n=1}^{\mathrm{}}`$ which converges to some $`x\overline{f\left(𝐑^\kappa \right)}`$. But then by the above, we must have that $`\left\{n_iA\left(v\right)\right\}_{i=1}^{\mathrm{}}`$ converges in $`𝐑^\kappa `$. And so if $`w_{\mathrm{}}`$ is any non-zero component of $`A\left(v\right)=w_i_{i=1}^\kappa `$ , we would then have that $`\left\{n_iw_{\mathrm{}}\right\}_{i=1}^{\mathrm{}}`$ converges in $`𝐑`$, which is impossible since $`\left\{n_i\right\}_{i=1}^{\mathrm{}}`$ is unbounded. Thus, we must have $`\mathrm{ker}A=𝐑^\kappa `$ and $`A`$ is the zero map. That $`_f`$ is countable then follows from the fact that $`[\overline{f\left(𝐑^\kappa \right)};𝐓^\kappa ]`$ is countable.
$`\left(4\right)`$ If $`\kappa <\mathrm{}`$ and if $`\left\{t_i\right\}_{i=1}^{\mathrm{}}`$ is an unbounded $`f`$–sequence, then $`\left\{\left(AB\right)\left(t_i\right)\right\}_{i=1}^{\mathrm{}}`$ would be unbounded if $`\left(AB\right)`$ were invertible. $`\mathrm{}`$
We know that when $`f`$ has an image which is not compact, $`_f`$ may not be countable and $`\iota `$ may not be an embedding. For example, for $`f:𝐑^2𝐑\times 𝐓^1;`$ $`t_{1,}t_2t_1,\pi ^1\left(t_2\right)`$, making the identification of $`Hom(𝐑^2,𝐑^2)`$ with $`Mat(2,𝐑)`$ (the group of $`2\times 2`$ matrices with real entries under component–wise addition), we have
$$\{\left(\begin{array}{cc}c\hfill & 0\hfill \\ 0\hfill & z\hfill \end{array}\right)Mat(2,𝐑)c𝐑\text{ and }z𝐙\}_f$$
and $`\iota \left(\left(\begin{array}{cc}c\hfill & 0\hfill \\ 0\hfill & z\hfill \end{array}\right)\right)=\iota \left(\left(\begin{array}{cc}c^{}\hfill & 0\hfill \\ 0\hfill & z\hfill \end{array}\right)\right)`$ for any $`c`$ and $`c^{}`$ in $`𝐑`$. And since there are unbounded $`f`$–sequences, we know that $`\left(\begin{array}{cc}c\hfill & 0\hfill \\ 0\hfill & z\hfill \end{array}\right)\left(\begin{array}{cc}c^{}\hfill & 0\hfill \\ 0\hfill & z\hfill \end{array}\right)`$ is not invertible.
###### Theorem 3.5.
If $`f:𝐑^\kappa X`$ has exponent group $`_f=\{A_\lambda \lambda \mathrm{\Lambda }\}`$, we have the map
$$h_f:\overline{f\left(𝐑^\kappa \right)}\underset{\lambda \mathrm{\Lambda }}{}𝐓^\kappa ;\text{ }xf_{A_\lambda }\left(x\right)_{\lambda \mathrm{\Lambda }},$$
and $`h_ff`$ is a homomorphism of $`𝐑^\kappa `$ into $`\underset{\lambda \mathrm{\Lambda }}{}𝐓^\kappa `$. And if $`\overline{f\left(𝐑^\kappa \right)}`$ is compact, then $`h_f\left(\overline{f\left(𝐑^\kappa \right)}\right)`$ is a subgroup of $`\underset{\lambda \mathrm{\Lambda }}{}𝐓^\kappa `$.
Proof: Since each $`f_{A_\lambda }`$ is continuous, we know that $`h_f`$ is also continuous. When we give $`\underset{\lambda \mathrm{\Lambda }}{}𝐓^\kappa `$ the group structure of component–wise addition, it is a topological group with the product (Tychonoff) topology and
$`h_ff\left(s+t\right)`$ $`=`$ $`\pi ^\kappa A_\lambda \left(s+t\right)_{\lambda \mathrm{\Lambda }}=\pi ^\kappa A_\lambda \left(s\right)_{\lambda \mathrm{\Lambda }}+\pi ^\kappa A_\lambda \left(t\right)_{\lambda \mathrm{\Lambda }}`$
$`=`$ $`h_ff\left(s\right)+h_ff\left(t\right),`$
and so $`h_ff`$ is a homomorphism. When $`\overline{f\left(𝐑^\kappa \right)}`$ is compact, we have that $`h_f\left(\overline{f\left(𝐑^\kappa \right)}\right)=\overline{h_ff\left(𝐑^\kappa \right)}`$ (recall that $`\mathrm{\Lambda }`$ is countable when $`\overline{f\left(𝐑^\kappa \right)}`$ is compact). We now know that $`h_ff\left(𝐑^\kappa \right)`$ is a subgroup of $`\underset{\lambda \mathrm{\Lambda }}{}𝐓^\kappa `$, and so $`\overline{h_ff\left(𝐑^\kappa \right)}`$ is also a subgroup. $`\mathrm{}`$
We are now in a position to give a generalization of almost periodic maps $`𝐑X`$ to almost periodic maps $`𝐑^\kappa X`$.
###### Definition 3.6.
*A map* $`f:𝐑^\kappa X`$ *is* almost periodic *if* $`_f`$ *is countable and if the following condition* *is satisfied*:
$$()\left[\left\{t_i\right\}_{i=1}^{\mathrm{}}\text{ is an}\text{ }f\text{}\text{sequence}\right]\left[\left\{\pi ^\kappa \left(A\left(t_i\right)\right)\right\}_{i=1}^{\mathrm{}}\text{ converges for all}\text{ }A_f\right].$$
It is known that any classically defined almost periodic function satisfies $`()`$ when $`\kappa =1`$ and $`X`$ is complete, see \[Cl\] . And below we shall see that any map satisfying $`()`$ is indeed almost periodic in the classical sense when $`\kappa =1`$ and similar techniques may be used to show that generally the above definition coincides with Bochner’s \[B\].
###### Theorem 3.7.
If $`f`$ is almost periodic, then $`\overline{f\left(𝐑^\kappa \right)}`$ is compact, and $`h_f`$ as in Theorem 3.5 is a homeomorphism of $`\overline{f\left(𝐑^\kappa \right)}`$ onto the subgroup $`h_f\left(\overline{f\left(𝐑^\kappa \right)}\right)`$ of $`\underset{n=1}{\overset{\mathrm{}}{}}𝐓^\kappa `$.
Proof: Let $`\left\{x^j\right\}_{j=1}^{\mathrm{}}`$ be a sequence in $`\overline{f\left(𝐑^\kappa \right)}`$. Choose a subsequence of $`\left\{h_f\left(x^j\right)\right\}_{j=1}^{\mathrm{}}`$ which converges in $`\underset{n=1}{\overset{\mathrm{}}{}}𝐓^\kappa `$ and which we label $`\left\{h_f\left(y^j\right)\right\}_{j=1}^{\mathrm{}}`$ for convenience. With $`y^j=\underset{i}{lim}\left\{f\left(t_i^j\right)\right\}_{i=1}^{\mathrm{}}`$ and with $`_f=\{A_1,A_2,\mathrm{}\},`$ for each $`j\{1,2,\mathrm{}\}`$ choose $`t_{i_j}^j`$ such that
$$d(f\left(t_{i_j}^j\right),y^j)<\frac{1}{j}\text{ and }\mathrm{max}\left\{d_{𝐓^\kappa }(\pi ^\kappa \left(A_n\left(t_{i_j}^j\right)\right),f_{A_n}\left(y^j\right))n\{1,\mathrm{},j\}\right\}<\frac{1}{j}\text{.}$$
Since we have that $`\left\{f_{A_n}\left(y^j\right)\right\}_{j=1}^{\mathrm{}}`$ converges for each $`n\{1,2,\mathrm{}\}`$, we must then have that $`\left\{\pi ^\kappa \left(A\left(t_{i_j}^j\right)\right)\right\}_{j=1}^{\mathrm{}}`$ converges for all $`A_f`$, and so $`\left\{t_{i_j}^j\right\}_{j=1}^{\mathrm{}}`$ is an $`f`$–sequence by $`()`$, say $`\left\{f\left(t_{i_j}^j\right)\right\}_{j=1}^{\mathrm{}}yX`$. And so we also have $`\left\{f\left(y^j\right)\right\}_{j=1}^{\mathrm{}}y`$, demonstrating that $`\overline{f\left(𝐑^\kappa \right)}`$ is compact. Now suppose that $`h_f\left(x\right)=h_f\left(y\right)`$ with $`x=\underset{i}{lim}\left\{f\left(s_i\right)\right\}`$ and $`y=\underset{i}{lim}\left\{f\left(t_i\right)\right\}`$. We then have that for each $`n\{1,2,\mathrm{}\}`$
$$f_{A_n}\left(x\right)=\underset{i}{lim}\left\{\pi ^\kappa \left(A_n\left(s_i\right)\right)\right\}_{i=1}^{\mathrm{}}=\underset{i}{lim}\left\{\pi ^\kappa \left(A_n\left(t_i\right)\right)\right\}_{i=1}^{\mathrm{}}=f_{A_n}\left(y\right),$$
and so with $`u_{2i1}=s_i`$ and $`u_{2i}=t_i`$ for $`i\{1,2,\mathrm{}\},`$ we have that $`\underset{i}{lim}\left\{\pi ^\kappa \left(A_n\left(u_i\right)\right)\right\}_{i=1}^{\mathrm{}}=f_{A_n}\left(x\right)`$ for each $`n\{1,2,\mathrm{}\}.`$ Then by $`()`$ we have that $`\left\{f\left(u_i\right)\right\}_{i=1}^{\mathrm{}}`$ converges in $`X`$, which is only possible if $`x=y`$. Thus, $`h_f`$ is one–to–one and so is a homeomorphism onto its image. $`\mathrm{}`$
When $`f`$ is almost periodic, there is a topological isomorphism $`𝔦_f`$ from $`h_f\left(\overline{f\left(𝐑^\kappa \right)}\right)`$ onto some $`\lambda `$–solenoid $`_{\overline{M}}`$ $`\left(\lambda \mathrm{}\right)`$ since $`h_f\left(\overline{f\left(𝐑^\kappa \right)}\right)`$ is a compact connected abelian group; see \[P\] Thm 68, \[NS\] V,8.16 \[C\] and \[Cl\]. Then, using the notation of \[C\], we have a map $`𝔥:(𝐑^\kappa ,\mathrm{𝟎})(𝐑^\lambda ,\mathrm{𝟎})`$ making the following diagram commute:
$$\begin{array}{ccc}& & 𝐑^\lambda \\ & 𝔥& \pi _{\overline{M}}\\ 𝐑^\kappa & \stackrel{𝔦_fh_ff}{}& _{\overline{M}}\end{array}$$
since $`\pi _{\overline{M}}`$ is a fibration with unique path lifting. Then for $`s,t𝐑^\kappa `$, since $`𝔦_fh_ff`$ and $`\pi _{\overline{M}}`$ are homomorphisms, we have
$`\pi _{\overline{M}}\left(𝔥\left(s\right)+𝔥\left(t\right)\right)`$ $`=`$ $`\pi _{\overline{M}}𝔥\left(s\right)+\pi _{\overline{M}}𝔥\left(t\right)=𝔦_fh_ff\left(s\right)+𝔦_fh_ff\left(t\right)`$
$`=`$ $`𝔦_fh_ff\left(s+t\right)=\pi _{\overline{M}}𝔥\left(s+t\right),`$
and so $`𝔥\left(s\right)+𝔥\left(t\right)𝔥\left(s+t\right)\mathrm{ker}\pi _{\overline{M}}`$, implying that $`𝔥\left(s\right)+𝔥\left(t\right)=𝔥\left(s+t\right)`$ for all $`s,t𝐑^\kappa `$, see \[C\] 3.4. Hence, $`𝔥\left(𝐑^\kappa \right)`$ is a subgroup of $`𝐑^\lambda `$. And so we may think of $`\pi _{\overline{M}}𝔥\left(𝐑^\kappa \right)`$ as a “linear subspace” of $`_{\overline{M}}`$.
###### Theorem 3.8.
If $`f:𝐑^\kappa X`$ is almost periodic, then $`f`$ may be extended to a continuous group action of $`(𝐑^\kappa ,+)`$ on all of $`\overline{f\left(𝐑^\kappa \right)}:`$
$$\alpha _f:𝐑^\kappa \times \overline{f\left(𝐑^\kappa \right)}\overline{f\left(𝐑^\kappa \right)},\text{ where }\alpha _f(t,f\left(0\right))=f\left(t\right).$$
Proof: We have the group action
$$\varphi _f:𝐑^\kappa \times _{\overline{M}}_{\overline{M}};(t,x)\pi _{\overline{M}}𝔥\left(t\right)+x.$$
For $`x_{\overline{M}}`$, we have $`\varphi _f(\mathrm{𝟎},x)=\pi _{\overline{M}}𝔥\left(\mathrm{𝟎}\right)+x=e_{\overline{M}}+x=x`$ and for $`s,t𝐑^\kappa `$ we have
$`\varphi _f(s+t,x)`$ $`=`$ $`\pi _{\overline{M}}𝔥\left(s+t\right)+x=\pi _{\overline{M}}𝔥\left(s\right)+\pi _{\overline{M}}𝔥\left(t\right)+x`$
$`=`$ $`\varphi _f(s,\varphi _f(t,x)),`$
and so $`\varphi _f`$ is indeed a group action. We then define
$$\alpha _f(s,x)\stackrel{def}{=}\left(𝔦_fh_f\right)^1\left(\varphi _f(s,𝔦_fh_f\left(x\right))\right),$$
i.e., the action conjugate via $`𝔦_fh_f`$ to $`\varphi _f`$. It then follows directly that $`\alpha _f`$ is a group action, and
$`\alpha _f(t,f\left(0\right))`$ $`=`$ $`\left(𝔦_fh_f\right)^1\left(\varphi _f(t,𝔦_fh_f\left(f\left(0\right)\right))\right)=\left(𝔦_fh_f\right)^1\left(\pi _{\overline{M}}𝔥\left(t\right)+e_{\overline{M}}\right)`$
$`=`$ $`\left(𝔦_fh_f\right)^1\left(𝔦_fh_ff\left(t\right)\right)=f\left(t\right).\mathrm{}`$
And when $`\kappa =1`$, $`\alpha _f`$ is a flow on $`\overline{f\left(𝐑^\kappa \right)}`$ and $`𝔦_fh_f`$ provides an equivalence of $`\alpha _f`$ and an almost periodic linear flow as described in \[C\] , implying that $`f`$ is itself almost periodic in the classical sense since it is an orbit of this flow.
Generally, it is now clear that all $`\varphi _f`$–orbits are translates of $`\pi _{\overline{M}}𝔥\left(𝐑^\kappa \right)`$, and so we consider the decomposition of $`_{\overline{M}}`$ determined by $`\varphi _f`$ as a decomposition into “linear subspaces.” Thus, $`f`$ determines a group action on $`\overline{f\left(𝐑^\kappa \right)}`$ which is equivalent to a “linear foliation” of $`_{\overline{M}}`$. It then follows that if $`f:𝐑^nM`$ is a smooth almost periodic map onto a leaf of an $`n`$–dimensional foliation $``$ of the manifold $`M`$ that $`𝔦_fh_f`$ then provides an equivalence of $`_{\overline{f\left(𝐑^n\right)}}`$ with the linear foliation $`\varphi _f`$ since the leaves of $`_{\overline{f\left(𝐑^n\right)}}`$ will coincide with the orbits of $`\alpha _f`$ due to the density of $`f\left(𝐑^n\right)`$ in $`\overline{f\left(𝐑^n\right)}`$.
Given $`𝐯𝐑^\kappa `$, any map $`f:𝐑^\kappa X`$ induces the map
$$𝐑t𝐯𝐑^\kappa \stackrel{f}{}X$$
and the one–dimensional exponent group of this map $`𝐑X`$ is frequently sufficient to determine as much about $`\stackrel{ˇ}{H}^1\left(\overline{f\left(𝐑^\kappa \right)}\right)[\overline{f\left(𝐑^\kappa \right)};𝐓^1]`$ as $`_f`$ reveals. But when $`f`$ extends to a group action, $`f_A`$ will be a fiber bundle projection, and some of the bundle structure is lost in considering only the one–dimensional maps. Thus, the higher dimensional groups have content not captured in the one-dimensional exponent group.
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# 1 Graviton profile on the Schwarzschild-AdS space-time. For large r the graviton behaves like in ordinary AdS space and is not normalizable. A cutoff in form of a brane is needed at r=r₀. For small values of r, the graviton also diverges, but is hidden behind a horizon.
IFT-UAM/CSIC-00-11
hep-th/0003002
Brane World with Bulk Horizons
César Gómez,<sup>1</sup><sup>1</sup>1E-mail address: cesar.gomez@uam.es Bert Janssen<sup>2</sup><sup>2</sup>2E-mail address: bert.janssen@uam.es and Pedro J. Silva<sup>3</sup><sup>3</sup>3E-mail address: psilva@delta.ft.uam.es
Instituto de Física Teórica, C-XVI,
Departamento de Física Teórica, C-XI,
Universidad Autónoma de Madrid
E-28006 Madrid, Spain
ABSTRACT
A brane world in the presence of a bulk black hole is constructed. The brane tension is fine tuned in terms of the black hole mass and cosmological constant. Gravitational perturbations localized on the brane world are discussed.
1. Introduction
A brane world with induced four dimensional gravity was first introduced in on the basis of a $`AdS_5`$ bulk geometry. In this scheme normalizable gravitational zero modes are allowed due to the ultraviolet cutoff induced by the brane wall. The dilaton field is constant and the holographic degrees of freedom on the wall define a conformal field theory coupled to gravity . In a series of recent papers this framework was extended to the non conformal case i.e to dilatonic domain walls. In these cases, both with vanishing and non vanishing cosmological constant, we observe the phenomena of induced four dimensional gravity, however a naked curvature singularity is induced by the non constant dilaton in the bulk at a finite proper distance from the brane wall. The physics interpretation of such a singularity from the four dimensional point of view is still an open problem.
In this letter we will look for a Randall-Sundrum scenario but this time in the bulk geometry of a real black hole with the singularity inside a trapped surface. A similar analysis was with a non static Ansatz was first by in the framework of cosmological models. We will work out the static case in a Schwarzschild-AdS bulk metric. This will correspond to a brane world in a thermal bath at the Hawking temperature. The first question we will address would be the fine tuning relations between the brane wall tension and the parameters $`\mathrm{\Lambda }`$ and $`M`$ characterizing the Schwarzschild-AdS metric. These fine tuning relations would be obtained by solving the corresponding jump equations once we introduce the wall as an ultraviolet cutoff analogously to the AdS case. Since our space is asymptotically AdS this cutoff could be enough to induce four dimensional gravity on the wall in terms of normalizable graviton zero modes.
2. Construction of the solution
Our starting point is the following five-dimensional action of gravity in the presence of a cosmological constant $`\mathrm{\Lambda }`$ with a domain wall source term given by:
$$S=\frac{1}{\kappa }d^4x𝑑y\sqrt{|g|}\left[\mathrm{\Lambda }\right]+d^4x\sqrt{|\stackrel{~}{g}|}V_0,$$
(1)
where $`V_0`$ is the tension of the brane and $`\stackrel{~}{g}_{mn}=g_{\mu \nu }\delta _m^\mu \delta _n^\nu `$ the induced metric on the brane.
We are interested in a solution resembling the Schwarzschild-AdS solution, which is given by :
$$ds^2=(1+R^2r^2\frac{2M}{r^2})dt^2\frac{1}{(1+R^2r^2\frac{2M}{r^2})}dr^2r^2d\mathrm{\Omega }_3^2.$$
(2)
where $`R=12\mathrm{\Lambda }^1`$ and $`M`$ is basically the black hole mass. In order to apply the Randall-Sundrum program to this type of metrics, we consider a new set of coordinates defined by,
$$dz=\frac{1}{\sqrt{1+R^2r^2\frac{2M}{r^2}}}dr,$$
(3)
therefore ending up in a holographic-like frame
$$ds^2=A^2(z)dt^2+B^2(z)d\mathrm{\Omega }_3^2dz^2,$$
(4)
where $`A(z)`$ and $`B(z)`$ are function of the holografic coordinate $`z`$, implicitly given by
$$A(r)=\sqrt{1+R^2r^2\frac{2M}{r^2}},B(r)=r,$$
(5)
with $`r`$ a function of $`z`$, defined by the relation (3). Using this Ansatz (4) in the corresponding equations of motion of (1), we get the following system of equations:
$`B^2\mathrm{\Lambda }B^2(B^{})^2A^1A^{}B^1B^{}=0,`$
$`B^2+B^1B^{\prime \prime }+B^2(B^{})^2+\mathrm{\Lambda }+\frac{1}{6}\kappa V_0\delta (z)=0,`$ (6)
$`1+\frac{1}{2}B^2\mathrm{\Lambda }+A^1B^2A^{\prime \prime }+2BB^{\prime \prime }+(B^{})^2+2A^1A^{}BB^{}+\frac{1}{2}\kappa V_0\delta (z)=0,`$
where “ ” means derivatives with respect to $`z`$ and have used the fact that $`d\mathrm{\Omega }_3^2`$ is a three-dimensional sphere of radius one.
To obtain the desired solution, we use the fact that (5) is a solution of the action without source term, therefore in our case the solution to the full equations (S0.Ex2) is given by the functions (5) with a modification of the relation between $`z`$ and $`r`$, defined as follows
$$|\overline{z}|=z_0\frac{R}{2}\mathrm{log}\left|\frac{2R^1rA+2R^2r^2+1}{\sqrt{1+8MR^2}}\right|$$
(7)
where $`\overline{z}=z_0z`$, which translates into:
$$d|\overline{z}|=\frac{1}{\sqrt{1+R^2r^2\frac{2M}{r^2}}}dr.$$
(8)
The modulus of the radial coordinate $`\overline{z}`$ runs between the the position of the brane at $`\overline{z}=0`$ and the black hole horizons at $`|\overline{z}|=z_0`$. In this coordinates this means that the brane is located somewhere between the horizon and the boundary.
By solving the jump equations we find that the brane tension $`V_0`$ is given in terms of the cosmological constant $`\mathrm{\Lambda }`$, the black hole mass $`M`$ and $`r(\overline{z}=0)`$, provided the brane is located at the origin in $`\overline{z}`$ coordinates, by the following relation:
$$V_0=\frac{6}{\kappa r(0)}\sqrt{1+R^2r(0)^2\frac{2M}{r(0)^2}}$$
(9)
If the black hole horizon is smaller than the AdS radius $`M<R^2`$, we could choose to do the cutoff at $`r(\overline{z}=0)=R`$ and the above formula reduces to:
$$V_0=\frac{6}{\kappa R^2}\sqrt{2R^22M}=\frac{1}{\kappa }\sqrt{\frac{\mathrm{\Lambda }}{6}(1\frac{1}{12}M\mathrm{\Lambda })}$$
(10)
Note that in the limit $`M0`$, we recover the Randall-Sundrum relation between $`V_0`$ and $`\mathrm{\Lambda }`$ <sup>4</sup><sup>4</sup>4Also the limit $`R\mathrm{}`$ is well defined, recovering the standard Schwarzschild solution in Minkowski space, where the relation between $`r`$ and $`\overline{z}`$ is: $`\overline{z}=z_0\sqrt{r^22M}`$. .
In summary what we have done is basically to consider Schwarzschild-AdS space time with a brane located at a given distance $`r(0)`$ from the event horizon. Then replace the part of the space time outside the brane ($`r>r(0)`$) with a copy of the inner part, ending up with a finite range for the radial variable. It is important to note that this space time comes with two space-like singularities hidden inside the event horizons. Comparison with the dilatonic solution found on previous work , shows that the role of the singularity on those solutions is replaced by the event horizon in this new model. Nevertheless we also have a non isotropic worldbrane, the time direction scales differently than the space directions under radial flow.
3. Gravitational perturbations
To calculate the behavior of the graviton we add small fluctuations $`h_{\mu \nu }`$ to the above background, choosing the following gauge:
$$\widehat{g}_{\mu \nu }=g_{\mu \nu }+h_{\alpha \beta }\delta _\mu ^\alpha \delta _\nu ^\beta ,$$
(11)
where $`\alpha ,\beta `$ run over the coordinates $`t`$ and the angular coordinates $`x^m`$. Furthermore we have $`h=g^{\mu \nu }h_{\mu \nu }=A^2h_{tt}+B^2\overline{h}`$ and $`h=0`$, $`\overline{}^mh_{m\mu }=0`$, where $`\overline{}_m`$ stands for the covariant derivative of the angular coordinates. Notice that this is not the usual de Donder gauge since the perturbation is not traceless. Nevertheless if we are interested in a real graviton with two helicity states more constraints should be added.
The equation of motion, on this gauge for the fluctuation $`h_{\mu \nu }`$ is:
$$^2h_{\mu \nu }2^\rho _{(\mu }h_{\nu )\rho }+\frac{1}{3}\mathrm{\Lambda }h_{\mu \nu }\frac{1}{2}\mathrm{\Lambda }g_{\mu \nu }h=0,$$
(12)
Introducing the background (4), the components $`\{tt\}`$ and $`\{zz\}`$ of the above equation, reduce to:
$`A^3A^{}_zh_{tt}2A^4(A^{})^2h_{tt}+B^3B^{}_z\overline{h}2B^4(B^{})^2\overline{h}=0`$
$`_z^2h_{tt}+A^2_t^2h_{tt}B^2\overline{}^2h_{tt}3(B^1B^{}A^1A^{})_zh_{tt}`$
$`+(4A^2(A^{})^2\frac{1}{3}\mathrm{\Lambda })h_{tt}+\frac{1}{2}A^2h=0`$ (13)
To describe four-dimensional zero modes, we consider eigenfunction of the world brane variables $`x^\alpha `$, satisfying
$$A^2_t^2h_{tt}B^2\overline{}^2h_{tt}=0.$$
(14)
Under these conditions we find a very simple solution:
$$h_{tt}=A^2,\overline{h}=B^2$$
(15)
In principle we could turn on more degrees of freedom, to determine a more realistic graviton, nevertheless this mode shows the correct behavior to illustrate the location of the perturbation, and its normalizability. Note that the specific form of our perturbation reproduces the desired localization on the brane (see fig. 1,2) as well as the normalizability condition.
To end this letter we would like to relate the world brane Newton constant with the five dimensional Newton constant. To proceed on this direction we note that a straightforward definition is not possible since the obvious Kaluza-Klein reduction gives no terms in the effective action that could be related to the Einstein term. This is a consequence of the anisotropy of the world brane. Fortunately far from the event horizon this space time looks like AdS, therefore our brane becomes isotropic with warp factor $`A^2`$. Then we can proceed as usual to define the Newton constant.
The Newton constant, far from the horizon is essentially AdS in static coordinates plus a correction coming from the black hole:
$$M_4^2=M_5^3_0^{z_0}𝑑zA^2(r(\overline{z}))=r_0\left(1+\frac{r_0^2}{R^2}+\frac{2M}{r_0^2}\right)$$
(16)
For $`r_0=R`$ and $`MR`$
$$M_4^2=M_5^3\sqrt{\frac{64}{3\mathrm{\Lambda }}}\left(1\frac{1}{8}M\mathrm{\Lambda }\right)$$
(17)
Again the warp factor defines the hierarchy and goes like $`A^2`$.‘
Acknowledgments
We thank E. Alvarez for helpful conversations. The work of C.G. and B.J. has been supported by the TMR program FMRX-CT96-0012 on Integrability, non-perturbative effects, and symmetry in quantum field theory. The work of P.S. was partially supported by the government of Venezuela.
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# R-parity violating anomaly mediated supersymmetry breaking
## 1 Introduction
Low energy supersymmetry remains the most promising known perturbative solution to the gauge hierarchy problem that afflicts the Standard Model. It is clear from current data however, that for supersymmetry to be present in nature, it must be broken. Three phenomenologically distinct mechanisms for translating supersymmetry breaking from a hidden sector to the observable sector are currently recognised: tree-level gravity, gauge or anomaly mediation. The last mediator has received relatively little attention, and it is upon this mechanism that we focus the attention of this letter.
Anomaly mediated supersymmetry breaking (AMSB) in the minimal supersymmetric standard model (MSSM) provides a potential solution to the supersymmetric (SUSY) flavour problem . This is a problem of many supergravity theories in which squarks and sleptons typically acquire unacceptably large flavour-changing neutral currents (FCNCs) through flavour mixings in their mass matrices. In AMSB, SUSY breaking is assumed to take place in a hidden (“sequestered”) sector. A re-scaling anomaly in the super-Weyl conformal transformation transmits the SUSY breaking to the observable sector. It was suggested that the MSSM superfields be confined to a 3-brane in a higher dimensional bulk space-time, with the sequestered sector separated by the extra dimension from the visible sector brane. If direct Kahler couplings between the sequestered and visible sectors are suppressed (as is the case in the above geometrical set-up), these SUSY breaking terms can be the dominant forms of SUSY breaking in the visible sector. This scenario produces a supersymmetric spectrum dependent upon only three unknown parameters, an overall supersymmetric breaking mass scale and the MSSM Higgs potential parameters $`\mu `$ and $`B`$. For example, the AMSB mass squared values for scalar components of chiral matter supermultiplets are given by ,
$$(m^2)_{\mathrm{\Phi }_i}^{\mathrm{\Phi }_j}|_{AM}=\frac{1}{4}M_{\mathrm{aux}}^2\left[\mu \frac{d^2}{d\mu ^2}\mathrm{ln}Z_i^j\right],$$
(1)
where $`\mu `$ denotes the t’Hooft renormalization scale and $`Z_j^i`$ is the matter field wave function of the superfield $`\mathrm{\Phi }_i`$. $`M_{\mathrm{aux}}`$ is the vacuum expectation value of a compensator superfield , and sets the overall mass scale for visible sector SUSY breaking. Defining
$$\mathrm{\Gamma }_{\mathrm{\Phi }_i}^{\mathrm{\Phi }_j}1/2\frac{d(\mathrm{ln}Z_j^i)}{d\mathrm{ln}\mu },$$
(2)
we may write
$$(m^2)_{\mathrm{\Phi }_i}^{\mathrm{\Phi }_j}|_{AM}=\frac{1}{2}M_{\mathrm{aux}}^2\left[\beta (Y)\frac{}{Y}\mathrm{\Gamma }_{\mathrm{\Phi }_i}^{\mathrm{\Phi }_j}+\beta (g)\frac{}{g}\mathrm{\Gamma }_{\mathrm{\Phi }_i}^{\mathrm{\Phi }_j}\right]$$
(3)
summed over all Yukawa couplings $`Y`$ and gauge couplings $`g`$. $`\beta (x)`$ represents the beta function $`dx/d\mathrm{ln}\mu `$ of coupling $`x`$. An interesting fact is that the AMSB soft terms are valid to all orders in perturbation theory . For scalars of the first two families, Yukawa couplings can be safely neglected and so the dominant terms in eq. (3) are those proportional to the gauge couplings. These, being family universal, highly suppress the most problematic FCNC processes and thus solve the SUSY flavour problem.
The trilinear soft term $`A_Y`$ corresponding to Yukawa coupling $`Y`$ is given by
$$YA_Y=\beta (Y)M_{\mathrm{aux}},$$
(4)
and the gaugino mass $`M_g`$ associated with each gauge group of coupling $`g`$ is
$$gM_g=\beta (g)M_{\mathrm{aux}}.$$
(5)
One particularly desirable feature is that eqs. (3-5) are renormalisation invariant , provided that there are no massive parameters present in the superpotential (for example from bilinear terms). This means that AMSB has the advantage over gravity and gauge mediation that the low energy phenomenology does not depend upon the renormalisation of parameters at high scales, where unknown corrections from new physics apply. Unfortunately in the minimal version of the AMSB MSSM, the sleptons have negative mass squared values, indicating that the true vacuum state of the model is not the desired electroweak one. There have been several successful attempts to fix this problem, for example positive bulk Standard Model singlet contributions to the scalar masses , non-decoupling effects , heavy vector-like messengers coupled to light modulus fields have all been proposed. All of the above models have spoiled the desirable renormalisation group invariant feature, rendering the theories potentially sensitive to unknown ultra-violet effects.
In a recent paper , a model which solves the slepton problem in AMSB with SUSY breaking Fayet-Iliopoulos D-terms of an additional U(1) gauge symmetry was presented. Extra Standard-Model gauge singlet chiral superfields are added to the MSSM. The model has the advantage of soft terms that do not depend upon the ultra-violet physics ,. The model of ref. extends the MSSM by 3 extra Higgs doublets, a vector-like pair of extra triplets and 4 new singlets near the weak scale. Large Yukawa couplings between the extra Higgs and MSSM leptons provide additional positive contributions to the slepton mass squareds in eq. (3), while preserving renormalisation group invariance of the mass relations. We note that the above attempts to solve the AMSB slepton problem have all had $`R`$-parity invariance as a feature.
Here, we make a new proposal which preserves the renormalisation invariance of the AMSB supersymmetry breaking mass relations and requires no superfields additional to those in the MSSM coupling directly to the visible sector. By considering a subset of trilinear R-parity violating ($`\overline{)}R_p`$) operators in the superpotential, we change the wave-function renormalisation of the sparticles, in particular providing new positive contributions to the slepton mass squared values. Throughout this work we will assume the dogma of minimality with respect to solving the slepton problem in AMSB, for brevity and simplicity.
First, in section 2, we will classify models that simultaneously solve the MSSM AMSB tachyonic slepton problem while not generating dangerous lepton flavour violating operators. There emerges a scenario with 3 non-negligible lepton number violating ($`\overline{)}L`$)-operators only. In section 3, we then present the slepton masses in terms of the supersymmetric couplings explicitly and briefly discuss the other soft breaking terms. All other soft masses are equivalent to the $`R_p`$ conserving AMSB MSSM to one-loop order. In section 4, we impose constraints upon the model, the most restrictive being from lepton non-universality. For the $`\overline{)}L`$-contributions to be sufficient to raise all slepton mass squared values above zero, we require some $`\overline{)}L`$-couplings of order 1. We demonstrate that this is not in conflict with current data if the scalar sparticles are rather heavy, above 1.2 TeV. We also present the sparticle spectra. In section 5, some implications for collider searches at the Tevatron and LHC are presented. Finally, we summarise the main features of the model, reviewing its successful features and noting possible future work in section 6.
## 2 Rescuing the AMSB MSSM with R-parity violation
We begin with the AMSB MSSM including general trilinear $`\overline{)}R_p`$-operators. We then identify the subset of operators which are useful in solving the AMSB slepton problem. In the notation of ref. , the general trilinear $`\overline{)}R_p`$-MSSM superpotential is written
$`W_3`$ $`=`$ $`(Y_E)_{ij}L_iH_1\overline{E}_j+(Y_D)_{ij}Q_iH_1\overline{D}_j+(Y_U)_{ij}Q_iH_2\overline{U}_j+`$ (6)
$`{\displaystyle \frac{1}{2}}\lambda _{ijk}L_iL_j\overline{E}_k+\lambda _{ijk}^{}L_iQ_j\overline{D}_k+{\displaystyle \frac{1}{2}}\lambda _{ijk}^{\prime \prime }\overline{U}_i\overline{D}_j\overline{D}_k,`$
where we have suppressed gauge indices, and $`i,j,k,\mathrm{}=1,2,3`$ are family indices. The anomalous dimensions $`\mathrm{\Gamma }_{\mathrm{\Phi }_i}^{\mathrm{\Phi }_j}`$ relevant for substitution into eq. (3) for the superpotential $`W_3`$ have been presented in ref. to two-loop order. Their one-loop truncation is annexed here to Appendix A for ease of reference.
For example, substituting $`\mathrm{\Gamma }_{E_{R}^{}{}_{i}{}^{}}^{E_{R}^{}{}_{j}{}^{}}`$ into eq. (3), we obtain the most problematic soft mass squared, that of the right handed sleptons:
$$16\pi ^2(m_{E_R}^2)_i^j=\frac{1}{2}M_{\mathrm{aux}}^2[2(Y_E^{})_{jk}\beta (Y_E)_{ki}+\lambda _{kmi}\beta (\lambda _{kmj})+(ij)G_E],$$
(7)
where $`G_E396g_1^4/(25\times 16\pi ^2)`$. The first term on the right hand side is negligible for selectrons and smuons because it is proportional to the electron and muon mass respectively. In the R-parity conserving scenario where all $`\lambda _{ijk}=0`$, the last (negative) term therefore forces the right handed selectrons and smuons to have negative mass squared values. If $`\mathrm{tan}\beta >40`$, the positive contribution from $`(Y_E)_{33}`$ becomes non-negligible and the right handed stau mass squared may be raised above zero. However, we immediately see that a positive contribution may be obtained to $`(m_{E_R})_k^k`$ from $`\lambda _{ijk}0`$, and it is this possibility that we exploit<sup>1</sup><sup>1</sup>1Following the minimality dogma, we assume real $`\overline{)}R_p`$-couplings.. So far, the condition of minimality identifies<sup>2</sup><sup>2</sup>2We set all $`\overline{)}R_p`$operators not explicitly mentioned to zero. the combination
$$\lambda _{jk1},\lambda _{lm2},\lambda _{nq3}0,$$
(8)
which provides positive contributions to all three right handed slepton masses. We also observe that $`\lambda _{ijk}^{}`$, $`\lambda _{ijk}^{\prime \prime }`$ cannot solve the problem of negative right handed slepton masses and so we drop them from the discussion. We will assume in the present models that they are zero, and indeed this assumption will make it simpler to satisfy stringent empirical limits upon successful scenarios.
The left-handed sleptons also have negative mass squared values in the usual $`R_p`$-conserving AMSB MSSM scenario. Including LLE operators by substituting $`\mathrm{\Gamma }_{L_i}^{L_j}`$ into eq. (3) and setting $`\lambda _{ijk}^{},\lambda _{ijk}^{\prime \prime }=0`$ for all $`i,j,k`$,
$$16\pi ^2(m_L^2)_i^j=\frac{1}{2}M_{\mathrm{aux}}^2[(Y_E^{})_{jk}\beta (Y_E)_{ki}+\lambda _{ikq}\beta (\lambda _{jkq})+(ij)G_L],$$
(9)
where $`G_L(99g_1^2+75g_2^2)/(25\times 16\pi ^2)`$. Thus the R-parity conserving scenario (all $`\lambda _{ijk}=0`$) suffers from negative mass squared values for the left-handed selectron and smuon (and the stau if $`\mathrm{tan}\beta `$ is not large), analogous to the right handed sleptons. A positive contribution to $`(m_L^2)_i^i`$ results if $`\lambda _{ijk}0`$, i.e. we require
$$\lambda _{1jk},\lambda _{2lm},\lambda _{3nq}0,$$
(10)
to provide additional positive contributions to all left-handed slepton mass squareds.
In order to keep the number of couplings to a minimum, we require that the same non-zero couplings that render the right-handed slepton masses squared values positive in eq. (8) also provide us with positive left-handed slepton mass squared values in eq. (10).
We now identify a further constraint upon $`\lambda _{ijk}`$
$$\lambda _{imm}=0,i$$
(11)
(no sum on $`m`$) to avoid the generation of large off-diagonal slepton mass terms. Such terms would generate an empirically unacceptable amount of lepton flavour violation , such as $`\mu e\gamma `$. Eq. (11) also forbids the generation of lepton-Higgs mixing, as can be seen from eq. (39). Simultaneously applying the above constraints in eqs. (8-11) leads to the unique combination
$$\lambda _{123},\lambda _{132},\lambda _{231}0.$$
(12)
Thus we have identified the $`\overline{)}L`$-couplings that will solve the AMSB MSSM slepton problem without generating lepton flavour violating effects that are too large. We note that for high $`\mathrm{tan}\beta >40`$, we could set $`\lambda _{123}=0`$ and still have positive mass squared values for the sleptons. For the moment, we include all three couplings for generality, and indeed below, we focus on a particular model for which $`\mathrm{tan}\beta =4.2`$, requiring us to include $`\lambda _{123}0`$.
## 3 Soft supersymmetry breaking terms
We now discuss the (one loop) equations for the soft supersymmetry breaking terms. We work in a basis where Yukawa couplings apart from $`Y_\tau (Y_E)_{33}`$, $`Y_t(Y_U)_{33}`$, $`Y_b(Y_D)_{33}`$ (the tau, top and bottom Yukawa couplings respectively) and $`\overline{)}R_p`$-couplings not discussed above are sub-leading and are therefore neglected.
### 3.1 Slepton masses
The slepton soft masses are given by eqs. (9),(7) as
$`(m^2)_{L_1}^{L_1}`$ $`=`$ $`{\displaystyle \frac{M_{aux}^2}{(16\pi ^2)}}\left[\lambda _{123}\beta (\lambda _{123})+\lambda _{132}\beta (\lambda _{132})\left({\displaystyle \frac{3}{10}}g_1\beta (g_1)+{\displaystyle \frac{3}{2}}g_2\beta (g_2)\right)\right]`$ (13)
$`(m^2)_{L_2}^{L_2}`$ $`=`$ $`{\displaystyle \frac{M_{aux}^2}{(16\pi ^2)}}\left[\lambda _{231}\beta (\lambda _{231})+\lambda _{123}\beta (\lambda _{123})\left({\displaystyle \frac{3}{10}}g_1\beta (g_1)+{\displaystyle \frac{3}{2}}g_2\beta (g_2)\right)\right]`$ (14)
$`(m^2)_{L_3}^{L_3}`$ $`=`$ $`{\displaystyle \frac{M_{aux}^2}{(16\pi ^2)}}\left[Y_\tau \beta (Y_\tau )+\lambda _{132}\beta (\lambda _{132})+\lambda _{231}\beta (\lambda _{231})\left({\displaystyle \frac{3}{10}}g_1\beta (g_1)+{\displaystyle \frac{3}{2}}g_2\beta (g_2)\right)\right]`$
$`(m^2)_{E_1}^{E_1}`$ $`=`$ $`{\displaystyle \frac{M_{aux}^2}{(16\pi ^2)}}\left[2\lambda _{231}\beta (\lambda _{231}){\displaystyle \frac{6}{5}}g_1\beta (g_1)\right]`$ (15)
$`(m^2)_{E_2}^{E_2}`$ $`=`$ $`{\displaystyle \frac{M_{aux}^2}{(16\pi ^2)}}\left[2\lambda _{132}\beta (\lambda _{132}){\displaystyle \frac{6}{5}}g_1\beta (g_1)\right]`$ (16)
$`(m^2)_{E_3}^{E_3}`$ $`=`$ $`{\displaystyle \frac{M_{aux}^2}{(16\pi ^2)}}\left[2Y_\tau \beta (Y_\tau )+2\lambda _{123}\beta (\lambda _{123}){\displaystyle \frac{6}{5}}g_1\beta (g_1)\right]`$ (17)
where the $`\beta `$-functions are given by
$`\beta (Y_\tau )`$ $`=`$ $`{\displaystyle \frac{Y_\tau }{16\pi ^2}}\left[4Y_\tau ^2+3Y_b^2+2\lambda _{123}^2+\lambda _{132}^2+\lambda _{231}^2\left({\displaystyle \frac{9}{5}}g_1^2+3g_2^2\right)\right]`$ (18)
$`\beta (\lambda _{123})`$ $`=`$ $`{\displaystyle \frac{\lambda _{123}}{16\pi ^2}}\left[2Y_\tau ^2+4\lambda _{123}^2+\lambda _{231}^2+\lambda _{132}^2\left({\displaystyle \frac{9}{5}}g_1^2+3g_2^2\right)\right]`$ (19)
$`\beta (\lambda _{231})`$ $`=`$ $`{\displaystyle \frac{\lambda _{231}}{16\pi ^2}}\left[Y_\tau ^2+4\lambda _{231}^2+\lambda _{123}^2+\lambda _{132}^2\left({\displaystyle \frac{9}{5}}g_1^2+3g_2^2\right)\right]`$ (20)
$`\beta (\lambda _{132})`$ $`=`$ $`{\displaystyle \frac{\lambda _{132}}{16\pi ^2}}\left[Y_\tau ^2+4\lambda _{132}^2+\lambda _{123}^2+\lambda _{231}^2\left({\displaystyle \frac{9}{5}}g_1^2+3g_2^2\right)\right]`$ (21)
### 3.2 Other soft terms
The soft terms for squark masses and trilinear couplings can be derived from the general formulae in the Appendix. To one-loop accuracy, the rest are equivalent to the $`R_p`$ conserving MSSM soft terms , except for the trilinear slepton couplings and $`m_{H_1}`$:
$$m_{H_1}^2=\frac{M_{aux}^2}{16\pi ^2}\left[3Y_b\beta (Y_b)+Y_\tau \beta (Y_\tau )\left(\frac{3}{10}g_1\beta (g_1)+\frac{3}{2}g_2\beta (g_2)\right)\right].$$
(22)
From eq. (18), we see that $`m_{H_1}`$ depends upon the combination $`\kappa 2\lambda _{123}^2+\lambda _{132}^2+\lambda _{231}^2`$. Because $`\mu `$ is fixed partly by $`m_{H_1}`$ in the electroweak symmetry breaking condition , it is altered from the $`R_p`$-conserving scenario by $`\kappa 0`$. We note in particular the anomaly-mediated contribution to the $`B`$-term realised in a specific model with a bulk contribution :
$`B={\displaystyle \frac{M_{aux}}{16\pi ^2}}\left[3Y_t^2+3Y_b^2+Y_\tau ^2\left({\displaystyle \frac{3}{5}}g_1^2+3g_2^2\right)\right].`$ (23)
We shall utilise this model in order to cut the parameter space down. The prediction of $`B`$ (for a given value of $`M_{\mathrm{aux}}`$) results in a prediction of $`\mathrm{tan}\beta `$ from the potential minimisation conditions . We note that in the AMSB MSSM, a term $`\mu H_1H_2`$ in the superpotential spoils the conformal invariance. However, $`\mu `$ can be viewed as a result of supersymmetry breaking , providing a natural explanation the size of $`\mu `$ necessary to obtain $`M_Z=91.18`$ GeV. In our convention, $`B`$ and $`\mu `$ have opposite signs in successful minima, so the $`B`$ term predicted also constrains the sign of $`\mu `$ to be positive.
## 4 Spectra and constraints
Some products of the $`\overline{)}L`$-couplings are constrained to be tiny and practically useless in solving the AMSB slepton mass problem ,
$`\lambda _{mni}`$ $`/`$ $`\lambda _{mnj}ij,`$ (24)
$`\lambda _{imn}`$ $`/`$ $`\lambda _{jmn}ij,`$ (25)
where $`/`$ stands for ‘not non-zero with’. The combination of couplings $`\lambda _{123},\lambda _{132},\lambda _{231}0`$ respects this constraint. In addition, the most recent bounds upon the individual couplings are :
$`\lambda _{123}`$ $`\stackrel{<}{}`$ $`0.49\times {\displaystyle \frac{m_{\stackrel{~}{\tau }_R}}{1\mathrm{T}eV}}`$ (26)
$`\lambda _{132}`$ $`\stackrel{<}{}`$ $`0.62\times {\displaystyle \frac{m_{\stackrel{~}{\mu }_R}}{1\mathrm{T}eV}}`$ (27)
$`\lambda _{231}`$ $`\stackrel{<}{}`$ $`0.70\times {\displaystyle \frac{m_{\stackrel{~}{e}_R}}{1\mathrm{T}eV}},`$ (28)
from charged lepton universality . As we show below, these provide the most severe constraints upon the model. The couplings also pass the $`\mu e\gamma `$ conversion limits .
### 4.1 Sparticle Spectra
We now perform a one-loop accuracy numerical analysis to determine the sparticle and Higgs spectrum. The full one-loop Higgs potential is minimised at a scale $`Q`$2 TeV, where the radiative corrections to the potential are small, to determine $`\mu `$. This choice of scale can change the $`\mathrm{tan}\beta `$ prediction a little. For experimental inputs on the gauge couplings we use $`\alpha _s(M_Z)=0.119`$, $`\mathrm{sin}^2\theta _w=0.2312`$, $`\alpha (M_Z)=1/127.9`$ in the $`\overline{MS}`$ scheme. 3 loop QCD$``$2 loop QED is used as an effective field theory below $`M_Z`$. Central values of the top and bottom pole masses $`M_t=174.3,M_b=4.9`$ GeV are taken. For simplicity we set $`\lambda \lambda _{123}=\lambda _{231}=\lambda _{312}`$, a renormalisation group invariant choice in the limit that $`Y_\tau =0`$. In fact, all soft masses and couplings except those of the sleptons and one of the Higgs are independent of the choice of these three couplings. Here, we choose $`\lambda (M_Z)=0.73`$, which is sufficient to render all slepton mass squared values positive ($`\lambda >0.66`$). If $`\lambda (M_Z)`$ is set too large, then $`M_{aux}`$ must be set very large in order to produce sleptons heavy enough to evade eqs. (26-28). Remarkably, $`\lambda (M_Z)=0.73`$ is near a common quasi-fixed point value for all three couplings if we assume that they are set large at a scale $`M_{GUT}2\times 10^{16}`$ GeV, where the gauge couplings unify. The quasi-fixed behaviour is exhibited by displaying insensitivity to the ultra-violet boundary condition . This behaviour is exhibited in Fig. 1
As can be seen from the figure, $`\lambda _{132}`$ and $`\lambda _{231}`$ both approach the 0.7-0.8 region, roughly independent of the values assumed for them at $`M_{GUT}`$. We have checked that $`\lambda _{123}`$ approaches the 0.7-0.8 region also.
Fig. 2 shows the supersymmetric particle spectrum in the AMSB scenario for different values of $`M_{aux}`$. The value of $`\mathrm{tan}\beta 4.2`$ is predicted where the minimisation conditions of the Higgs potential satisfy eq. (23). In fact, this value has a small (neglected) dependence upon $`M_{aux}`$ and the scale at which the potential is minimised, and thus has an uncertainty of about $`\pm 1.0`$. The dashed region is excluded from the experimental bounds derived from charged lepton universality, the most stringent being eq. (26). An improvement of these bounds by a factor of two would test up to $`M_{aux}=400`$ TeV. The LSP neutralino is quasi-degenerate with the lightest chargino, as usual in AMSB . The Higgs mass determinations are performed using state-of-the-art two loop corrections . The lightest CP-even Higgs mass $`m_{h^0}`$ is insensitive to $`M_{\mathrm{aux}}`$ and $`\lambda (M_Z)`$, but has the usual large dependence upon $`M_t`$. For $`M_t=174.3`$ GeV however, $`M_{h^0}=117.5\pm 0.5`$ GeV, with an estimated 2 GeV uncertainty coming from higher order corrections. As can be seen from the figure, the excluded region forces all sparticles to be rather heavy - the lightest chargino and neutralino can be as light as 500 GeV, but all other sparticles and heavy Higgs must be heavier than 1.1 TeV.
## 5 Collider phenomenology
The phenomenology and search prospects of the AMSB $`R_p`$-conserving scenarios have been considered in refs. ,,,,. The present scenario differs in two main respects. Firstly, the $`\overline{)}R_p`$-coupling exclusion limits shown in fig. 2 force superpartners to be heavier than was previously considered with $`R_p`$-conserving AMSB. Secondly, the decays of $`\chi _1^\pm `$, $`\chi _1^0`$ are qualitatively different. In the $`R_p`$ conserving case, the quasi-degeneracy of $`\chi _1^\pm `$ and $`\chi _1^0`$ means that the $`\chi _1^\pm `$ is quasi-stable , and of course the $`\chi _1^0`$ is undetected, except as missing energy. A classic signature for the lightest chargino is then the presence of heavily ionising tracks, with possible slow decays into pions/leptons. In the present scenario however, the lightest chargino and neutralinos decay almost immediately into 3 leptons.
To illustrate the decays and cross sections of the model, we pick a particular value of $`M_{\mathrm{aux}}=220`$ TeV. The detailed spectrum and parameters are displayed in Table 1. HERWIG6.1 was utilised to estimate sparticle discovery prospects of this spectrum at the Tevatron and LHC. We display the cross-sections of the hard sub-process of sparticle production in Table LABEL:tab:cross. Also shown is the number of expected events for luminosities of $`=10,30`$ fb<sup>-1</sup> at the Tevatron and LHC respectively<sup>3</sup><sup>3</sup>3Equivalent to approximately 3 years running at low luminosity at the LHC. The table shows that charginos and neutralinos are produced at the LHC at a detectable rate, but the Tevatron should see no superparticles. However, we note that the lightest CP-even Higgs of mass $`m_{h^0}118`$ GeV should be discovered at the Tevatron .
We ran the weak-scale spectrum through a version of ISASUSY modified to take $`\overline{)}R_p`$-interactions into account . This then calculated the relevant decays. The lightest neutralino decays through twelve channels with equal branching ratios of 1/12 and partial widths of 2.25$`\times 10^5`$ GeV:
$`\chi _1^0`$ $``$ $`e^+\overline{\nu }_\mu \tau ^{},e^{}\nu _\mu \tau ^+,e^+\overline{\nu }_\tau \mu ^{},e^{}\nu _\tau \mu ^+,\mu ^+\overline{\nu }_e\tau ^{},\mu ^{}\nu _e\tau ^+,`$ (29)
$`\mu ^+\overline{\nu }_\tau e^{},\mu ^{}\nu _\tau e^+,\tau ^+\overline{\nu }_e\mu ^{},\tau ^{}\nu _e\mu ^+,\tau ^+\overline{\nu }_\mu e^{},\tau ^{}\nu _\mu e^+,`$
whereas $`\chi _1^+`$ has six decay channels
$$\chi _1^+\nu _\mu \nu _e\tau ^+,\nu _\tau \nu _e\mu ^+,\nu _\tau \nu _\mu e^+,\mu ^+e^+\tau ^{},\tau ^+e^+\mu ^{},\tau ^+\mu ^+e^{}$$
(30)
again with equal branching ratios of 1/6 and partial widths of 4.5$`\times 10^5`$ GeV. These decays should be easy to find with low backgrounds at the LHC. Double chargino production can be found by decays into six charged leptons, or chargino/neutralino production via a five charged lepton channel, with distinctive flavour structure. Lepton flavour violation is usually explicit in the final state. We have obtained approximately equal branching ratios here mainly because we have assumed $`\lambda _{123}=\lambda _{231}=\lambda _{132}`$. In the case they are non-degenerate, this will change and the relative branching ratios into different final states will help measure the $`\overline{)}R_p`$-couplings.
## 6 Conclusions
We have proposed a new solution to the problem of slepton negative mass squared values in the AMSB MSSM. It involves including 3 $`\overline{)}L`$-operators in the superpotential which were previously assumed to be absent. This leads to positive mass squared values for all of the sleptons and renormalisation-group invariant relations between supersymmetry breaking terms and the measured supersymmetric couplings. This has the advantage of rendering the model insensitive to unknown ultra-violet effects. The $`\mu `$ problem has a natural solution , indeed the prediction of the $`B`$term in a particular model results in a prediction for $`\mathrm{tan}\beta `$. All of the sparticle spectrum except for the slepton masses are then given in terms of two free parameters: $`M_{\mathrm{aux}}`$ and $`\kappa `$.
We have therefore assumed the MSSM spectrum near the weak scale, and that the dominant source of supersymmetry breaking terms in the observable sector is from anomaly mediation. Experimental limits on the $`\overline{)}L`$-operators provide stringent constraints upon the model, meaning that sparticle masses must be rather high. Whereas the lightest charginos and neutralinos can be as light as 500 GeV, the other sparticles must be heavier than 1.1 TeV. The Tevatron therefore sees no new particles except the lightest Higgs of mass 118 GeV, and the LHC can detect the lightest charginos and neutralinos via distinctive leptonic decays. Charged lepton universality violation is predicted to be close to its experimental bound, within a factor of two.
Neutrino masses and mixings are beyond the scope of this paper, but it is well known that the $`\overline{)}L`$-operators we have introduced can generate them at the loop level . We intend to pursue them in future work , and the small number of free parameters should allow a strict correlation with lepton flavour violating predictions. We believe the present model of supersymmetry breaking in the observable sector to warrant several future works. For example, a more accurate calculation of the spectrum and a determination of the LHC reach in parameter space would be useful. It would be desirable to find symmetries to ban or suppress the other $`\overline{)}R_p`$-couplings. Aside from these, the usual calculations in the MSSM (quark FCNCs, charge and colour breaking constraints etc) could be performed. It will be interesting to investigate the present idea in a more general framework, for example when $`\mathrm{tan}\beta `$ is large (the prediction of the lightest MSSM Higgs mass is likely to change from the one presented here) and there are only two LLE couplings, or where splittings between the $`\overline{)}R_p`$couplings occur. Relaxing the assumption that $`\lambda _{ijk}^{}=0`$ might lead to the possible observation of a single slepton at the LHC via slepton-strahlung .
To summarise, our scenario is a predictive scheme of supersymmetry breaking, containing natural solutions to the $`\mu `$ problem and supersymmetric flavour problem. The spectrum depends upon only two parameters apart from the slepton masses. In the case that the $`\overline{)}R_p`$-couplings are at their quasi-fixed point values, the slepton masses approximately only depend upon these same two parameters. If one assumes a high-energy cut-off scale, such as the GUT scale for example, we note that the weak-scale values of the couplings are approximately predicted by the quasi-fixed point and there is only one free parameter on which the whole sparticle spectrum depends. $`\mathrm{tan}\beta `$ is predicted in a specific model and the soft masses are given by renormalisation group invariant relations with the measured SUSY couplings. The phenomenology is rather distinctive and should be easily disentangled from other possibilities at the LHC, after being tested at the Tevatron by the Higgs mass prediction. The present model is the only example of a model with both insensitivity of the soft terms to unknown ultra-violet physics and the MSSM spectrum near the weak scale, and as such is important to investigate further.
## Appendix A One-loop anomalous dimensions and beta functions in the $`\overline{)}R_p`$-MSSM
$`\mathrm{\Lambda }_{U^i},\mathrm{\Lambda }_{D^i},\mathrm{\Lambda }_{E^i}`$ were written in a matrix notation in as
$`(\mathrm{\Lambda }_{E^k})_{ij}=\lambda _{ijk},(\mathrm{\Lambda }_{D^k})_{ij}=\lambda _{ijk}^{},(\mathrm{\Lambda }_{U^k})_{ij}=\lambda _{ijk}^{\prime \prime }`$ (31)
and we adopt this notation for presenting results with general family structures. The one-loop anomalous dimensions of the $`\overline{)}R_p`$-MSSM are
$`16\pi ^2\mathrm{\Gamma }_{L_i}^{(1)L_j}`$ $`=`$ $`\left(Y_EY_E^{}\right)_{ji}+(\mathrm{\Lambda }_{E^q}\mathrm{\Lambda }_{E^q}^{})_{ji}+3(\mathrm{\Lambda }_{D^q}\mathrm{\Lambda }_{D^q}^{})_{ji}\delta _i^j({\displaystyle \frac{3}{10}}g_1^2+{\displaystyle \frac{3}{2}}g_2^2),`$ (32)
$`16\pi ^2\mathrm{\Gamma }_{E_i}^{(1)E_j}`$ $`=`$ $`2\left(Y_E^{}Y_E\right)_{ji}+\text{Tr}(\mathrm{\Lambda }_{E^j}\mathrm{\Lambda }_{E^i}^{})\delta _i^j({\displaystyle \frac{6}{5}}g_1^2),`$ (33)
$`16\pi ^2\mathrm{\Gamma }_{Q_i}^{(1)Q_j}`$ $`=`$ $`\left(Y_DY_D^{}\right)_{ji}+\left(Y_UY_U^{}\right)_{ji}+(\mathrm{\Lambda }_{D^q}^{}\mathrm{\Lambda }_{D^q})_{ij}`$ (34)
$`\delta _i^j({\displaystyle \frac{1}{30}}g_1^2+{\displaystyle \frac{3}{2}}g_2^2+{\displaystyle \frac{8}{3}}g_3^2),`$
$`16\pi ^2\mathrm{\Gamma }_{D_i}^{(1)D_j}`$ $`=`$ $`2\left(Y_D^{}Y_D\right)_{ij}+2\text{Tr}(\mathrm{\Lambda }_{D^i}^{}\mathrm{\Lambda }_{D^j})+2(\mathrm{\Lambda }_{U^q}\mathrm{\Lambda }_{U^q}^{})_{ji}`$ (35)
$`\delta _i^j({\displaystyle \frac{2}{15}}g_1^2+{\displaystyle \frac{8}{3}}g_3^2)),`$
$`16\pi ^2\mathrm{\Gamma }_{U_i}^{(1)U_j}`$ $`=`$ $`2\left(Y_U^{}Y_U\right)_{ij}+\text{Tr}(\mathrm{\Lambda }_{U^j}\mathrm{\Lambda }_{U^i}^{})\delta _i^j({\displaystyle \frac{8}{15}}g_1^2+{\displaystyle \frac{8}{3}}g_3^2)),`$ (36)
$`16\pi ^2\mathrm{\Gamma }_{H_1}^{(1)H_1}`$ $`=`$ $`\text{Tr}\left(3Y_DY_D^{}+Y_EY_E^{}\right)({\displaystyle \frac{3}{10}}g_1^2+{\displaystyle \frac{3}{2}}g_2^2),`$ (37)
$`16\pi ^2\mathrm{\Gamma }_{H_2}^{(1)H_2}`$ $`=`$ $`3\text{Tr}\left(Y_UY_U^{}\right)({\displaystyle \frac{3}{10}}g_1^2+{\displaystyle \frac{3}{2}}g_2^2),`$ (38)
$`16\pi ^2\mathrm{\Gamma }_{L_i}^{(1)H_1}`$ $`=`$ $`16\pi ^2\mathrm{\Gamma }_{H_1}^{(1)L_i}{}_{}{}^{}=3(\mathrm{\Lambda }_{D^q}^{}Y_D)_{iq}(\mathrm{\Lambda }_{E^q}^{}Y_E)_{iq}.`$ (39)
The beta functions of the couplings appearing in the superpotential in eq. (6) are:
$`\beta (Y_E)_{ij}`$ $`=`$ $`(Y_E)_{ik}\mathrm{\Gamma }_{E_k}^{E_j}+(Y_E)_{ij}\mathrm{\Gamma }_{H_1}^{H_1}(\mathrm{\Lambda }_{E^j})_{ki}\mathrm{\Gamma }_{L_k}^{H_1}+(Y_E)_{kj}\mathrm{\Gamma }_{L_k}^{L_i},`$ (40)
$`\beta (Y_D)_{ij}`$ $`=`$ $`(Y_D)_{ik}\mathrm{\Gamma }_{D_k}^{D_j}+(Y_D)_{ij}\mathrm{\Gamma }_{H_1}^{H_1}(\mathrm{\Lambda }_{D^j})_{ki}\mathrm{\Gamma }_{L_k}^{H_1}+(Y_D)_{kj}\mathrm{\Gamma }_{Q_k}^{Q_i},`$ (41)
$`\beta (Y_U)_{ij}`$ $`=`$ $`(Y_U)_{ik}\mathrm{\Gamma }_{U_k}^{U_j}+(Y_U)_{ij}\mathrm{\Gamma }_{H_2}^{H_2}+(Y_U)_{kj}\mathrm{\Gamma }_{Q_k}^{Q_i},`$ (42)
$`\beta (\mathrm{\Lambda }_{E^k})_{ij}`$ $`=`$ $`(\mathrm{\Lambda }_{E^l})_{ij}\mathrm{\Gamma }_{E_l}^{E_k}+(\mathrm{\Lambda }_{E^k})_{il}\mathrm{\Gamma }_{L_l}^{L_j}+(Y_E)_{ik}\mathrm{\Gamma }_{H_1}^{L_j}`$ (43)
$`(\mathrm{\Lambda }_{E^k})_{jl}\mathrm{\Gamma }_{L_l}^{L_i}(Y_E)_{jk}\mathrm{\Gamma }_{H_1}^{L_i},`$
$`\beta (\mathrm{\Lambda }_{D^k})_{ij}`$ $`=`$ $`(\mathrm{\Lambda }_{D^l})_{ij}\mathrm{\Gamma }_{D_l}^{D_k}+(\mathrm{\Lambda }_{D^k})_{il}\mathrm{\Gamma }_{Q_l}^{Q_j}+(\mathrm{\Lambda }_{D^k})_{lj}\mathrm{\Gamma }_{L_l}^{L_i}(Y_D)_{jk}\mathrm{\Gamma }_{H_1}^{L_i},`$ (44)
$`\beta (\mathrm{\Lambda }_{U^i})_{jk}`$ $`=`$ $`(\mathrm{\Lambda }_{U^i})_{jl}\mathrm{\Gamma }_{D_l}^{D_k}+(\mathrm{\Lambda }_{U^i})_{lk}\mathrm{\Gamma }_{D_l}^{D_j}+(\mathrm{\Lambda }_{U^l})_{jk}\mathrm{\Gamma }_{U_l}^{U_i}.`$ (45)
###### Acknowledgments.
BCA would like to thank D.R.T. Jones, P. Richardson and M.A. Parker for useful discussions and the IFIC, University of Valencia for hospitality offered while some of this work was carried out. We thank H. Dreiner for useful discussions and comments on the draft. AD is supported from Marie Curie Research Training Grants ERB-FMBI-CT98-3438 and thanks DAMTP, University of Cambridge for hospitality. This work has been partially supported by PPARC.
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# Motional sidebands and direct measurement of the cooling rate in the resonance fluorescence of a single trapped ion
## Abstract
Resonance fluorescence of a single trapped ion is spectrally analyzed using a heterodyne technique. Motional sidebands due to the oscillation of the ion in the harmonic trap potential are observed in the fluorescence spectrum. From the width of the sidebands the cooling rate is obtained and found to be in agreement with the theoretical prediction.
Since the first preparation of a single atom in a Paul trap and observation of its resonance fluorescence , investigation of this light has revealed a range of unique properties. Examples are its nonclassical nature and the highly nonlinear response, in the form of sudden intensity jumps, of a multi-level atom to continuous laser excitation . The fluorescence is, at the same time, a unique tool for determining the state of the atom. This is particularly obvious for a single particle where each photon emission marks the respective projection of the atomic wave function into the final state of the corresponding transition. It is also of great interest to study, through its resonance fluorescence, the motion of a single laser-excited particle, e.g. for investigating laser cooling schemes or in connection with proposals for quantum state manipulation or quantum information processing with trapped particles .
The spectrum of fluorescence of a motionless two-level atom exhibits an elastic part (Rayleigh peak) which is $`\delta `$-correlated with the exciting light, and an inelastic contribution, the Mollow triplet . While the latter marks spontaneous transitions between the dressed states of the combined atom-light quantum system and appears when the light intensity approaches saturation, the elastic contribution dominates for light levels below saturation. For a free atom this spectrum is modified by the recoil shift. For a trapped atom the elastic peak is unshifted, and sidebands appear at the characteristic frequencies of the motion in the trap , with sizes that depend on the amplitude of this motion. The spectral width of these sidebands reflects the effective decay of the motional states of an atom in the trap, i.e. the rate of transitions which change this state. In particular, if laser excitation provides optical cooling of the trapped atom, the sideband width reflects the equilibrium of heating and cooling transitions in the steady state. While such sidebands have been observed in the fluorescence spectrum of ensembles of trapped neutral atoms , their investigation for a single particle and analysis of their width has not been done so far . The measurement of cooling rates is highly interesting in experiments with cold atoms or ions, in particular when many levels or several light fields are involved such that an optimal set of parameters is hard to find solely from theoretical arguments.
In this paper we report on a measurement of the resonance fluorescence of a single trapped Barium ion which reveals sidebands of the elastically scattered light due to the various components of the ion’s motion in the trapping potential. From the width of the sidebands which correspond to one of its vibrational modes in the Paul trap quasi-potential we deduce the cooling rate induced by the exciting laser. This method will enable us to perform detailed studies of motional effects of laser radiation in situ, i.e. during the laser excitation without further analysis tools.
In the experiment, a single Ba<sup>+</sup> ion is trapped in a 1 mm diameter Paul trap. The ion is generated by impact ionization of a weak thermal Ba atomic beam with an electron beam inside the trap. The trap is suspended in UHV and driven with a 500 V<sub>pp</sub> radio frequency signal at $`f_{Paul}19`$ MHz. The ion is laser-cooled by simultaneous excitation on its S<sub>1/2</sub> $``$ P<sub>1/2</sub> and P<sub>1/2</sub> $``$ D<sub>3/2</sub> resonance lines at 493.4 nm and 649.7 nm, respectively . See Fig. 1 for the relevant levels of Ba<sup>+</sup>. The 493 nm light is produced by a frequency doubled diode laser with external grating resonator described in Ref. . This laser is frequency stabilized to a Te<sub>2</sub> resonance line 666 MHz away from the Ba<sup>+</sup> line, by modulation transfer spectroscopy . The 650 nm light is generated by a diode laser with an external grating resonator, stabilized to an optical resonator. Both lasers have linewidths well below 100 kHz. The laser beams are combined on a dichroic beamsplitter before they are focused into the trap, and both light fields are linearly polarized. The laser intensities at the position of the ion are in the range of 200 mW/cm<sup>2</sup> (493 nm) and 100 mW/cm<sup>2</sup> (650 nm). The 650 nm laser is close to resonance, the 493 nm laser is red-detuned by about the transition linewidth ($`\mathrm{\Gamma }=15.1`$ MHz) for Doppler cooling. A 2.8 Gauss magnetic field which is orthogonal to both the laser wave vector and the laser polarization defines a quantization axis and lifts the degeneracy of the Zeeman sublevels. The precise parameters are determined by fitting an 8-level Bloch equation calculation to a scan of the fluorescence intensity vs. the detuning of the 650 nm laser .
The 493 nm resonance fluorescence of the ion is analyzed with a heterodyne detection setup shown in Fig. 1 . Using a slightly simpler scheme, with a frequency shift in only one arm of the heterodyne setup and with a frequency selective photo detector, Höffges et al. have observed the elastic part of an ion’s resonance fluorescence, but no motional sidebands. In our setup, the fluorescence at right angle to the direction of the laser beams, in the direction of the magnetic field, is collimated with an f/0.7 lens. This light consists of the two $`\sigma `$-polarized components corresponding to the P$`{}_{1/2}{}^{}(m=\pm 1/2)`$ to S$`{}_{1/2}{}^{}(m=1/2)`$ transitions, and its elastic part is linearly polarized, because the two $`\sigma `$-polarized components are coherently superimposed. For creating a local oscillator beam, the green laser light is divided into two beams which are frequency shifted by 112.5 MHz (beam 1) and 80 MHz (beam 2) using acousto-optical modulators (AOMs). Beam 1 excites the ion while beam 2 is superimposed with the collimated fluorescence on a polarizing beam splitter PBS1. The beam sizes of the collimated fluorescence and the local oscillator (beam 2) are adjusted for optimum overlap, using telescopes. Their orthogonal polarizations after PBS1 are mixed with a $`\lambda `$/2 plate and a second polarizing beam splitter PBS2 to create an interference signal. The two output signals of PBS2 are detected on two fast photodiodes (50 MHz bandwidth, 80$`\%`$ quantum efficiency). By subtracting their photodiode currents, the interference term $`SE_{lo}E_{fluor}`$ between the fluorescence and the local oscillator is filtered out. Due to the frequency shifts in beams 1 and 2, this interference product (i.e. its elastic part) is expected at a frequency of 32.5 MHz, the difference frequency of the two AOM drives, while any motional sideband due to an oscillatory motion with frequency $`f_s`$ would appear at 32.5 MHz $`\pm nf_s,n=1,2,..`$. The reason for using two AOMs is that their difference frequency is not generated elsewhere in the setup, such that no electrical crosstalk perturbs the final signal, neither can any residual amplitude modulation in one of the beams create a 32.5 MHz signal on the photodiodes. The interference signal $`S`$ is mixed with an rf reference signal at variable frequency $`f_{mix}`$ around 32.5 MHz, low-pass filtered, and finally analyzed on a Fourier spectrum analyzer with 100 kHz maximum bandwidth. All rf sources, i.e. the trap drive, the AOM supplies, and the rf reference, are phase locked to the same 10 MHz master oscillator.
With the mixer frequency $`f_{mix}`$ set to $`f_0=32.45`$ MHz, the elastically scattered light produces a signal at 50 kHz on the spectrum analyzer. Such a signal, with the resolution bandwidth set to 61 mHz, is shown in Fig. 2. Clearly, only one data point is significantly above the background noise, thus verifying the $`\delta `$-correlation between exciting laser and fluorescence. The signal-to-noise ratio (SNR) is 17 dB (at unit bandwidth). The maximum SNR which can be achieved depends on the rate $`N`$ of fluorescence photons that are detected. It is derived in the following way: The spectral power of the signal $`S`$ is $`P_SS^2/\mathrm{\Delta }\nu `$, where $`\mathrm{\Delta }\nu `$ is the detection bandwidth (or signal bandwidth, whichever dominates), while the noise power is $`P_NE_{lo}^2`$, such that their ratio is $`\mathrm{SNR}_{max}=E_{fluor}^2/\mathrm{\Delta }\nu =N/\mathrm{\Delta }\nu `$. For this to hold, the noise created by $`E_{lo}`$ with no fluorescence present has to be the dominant noise in the photodiode signal $`S`$. By varying the local oscillator power without a fluorescence signal we confirmed that this is the case in our experiment and that the noise is close the the quantum limit. With a typical scattering rate of $`2.55\times 10^4`$ photons/s into the solid angle that is collimated, and with $`80\%`$ photodiode quantum efficiency, the maximum possible SNR turns out to be 35-40 dB. The comparatively low value of 17 dB which we find is due to phase front distortions induced by the collimating lens and a resulting low degree of mode matching between the collimated fluorescence and the local oscillator.
Due to the quadrupole radiofrequency field in a Paul trap, the ion undergoes a driven oscillation at the frequency $`f_{Paul}=18.53`$ MHz. This so-called micromotion is in phase with the driving field, and its amplitude is proportional to the ion’s distance from the trap center. The sidebands to the elastic peak which this oscillation generates can be observed on the spectrum analyser, in the same manner as the carrier, by setting $`f_{mix}`$ to $`f_0\pm nf_{Paul},n=1,2,..`$. In Fig. 3 we show the results of such measurements. Due to the fact that the phase of the micromotion is well-defined, and because all rf sources are phase locked, the width of the micromotion sidebands is limited by the resolution bandwidth of the spectrum analyser, just like the width of the elastic peak in Fig. 2.
The sizes of the micromotion sidebands are expected to be proportional to $`|J_n(m)|^2`$ where $`J`$ is a Bessel function, $`n=0,\pm 1,\pm 2,\mathrm{}`$ is the number of the sideband, and $`m`$ is the modulation index corresponding to the ion’s oscillation. With the micromotion described by $`\stackrel{}{a}\mathrm{sin}\mathrm{\Omega }t`$, ($`\mathrm{\Omega }=2\pi f_{Paul}`$) and with $`\stackrel{}{k_l}`$ and $`\stackrel{}{k_d}`$ describing the laser and fluorescence wave vectors, respectively, $`m`$ is given by $`m=\stackrel{}{a}(\stackrel{}{k_d}\stackrel{}{k_l})`$ . The particular value of $`m`$ found from the sidebands in Fig. 3 is 0.47. This corresponds to a micromotion amplitude in the direction of $`\stackrel{}{k_l}\stackrel{}{k_d}`$ of about 26 nm. By minimising the micromotion, i.e. the value of $`m`$, with the aid of such a measurement, the ion can be placed in the trap center where the trapping conditions are optimal . At optimum conditions, i.e. if the SNR reaches 40 dB, and $`\stackrel{}{a}\stackrel{}{k_d}\stackrel{}{k_l}`$, this measurement is sensitive to micromotion amplitudes of about 1 nm.
Apart from the forced micromotion oscillation, the three-dimensional trapping in the rf-generated pseudo-potential of a Paul trap results in three independent modes of free vibration along orthogonal axes at three (in general different) frequencies $`f_{macro}`$. In our case, these frequencies are 620.5, 670, 1301 kHz. Investigation of this so-called macromotion is our main purpose, because these motional degrees of freedom interact with the laser in a cooling process, and the corresponding sidebands contain information about the motional state of the ion, the efficiency of the cooling, and their dependences on the laser parameters. It is also the macromotion which is used in experiments on quantum state manipulation and entanglement, in particular in the framework of quantum information and quantum computation .
The macromotion is harder to detect in the heterodyne signal than the micromotion because it is not correlated with any applied rf source. In view of the noise considerations above, detection of a macromotion sideband requires that the rate of photons which contribute to the sideband heterodyne signal is larger than the spectral width of that sideband. Due to the inefficient mode matching described above, this situation is not realized in our present setup. However, the macromotion corresponds to a damped harmonic oscillator (laser cooling being the damping mechanism) which can be driven with an additional external field at some frequency $`f_{drive}`$ close to the macromotion frequency $`f_{macro}`$. Then, with $`f_{mix}`$ set to $`f_0\pm nf_{drive}`$, the excited macromotion is detectable as a $`\delta `$-signal on the spectrum analyser, and a scan of $`f_{drive}`$ over one of the macromotion resonances reveals the response of that oscillator mode to the drive, in particular the broadening of the resonance due to laser cooling. The result of such a measurement is shown in Fig. 4. Here, the upper trace corresponds to the height of the elastic peak ($`f_{mix}=f_0`$) as in Fig. 2, while the lower trace corresponds to the height of the sideband at $`f_{drive}`$ ($`f_{mix}=f_0+f_{drive}`$), both as functions of $`f_{drive}`$ around the 620.5 kHz macromotion mode frequency. The driving voltage was applied to an electrode about 1 mm away from the ion, its power was -60 dBm. The axis of the excited vibrational mode is at 45<sup>o</sup> between the direction of the laser and the observation, as was found from the image of the ion at much higher drive power.
Fig. 4 shows how the elastic scattering decreases while the sideband scattering increases around the macromotion resonance. The two traces were fitted with functions $`|A_nJ_n(m(f_{drive}))|^2,n=0`$ for the carrier, $`n=1`$ for the sideband, where $`m(f_{drive})`$ was assumed as
$$m(f_{drive})=\frac{m_{max}}{1+(\frac{f_{drive}f_{macro}}{\mathrm{\Delta }f/2})^2}$$
(1)
as expected for a damped harmonic oscillator. We find a maximum modulation index $`m_{max}=1.5`$ and a width of the macromotion resonance $`\mathrm{\Delta }f=750`$ Hz. This width corresponds to the effective linewidth that the ion’s oscillator eigenstates in the trap acquire due to the ongoing laser excitation which makes the ion change its motional state .
To compare this cooling rate with the one derived from a simple model, we describe the motion of the ion by a driven and damped harmonic oscillator. Some limiting conditions are fulfilled in our experiment: The cooling rate is much smaller than the linewidth of the transition $`\mathrm{\Gamma }`$, i.e. the velocity changes only negligibly during one lifetime, and both the recoil frequency $`\mathrm{}k^2/2M=2\pi \times 5.9`$ kHz, $`M`$ being the ion mass, and the maximum Doppler shift $`kv_{max}`$ are smaller than the oscillation frequency (the latter being the Lamb-Dicke condition). Therefore we can calculate the cooling rate, i.e. the damping coefficient, from the radiation pressure that acts on the ion during its oscillation
$$F(v)=\mathrm{}k\mathrm{\Gamma }P_\mathrm{P}(\mathrm{\Delta }kv),$$
(2)
where $`P_\mathrm{P}`$ is the probability of finding the ion in the P<sub>1/2</sub> state and $`\mathrm{\Delta }`$ is the detuning of the 493 nm laser from the atomic resonance frequency. With $`\mathrm{\Delta }`$ negative and of the order of $`\mathrm{\Gamma }`$, and for $`kv\mathrm{\Gamma },|\mathrm{\Delta }|`$ the velocity-dependent part of $`F(v)`$ amounts, in first order, to a friction force $`F_f=\alpha Mv`$ that the ion experiences , and
$$\alpha =2\frac{\mathrm{}k^2}{2M}\mathrm{\Gamma }\frac{dP_\mathrm{P}}{d\mathrm{\Delta }}$$
(3)
corresponds to the linewidth induced by the laser cooling. All parameters in Eq.(3) refer to the 493 nm laser. We can neglect the contribution of the 650 nm laser to the cooling because this laser is set to yield maximum fluorescence ($`\mathrm{\Delta }_{650}=2\pi \times 5`$ MHz), such that $`dP_\mathrm{P}/d\mathrm{\Delta }_{650}0`$. By calculating $`dP_\mathrm{P}/d\mathrm{\Delta }`$ from the same 8-level Bloch equations that determine our experimental parameters $`\mathrm{\Delta }_{493}=2\pi \times 19`$ MHz, $`I_{493}=189\mathrm{mW}/\mathrm{cm}^2`$, $`I_{650}=107\mathrm{mW}/\mathrm{cm}^2`$, we get $`\alpha =2\pi \times 640`$ Hz which agrees well with the measurement.
In conclusion, the spectrum of resonance fluorescence from a single harmonically confined ion in a Paul trap was observed using heterodyne detection. Aside from the elastic peak, the spectrum exhibits sidebands due to the weakly excited macromotion of the ion in the trap. Investigation of these sidebands allows for an analysis of the cooling rate, without affecting the motional state of the ion by, e.g., probing pulses. The measured cooling rate in our experiment compares well with a simple model calculation. Moreover, micromotion sidebands are observed which, by minimizing their amplitude, can be used to actively compensate for residual micromotion down to the nm-level. These techniques will prove useful in the context of precision spectroscopy with ion traps and in particular for the preparation and manipulation of vibrational quantum states of motion as they are required for quantum information experiments with trapped ions.
We acknowledge support by the Fonds zur Förderung der wissenschaftlichen Forschung (FWF) (project P11467-PHY) and by the EC (TMR network ”Quantum Structures”, ERB-FMRX-CT96-0077).
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# Spontaneous deformation of the Fermi surface due to strong correlation in the two-dimensional 𝑡-𝐽 model
## Abstract
Fermi surface of the two-dimensional $`t`$-$`J`$ model is studied using the variational Monte Carlo method. We study the Gutzwiller projected $`d`$-wave superconducting state with an additional variational parameter $`t_\mathrm{v}^{}`$ corresponding to the next-nearest neighbor hopping term. It is found that the finite $`t_\mathrm{v}^{}<0`$ gives the lowest variational energy in the wide range of hole-doping rates. The obtained momentum distribution function shows that the Fermi surface deforms spontaneously. It is also shown that the van Hove singularity is always located very close to the Fermi energy. Using the Gutzwiller approximation, we show that this spontaneous deformation is due to the Gutzwiller projection operator or the strong correlation.
71.10.Fd, 71.10.Pm, 79.60.-i
The effect of strong correlation is one of the most important issues for understanding the high-$`T_c`$ superconductivity (SC). Among various anomalous electronic properties, the experiments of angle resolved photoemission spectroscopy (ARPES) have revealed that a flat band around ($`\pi `$, 0) and (0, $`\pi `$) is pinned just below the Fermi energy. This phenomenon is unexpected in the band calculations and it is considered to be closely related to the opening of the pseudogap on the Fermi surface (FS) , which is also an extraordinary feature in high-$`T_c`$ cuprates. This anomalous nature of the FS will be the direct evidence for the non-Fermi liquid behavior. It is thus an interesting issue to study the FS in the presence of strong correlation.
The effect of the flat band and the geometry of the FS can be taken into account by using the $`t`$-$`t^{}`$-$`J`$ model or the $`t`$-$`t^{}`$-$`U`$ Hubbard model in which the next-nearest neighbor hopping term $`t^{}`$ is introduced as a fitting parameter. If one chooses $`t^{}<0`$, the FS centered at ($`\pi `$, $`\pi `$) observed experimentally can be reproduced in the tight-binding model. However high temperature expansion studies on the momentum distribution function for the $`t`$-$`J`$ model have shown that the FS is similar to that with $`t^{}<0`$ even though the $`t^{}`$-term is absent in the Hamiltonian. On the other hand, the conventional mean-field theories, such as slave-boson theory, simply give the FS with $`t^{}=0`$. Therefore the strong correlation which is not included in the mean-field theories will be the origin of the change of the FS geometry.
Here we study this problem from a different point of view. Since the calculation in Ref. is carried out in the high temperature region, it is not clear whether or not the FS deforms down to zero temperature. To study the FS of the ground state is generally very difficult. The exact diagonalization study of small clusters does not give enough resolution in the $`𝒌`$ space. The quantum Monte Carlo simulations have been often useless for the two-dimensional $`t`$-$`J`$ model because of the minus sign problem. Therefore we use the variational Monte Carlo (VMC) method in this paper, which is free from the limitation of the system size as well as from the sign problem. The VMC method treats exactly the constraints of no doubly occupied sites and gives accurate estimates of the expectation values such as the variational energies and the momentum distribution functions.
Although it is a variational theory, the VMC method is powerful to see whether some kind of symmetry breaking takes place or not. In this paper we examine the Gutzwiller-projected $`d`$-wave superconducting state which contains an additional variational parameter $`t_\mathrm{v}^{}`$ corresponding to the next-nearest neighbor hopping term. We can safely discuss the relative energy difference between the variational states with and without $`t_\mathrm{v}^{}`$, although the absolute values of the variational energies can still be lowered. We find that the wave function with $`t_\mathrm{v}^{}0.1`$ has the lowest variational energy even though the Hamiltonian does not contain $`t^{}`$-term. This means that the FS deforms spontaneously. The momentum distribution function $`n(𝒌)`$ calculated in the optimized wave function is consistent with that in the high temperature expansion. Our method gives an independent and complementary support of the result that the deformation of the FS is a distinctive feature of strongly correlated electron systems.
In addition to this, we can identify the physical origin of the FS deformation in our variational approach. We show the relation between the energy gain and the van Hove singularity. It has been argued that a remarkable enhancement of SC correlation is achieved if the van Hove singularity is close to the Fermi energy. Our results show some similarity to this picture. Furthermore, by comparing the obtained results with the Gutzwiller approximation, we can see that the finite $`t_\mathrm{v}^{}`$ is caused solely by the Gutzwiller projection.
We use the two-dimensional $`t`$-$`J`$ model on a square lattice,
$$H=t\underset{ij\sigma }{}P_\mathrm{G}(c_{i\sigma }^{}c_{j\sigma }+h.c.)P_\mathrm{G}+J\underset{ij}{}𝑺_i𝑺_j,$$
(1)
where $`ij`$ represents the sum over the nearest-neighbor sites. $`c_{i\sigma }^{}`$ ($`c_{i\sigma }`$) is a creation (annihilation) operator of $`\sigma `$ ($``$ or $``$) electron at $`i`$-site and $`𝑺_i=c_{i\alpha }^{}(\frac{1}{2}𝝈)_{\alpha \beta }c_{i\beta }`$. The Gutzwiller’s projection operator $`P_\mathrm{G}`$ is defined as $`P_\mathrm{G}=\mathrm{\Pi }_i(1\widehat{n}_i\widehat{n}_i)`$, which prohibits the doubly occupied sites. We set $`J/t=0.3`$.
We use a Gutzwiller-projected mean-field type wave function $`P_\mathrm{G}P_{N_\mathrm{e}}|\varphi _0`$ as a trial state with fixing the number of electrons $`N_\mathrm{e}`$ through $`P_{N_\mathrm{e}}`$. The state is written as
$`P_\mathrm{G}P_{N_\mathrm{e}}|\varphi _0`$ $`=`$ $`P_\mathrm{G}P_{N_\mathrm{e}}{\displaystyle \underset{k}{}}(u_k+v_kc_k^{}c_k^{})|0`$ (2)
$`=`$ $`P_\mathrm{G}P_{N_\mathrm{e}}{\displaystyle \underset{k}{}}u_k\mathrm{exp}\left[{\displaystyle \underset{k}{}}{\displaystyle \frac{v_k}{u_k}}c_k^{}c_k^{}\right]|0`$ (3)
$`=`$ $`P_\mathrm{G}P_{N_\mathrm{e}}{\displaystyle \underset{k}{}}u_k\mathrm{exp}\left[{\displaystyle \underset{ij}{}}a_{ij}c_i^{}c_j^{}\right]|0`$ (4)
$`=`$ $`P_\mathrm{G}{\displaystyle \underset{k}{}}u_k{\displaystyle \frac{1}{(N_\mathrm{e}/2)!}}\left({\displaystyle \underset{ij}{}}a_{ij}c_i^{}c_i^{}\right)^{N_\mathrm{e}/2}|0,`$ (5)
where $`v_k/u_k=\mathrm{\Delta }_k/(ϵ_k\mu +\sqrt{(ϵ_k\mu )^2+\mathrm{\Delta }_k^2})`$ and $`a_{ij}`$ is a Fourier transform of $`v_k/u_k`$.
Usually $`ϵ_k`$ is chosen to be $`ϵ_k=2(\mathrm{cos}k_x+\mathrm{cos}k_y)`$, which is in accordance with the Hamiltonian (1). However in this paper, we introduce an additional variational parameter $`t_\mathrm{v}^{}`$ which changes the FS of the variational state. We assume $`ϵ_k`$ and $`\mathrm{\Delta }_k`$ as
$`ϵ_k`$ $`=`$ $`2(\mathrm{cos}k_x+\mathrm{cos}k_y)4t_\mathrm{v}^{}\mathrm{cos}k_x\mathrm{cos}k_y,`$ (6)
$`\mathrm{\Delta }_k`$ $`=`$ $`2\mathrm{\Delta }_d(\mathrm{cos}k_x\mathrm{cos}k_y).`$ (7)
The present wave function contains three variational parameters $`t_\mathrm{v}^{}`$, $`\mu `$ and $`\mathrm{\Delta }_d`$.
At first, using the above wave function we calculate the variational energy $`E_{\mathrm{var}}`$ of the Hamiltonian (1)
$$E_{\mathrm{var}}=\frac{\varphi _0|P_{N_\mathrm{e}}P_\mathrm{G}HP_\mathrm{G}P_{N_\mathrm{e}}|\varphi _0}{\varphi _0|P_{N_\mathrm{e}}P_\mathrm{G}P_\mathrm{G}P_{N_\mathrm{e}}|\varphi _0},$$
(8)
by means of the VMC method. The distribution of the wave vectors $`𝒌`$ is determined in the periodic boundary conditions in the $`x`$ direction and in the antiperiodic ones in the $`y`$ direction so as to avoid the gap node of $`d`$-wave superconductivity. Although the results in the 10$`\times `$10 square lattice are mainly shown in the following, we also calculate larger sizes up to $`20\times 20`$ to investigate the size dependence.
Figure 1 shows the $`\mathrm{\Delta }_d`$ dependence of $`E_{\mathrm{var}}`$ for various values of $`t_\mathrm{v}^{}`$ at the doping rate $`\delta =0.12`$. Apparently $`t_\mathrm{v}^{}0.1`$ gives the lowest variational energy. Since the Hamiltonian does not contain next-nearest neighbor hopping terms, the present result means that the shape of the FS of the ground state is different from that of the non-interacting Hamiltonian. We have also checked that, if the Hamiltonian has the next-nearest neighbor hopping term $`t^{}`$, the optimized variational state has $`t_\mathrm{v}^{}`$, whose amplitude is larger than $`t^{}`$.
The most significant effect of this result appears in the shape of momentum distribution functions. Figure 2 shows a contour map of the gradient of the momentum distribution function $`|_kn(𝒌)|`$ for $`t_\mathrm{v}^{}=0.1`$ calculated on the 20$`\times `$20 square lattice. Although we have used the optimized variational parameters $`\mathrm{\Delta }_d`$ and $`\mu `$ on the 10$`\times `$10 lattice, it is justified because their size dependeces are negligible. Brighter areas in Fig. 2 correspond to the momentum $`𝒌`$ with larger values of $`|_kn(𝒌)|`$. Although we cannot specify exactly the location of the FS due to the $`d`$-wave SC gap, we expect that the FS lies close to the area where $`|_kn(𝒌)|`$ is large.
Our result of momentum distribution function is similar to that obtained in high temperature expansion by Putikka et al. Since we take an opposite approach to high temperature studies, i.e., in the zero temperature, it is confirmed that the FS shown in Fig. 2 is an intrinsic feature of the $`t`$-$`J`$ model. Note here that the smearing of the FS around ($`\pi `$, 0) in our calculation is due to the $`d`$-wave SC gap. This suggests that the similar smearing observed in Ref. may be due to the pseudogap with $`d`$-wave symmetry, in addition to the smearing due to finite temperature.
For the wide range of doping $`\delta =0.040.20`$, we find that $`E_{\mathrm{var}}`$ is minimized around $`t_\mathrm{v}^{}0.1`$ and the chemical potential $`\mu 0.5\pm 0.05`$. Because of the insensitiveness of $`\mu `$ as a function of doping, the area of the momentum space enclosed by the FS is also insensitive to the doping rate. This result supports the violation of the Luttinger theorem suggested by Putikka et al. Actually $`n(𝒌)`$ at the doping $`\delta =0.2`$ in Ref. is very close to our results in Fig. 2 at $`\delta =0.12`$.
Figure 3 shows the energy difference between the value at $`t_\mathrm{v}^{}=0`$ and at $`t_\mathrm{v}^{}=0.1`$, i.e. the energy gain due to the finite $`t_\mathrm{v}^{}`$, for several system sizes. Although the Monte Carlo results scatter a little, there is apparently a tendency that the energy gain due to the finite $`t_\mathrm{v}^{}`$ becomes maximum around $`\delta =0.12`$. As we increase the system size, the energy gain slightly decreases, but it will remain finite in the thermodynamic limit.
Let us discuss here the relation between the energy gain and the van Hove singularity. For the optimized value $`t_\mathrm{v}^{}0.1`$, $`ϵ_k`$ at $`𝒌=(\pi ,0)`$ becomes $`0.4`$. On the other hand, the optimized chemical potential $`\mu `$ is around $`0.5`$. This means that the suddle point or the position of the flat band near $`𝒌=(\pi ,0)`$ is very close to the chemical potential. Since the $`d`$-wave SC gap has a maximum at $`(\pi ,0)`$, the enhancement of the density of states near the Fermi energy due to the van Hove singularity causes the energy gain. Actually, if we assume $`\mathrm{\Delta }_d=0`$, the lowest energy is achieved at $`t_\mathrm{v}^{}=0`$. This indicates that the FS deforms itself so as to fix the van Hove singularity to the chemical potential in the presence of the $`d`$-wave SC gap. This looks consistent with the mechanism of SC due to the van Hove singularity.
If we use a Hamiltonian $`\stackrel{~}{H}_{t\text{-}J}`$ without projection operator and the mean-field wave function $`|\varphi _0`$, it is apparent that the variational energy $`\varphi _0|\stackrel{~}{H}_{t\text{-}J}|\varphi _0/\varphi _0|\varphi _0`$ is minimized at $`t_\mathrm{v}^{}=0`$. Therefore the energy gain due to the non-zero value of $`t_\mathrm{v}^{}`$ is solely from the Gutzwiller’s projection operator. In order to clarify the effect of the projection, we examine the Gutzwiller approximation, in which the effect of constraints are taken into account by statistical weighting factors. For the $`t`$-$`J`$ model, we have
$$\frac{\varphi _0|P_\mathrm{G}c_{i\sigma }^{}c_{j\sigma }P_\mathrm{G}|\varphi _0}{\varphi _0|P_\mathrm{G}P_\mathrm{G}|\varphi _0}=g_t\varphi _0|c_{i\sigma }^{}c_{j\sigma }|\varphi _0=g_tc_{i\sigma }^{}c_{j\sigma }_0$$
(9)
and
$$\frac{\varphi _0|P_\mathrm{G}𝑺_i𝑺_jP_\mathrm{G}|\varphi _0}{\varphi _0|P_\mathrm{G}P_\mathrm{G}|\varphi _0}=g_s\varphi _0|𝑺_i𝑺_j|\varphi _0=g_s𝑺_i𝑺_j_0,$$
(10)
where $`g_t`$ and $`g_s`$ are the renormalization factors due to the projection. In the simplest Gutzwiller approximation, $`g_t`$ and $`g_s`$ are constant, i.e., $`g_t=2\delta /(1+\delta )`$ and $`g_s=4/(1+\delta )^2`$. In this case, the Gutzwiller projection does not alter the mean-field results. However it was recently shown that the dependence of the renormalization factors on the expectation values, such as $`\chi =c_{i\sigma }^{}c_{j\sigma }_0`$ and $`\mathrm{\Delta }=c_i^{}c_j^{}_0`$, plays a crucial role in evaluating the variational energies. If we use this Gutzwiller approximation, we can show that
$$\delta E_{\mathrm{var}}\frac{H}{\chi ^{}}\delta t_\mathrm{v}^{}+8Nt_\mathrm{v}^{}\delta t_\mathrm{v}^{}$$
(11)
where $`\chi ^{}=c_{i\sigma }^{}c_{j\sigma }_0`$ with $`(ij)`$ being the next-nearest neighbor sites and
$$\frac{H}{\chi ^{}}t\frac{g_t}{\chi ^{}}\underset{ij\sigma }{}(c_{i\sigma }^{}c_{j\sigma }_0+c.c.)+J\frac{g_s}{\chi ^{}}\underset{ij}{}𝑺_i𝑺_j_0.$$
(12)
The first term on the r.h.s. of eq. (11) is linear with respect to $`t_\mathrm{v}^{}`$ so that $`E_{\mathrm{var}}`$ is minimized at a finite value of $`t_\mathrm{v}^{}`$ which satisfies
$$t_\mathrm{v}^{}=\frac{1}{8N}\frac{H}{\chi ^{}}.$$
(13)
Apparently the renormalization factors $`g_t`$ and $`g_s`$ due to the projection operator and their nonlinear dependence on $`\chi ^{}`$ are the origin of the spontaneous deformation of the FS. These phenomena cannot be found in the mean-field theories. The explicit calculations will be published elsewhere.
In summary, we investigated the shape of the FS in the two-dimensional $`t`$-$`J`$ model by means of the VMC calculation introducing an additional variational parameter $`t_\mathrm{v}^{}`$. We found that the variational energy is minimized around $`t_\mathrm{v}^{}0.1`$ for various doping rates. The system size dependence indicate this effect is realized even in the thermodynamic limit. The magnitude of the energy gain is large enough compared with other VMC studies. For example, the energy difference between the pure $`d`$-wave SC phase and the coexistent phase of AF and $`d`$-wave SC is comparable to the present energy gain at the doping rate $`\delta 0.08`$. Then we have clarified the origin of the energy gain by examining the van Hove singularity and the effect of the projection using the Gutzwiller approximation. Combining our results at zero temperature and those in high temperature expansion, we consider that the FS deformation is the most significant phenomenon in the presence of strong correlation.
Since the nesting property of the FS becomes worse in the present wave function than in the original $`t`$-$`J`$ model, the coexistence of AF and $`d`$-wave SC near half-filling will be suppressed. Alternatively we expect some incommensurate AF correlations which were unexpected in the $`t`$-$`J`$ model. This is presumably related to the stripe state observed experimentally.
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# The Anderson prescription for surfaces and impurities
\[
## Abstract
We test the Anderson prescription , a BCS formalism for describing superconductivity in inhomogeneous systems, and compare results with those obtained from the Bogoliubov-de Gennes formalism, using the attractive Hubbard model with surfaces and nonmagnetic impurities. The Anderson approach captures the essential features of the spatial variation of the gap parameter and electron density around a surface or an impurity over a wide range of parameters. It breaks down, however, in the strong-coupling regime for a weak impurity potential.
\]
In microscopic treatments of inhomogeneity effects in superconductors, impurities are often averaged over in some manner . In the last decade, partly because of the increased computational resources now available, and partly because of the technical advances which allow small particle fabrication and single atom manipulation , the role of inhomogeneous effects in superconductors has received more wide-spread attention . One of the theoretical frameworks for addressing these questions is the Bogoliubov-de Gennes (BdG) formalism . This formalism allows one to answer questions regarding surfaces, interfaces, and impurity effects at a level of detail not previously addressed. At present, however, the scope of problems for which one can compute results accurately is limited by computer resources, since matrices whose dimension grows with system size require full diagonalization and should be solved self-consistently. On the other hand, the Anderson prescription , first presented to examine the impact of impurities on superconductivity, is a BCS formalism for an inhomogeneous system and requires one diagonalization of the single-particle problem. The purpose of this paper is to examine the limits of applicability of the Anderson approach, as compared to the BdG formalism. We first summarize the two approaches, and then follow with some concrete examples, utilizing surfaces and impurities as sources of inhomogeneities.
We find that the Anderson prescription works very well for surfaces and in some regimes for a single impurity. It breaks down for strong coupling with weak impurity scattering.
For our purposes we adopt a convenient model to describe s-wave superconductivity, the attractive Hubbard Hamiltonian (the results of our study will presumably apply to d-wave or other symmetry states):
$`H\mu N_e`$ $`=`$ $`{\displaystyle \underset{\genfrac{}{}{0pt}{}{i,\delta }{\sigma }}{}}t_\delta (a_{i+\delta ,\sigma }^{}a_{i\sigma }+\mathrm{h}.\mathrm{c}.)`$ (2)
$`{\displaystyle \underset{i,\sigma }{}}(\mu ϵ_i)n_{i\sigma }|U|{\displaystyle \underset{i}{}}n_in_i.`$
Here, $`a_{i\sigma }^{}`$ ($`a_{i\sigma }`$) creates (annihilates) an electron with spin $`\sigma `$ at site $`i`$ and $`n_{i\sigma }`$ is the number operator for an electron with spin $`\sigma `$ at site $`i`$. The $`t_\delta `$ is the hopping rate of electrons from one site to a neighbouring site (often nearest neighbours only are included, and we will adopt this model here), and $`|U|`$ is the attractive coupling strength between electrons on the same site. As usual, this attraction is justified in terms of an electron-phonon coupling, where retardation effects are unimportant, as is the case with many conventional superconductors. The second term includes the chemical potential $`\mu `$ and the impurity potential at site $`i`$, $`ϵ_i`$. We assume that impurities act to raise or lower single site energy levels. It is worth noting that impurity effects can certainly enter in other ways. For example if an impurity occupies one of the sites, undoubtedly the hopping amplitude to and from that site will also be altered, as will the interaction between two electrons occupying the same orbital on that site. In many studies (see, for example, Ref. ), the $`ϵ_i`$ are randomly distributed with some probability distribution, and then the results are averaged, to reflect the fact that we generally have no control over the precise distribution of impurities in the bulk. However, in systems where a single impurity can be added to the surface, for example (see, eg., Ref. ), we would want to study this model with only one impurity, at a particular site (in this case, on the surface).
Equation (2) also allows us the freedom to choose periodic boundary conditions (PBC) (to recover well-known results) or ‘open’ boundary conditions (OBC). The latter are natural in a tight-binding context; they require no assumptions about the order parameter, for example. ‘Open’ here simply means that electrons cannot hop beyond the surface. Here again more sophisticated boundary effects could be included — for example, in a real system the hopping integral at the surface will no doubt differ from that in the bulk, but we leave aside these finer points.
The BdG equations are obtained by defining an effective Hamiltonian, with effective potentials . By diagonalizing this effective Hamiltonian through the generalized Bogoliubov-Valatin transformation , one arrives at the two BdG equations :
$`E_nu_n(i)={\displaystyle \underset{i^{}}{}}A_{ii^{}}u_n(i^{})+V_iu_n(i)+\mathrm{\Delta }_iv_n(i)`$ (3)
$`E_nv_n(i)={\displaystyle \underset{i^{}}{}}A_{ii^{}}v_n(i^{})V_iv_n(i)+\mathrm{\Delta }_i^{}u_n(i)`$ (4)
where
$$A_{ii^{}}=t\underset{\delta }{}\left(\delta _{i^{},i\delta }+\delta _{i^{},i+\delta }\right)\delta _{ii^{}}\left(\mu ϵ_i\right).$$
(5)
The self-consistent potentials, $`V_i`$, and $`\mathrm{\Delta }_i`$, are given by
$`\mathrm{\Delta }_i=|U|{\displaystyle \underset{n}{}}u_n(i)v_n^{}(i)(12f_n)`$ (6)
$`V_i=|U|{\displaystyle \underset{n}{}}\left[|u_n(i)|^2f_n+|v_n(i)|^2(1f_n)\right].`$ (7)
We use the index $`n`$ to label the eigenvalues (there are $`2N`$ of them), the index $`i`$ to label the sites (1 through N), and the composite eigenvector is given by $`\left(\genfrac{}{}{0pt}{}{u_n}{v_n}\right)`$, of total length $`2N`$. The sums in Eqs. (6,7) are over positive eigenvalues only. The $`f_n`$ is the Fermi function, with argument $`\beta E_n`$, where $`\beta 1/k_BT`$, with $`T`$ the temperature. The single site electron density, $`n_i`$, is given, through Eq. (7), by $`V_i=|U|n_i/2`$.
The equations (6,7) for the effective potentials were determined through a variational principle so that the effective Hamiltonian allows fluctuations in any number of mean fields . As written, there are two possible mean field potentials, the Hartree potential, $`V_i`$, and the pair potential, $`\mathrm{\Delta }_i`$, from which the ground state energy and other properties may be obtained.
An alternate prescription was originally proposed by Anderson , whereby one first solves for the eigenvalues and eigenstates of the ‘non-interacting’ problem, i.e.,
$$E_n^0w_n(i)=\underset{i^{}}{}A_{ii^{}}w_n(i^{}).$$
(8)
Using the unitary matrix, $`U_{in}`$, for a basis which diagonalizes the single-particle Hamiltonian, one can determine the transformed electron-electron interaction:
$$V_{nm,n^{}m^{}}=|U|\underset{i}{}U_{in}^{}U_{im}^{}U_{in^{}}U_{im^{}},$$
(9)
which now mediates the (generally off-diagonal) electron-electron interaction. The gap and number equations are obtained just as in BCS theory, except that now the label is not the wave vector k, but rather some quantum number $`n`$, which simply enumerates the single particle eigenvalues . From the solution, one can obtain the ground state energy and, by transforming back to space coordinates, site-dependent quantities.
The simplest origin of gap inhomogeneity in a superconductor is the surface. The presence of surfaces beyond which electrons are unable to move yields a gap parameter (i.e., pair potential) which can exhibit a variety of behaviour near the surface. Traditionally in Ginzburg-Landau treatments the gap function is given a priori a boundary condition ; here the behaviour near a boundary (or impurity) is a derived quantity, i.e. as the solution to the BdG (or Anderson) equations.
The advantage (for the Anderson approach) of examining the impact of surfaces on the gap parameter is that an analytical solution exists for a simple tight-binding model . The eigenstates for a chain of length N, with lattice spacing $`a`$ and nearest-neighbour hopping $`t`$, are
$$a_{k\sigma }=\sqrt{\frac{2}{N+1}}\underset{i}{}\mathrm{sin}(kR_i)a_{i\sigma },$$
(10)
and the eigenenergies are
$$E_k^{(0)}=2t\mathrm{cos}(ka),ka=\frac{\pi n}{N+1},$$
(11)
where $`n=1,2,\mathrm{},N`$. We use these analytical results in the Anderson approach (making it not significantly more difficult than BCS theory), while in the BdG approach these analytical solutions are not particularly helpful.
In Fig. 1 we show the gap parameter $`\mathrm{\Delta }_i`$ as a function of site number $`i`$, for electron density $`n`$ ranging from half-filling to zero. The chain length is $`64`$ sites, and we have used OBC and the coupling strength $`|U|=1.5t`$. Here $`\mathrm{\Delta }_i`$ is shown for half the chain length (from $`i=1`$ to $`32`$): the gap parameter is symmetric about the middle. In Fig. 1(a) we plot the result from the BdG equations, while in Fig. 1(b) we show the corresponding results from the Anderson prescription. The first thing to note is that in either case the behaviour near the surface is markedly different as a function of electron density. At half-filling the gap parameter actually peaks at the surface, with several ‘Friedel-like’ oscillations ensuing towards the center of the sample, while at low fillings the gap parameter is much smoother by comparison. A comparison of the two figures shows quantitative differences, but overall, qualitatively they are very similar. It is evident that the Anderson prescription captures the essence of the BdG results remarkably well.
The accuracy of the Anderson results can be seen more closely in Fig. 2, where the cross sections of $`\mathrm{\Delta }_i`$ versus $`i`$ in Fig. 1 for (a) $`n=0.2`$ and (b) $`n=1.0`$ are shown. It is clear that the essential features of the BdG results are reproduced in the Anderson approach.
To examine the effect of impurities, we show in Fig. 3 $`\mathrm{\Delta }_i`$ as a function of $`i`$, for $`N=32`$ with an impurity at the central site with varying energy (both negative and positive). We have used PBC and an intermediate coupling strength, $`|U|=2t`$, at electron density $`n=0.9`$. We have intentionally stayed away from half-filling, at which the ground state is not a superconducting state, but a charge density wave. One may recall that with periodic boundary conditions and with no impurities, the ground state for the attractive Hubbard model at half-filling is doubly degenerate: both superconducting and charge density wave solutions coexist at this point. However, the presence of an impurity tilts the balance in favour of the charge density wave, and the BdG equations converge to a solution in which the pair potential, $`\mathrm{\Delta }_i`$, is identically zero at all sites. The Hartree potential, $`V_i`$, on the other hand, oscillates as a function of site position. The Anderson prescription is unable to reproduce this (correct) feature at half filling, and gives a superconducting solution with nonzero gap parameters. Also if the self-consistency of the Hartree potential is neglected in the BdG equations, this physics is missed, and the ensuing BdG result is similar to the Anderson solution.
Returning to Fig. 3 we have plotted the BdG results in Fig. 3(a), while in Fig. 3(b) we show the results from the Anderson equations. The gap is suppressed at the site with a positive impurity potential (as is the site density $`n_i`$) and exhibits ‘Friedel-like’ oscillations around it. The Anderson prescription captures this behaviour qualitatively, while it tends to underestimate the amplitudes of the oscillations (compare the solid curves in Fig. 3(a) and (b) for $`ϵ_{16}=0.5t`$, and note the magnified scale in the latter graph). The Anderson results become better for stronger impurity potentials and for electron density $`n`$ further away from half filling. For a negative impurity potential, when the potential strength is very weak, $`\mathrm{\Delta }_i`$ (and $`n_i`$) has a peak at the impurity site, as can be seen in Fig. 3(a) for $`ϵ_{16}=\mathrm{\hspace{0.17em}0.1}t`$ (the dotted curve). Though smaller in scale, the Anderson result in Fig. 3(b) has similar behaviour. As the potential strength increases, however, an attractive impurity tends to break the pairing, and suppresses the gap not only at the impurity site but also at surrounding sites (see the dashed curve for $`ϵ_{16}=\mathrm{\hspace{0.17em}0.5}t`$ in Fig. 3(a)). In such cases, the Anderson method overestimates the gap parameter around the impurity site. This can be seen in Fig. 3(b), where the gap has the correct oscillating pattern, but with much smaller amplitudes.
We study the attractive impurity case further in Fig. 4, where we show $`\mathrm{\Delta }_i`$ and $`n_i`$ versus $`i`$ for $`N=32`$ with PBC and $`ϵ_{16}=\mathrm{\hspace{0.17em}0.5}t`$, for (a) $`|U|=1t`$ and (b) $`|U|=3t`$. When the coupling is weaker than or comparable to the impurity potential, the Anderson approach captures the main features of the gap parameter around the impurity. This is the case for $`|U|=1t`$ in Fig. 4(a), and indeed the Anderson results show excellent agreement with the BdG results. As the coupling becomes stronger compared to the impurity potential, the impact of the impurity becomes more drastic, even with relatively weak strength. The gap is more suppressed around the impurity, and the density distribution exhibits ‘Friedel-like’ oscillations more enhanced. This can be seen for $`|U|=3t`$ in Fig. 4(b) (solid curves). On the other hand, the Anderson method yields a gap parameter that is more uniform as a function of site position, and does not reproduce the correct oscillations in the site densities.
In conclusion, we have formulated the BdG equations for a tight-binding model with an on-site attractive interaction. We have retained the self-consistent Hartree potential in the BdG equations, and found that a single impurity breaks the superconducting/charge density wave degeneracy which would otherwise exist at half-filling in this model. We have also formulated the prescription set out by Anderson, without impurity-averaging, and found good qualitative agreement with the BdG results.
To our knowledge, the spatial dependence of the order parameter and the electron density has not been previously explored in detail within the Anderson prescription. We have found, somewhat to our surprise, good agreement with results from the BdG formalism, for surfaces and single impurities (aside from the weak scattering limit). Throughout this study, it is important to note that in the vicinity of surfaces or impurities, “charge ordered” states and superconductivity in general coexist (at the mean field level). In future work we will examine various correlation functions and the local density of states, the latter of which has been and will continue to be measured using scanning tunneling microscopy.
Acknowledgements We thank Jorge Hirsch for suggesting the weak impurity potential regime as one in which the Anderson prescription should break down. We also thank Kamran Kaveh for stimulating discussions. Calculations were performed on the 64-node SGI parallel processor at the University of Alberta. This research was supported by the Avadh Bhatia Fellowship and by the Natural Sciences and Engineering Research Council of Canada and the Canadian Institute for Advanced Research. One of us (F.M.) acknowledges the hospitality of the Aspen Center for Physics, where some of this work was performed.
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# The regular representations and the 𝐴_𝑛(𝑉)-algebras
## 1. Introduction
In \[Li2\], for a vertex operator algebra $`V`$ and a nonzero complex number $`z`$, a weak $`VV`$-module $`𝒟_{P(z)}(V)`$ was constructed out of the dual space $`V^{}`$, and certain results of Peter-Weyl type were obtained. The weak $`VV`$-modules $`𝒟_{P(z)}(V)`$ were referred as regular representations. In \[Li3\], as a generalization, weak $`VV`$-modules $`𝒟_{P(z)}(V,U)`$ were constructed for any vector space $`U`$. Furthermore, Zhu’s $`A(V)`$-theory (\[Z1\], \[FZ\]) was related to the regular representations in the spirit of the induced module theory for a Lie group (cf. \[Ki\]), and a notion of an induced $`V`$-module from an $`A(V)`$-module was formulated in terms of the regular representations. The induced $`V`$-module from an $`A(V)`$-module $`U`$ was defined in \[Li3\] as follows: First consider linear functions from $`V`$ to $`U`$, which are lifted from linear functions from $`A(V)`$ to $`U`$, or simply just linear functions from $`A(V)`$ to $`U`$. Second, it was shown that $`\mathrm{Hom}(A(V),U)`$ is a subspace of $`𝒟_{P(1)}(V,U)`$, and what is more, $`\mathrm{Hom}(A(V),U)`$ and $`\mathrm{\Omega }(𝒟_{P(1)}(V,U))`$ $`(\mathrm{Hom}(V,U))`$ coincide as natural $`A(V)A(V)`$-modules. Meanwhile, all the (left) $`A(V)`$-invariant functions from $`A(V)`$ to $`U`$ give us the space $`\mathrm{Hom}_{A(V)}(A(V),U)`$, which is canonically isomorphic to $`U`$ as an $`A(V)`$-module. Third, the induced module $`\mathrm{Ind}_{A(V)}^VU`$ was defined to be the submodule of $`𝒟_{P(1)}(V,U)`$, generated by $`\mathrm{Hom}_{A(V)}(A(V),U)`$ $`(=U)`$ under the action of $`V`$.
In \[DLM2\], as a generalization of Zhu’s $`A(V)`$-theory, a family of associative algebras $`A_n(V)`$ were constructed and a family of functors $`\mathrm{\Omega }_n`$ from the category of weak $`V`$-modules to the category of $`A_n(V)`$-modules and a family of functors $`M_n`$ (with certain properties) from the category of $`A_n(V)`$-modules to the category of $``$-graded weak $`V`$-modules were constructed. By definition, $`\mathrm{\Omega }_n(W)`$ consists of each $`w`$ such that $`v_mw=0`$ for homogeneous $`vV`$ and for $`m\mathrm{wt}v+n`$. (Of course, $`\mathrm{\Omega }_n(W)`$ can also be considered as the invariant space with respect to a certain Lie algebra.) In the case that $`W`$ is a lowest weight generalized irreducible $`V`$-module, $`\mathrm{\Omega }_n(W)`$ is the sum of the first $`n`$ lowest weight subspaces.
In 1993, Zhu \[Z2\] gave a general construction of associative algebras from a vertex operator algebra for a certain purpose. The algebras $`A_n(V)`$ might be related to those algebras in a certain way.
In this paper, we shall relate $`A_n(V)`$-theory to the (generalized) regular representations of $`V`$ on $`𝒟_{P(1)}(V,U)`$. When $`n1`$, unlike the $`n=0`$ case \[Li3\], there are certain complicated factors. It is proved (Propositions 3.11, 3.15, and Corollary 3.17) that as vector spaces, $`\mathrm{Hom}(A_n(V),U)`$ is a subspace of $`\mathrm{\Omega }_n(𝒟_{P(1)}(V,U))`$. However, as natural $`A_n(V)A_n(V)`$-modules, $`\mathrm{Hom}(A_n(V),U)`$ is not a submodule. It turns out that the $`A_nA_n(V)`$-module structure on $`\mathrm{Hom}(A_n(V),U)`$ coincides with a twisted or deformed $`A_nA_n(V)`$-module structure on $`\mathrm{\Omega }_n(𝒟_{P(1)}(V,U))`$ with respect to a certain linear automorphism on $`\mathrm{\Omega }_n(𝒟_{P(1)}(V,U))`$ (Theorem 3.18). Using this connection we formulate a notion of induced $`V`$-module from an $`A_n(V)`$-module and we show that the induced modules are lowest weight generalized $`V`$-modules if the given $`A_n(V)`$-modules are irreducible.
An induced module theory from modules for a vertex operator subalgebra was established in \[DLin\]. As mentioned in \[Li3\], the notion of induced module defined here and the notion of induced module defined in \[DLin\] are different in nature.
This paper is organized as follows: In Section 2, we review the construction of the weak $`VV`$-module $`𝒟_{P(z)}(W,U)`$. In Section 3, we relate $`A_n(V)`$-modules $`\mathrm{Hom}(A_n(W),U)`$ with $`\mathrm{\Omega }_n(𝒟_{P(1)}(W,U))`$, and we define the induced $`V`$-module $`\mathrm{Ind}_{A_n(V)}^VU`$ for a given $`A_n(V)`$-module $`U`$.
## 2. The weak $`VV`$-module $`𝒟_{P(z)}(W,U)`$
In this section we shall recall from \[Li3\] the construction of the weak $`VV`$-module $`𝒟_{P(z)}(W,U)`$ and there are nothing new.
We use standard definitions and notations as given in \[FLM\] and \[FHL\]. A vertex operator algebra is denoted by $`(V,Y,\mathrm{𝟏},\omega )`$, where $`\mathrm{𝟏}`$ is the vacuum vector and $`\omega `$ is the Virasoro element, or simply by $`V`$. We also use the notion of weak module as defined in \[DLM2\]—A weak module satisfies all the axioms given in \[FLM\] and \[FHL\] for the notion of a module except that no grading is required.
We typically use letters $`x,y,x_1,x_2,\mathrm{}`$ for mutually commuting formal variables and $`z,z_0,\mathrm{}`$ for complex numbers. For a vector space $`U`$, $`U[[x,x^1]]`$ is the vector space of all (doubly infinite) formal series with coefficients in $`U`$, $`U((x))`$ is the space of formal Laurent series in $`x`$, and $`U((x^1))`$ is the space of formal Laurent series in $`x^1`$. We emphasize the following standard formal variable convention:
(2.1) $`(x_1x_2)^n={\displaystyle \underset{i0}{}}(1)^i\left({\displaystyle \genfrac{}{}{0pt}{}{n}{i}}\right)x_1^{ni}x_2^i,`$
(2.2) $`(xz)^n={\displaystyle \underset{i0}{}}(z)^i\left({\displaystyle \genfrac{}{}{0pt}{}{n}{i}}\right)x^{ni},`$
(2.3) $`(zx)^n={\displaystyle \underset{i0}{}}(1)^iz^{ni}\left({\displaystyle \genfrac{}{}{0pt}{}{n}{i}}\right)x^i`$
for $`n,z^\times `$.
For vector spaces $`U_1,U_2`$, a linear map $`f\mathrm{Hom}(U_1,U_2)`$ extends canonically to a linear map from $`U_1[[x,x^1]]`$ to $`U_2[[x,x^1]]`$. We shall use this canonical extension without any comments.
Let $`V`$ be a vertex operator algebra. For $`vV`$, we set (cf. \[FHL\], \[HL1\])
(2.4) $`Y^o(v,x)=Y(e^{xL(1)}(x^2)^{L(0)}v,x^1).`$
For a weak $`V`$-module $`W`$, because $`e^{xL(1)}(x^2)^{L(0)}vV[x,x^1]`$ and $`Y(u,x^1)wW((x^1))`$ for $`uV,wW`$, $`Y^o(v,x)`$ lies in $`\mathrm{Hom}(W,W((x^1)))`$. More generally, for any complex number $`z_0`$, $`Y^o(v,x+z_0)`$ lies in $`\mathrm{Hom}(W,W((x^1)))`$, where by definition
(2.5) $`Y^o(v,x+z_0)w=(Y^o(v,y)w)|_{y=x+z_0}`$
for $`wW`$. Let $`W`$ be a weak $`V`$-module and let $`U`$ be a vector space, e.g., $`U=`$. For $`vV,f\mathrm{Hom}(W,U)`$, the compositions $`fY^o(v,x)`$ and $`fY^o(v,x+z_0)`$ for any complex number $`z_0`$ are elements of $`(\mathrm{Hom}(W,U))[[x,x^1]]`$.
Let $`(x)`$ be the algebra of rational functions of $`x`$ (and $`[[x,x^1]]`$ be the vector space of all doubly infinite formal series in $`x`$ with complex coefficients). The $`\iota `$-maps $`\iota _{x;0}`$ and $`\iota _{x;\mathrm{}}`$ from $`(x)`$ to $`[[x,x^1]]`$ are defined as follows: for any rational function $`f(x)`$, $`\iota _{x;0}f(x)`$ is the Laurent series expansion of $`f(x)`$ at $`x=0`$ and $`\iota _{x;\mathrm{}}f(x)`$ is the Laurent series expansion of $`f(x)`$ at $`x=\mathrm{}`$. These are injective $`[x,x^1]`$-linear maps. In terms of the formal variable convention, we have
(2.6) $`\iota _{x;0}\left((xz)^nf(x)\right)=(z+x)^n\iota _{x;0}f(x),`$
(2.7) $`\iota _{x;\mathrm{}}\left((xz)^nf(x)\right)=(xz)^n\iota _{x;\mathrm{}}f(x)`$
for $`n,z^\times ,f(x)(x)`$.
###### Definition 2.1.
Let $`W`$ be a weak $`V`$-module, $`U`$ a vector space and $`z`$ a nonzero complex number. Define $`𝒟_{P(z)}(W,U)`$ to be the subspace of $`\mathrm{Hom}(W,U)`$, consisting of each $`f`$ such that for $`vV`$, there exist $`k,l`$ such that
(2.8) $`(xz)^kx^lu^{},fY^o(v,x)w[x]`$
for all $`u^{}U^{},wW`$, or what is equivalent, for all $`u^{}U^{},wW`$, the formal series
$$u^{},fY^o(v,x)w,$$
an element of $`((x^1))`$, absolutely converges in the domain $`|x|>|z|`$ to a rational function of the form $`x^l(xz)^kg(x)`$ for $`g(x)[x]`$.
The following are equivalent definitions of $`𝒟_{P(z)}(W,U)`$ in terms of formal series:
###### Lemma 2.2.
Let $`f\mathrm{Hom}(W,U)`$. Then the following statements are equivalent:
(a) $`f𝒟_{P(z)}(W,U)`$.
(b) For $`vV`$, there exist $`k,l`$ such that
(2.9) $`(xz)^kx^lfY^o(v,x)(\mathrm{Hom}(W,U))[[x]].`$
(c) For $`vV`$, there exist $`k,l`$ such that for each $`wW`$,
(2.10) $`(xz)^kx^lfY^o(v,x)wU[x].`$
Let $`vV,f𝒟_{P(z)}(W,U)`$ and let $`k,l`$ be such that (2.10) holds. Then by changing variable we get
(2.11) $`x^k(x+z)^lfY^o(v,x+z)wU[x]`$
for $`wW`$.
###### Definition 2.3.
Let $`W,U`$ and $`z`$ be given as before. For
$$vV,f𝒟_{P(z)}(W,U),$$
we define two elements $`Y_{P(z)}^L(v,x)f`$ and $`Y_{P(z)}^R(v,x)f`$ of $`(\mathrm{Hom}(W,U))[[x,x^1]]`$ by
(2.12) $`(Y_{P(z)}^L(v,x)f)(w)`$ $`=`$ $`(z+x)^l\left((x+z)^lf(Y^o(v,x+z)w)\right)`$
(2.13) $`(Y_{P(z)}^R(v,x)f)(w)`$ $`=`$ $`(z+x)^k\left((xz)^kf(Y^o(v,x)w)\right)`$
for $`wW`$, where $`k,l`$ are any pair of (possibly negative) integers such that (2.9) holds.
First, in view of (2.10) and (2.11), both $`(z+x)^l\left((x+z)^lf(Y^o(v,x+z)w)\right)`$ and $`(z+x)^k\left((xz)^kf(Y^o(v,x)w)\right)`$ lie in $`U((x))`$, so that $`Y_{P(z)}^L(v,x)f`$ and $`Y_{P(z)}^R(v,x)f`$ make sense. However, we are not allowed to remove the left-right brackets to cancel $`(xz)^k`$ or $`(x+z)^l`$ because of the nonexistence of $`(z+x)^lf(Y^o(v,x+z)w)`$ and $`(z+x)^kf(Y^o(v,x)w)`$. Second, they are also well defined, i.e., they are independent of the choice of the pair of integers $`k,l`$. Indeed, if $`k^{},l^{}`$ are another pair of integers such that (2.9) holds, say for example, $`kk^{}`$, then
$`(z+x)^k\left((xz)^kfY^o(v,x)w\right)`$
$`=`$ $`(z+x)^k\left((xz)^{kk^{}}(xz)^k^{}fY^o(v,x)w\right)`$
$`=`$ $`(z+x)^k(xz)^{kk^{}}\left((xz)^k^{}fY^o(v,x)w\right)`$
$`=`$ $`(z+x)^k^{}\left((xz)^k^{}fY^o(v,x)w\right).`$
From definition we immediately have:
###### Lemma 2.4.
For $`vV,f𝒟_{P(z)}(W,U)`$,
(2.15) $`(z+x)^lY_{P(z)}^L(v,x)f=(x+z)^lfY^o(v,x+z),`$
(2.16) $`(z+x)^kY_{P(z)}^R(v,x)f=(xz)^kfY^o(v,x),`$
where $`k,l`$ are any pair of (maybe negative) integers such that (2.9) holds.
From the definition, $`u^{},fY^o(v,x)w`$ lies in the range of $`\iota _{x;\mathrm{}}`$ for $`u^{}U^{},f𝒟_{P(z)}(W,U),vV,wW`$. Then $`\iota _{x;\mathrm{}}^1u^{},fY^o(v,x)w`$ is a well defined element of $`(x)`$. In terms of rational functions and the $`\iota `$-maps we immediately have:
###### Lemma 2.5.
For $`vV,f𝒟_{P(z)}(W,U),u^{}U^{},wW`$,
(2.17) $`u^{},(Y_{P(z)}^L(v,x)f)(w)=\iota _{x;0}\iota _{x;\mathrm{}}^1u^{},fY^o(v,x+z)w,`$
(2.18) $`u^{},(Y_{P(z)}^R(v,x)f)(w)=\iota _{x;0}\iota _{x;\mathrm{}}^1u^{},fY^o(v,x)w.`$
###### Theorem 2.6.
Let $`W`$ be a weak $`V`$-module, $`U`$ a vector space and $`z`$ a nonzero complex number. Then the pairs $`(𝒟_{P(z)}(W,U),Y_{P(z)}^L)`$ and $`(𝒟_{P(z)}(W,U),Y_{P(z)}^R)`$ carry the structure of a weak $`V`$-module and the actions $`Y_{P(z)}^L`$ and $`Y_{P(z)}^R`$ of $`V`$ on $`𝒟_{P(z)}(W,U)`$ commute. Furthermore, set
(2.19) $`Y_{P(z)}=Y_{P(z)}^LY_{P(z)}^R.`$
Then the pair $`(𝒟_{P(z)}(W,U),Y_{P(z)})`$ carries the structure of a weak $`VV`$-module.
The following relation among $`fY^o(v,x),Y^L(v,x)f`$ and $`Y^R(v,x)f`$ holds \[Li3\]:
###### Proposition 2.7.
Let $`vV,f𝒟_{P(z)}(W,U)`$. Then
$`x_0^1\delta \left({\displaystyle \frac{xz}{x_0}}\right)fY^o(v,x)x_0^1\delta \left({\displaystyle \frac{zx}{x_0}}\right)Y_{P(z)}^R(v,x)f`$
$`=`$ $`z^1\delta \left({\displaystyle \frac{xx_0}{z}}\right)Y_{P(z)}^L(v,x_0)f.`$
## 3. The associative algebras $`A_n(V)`$ and induced modules $`\mathrm{Ind}_{A_n(V)}^VU`$
In this section, the nonzero complex number $`z`$ in the notion of weak $`VV`$-module $`𝒟_{P(z)}(W,U)`$ will be specified as $`1`$. We shall simply use $`Y^L`$ and $`Y^R`$ for $`Y_{P(1)}^L`$ and $`Y_{P(1)}^R`$. Throughout this section, $`n`$ will represent a nonnegative integer.
We shall need the following notions. A generalized $`V`$-module \[HL1\] is a weak $`V`$-module on which $`L(0)`$ semisimply acts. Then for a generalized $`V`$-module $`W`$ we have the $`L(0)`$-eigenspace decomposition: $`W=_hW_{(h)}`$. Thus, a generalized $`V`$-module satisfies all the axioms defining the notion of a $`V`$-module (\[FLM\], \[FHL\]) except the two grading restrictions on homogeneous subspaces. If a generalized $`V`$-module furthermore satisfies the lower truncation condition (one of the two grading restrictions), it is called a lower truncated generalized module \[H1\].
A lowest weight generalized $`V`$-module is a generalized $`V`$-module such that $`W=_nW_{(h+n)}`$ for some $`h`$ and $`W_{(h)}`$ generates $`W`$ as a weak $`V`$-module. Furthermore, if $`W0`$, we call the unique $`h`$ the lowest weight of $`W`$.
Now we recall the construction of $`A_n(V)`$ algebra and some basic results from \[DLM2\].
###### Definition 3.1.
Let $`V`$ be a vertex operator algebra and let $`n`$. Define a subspace $`O_n(V)`$ of $`V`$, linearly spanned by elements
(3.1) $`(L(1)+L(0))v,`$
(3.2) $`\mathrm{Res}_x{\displaystyle \frac{(1+x)^{\mathrm{wt}u+n}}{x^{2n+2}}}Y(u,x)v`$
for homogeneous $`u,vV`$. Define
$`u_nv`$ $`=`$ $`{\displaystyle \underset{m=0}{\overset{n}{}}}(1)^m\left({\displaystyle \genfrac{}{}{0pt}{}{m+n}{n}}\right)\mathrm{Res}_x{\displaystyle \frac{(1+x)^{\mathrm{wt}u+n}}{x^{n+m+1}}}Y(u,x)v`$
$`\left(={\displaystyle \underset{m=0}{\overset{n}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n1}{m}}\right)\mathrm{Res}_x{\displaystyle \frac{(1+x)^{\mathrm{wt}u+n}}{x^{n+m+1}}}Y(u,x)v\right).`$
We have:
###### Lemma 3.2.
\[DLM2\] Let $`u,vV`$ be homogeneous. Then
$`u_nv{\displaystyle \underset{m=0}{\overset{n}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n1}{m}}\right)(1)^{nm}\mathrm{Res}_xx^{nm1}(1+x)^{\mathrm{wt}v+m1}Y(v,x)uO_n(V).`$
Set $`A_n(V)=V/O_n(V)`$.
###### Proposition 3.3.
\[DLM2\] Let $`(V,Y,\mathrm{𝟏},\omega )`$ be a vertex operator algebra. Then
(a) $`O_n(V)`$ is a two-sided ideal of the nonassociative algebra $`(V,_n)`$ and the quotient algebra $`A_n(V)`$ is an associative algebra with $`\mathrm{𝟏}+O_n(V)`$ as its identity element, with $`\omega +O_n(V)`$ being central and with an involution (anti-automorphism)
(3.4) $`\theta :v+O_n(V)e^{L(1)}(1)^{L(0)}v+O_n(V).`$
(b) For each $`n0`$, the identity map of $`V`$ gives rise to an algebra homomorphism $`\psi _n`$ from $`A_{n+1}(V)`$ onto $`A_n(V)`$.
For any weak $`V`$-module $`W`$ and any $`n`$, we define \[DLM2\]
(3.5) $`\mathrm{\Omega }_n(W)=\{wW|v_{\mathrm{wt}v+m}w=0\text{ for homogeneous }vV,mn\}.`$
For $`w\mathrm{\Omega }_n(W)`$ and for homogeneous $`vV`$, we have $`x^{\mathrm{wt}v+n}Y(v,x)wW[[x]]`$. When $`n=0`$, we have $`\mathrm{\Omega }(W)=\mathrm{\Omega }_0(W)`$. Clearly,
(3.6) $`\mathrm{\Omega }_0(W)\mathrm{\Omega }_1(W)\mathrm{}.`$
We similarly define $`\mathrm{\Omega }_1(W),\mathrm{\Omega }_2(W),\mathrm{}`$. Since $`\mathrm{wt}\mathrm{𝟏}=0`$ and $`\mathrm{𝟏}_rw=\delta _{r,1}w`$, we have $`\mathrm{\Omega }_n(W)=0`$ for $`n1`$.
The following result was proved in \[DLM2\]:
###### Proposition 3.4.
Let $`W`$ be a weak $`V`$-module and let $`n0`$. Then $`\mathrm{\Omega }_n(W)`$ is an $`A_n(V)`$-module where $`v+O_n(V)`$ acts as $`v_{\mathrm{wt}v1}`$ for homogeneous $`vV`$.
Let $`W_1,W_2`$ be weak $`V`$-modules and let $`\psi `$ be a $`V`$-homomorphism from $`W_1`$ to $`W_2`$. It is clear that $`\psi (\mathrm{\Omega }_n(W_1))\mathrm{\Omega }_n(W_2)`$ and the restriction $`\mathrm{\Omega }_n(\psi ):=\psi |_{\mathrm{\Omega }_n(W_1)}`$ is an $`A_n(V)`$-homomorphism. It is routine to check that we have obtained a functor $`\mathrm{\Omega }_n`$ from the category of weak $`V`$-modules to the category of $`A_n(V)`$-modules.
###### Lemma 3.5.
Let $`W`$ be a weak $`V`$-module and set
(3.7) $`𝒮(W)=_{n0}\mathrm{\Omega }_n(W).`$
Let $`uV`$ be homogeneous and let $`r`$. Then
(3.8) $`u_r\mathrm{\Omega }_n(W)\mathrm{\Omega }_n(W)`$
if $`r\mathrm{wt}u1`$, and
(3.9) $`u_r\mathrm{\Omega }_n(W)\mathrm{\Omega }_{n+\mathrm{wt}ur1}(W)`$
if $`r<\mathrm{wt}u1`$. In particular, $`𝒮(W)`$ is a sub-weak-module of $`W`$. Furthermore,
(3.10) $`\mathrm{\Omega }_n(𝒮(W))=\mathrm{\Omega }_n(W).`$
###### Proof.
Let $`w\mathrm{\Omega }_n(W)`$, let $`vV`$ be homogeneous and let $`m`$. By Borcherds commutator formula,
$`v_{\mathrm{wt}v+m}u_rw`$
$`=`$ $`u_rv_{\mathrm{wt}v+m}w+{\displaystyle \underset{i0}{}}\left({\displaystyle \genfrac{}{}{0pt}{}{\mathrm{wt}v+m}{i}}\right)(v_iu)_{\mathrm{wt}v+m+ri}w`$
$`=`$ $`u_rv_{\mathrm{wt}v+m}w+{\displaystyle \underset{i0}{}}\left({\displaystyle \genfrac{}{}{0pt}{}{\mathrm{wt}v+m}{i}}\right)(v_iu)_{\mathrm{wt}(v_iu)+m+r\mathrm{wt}u+1}w.`$
Then the first part follows immediately. Since $`𝒮(W)`$ is a submodule of $`W`$, we have $`\mathrm{\Omega }_n(𝒮(W))\mathrm{\Omega }_n(W)`$. It is easy to see that $`\mathrm{\Omega }_n(W)\mathrm{\Omega }_n(𝒮(W))`$. This completes the proof. ∎
We shall need the result that $`L(1)`$ is locally nilpotent on $`𝒮(W)`$ for any weak $`V`$-module $`W`$. To prove this result, we recall from \[Li3\] the following result, which is a reformulation of a result in \[DLM2\] (Remark 3.3):
###### Lemma 3.6.
Let $`W`$ be a weak $`V`$-module, $`wW`$. Let $`u,vV`$ and let $`k`$ be such that
(3.12) $`x^kY(u,x)wW[[x]],`$
or equivalently,
(3.13) $`u_{k+m}w=0\text{ for }m0.`$
Then for $`p,q`$,
(3.14) $`u_pv_qw={\displaystyle \underset{i=0}{\overset{s}{}}}{\displaystyle \underset{j0}{}}\left({\displaystyle \genfrac{}{}{0pt}{}{pk}{i}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{k}{j}}\right)(u_{pki+j}v)_{q+k+ij}w.`$
where $`s`$ is any nonnegative integer such that $`x^{s+1+q}Y(v,x)wW[[x]]`$.
As an immediate consequence we have (\[DM\] and \[Li1\]):
###### Corollary 3.7.
Let $`W`$ be a weak $`V`$-module and let $`wW`$. Set
(3.15) $`w=\text{linear span }\{v_mw|vV,m\}.`$
Then $`w`$ is the sub-weak-module of $`W`$, generated by $`w`$.
###### Lemma 3.8.
Let $`W`$ be a weak $`V`$-module. Then for any $`r`$ homogeneous vectors $`v^1,\mathrm{},v^rV`$,
(3.16) $`v_{m_1}^1\mathrm{}v_{m_r}^r\mathrm{\Omega }_n(W)=0`$
for $`m_i`$ with
$$m_1+\mathrm{}+m_r\mathrm{wt}v^1+\mathrm{}+\mathrm{wt}v^rr+n.$$
In particular, for homogeneous $`vV`$ and for $`m\mathrm{wt}v`$,
(3.17) $`(v_m)^n\mathrm{\Omega }_n(W)=0.`$
###### Proof.
We shall prove the first part by induction on $`r`$. From the definition of $`\mathrm{\Omega }_n(W)`$, the lemma is true for $`r=1`$. Assume it is true for any $`r`$ homogeneous vectors in $`V`$. Now let $`v^1,\mathrm{},v^r,v^{r+1}V`$ be homogeneous and let $`m_i`$ with
(3.18) $`m_1+\mathrm{}+m_r+m_{r+1}\mathrm{wt}v^1+\mathrm{}+\mathrm{wt}v^{r+1}(r+1)+n.`$
Set
$$u=v^r,v=v^{r+1},p=m_r,q=m_{r+1}.$$
Since $`w\mathrm{\Omega }_n(W)`$, in Lemma 3.6, we may take $`k=\mathrm{wt}u+n=\mathrm{wt}v^r+n`$. Let $`s`$ be any nonnegative integer such that $`x^{s+1+q}Y(v,x)wW[[x]]`$. By Lemma 3.6, we have
$`u_pv_qw`$
$`=`$ $`{\displaystyle \underset{i=0}{\overset{s}{}}}{\displaystyle \underset{j0}{}}\left({\displaystyle \genfrac{}{}{0pt}{}{p\mathrm{wt}un}{i}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{\mathrm{wt}u+n}{j}}\right)(u_{p\mathrm{wt}uni+j}v)_{q+\mathrm{wt}u+n+ij}w.`$
Notice that
$`\mathrm{wt}(u_{p\mathrm{wt}uni+j}v)`$ $`=`$ $`\mathrm{wt}u+\mathrm{wt}v+\mathrm{wt}u+n+ij1p`$
$`=`$ $`2\mathrm{w}\mathrm{t}v^r+\mathrm{wt}v^{r+1}+n+ij1p.`$
Thus
$`m_1+\mathrm{}+m_{r1}+(q+\mathrm{wt}u+n+ij)`$
$``$ $`\mathrm{wt}v^1+\mathrm{}+\mathrm{wt}v^{r+1}(r+1)+n+(q+\mathrm{wt}u+n+ij)m_rm_{r+1}`$
$`=`$ $`\mathrm{wt}v^1+\mathrm{}+\mathrm{wt}v^{r1}+\mathrm{wt}(u_{p\mathrm{wt}uni+j}v)r+n.`$
Then it follows from the inductive hypothesis that
$`v_{m_1}^1\mathrm{}v_{m_{k+1}}^{k+1}w`$
$`=`$ $`{\displaystyle \underset{i=0}{\overset{s}{}}}{\displaystyle \underset{j0}{}}\left({\displaystyle \genfrac{}{}{0pt}{}{pk}{i}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{k}{j}}\right)v_{m_1}^1\mathrm{}v_{m_{k1}}^{k1}(u_{pki+j}v)_{q+k+ij}w`$
$`=`$ $`0.`$
This finishes the induction and concludes the proof. ∎
In view of Lemma 3.8, noticing that $`L(1)=\omega _2`$ and $`\mathrm{wt}\omega =2`$, we immediately have:
###### Corollary 3.9.
Let $`W`$ be a weak $`V`$-module, let $`vV`$ be homogeneous and let $`m\mathrm{wt}v`$. Then $`v_m`$ is locally nilpotent on $`𝒮(W)`$. In particular, $`L(1)`$ is locally nilpotent on $`𝒮(W)`$.
Let $`W`$ be a weak $`V`$-module. We define $`O_n^{}(W)`$ to be the subspace of $`W`$, linearly spanned by elements of the form:
(3.22) $`v_nw:=\mathrm{Res}_xx^{2n2}(1+x)^{\mathrm{wt}v+n}Y(v,x)w`$
for $`wW`$ and for homogeneous $`vV`$. The proof of Lemma 2.1.2 of \[Z1\] with minor necessary changes directly gives:
###### Lemma 3.10.
Let $`W`$ be a weak $`V`$-module, let $`wW`$, and let $`vV`$ be homogeneous. Then
(3.23) $`\mathrm{Res}_xx^{2n2r}(1+x)^{\mathrm{wt}v+n+s}Y(v,x)wO_n^{}(W)`$
for $`rs0`$.
In the following there will be several module structures on a certain vector space. For this reason, we shall use $`\mathrm{\Omega }_n(W,Y_W)`$ including the vertex operator map $`Y_W`$ in the notation for $`\mathrm{\Omega }_n(W)`$. Since $`\mathrm{Hom}(,U)`$ is a contravariant functor for the category of vector spaces, for any vector spaces $`A,B`$ and any surjective linear map $`g\mathrm{Hom}(A,B)`$, we have an injective linear map $`\mathrm{Hom}(g,U)`$ from $`\mathrm{Hom}(B,U)`$ into $`\mathrm{Hom}(A,U)`$. In particular, if $`B`$ is a quotient space of $`A`$, we may naturally identify $`\mathrm{Hom}(B,U)`$ as a subspace of $`\mathrm{Hom}(A,U)`$.
###### Proposition 3.11.
Let $`W`$ be a weak $`V`$-module and let $`U`$ be a vector space. Set $`A_n^{}(W)=W/O_n^{}(W)`$. Then
(3.24) $`\mathrm{Hom}(A_n^{}(W),U)=\mathrm{\Omega }_n(𝒟_{P(1)}(W,U),Y^L)\mathrm{\Omega }_n(𝒟_{P(1)}(W,U),Y^R).`$
Furthermore, elements $`\alpha `$ of $`\mathrm{Hom}(A_n^{}(W),U)`$, a natural subspace of $`\mathrm{Hom}(W,U)`$, are characterized by the following property:
(3.25) $`x^{\mathrm{wt}v+n}(x+1)^{\mathrm{wt}v+n}\alpha Y^o(v,x)(\mathrm{Hom}(W,U))[[x]]`$
for homogeneous $`vV`$.
###### Proof.
Let $`T`$ be the set defined by the property (3.25). We shall prove
$`\mathrm{Hom}(A_n^{}(W),U)T\mathrm{\Omega }_n(𝒟_{P(1)}(W,U),Y^L)\mathrm{\Omega }_n(𝒟_{P(1)}(W,U),Y^R)`$
$`T\mathrm{Hom}(A_n^{}(W),U).`$
Let $`\alpha \mathrm{Hom}(A_n^{}(W),U)`$ and let $`vV`$ be homogeneous. Then for any $`m0`$,
$`\mathrm{Res}_xx^{\mathrm{wt}v+n+m}(x+1)^{\mathrm{wt}v+n}\alpha Y^o(v,x)w`$
$`=`$ $`\mathrm{Res}_xx^{\mathrm{wt}v+n+m}(x+1)^{\mathrm{wt}v+n}\alpha Y(e^{xL(1)}(x^2)^{L(0)}v,x^1)w`$
$`=`$ $`(1)^{\mathrm{wt}v}\mathrm{Res}_xx^{2n+m}(1+x^1)^{\mathrm{wt}v+n}\alpha Y(e^{xL(1)}v,x^1)w`$
$`=`$ $`(1)^{\mathrm{wt}v}\mathrm{Res}_xx^{2nm2}(1+x)^{\mathrm{wt}v+n}\alpha Y(e^{x^1L(1)}v,x)w`$
$`=`$ $`0`$
because (Lemma 3.10)
$`\mathrm{Res}_xx^{2nm2}(1+x)^{\mathrm{wt}v+n}Y(e^{x^1L(1)}v,x)w`$
$`=`$ $`{\displaystyle \underset{i0}{}}{\displaystyle \frac{1}{i!}}\mathrm{Res}_xx^{2nm2i}(1+x)^{\mathrm{wt}(L(1)^iv)+n+i}Y(L(1)^iv,x)w`$
$``$ $`O_n^{}(W).`$
This proves (3.25). Since $`Y^o(v,x)wW((x^1))`$ for $`wW`$, (3.25) implies
(3.28) $`x^{\mathrm{wt}v+n}(x+1)^{\mathrm{wt}v+n}\alpha Y^o(v,x)wU[x].`$
By changing variable we get
(3.29) $`(x1)^{\mathrm{wt}v+n}x^{\mathrm{wt}v+n}\alpha Y^o(v,x1)(\mathrm{Hom}(W,U))[[x]].`$
By Lemma 2.2, $`\alpha 𝒟_{P(1)}(W,U)`$ and by Lemma 2.4
(3.30) $`x^{\mathrm{wt}v+n}(1+x)^{\mathrm{wt}v+n}Y^R(v,x)\alpha =x^{\mathrm{wt}v+n}(x+1)^{\mathrm{wt}v+n}\alpha Y^o(v,x),`$
(3.31) $`(1+x)^{\mathrm{wt}v+n}x^{\mathrm{wt}v+n}Y^L(v,x)\alpha =(x1)^{\mathrm{wt}v+n}x^{\mathrm{wt}v+n}\alpha Y^o(v,x1).`$
Consequently,
(3.32) $`x^{\mathrm{wt}v+n}Y^R(v,x)\alpha ,x^{\mathrm{wt}v+n}Y^L(v,x)\alpha (\mathrm{Hom}(W,U))[[x]].`$
That is, $`\alpha \mathrm{\Omega }_n(𝒟_{P(1)}(W,U),Y^L)\mathrm{\Omega }_n(𝒟_{P(1)}(W,U),Y^R).`$
Conversely, let $`\alpha \mathrm{\Omega }_n(𝒟_{P(1)}(W,U),Y^L)\mathrm{\Omega }_n(𝒟_{P(1)}(W,U),Y^R).`$ Then (3.32) holds for each homogeneous $`vV`$. Recall (2.7) with $`f=\alpha ,z=1`$:
$`x_0^1\delta \left({\displaystyle \frac{x+1}{x_0}}\right)\alpha Y^o(v,x)x_0^1\delta \left({\displaystyle \frac{1+x}{x_0}}\right)Y_{P(z)}^R(v,x)\alpha `$
$`=`$ $`\delta (x+x_0)Y_{P(z)}^L(v,x_0)\alpha .`$
Applying $`\mathrm{Res}_{x_0}x^{\mathrm{wt}v+n}x_0^{\mathrm{wt}v+n}`$ to (3), then using (3.32) and the fundamental properties of delta function we get
$`x^{\mathrm{wt}v+n}(x+1)^{\mathrm{wt}v+n}\alpha Y^o(v,x)`$
$`=`$ $`x^{\mathrm{wt}v+n}(1+x)^{\mathrm{wt}v+n}Y^R(v,x)\alpha `$
$`\mathrm{Res}_{x_0}x^{\mathrm{wt}v+n}x_0^{\mathrm{wt}v+n}\delta (x+x_0)Y^L(v,x_0)\alpha `$
$`=`$ $`x^{\mathrm{wt}v+n}(1+x)^{\mathrm{wt}v+n}Y^R(v,x)\alpha `$
$``$ $`(\mathrm{Hom}(W,U))[[x]].`$
Furthermore, for any $`wW`$,
$`\mathrm{Res}_xx^{2n2}(1+x)^{\mathrm{wt}v+n}\alpha Y(v,x)w`$
$`=`$ $`\mathrm{Res}_xx^{2n2}(1+x)^{\mathrm{wt}v+n}\alpha Y^o(e^{xL(1)}(x^2)^{L(0)}v,x^1)w`$
$`=`$ $`\mathrm{Res}_xx^{2n}(1+x^1)^{\mathrm{wt}v+n}\alpha Y^o(e^{x^1L(1)}(x^2)^{L(0)}v,x)w`$
$`=`$ $`\mathrm{Res}_x(1)^{\mathrm{wt}v}x^{\mathrm{wt}v+n}(x+1)^{\mathrm{wt}v+n}\alpha Y^o(e^{x^1L(1)}v,x)w`$
$`=`$ $`{\displaystyle \underset{i0}{}}(1)^{\mathrm{wt}v}{\displaystyle \frac{1}{i!}}\mathrm{Res}_xx^{\mathrm{wt}vi+n}(x+1)^{\mathrm{wt}v+n}\alpha Y^o(L(1)^iv,x)w`$
$`=`$ $`{\displaystyle \underset{i0}{}}(1)^{\mathrm{wt}v}{\displaystyle \frac{1}{i!}}\mathrm{Res}_xx^{\mathrm{wt}(L(1)^iv)+n}(x+1)^{\mathrm{wt}(L(1)^iv)+n+i}\alpha Y^o(L(1)^iv,x)w`$
$`=`$ $`0.`$
Thus $`\alpha (O_n^{}(W))=0`$, hence $`\alpha \mathrm{Hom}(A_n^{}(W),U)`$. This completes the proof. ∎
It follows from Theorem 2.6 and Proposition 3.4 that $`\mathrm{\Omega }_n(𝒟_{P(1)}(W,U),Y^R)`$ is a natural $`A_n(V)`$-module. Since $`Y^L`$ and $`Y^R`$ commute, $`\mathrm{\Omega }_n(𝒟_{P(1)}(W,U),Y^R)`$ is also a weak $`V`$-module under the vertex operator map $`Y^L`$. Then it follows from Proposition 3.4 again that
$$\mathrm{\Omega }_n(\mathrm{\Omega }_n(𝒟_{P(1)}(W,U),Y^R),Y^L)$$
is an $`A_n(V)A_n(V)`$-module. Clearly,
$`\mathrm{\Omega }_n(\mathrm{\Omega }_n(𝒟_{P(1)}(W,U),Y^R),Y^L)`$
$`=`$ $`\mathrm{\Omega }_n(𝒟_{P(1)}(W,U),Y^L)\mathrm{\Omega }_n(𝒟_{P(1)}(W,U),Y^R).`$
Thus, $`\mathrm{\Omega }_n(𝒟_{P(1)}(W,U),Y^L)\mathrm{\Omega }_n(𝒟_{P(1)}(W,U),Y^R)`$ is an $`A_n(V)A_n(V)`$-module. For convenience, we refer to this $`A_n(V)A_n(V)`$-module structure as the canonical module structure. From definition, we have
(3.36) $`\mathrm{\Omega }_n(𝒟_{P(1)}(W,U))\mathrm{\Omega }_n(𝒟_{P(1)}(W,U),Y^L)\mathrm{\Omega }_n(𝒟_{P(1)}(W,U),Y^R).`$
(The equality of (3.36) holds when $`n=0`$, but the equality does not hold for $`n1`$.) It is easy to see that $`\mathrm{\Omega }_n(𝒟_{P(1)}(W,U))`$ is an $`A_n(V)A_n(V)`$-submodule.
Motivated by \[Li3\] for $`n=0`$, we should identify $`\mathrm{Hom}(A_n^{}(W),U)`$ with
$$\mathrm{\Omega }_n(𝒟_{P(1)}(W,U),Y^L)\mathrm{\Omega }_n(𝒟_{P(1)}(W,U),Y^R)$$
as natural $`A_n(V)A_n(V)`$-modules. We shall prove that $`A_n^{}(W)`$ just like $`A_n(V)`$ has a natural $`A_n(V)A_n(V)`$-module structure and so does $`\mathrm{Hom}(A_n^{}(W),U)`$. It turns out that the $`A_n(V)A_n(V)`$-module $`\mathrm{Hom}(A_n^{}(W),U)`$ is naturally isomorphic to
$$\mathrm{\Omega }_n(𝒟_{P(1)}(W,U),Y^L)\mathrm{\Omega }_n(𝒟_{P(1)}(W,U),Y^R)$$
with a deformed $`A_n(V)A_n(V)`$-module structure.
To achieve our goal we shall need the following result (cf. \[Li2\], Remark 2.10):
###### Proposition 3.12.
Let $`(E,Y_E)`$ be a weak $`V`$-module on which $`L(1)`$ is locally nilpotent, and let $`z_0`$ be any complex number. For $`vV`$, we define
(3.37) $`Y_E^{[z_0]}(v,x)=Y_E(e^{z_0(1+z_0x)L(1)}(1+z_0x)^{2L(0)}v,x/(1+z_0x)).`$
Then the pair $`(E,Y_E^{[z_0]})`$ carries the structure of a weak $`V`$-module and $`e^{z_0L(1)}`$ is a $`V`$-isomorphism from $`(E,Y_E)`$ to $`(E,Y_E^{[z_0]})`$. Furthermore, for homogeneous $`vV`$ and for $`m`$, we have
$`\mathrm{Res}_xx^mY_E^{[z_0]}(v,x)`$
$`=`$ $`\mathrm{Res}_xx^m(1z_0x)^{2\mathrm{w}\mathrm{t}vm2}Y_E(e^{z_0(1z_0x)^1L(1)}v,x).`$
In particular,
$`\mathrm{Res}_xx^{\mathrm{wt}v1}Y^{[z_0]}(v,x)w`$
$`=`$ $`\mathrm{Res}_xx^{\mathrm{wt}v1}(1z_0x)^{\mathrm{wt}v1}Y(e^{z_0(1z_0x)^1L(1)}v,x).`$
###### Proof.
Recall the conjugation formula (5.2.38) of \[FHL\]:
$`e^{x_1L(1)}Y(v,x)e^{x_1L(1)}`$
$`=`$ $`Y(e^{x_1(1+x_1x)L(1)}(1+x_1x)^{2L(0)}v,x/(1+x_1x)).`$
Because $`L(1)`$ is locally nilpotent on $`E`$, we may set $`x_1=z_0`$, so that we have
$`e^{z_0L(1)}Y_E(v,x)e^{z_0L(1)}`$
$`=`$ $`Y_E(e^{z_0(1+z_0x)L(1)}(1+z_0x)^{2L(0)}v,x/(1+z_0x))`$
$`=`$ $`Y_E^{[z_0]}(v,x).`$
Then the first part of the proposition follows immediately.
By changing variable $`x=y/(1z_0y)`$ we get
$`\mathrm{Res}_xx^mY_E^{[z_0]}(v,x)`$
$`=`$ $`\mathrm{Res}_xx^mY_E(e^{z_0(1+z_0x)L(1)}(1+z_0x)^{2L(0)}v,x/(1+z_0x))`$
$`=`$ $`\mathrm{Res}_yy^m(1z_0y)^{m2}Y_E(e^{z_0(1+z_0y)^1L(1)}(1z_0y)^{2L(0)}v,y)`$
$`=`$ $`\mathrm{Res}_yy^m(1z_0y)^{2\mathrm{w}\mathrm{t}vm2}Y_E(e^{z_0(1z_0y)^1L(1)}v,y).`$
This completes the proof. ∎
By definition we have
$`(Y^{[z_0]})^{[z_0]}(v,x)`$
$`=`$ $`Y^{[z_0]}(e^{z_0(1z_0x)L(1)}(1z_0x)^{2L(0)}v,x/(1z_0x))`$
$`=`$ $`Y(e^{z_0(1+z_0x)L(1)}(1+z_0x)^{2L(0)}e^{z_0(1z_0x)L(1)}(1z_0x)^{2L(0)}v,x).`$
Recall (5.3.3) of \[FHL\]:
(3.43) $`x_1^{L(0)}L(1)x_1^{L(0)}=x_1L(1).`$
From this we immediately get
(3.44) $`x_1^{L(0)}e^{xL(1)}x_1^{L(0)}=e^{xx_1L(1)}.`$
In view of (3.44) we have
(3.45) $`e^{z_0(1+z_0x)L(1)}(1+z_0x)^{2L(0)}e^{z_0(1z_0x)L(1)}(1z_0x)^{2L(0)}=1,`$
hence
(3.46) $`(Y^{[z_0]})^{[z_0]}(v,x)=Y(v,x).`$
Continuing with Proposition 3.12 we have:
###### Proposition 3.13.
Let $`(E,Y_E)`$ be a weak $`V`$-module on which $`L(1)`$ is locally nilpotent and let $`z_0`$ be any complex number. Then
(3.47) $`\mathrm{\Omega }_n(E,Y_E)=\mathrm{\Omega }_n(E,Y_E^{[z_0]}).`$
Furthermore, $`e^{z_0L(1)}`$ is an $`A_n(V)`$-isomorphism from $`\mathrm{\Omega }_n(E,Y)`$ to $`\mathrm{\Omega }_n(E,Y^{[z_0]})`$.
###### Proof.
From (3.12) we easily get
(3.48) $`\mathrm{\Omega }_n(E,Y_E)\mathrm{\Omega }_n(E,Y_E^{[z_0]}).`$
Using this and the fact that $`Y_E=(Y_E^{[z_0]})^{[z_0]}`$, we get
(3.49) $`\mathrm{\Omega }_n(E,Y_E^{[z_0]})\mathrm{\Omega }_n(E,Y_E).`$
This proves (3.47). The second part follows from Proposition 3.12 immediately. ∎
We shall use Proposition 3.13 for $`z_0=0,1,1`$. Let $`W`$ and $`U`$ be given as before. Set
(3.50) $`E=𝒮(𝒟_{P(1)}(W,U)).`$
In view of Lemma 3.5 we have
(3.51) $`\mathrm{\Omega }_n(𝒟_{P(1)}(W,U))=\mathrm{\Omega }_n(E)`$
and it follows from Corollary 3.9 that $`L(1)`$ is locally nilpotent on $`E`$, so that we can apply Propositions 3.12 and 3.13 to $`E`$.
Let $`W`$ be a weak $`V`$-module. For homogeneous $`vV`$ and for $`wW`$, we define
$`v_nw`$
$`=`$ $`{\displaystyle \underset{m=0}{\overset{n}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n1}{m}}\right)\mathrm{Res}_xx^{nm1}(1+x)^{\mathrm{wt}v+n}Y(v,x)w,`$
$`w_nv`$
$`=`$ $`{\displaystyle \underset{m=0}{\overset{n}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n1}{m}}\right)(1)^{nm}\mathrm{Res}_xx^{nm1}(1+x)^{\mathrm{wt}v+m1}Y(v,x)w.`$
Then extend the definition by linearity.
Now, we are in a position to prove our key result.
###### Proposition 3.14.
Let $`W`$ be a weak $`V`$-module and let $`U`$ be a vector space. Let
$$f\mathrm{Hom}(A_n^{}(W),U)=\mathrm{\Omega }_n(𝒟_{P(1)}(W,U),Y^L)\mathrm{\Omega }_n(𝒟_{P(1)}(W,U),Y^R)$$
and $`wW`$. Then
(3.54) $`\left(\mathrm{Res}_xx^{\mathrm{wt}v}(Y^L)^{[1]}(v,x)f\right)(w)=f(w_nv)`$
(3.55) $`\left(\mathrm{Res}_xx^{\mathrm{wt}v}(Y^R)^{[1]}(v,x)f\right)(w)=f(\theta (v)_nw)`$
for homogeneous $`vV`$, where
(3.56) $`\theta (v)=e^{L(1)}(1)^{L(0)}v`$
(cf. (3.4)).
###### Proof.
First, using (3.44) we get (\[FHL\], (5.3.1)):
(3.57) $`e^{xL(1)}(x^2)^{L(0)}e^{x^1L(1)}=(x^2)^{L(0)},`$
(3.58) $`e^{xL(1)}(x^2)^{L(0)}e^{(x+1)^1L(1)}=e^{x/(x+1)L(1)}(x^2)^{L(0)}.`$
Because
(3.59) $`\left({\displaystyle \underset{m=0}{\overset{n}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n1}{m}}\right)(1)^{n+1m}x^m\right)(1+x)^{n+1}1+x^{n+1}[[x]],`$
for $`kn`$,
(3.60) $`\mathrm{Res}_xx^{\mathrm{wt}v+k}Y^L(v,x)f=0.`$
Since for any homogeneous $`uV`$,
(3.61) $`(1+x)^{\mathrm{wt}u+n}Y^L(u,x)f=(x1)^{\mathrm{wt}u+n}fY^o(u,x1)`$
(Proposition 3.11), we have
$`(1+x)^{\mathrm{wt}v+n}Y^L(e^{(1+x)^1L(1)}v,x)f`$
$`=`$ $`(x1)^{\mathrm{wt}v+n}fY^o(e^{(x1)^1L(1)}v,x1),`$
noting that $`\mathrm{wt}L(1)^iv=\mathrm{wt}vi`$ for $`i0`$. Using (3.12) and all the above information we have
$`\left(\mathrm{Res}_xx^{\mathrm{wt}v1}(Y^L)^{[1]}(v,x)f\right)(w)`$
$`=`$ $`\mathrm{Res}_x(1)^{\mathrm{wt}v1}x^{\mathrm{wt}v1}(1+x)^{\mathrm{wt}v1}\left(Y^L(e^{(1+x)^1L(1)}v,x)f\right)(w)`$
$`=`$ $`\mathrm{Res}_x{\displaystyle \underset{m=0}{\overset{n}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n1}{m}}\right)(1)^{\mathrm{wt}v+nm}x^{m+\mathrm{wt}v1}(1+x)^{\mathrm{wt}v+n}`$
$`\left(Y^L(e^{(1+x)^1L(1)}v,x)f\right)(w)`$
$`=`$ $`\mathrm{Res}_x{\displaystyle \underset{m=0}{\overset{n}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n1}{m}}\right)(1)^{\mathrm{wt}v+nm}x^{m+\mathrm{wt}v1}(x1)^{\mathrm{wt}v+n}`$
$`f\left(Y^o(e^{(x1)^1L(1)}v,x1)w\right)`$
$`=`$ $`{\displaystyle \underset{m=0}{\overset{n}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n1}{m}}\right)(1)^{\mathrm{wt}v+nm}`$
$`\mathrm{Res}_x(x+1)^{m+\mathrm{wt}v1}x^{\mathrm{wt}v+n}f\left(Y^o(e^{x^1L(1)}v,x)w\right)`$
$`=`$ $`{\displaystyle \underset{m=0}{\overset{n}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n1}{m}}\right)(1)^{\mathrm{wt}v+nm}`$
$`\mathrm{Res}_x(x+1)^{m+\mathrm{wt}v1}x^{\mathrm{wt}v+n}f(Y(e^{xL(1)}(x^2)^{L(0)}e^{x^1L(1)}v,x^1)w)`$
$`=`$ $`{\displaystyle \underset{m=0}{\overset{n}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n1}{m}}\right)(1)^{\mathrm{wt}v+nm}`$
$`\mathrm{Res}_x(x+1)^{m+\mathrm{wt}v1}x^{\mathrm{wt}v+n}f(Y((x^2)^{L(0)}v,x^1)w)`$
$`=`$ $`\mathrm{Res}_x{\displaystyle \underset{m=0}{\overset{n}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n1}{m}}\right)(1)^{nm}(x+1)^{m+\mathrm{wt}v1}x^{\mathrm{wt}v+n}f(Y(v,x^1)w)`$
$`=`$ $`\mathrm{Res}_x{\displaystyle \underset{m=0}{\overset{n}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n1}{m}}\right)(1)^{nm}(x^1+1)^{m+\mathrm{wt}v1}x^{\mathrm{wt}vn2}f(Y(v,x)w)`$
$`=`$ $`\mathrm{Res}_x{\displaystyle \underset{m=0}{\overset{n}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n1}{m}}\right)(1)^{nm}x^{nm1}(1+x)^{\mathrm{wt}v+m1}f(Y(v,x)w)`$
$`=`$ $`f(w_nv).`$
Similarly, using the fact
(3.63) $`\left({\displaystyle \underset{m=0}{\overset{n}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n1}{m}}\right)x^m\right)(1+x)^{\mathrm{wt}v+n}1+x^{n+1}[[x]]`$
we have
$`\left(\mathrm{Res}_xx^{\mathrm{wt}v1}(Y^R)^{[1]}(v,x)f\right)(w)`$
$`=`$ $`\mathrm{Res}_xx^{\mathrm{wt}v1}(1+x)^{\mathrm{wt}v1}\left(Y^R(e^{(1+x)^1L(1)}v,x)f\right)(w)`$
$`=`$ $`\mathrm{Res}_x{\displaystyle \underset{m=0}{\overset{n}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n1}{m}}\right)x^{m+\mathrm{wt}v1}(1+x)^{\mathrm{wt}v+n}\left(Y^R(e^{(1+x)^1L(1)}v,x)f\right)(w)`$
$`=`$ $`\mathrm{Res}_x{\displaystyle \underset{m=0}{\overset{n}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n1}{m}}\right)x^{m+\mathrm{wt}v1}(x+1)^{\mathrm{wt}v+n}f\left(Y^o(e^{(x+1)^1L(1)}v,x)w\right)`$
$`=`$ $`\mathrm{Res}_x{\displaystyle \underset{m=0}{\overset{n}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n1}{m}}\right)x^{m+\mathrm{wt}v1}(x+1)^{\mathrm{wt}v+n}`$
$`f\left(Y(e^{xL(1)}(x^2)^{L(0)}e^{(x+1)^1L(1)}v,x^1)w\right)`$
$`=`$ $`\mathrm{Res}_x{\displaystyle \underset{m=0}{\overset{n}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n1}{m}}\right)(1)^{\mathrm{wt}v}x^{m\mathrm{wt}v1}(x+1)^{\mathrm{wt}v+n}`$
$`f\left(Y(e^{x/(x+1)L(1)}(x^2)^{L(0)}v,x^1)w\right)`$
$`=`$ $`\mathrm{Res}_x{\displaystyle \underset{m=0}{\overset{n}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n1}{m}}\right)(1)^{\mathrm{wt}v}x^{m+\mathrm{wt}v1}(x^1+1)^{\mathrm{wt}v+n}`$
$`f\left(Y(e^{(1+x)^1L(1)}(x^2)^{L(0)}v,x)w\right)`$
$`=`$ $`\mathrm{Res}_x{\displaystyle \underset{m=0}{\overset{n}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n1}{m}}\right)(1)^{\mathrm{wt}v}x^{nm1}(1+x)^{\mathrm{wt}v+n}f\left(Y(e^{(1+x)^1L(1)}v,x)w\right)`$
$`=`$ $`{\displaystyle \underset{m=0}{\overset{n}{}}}{\displaystyle \underset{i0}{}}{\displaystyle \frac{1}{i!}}\left({\displaystyle \genfrac{}{}{0pt}{}{n1}{m}}\right)(1)^{\mathrm{wt}v}`$
$`\mathrm{Res}_xx^{nm1}(1+x)^{\mathrm{wt}(L(1)^iv)+n}f\left(Y(L(1)^iv,x)w\right)`$
$`=`$ $`f(\theta (v)_nw).`$
This completes the proof. ∎
One can in principle use similar arguments to those in \[Z1\], \[FZ\] and \[DLM2\] to show that the left and right actions of $`V`$ on $`W`$, defined by (3) and (3), give rise to an $`A_n(V)`$-bimodule structure on $`A_n^{}(W)`$, or $`A_n(W)`$ defined below. (From the proof of Theorem 2.3 of \[DLM2\], to prove the associativity for the right action it seems that we need to prove at least one more combinatorial identity in addition to those proved in \[DLM2\].) As a matter of fact, this easily follows from Proposition 3.14 and the (canonical and deformed) $`A_n(V)A_n(V)`$-module structures on
$$\mathrm{\Omega }_n(𝒟_{P(1)}(W,U),Y^L)\mathrm{\Omega }_n(𝒟_{P(1)}(W,U),Y^R).$$
###### Proposition 3.15.
Let $`W`$ be a weak $`V`$-module. Then the left and right actions of $`V`$ on $`W`$, defined by (3) and (3), give rise to an $`A_n(V)`$-bimodule structure on $`A_n^{}(W)`$.
###### Proof.
Let $`U=`$. For homogeneous $`vV`$ we set
(3.64) $`o_L^{[1]}(v)=\mathrm{Res}_xx^{\mathrm{wt}v1}(Y^L)^{[1]}(v,x),`$
(3.65) $`o_R^{[1]}(v)=\mathrm{Res}_xx^{\mathrm{wt}v1}(Y^R)^{[1]}(v,x).`$
Then extend the definition by linearity. It follows from Theorem 2.6 and Propositions 3.4, 3.12 and 3.13 that $`o_L^{[1]}o_R^{[1]}`$ gives rise to an $`A_n(V)A_n(V)`$-structure on
$$\mathrm{\Omega }_n(𝒟_{P(1)}(W,),Y^L)\mathrm{\Omega }_n(𝒟_{P(1)}(W,),Y^R)=\mathrm{Hom}(A_n^{}(W),).$$
In particular,
(3.66) $`o_L^{[1]}(O_n(V))=o_R^{[1]}(O_n(V))=0.`$
The following arguments are classical and routine in nature. Let $`u,vV`$ be homogeneous and let $`wW`$. For any $`f\mathrm{Hom}(A_n^{}(W),)`$, because
$$o_L^{[1]}(v)f\mathrm{Hom}(A_n^{}(W),),$$
in view of Proposition 3.14 we have
(3.67) $`f((u_nw)_nv)=o_L^{[1]}(v)f,u_nw=0.`$
Since $`f`$ is arbitrary, we must have
(3.68) $`(u_nw)_nvO_n^{}(W).`$
Using Proposition 3.14 and (3.66) we have
(3.69) $`f,w_n(u_nv)=o_L^{[1]}(u_nv)f,w=0`$
for every $`f\mathrm{Hom}(A_n^{}(W),)`$. Consequently,
(3.70) $`w_n(u_nv)O_n^{}(W).`$
Similarly, using the fact that $`\theta `$ gives rise to the involution $`\theta `$ of $`A_n(V)`$ (Proposition 3.3) we have
(3.71) $`v_n(u_nw),(u_nv)_nwO_n^{}(W).`$
Then the left action and right action of $`V`$ on $`W`$ give rise to a left action and right action of $`A_n(V)`$ on $`A_n^{}(W)`$. The rest can be proved similarly. ∎
Motivated by the definition of $`O_n(V)`$ we define
(3.72) $`O_n(W)=O_n^{}(W)+(L(1)+L(0))W.`$
The proof of of Lemma 2.1 of \[DLM2\] directly gives:
###### Lemma 3.16.
Let $`W`$ be a weak $`V`$-module, let $`wW`$ and let $`vV`$ be homogeneous. Then
(3.73) $`v_nww_nv\mathrm{Res}_x(1+x)^{\mathrm{wt}v1}Y(v,x)w\mathrm{mod}O_n(W).`$
Set
(3.74) $`A_n(W)=W/O_n(W).`$
Then we have:
###### Corollary 3.17.
The subspace $`O_n(W)`$ of $`W`$ is stable under the left and right actions of $`V`$ on $`W`$, defined by (3) and (3), and the quotient space $`A_n(W)`$ is an $`A_n(V)`$-bimodule which is the quotient module of $`A_n^{}(W)`$ modulo $`O_n(W)/O_n^{}(W)`$.
###### Proof.
In view of Proposition 3.15 we only need to prove
(3.75) $`(L(1)w+L(0)w)_nvO_n(W),`$
(3.76) $`v_n(L(1)w+L(0)w)O_n(W)`$
for $`vV`$, $`wW`$. Let us assume $`v`$ is homogeneous. First, from the proof of Lemma 2.2 in \[DLM2\] we have
(3.77) $`(L(1)v+L(0)v)_nw=(1)^n(2n+1)\left({\displaystyle \genfrac{}{}{0pt}{}{2n+1}{n}}\right)(v_nw)O_n(W).`$
Then using the fact
$`[L(1)+L(0),Y(v,x)]`$ $`=`$ $`(1+x)Y(L(1)v,x)+Y(L(0)v,x)`$
$`=`$ $`(1+x)(d/dx)Y(v,x)+(\mathrm{wt}v)Y(v,x),`$
we get
$`v_n(L(1)+L(0))w`$
$`=`$ $`(L(1)+L(0))(v_nw)+(L(1)v+L(0)v)_nwO_n(W).`$
Using Lemma 3.16 and (3), we get
$`(L(1)w+L(0)w)_nv`$
$``$ $`v_n(L(1)w+L(0)w)`$
$`\mathrm{Res}_x(1+x)^{\mathrm{wt}v1}Y(v,x)(L(1)w+L(0)w)\mathrm{mod}O_n(W)`$
$``$ $`\mathrm{Res}_x(1+x)^{\mathrm{wt}v1}Y(v,x)(L(1)w+L(0)w)\mathrm{mod}O_n(W)`$
$`=`$ $`\mathrm{Res}_x(1+x)^{\mathrm{wt}v1}(L(1)+L(0))Y(v,x)w`$
$`+\mathrm{Res}_x(1+x)^{\mathrm{wt}v}(d/dx)Y(v,x)w+\mathrm{Res}_x(\mathrm{wt}v)(1+x)^{\mathrm{wt}v1}Y(v,x)w`$
$`=`$ $`\mathrm{Res}_x(1+x)^{\mathrm{wt}v1}(L(1)+L(0))Y(v,x)w`$
$``$ $`0\mathrm{mod}O_n(W).`$
(This argument is also similar to one in the proof Lemma 2.2 of \[DLM2\].) ∎
Because $`A_n^{}(W)`$ is an $`A_n(V)`$-bimodule and $`\theta `$ is an involution of $`A_n(V)`$, from the classical fact $`\mathrm{Hom}(A_n^{}(W),U)`$ becomes an $`A_n(V)A_n(V)`$-module with
(3.81) $`((a_1,a_2)f)(w)=f(\theta (a_2)wa_1)`$
for $`a_1,a_2A_n(V),f\mathrm{Hom}(A_n^{}(W),U),wA_n(V)`$. We refer to this $`A_n(V)A_n(V)`$-module structure as the canonical dual module structure.
Combining Propositions 3.14 with 3.13 we immediately have:
###### Theorem 3.18.
Let $`W`$ be a weak $`V`$-module and $`U`$ a vector space. Let
$$\eta :\mathrm{Hom}(A_n^{}(W),U)\mathrm{\Omega }_n(𝒟_{P(1)}(W,U),Y^L)\mathrm{\Omega }_n(𝒟_{P(1)}(W,U),Y^R)$$
be the natural identification map (Proposition 3.11). Then the linear map
$$\sigma :=e^{L^R(1)L^L(1)}\eta $$
is an $`A_n(V)A_n(V)`$-isomorphism where $`\mathrm{Hom}(A_n^{}(W),U)`$ is equipped with the canonical dual $`A_n(V)A_n(V)`$-module structure and
$$\mathrm{\Omega }_n(𝒟_{P(1)}(W,U),Y^L)\mathrm{\Omega }_n(𝒟_{P(1)}(W,U),Y^R)$$
is equipped with the canonical module structure.
Let $`U`$ be an $`A_n(V)`$-module. Then $`U=\mathrm{Hom}_{A_n(V)}(A_n(V),U)`$. We also have the following $`A_n`$-module inclusion relations:
$$\mathrm{Hom}_{A_n(V)}(A_n(V),U)\mathrm{Hom}_{A_n(V)}(A_n^{}(V),U)\mathrm{Hom}_{}(A_n^{}(V),U).$$
With Theorem 3.18 we may and we should identify $`U`$ as a submodule of the $`A_n(V)`$-module $`\mathrm{\Omega }_n(𝒟_{P(1)}(V,U),Y^L)`$.
###### Definition 3.19.
Let $`U`$ be an $`A_n(V)`$-module. We define $`\mathrm{Ind}_{A_n(V)}^VU`$ to be the submodule of $`(𝒟_{P(1)}(V,U),Y^L)`$, generated by $`U`$ $`(=\mathrm{Hom}_{A_n(V)}(A_n(V),U))`$.
Using the proof of Lemma 3.14 in \[Li3\] with some minor changes, we have:
###### Proposition 3.20.
Let $`W`$ be a weak $`V`$-module and let $`U`$ be an irreducible $`A_n(V)`$-submodule of $`\mathrm{\Omega }_n(W)`$. Then the weak submodule $`M`$ of $`W`$ generated by $`U`$ is a lowest weight generalized $`V`$-module such that $`M_{(h)}=U`$ for some $`h`$ and $`M_{(k+h)}=0`$ for $`k<n`$. In particular, if $`U`$ is an irreducible $`A_n(V)`$-module, then $`\mathrm{Ind}_{A_n(V)}^VU`$ is a lowest weight generalized $`V`$-module with $`U`$ being the homogeneous subspace of some weight $`h`$ such that the lowest weight is no smaller than $`hn`$.
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# Ly𝛼 Absorption Systems and the Intergalactic Medium
## 1 Introduction
The existence of a forest of absorption lines blueward of the Ly$`\alpha `$ emission line in quasar spectra has been known for over 30 years (Bahcall and Salpeter 1965; Lynds 1971). These lines arise from Ly$`\alpha `$ absorption by neutral hydrogen from intervening structure along the line-of-sight. Early theoretical models interpreted this structure as absorption caused by discrete gas clouds in the intergalactic medium (IGM), either pressure confined by a hot IGM (Sargent et al. 1980; Ostriker and Ikeuchi 1983) or confined by the gravity of dark matter ‘mini-halos’ (e.g. Rees 1986). Over the last few years our understanding of the Ly$`\alpha `$ forest has undergone a transformation for at least two reasons. Firstly, observations with the Keck telescope have produced almost noise-free spectra of quasars at high spectral resolution over the redshift range $`2\stackrel{<}{}z\stackrel{<}{}4`$. The exquisite quality of Keck spectra has allowed observers to resolve Ly$`\alpha `$ absorption lines at low column densities ($`10^{12.5}\mathrm{cm}^2`$) and to study their evolution. Secondly, hydrodynamic numerical simulations of structure formation in cold dark matter (CDM) universes with high spatial resolution are now possible and have proved remarkably successful in reproducing many observed properties of the Ly$`\alpha `$ forest (Cen et al. 1994; Zhang , Anninos and Norman 1995, 1997; Miralda-Escudé et al. 1996; Hernquist et al. 1996, Theuns et al. 1998a). These simulations have shown that most of the Ly$`\alpha `$ lines at column densities $`\stackrel{<}{}10^{14.5}\mathrm{cm}^2`$ arise from modest fluctuations in the baryon density in a space filling photoionized IGM, rather than from distinct clouds. The properties of the Ly$`\alpha `$ lines can therefore be used to probe the structure and thermal history of the diffuse IGM and of the background UV radiation that determines its ionization state. The key characteristics of the numerical simulations are described in the next Section. Section 3 summarizes a number of results from these simulations and describes how the Ly$`\alpha `$ lines can be used to study the IGM. In this paper we discuss only the properties of the low column density Ly$`\alpha `$ lines. For a discussion of damped Ly$`\alpha `$ systems and metal lines see Pettini’s contribution to these proceedings. For a recent review of observations and theoretical models of Ly$`\alpha `$ absorption lines see Rauch (1998).
## 2 Numerical Simulations of the IGM
The simplest cosmological hydrodynamical simulations follow the evolution of (optically thin) gas and dark matter assuming a uniform photoionizing background. The simulation is therefore specified by:
$``$ parameters defining the cosmological model and its matter content (e.g. $`\mathrm{\Omega }_m`$, $`\mathrm{\Omega }_b`$, $`\mathrm{\Omega }_\mathrm{\Lambda }`$, $`H_0100h\mathrm{kms}^1\mathrm{Mpc}^1`$);
$``$ the amplitude and spectral shape of the mass fluctuations;
$``$ a model for the background UV flux as a function of redshift.
Here we will review the evolution of CDM universes with initially scale-invariant adiabatic fluctuations. The linear power spectrum for these models (in the limit of small baryon content) is given by Bardeen et al. (1986). We adopt a reference model (Model S) with parameters $`\mathrm{\Omega }_m=1`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$, $`h=0.5`$ and $`\mathrm{\Omega }_b=0.05`$. The physical density in baryons in this model is $`\omega _b\mathrm{\Omega }_bh^2=0.0125`$. The model is normalized so that the rms density fluctuations in spheres of radius $`8h^1\mathrm{Mpc}`$ is $`\sigma _8=0.7`$. The photoionizing background radiation is assumed to originate from quasars according to the model of Haardt and Madau (1996, herafter HM). The photoionization rates for hydrogen and helium in this model are plotted in the left hand panel of Figure 1 as a function of redshift. With these photoionization rates, and assuming a uniform IGM, hydrogen is reionized at a redshift of $`z6`$ and HeII is reionized at $`z4.5`$.
The right hand panel of Figure 1 shows the distribution of gas elements in the temperature-density plane in a numerical simulation of our reference model S at a redshift of $`z=3`$. There is a plume of shock heated gas extending to temperatures $`\stackrel{>}{}10^5`$K, but most of the gas has a low overdensity and follows a power law-like ‘equation of state’
$$T=T_0\left(\frac{\rho _b}{\overline{\rho }_b}\right)^{\gamma 1},$$
(1)
(Hui and Gnedin 1997). At times long after reionization, the diffuse IGM will settle into a state in which adiabatic cooling is balanced by photoheating. The temperature of the IGM will therefore tend towards
$`T3.2\times 10^4\mathrm{K}\left[{\displaystyle \frac{\mathrm{\Omega }_bh(1+\delta )(1+z)^3}{(2+\alpha )E(z)}}\right]^{1/1.76},`$ (2)
$`E(z)={\displaystyle \frac{H_0}{H(z)}}=[\mathrm{\Omega }_m(1+z)^3+\mathrm{\Omega }_k(1+z)^2+\mathrm{\Omega }_\mathrm{\Lambda }]^{1/2}`$
where we have assumed a power-law photoionizing background
$$J_\nu =J_{\nu _L}\left(\frac{\nu }{\nu _L}\right)^\alpha .$$
(3)
In equation (2.2) $`\mathrm{\Omega }_k=1\mathrm{\Omega }_m\mathrm{\Omega }_\mathrm{\Lambda }`$ and the exponent $`1/1.76`$ arise from the temperature dependence of the HII and HeIII recombination coefficients.
Figure 2 shows the spatial distribution of the gas in the simulation at $`z=3`$. The upper figure shows the gas with temperature $`T>10^5`$K. This hot gas fills a small fraction of the volume and is located in the dense knots and filaments corresponding to collapsed structures. In contrast, the cool gas with $`T<10^5`$K (shown in the lower figure) fills most of the computational volume. It is this diffuse low density gas which we believe accounts for the vast majority of the observed Ly$`\alpha `$ lines. Furthermore, because most of this gas is at low overdensities, $`\delta \stackrel{<}{}10`$, it is relatively easy to model numerically. By investigating the Ly$`\alpha `$ forest we can therfore hope to learn about the properties of the diffuse IGM, which at typical quasar redshifts of $`z24`$ contains most of the baryonic matter in the Universe. In particular:
$``$ the optical depth of HI and HeII absorption can set constraints on the baryonic density of the Universe and on the evolution, amplitude and spectrum of the photoionizing background;
$``$ the fluctuating optical depth of HI absorption can be used to construct the power spectrum of the matter fluctuations at redshifts $`z24`$;
$``$ the Ly$`\alpha `$ absorption line-widths can be used to infer the temperature and equation of state of the diffuse IGM and its evolution.
These topics will be discussed in more detail in the next section.
## 3 The Ly $`\alpha `$ forest as a probe of the IGM
### 3.1 Mean optical depth of HI and HeII absorption
The optical depth for HI Ly$`\alpha `$ absorption from an IGM in ionization equilibrium with density $`\rho _b`$ is given by
$$\tau (z)=6.5\times 10^4\left(\frac{\omega _b}{0.019}\right)^2\left(\frac{h}{0.65}\right)^1\frac{(1+z)^6}{E(z)}\frac{T_4^{0.76}}{(\mathrm{\Gamma }_{\mathrm{HI}}/10^{12}\mathrm{s}^1)}\left(\frac{\rho _b(z)}{\overline{\rho }_b(z)}\right)^2,$$
(4)
(e.g. Peebles 1993 §23), where $`\mathrm{\Gamma }_{\mathrm{HI}}`$ is the photoionization rate for HI and $`T_4`$ is the temperature of the IGM in units of $`10^4`$ K. Variations in $`\rho _b(z)`$ along the line-of-sight will produce a ‘fluctuating Gunn-Peterson’ effect. An observed absorption line spectrum can therefore be inverted to infer the clustering of the baryon distribution as pioneered by Croft et al. (1998, see Section 3c below).
The mean HI and HeII optical depths as a function of redshift are plotted in Figure 3 for our fiducial model S. The open squares show the optical depths derived using half the amplitude of the HM UV background. With this choice of photoionizing background, the mean HI optical depth of the simulation matches observations quite well over the redshift range $`2`$$`4`$. Evidently, the amplitude of the photoionizing background can be balanced by variations in other cosmological parameters according to equation (4) to preserve the match to observations. The simulations thefore imply that
$$\left(\frac{\omega _b}{0.0125}\right)^2\left(\frac{h}{0.5}\right)^1\mathrm{\Omega }_m^{1/2}\left(\frac{0.5\mathrm{\Gamma }_{\mathrm{HM}}}{\mathrm{\Gamma }_{\mathrm{HI}}}\right)\left(\frac{6\times 10^3}{T}\right)^{0.76}1.$$
(5)
This type of criterion has been used by Weinberg et al. (1997) to set a crude lower limit to the baryon density. At redshifts $`2`$$`3`$ one can be reasonably confident of the HM model of the photoionizing background at around the Lyman edge, because the quasar luminosity function is quite well determined at these redshifts (see Section 3 of Madau’s article in these proceedings). The HM model provides a lower limit to the photoionizing flux because it ignores additional photoionizing radiation from star formation. The temperature of a photoionized IGM cannot exceed a few times $`10^4`$K, but its precise value depends on the past thermal and ionization history of the IGM.<sup>1</sup><sup>1</sup>1Although a highly ionized IGM is in ionization equilibrium, it is not in thermal equilibrium and so retains a memory of the way in which it was heated. Uncertainties in the temperature of the IGM lead to an additional source of uncertainty in using equation (5) to derive a bound on the baryon density. Nevertheless, equation (5) implies $`\omega _b\stackrel{>}{}0.02\mathrm{\Omega }_m^{1/2}`$, interestingly close to the baryon density of $`\omega _b=0.019`$ inferred from primorial nucleosynthesis and the deuterium abundances measured in quasar spectra (Burles and Tytler 1998, Burles, Kirkman and Tytler 1999). The observed HI optical depth therefore leads to a consistent picture in which most of the baryons in the Universe at $`z2`$$`4`$ belong to the diffuse photoionized IGM.
The HeII optical depth is shown in the lower panel of Figure 3. The open squares show $`\tau _{\mathrm{HeII}}`$ computed from the simulation using the same amplitude for the photoionizing background that provides a good match to $`\tau _{\mathrm{HI}}`$. The results from the simulation lie below the observations at all redshifts, suggesting that the photoinizing background has a softer spectrum than computed by HM. In photoionization equilibrium, the optical depth in HeII is related to the optical depth in HI by $`\tau _{\mathrm{HeII}}/\tau _{\mathrm{HI}}\mathrm{\Gamma }_{\mathrm{HI}}/\mathrm{\Gamma }_{\mathrm{HeII}}J_{\mathrm{HI}}/J_{\mathrm{HeII}}`$, and so $`\tau _{\mathrm{HeII}}`$ can be raised by softening the photoionizing spectrum appropriately. By lowering $`\mathrm{\Gamma }_{\mathrm{HeII}}/\mathrm{\Gamma }_{\mathrm{HI}}`$ by a factor of two compared to the HM model of Figure 1, model S can match the observations at $`z\stackrel{<}{}2.8`$, but cannot match the high HeII optical depths at $`z\stackrel{>}{}2.9`$ found in the recent HST-STIS observations of the quasar Q0302-003 by Heap et al. (2000).
One possible interpretation of these results is that the typical quasar spectrum adopted by HM is too hard and that HeII reionization is delayed until a redshift $`z3`$, corresponding to the abrupt change in HeII opacity observed by Heap et al. (2000). Additional arguments to support this picture are discussed in Section (3d). If this interpretation is correct, then AGN and ‘mini-quasars’ cannot produce much photoionizing radiation capable of doubly-ionizing helium at redshifts $`z\stackrel{>}{}3`$ (see the article by Rees in this volume).
### 3.2 Evolution of the column density distribution
It has been known for many years that Ly$`\alpha `$ forest shows strong cosmological evolution (e.g. Sargent et al. 1980). Over a narrow range in redshift the evolution can be approximated by a power law
$$\frac{dN}{dz}(1+z)^ϵ,$$
(6)
where $`N`$ is the number of lines above a threshold rest frame equivalent width (typically $`W>0.32\AA `$). From high resolution Keck observations, Kim et al. (1997) find $`ϵ=2.78\pm 0.71`$ in the redshift range $`2<z<3.5`$ . Williger et al. (1994) find that the evolution is still stronger at higher redshifts with $`ϵ>4`$ at $`z>4`$. In contrast, observations with the HST indicate much weaker evolution at low redshifts with $`ϵ=0.48\pm 0.62`$ for $`z<1`$ (Morris et al. 1991; Bahcall et al. 1991, 1993; Impey et al. 1996).
The evolution of the Ly$`\alpha `$ forest, including the low rates of evolution at redshifts $`z\stackrel{<}{}2`$ can be reproduced quite simply in CDM models (Theuns, Leonard and Efstathiou 1998b; Davé et al. 1999). This is illustrated in Figure 4 which shows the evolution of the number of Ly$`\alpha `$ lines within a given range of column density for our reference model S. (The lines are identified by fitting Voigt profiles to simulated Keck spectra using the line-fitting program VPFIT, Webb 1987). As in Section (3a) we adopt the HM model of the photoionizing background with an amplitude divided by a factor of $`2`$ to match the observed optical depth in HI absorption. This model reproduces the observed column density distribution accurately over the column density range $`10^{12.5}\mathrm{cm}^2\stackrel{<}{}N_{\mathrm{HI}}\stackrel{<}{}10^{15}\mathrm{cm}^2`$ (see Figure 2 of Theuns et al. 1998b) and, as Figure 4 shows, also reproduces the observed rates of evolution as a function of column density. In particular, the decrease in the rate of evolution of the Ly$`\alpha `$ lines found from HST observations arises from the steep decline in the photoionizing background at $`z\stackrel{<}{}2`$ caused by the rapid drop in quasar numbers at low redshift.
### 3.3 Reconstruction of the matter power spectrum
Equations (2.1) and (3.1) can be combined to write the observed transmitted flux in terms of fluctuations in the baryon density,
$$F=\mathrm{exp}\left[A(\rho _b/\overline{\rho }_b)^\beta \right],\beta (2.760.76\gamma ).$$
(7)
Croft et al. (1998) have used this relation to infer the 1-dimensional power spectrum $`P_{1D}(k)`$ of the baryon fluctuations, from which the three-dimensional power spectrum can be recovered by differentiation,
$$P(k)=\frac{2\pi }{k}\frac{d}{dk}P_{1D}(k).$$
(8)
Croft et al. (1998) calibrate the amplitude of the matter power spectrum by comparing against numerical simulations. The procedure is not completely straightforward and we refer the reader to Croft et al. (1998) for details. By testing their inversion algorithm against numerical simulations these authors find that the amplitude and shape of the underlying matter power spectrum can be recovered accurately and that the recovery is insensitive to uncertainties in the equation of state. A variant of this technique, that incorporates a correction for the distortion of the clustering pattern by peculiar velocities, is described by Nusser and Haehnelt (1999).
Croft et al. (1999) describe an application of their method to a sample of $`19`$ quasar spectra spanning the redshift range $`2.08`$$`3.23`$. They recover $`P(k)`$ at $`z2.5`$ over the (comoving) wavenumber range $`2\pi /k2`$$`12h^1\mathrm{Mpc}`$ and find that it is well fitted by a power law $`P(k)k^n`$ with $`n=2.25\pm 0.28`$, consistent with what is expected from CDM models. This important result is the first attempted determination of the matter power spectrum at high redshift. The amplitude of the inferred power spectrum is high compared to that expected for spatially flat CDM models with $`\mathrm{\Omega }_m=1`$ normalized to reproduce the abundance of rich clusters at the present day. The best fitting CDM models have $`\mathrm{\Omega }_m+0.2\mathrm{\Omega }_\mathrm{\Lambda }0.46`$ (Weinberg et al. 1999; Phillips et al. 2000), in agreement with the parameters $`\mathrm{\Omega }_m0.3`$, $`\mathrm{\Omega }_\mathrm{\Lambda }0.7`$ derived from combining observations of anisotropies in the cosmic microwave background radiation and distant Type Ia supernovae (e.g. Efstathiou et al. 1999 and references therein). Constraints on neutrino masses from estimates of $`P(k)`$ at high redshift are discussed by Croft, Hu and Davé (1999).
### 3.4 Widths of the Ly$`\alpha `$ lines
The widths of the Ly$`\alpha `$ lines are usually characterised by the broadening parameter $`b`$ determined by fitting Voigt profiles. If the line width is caused by thermal broadening, the parameter $`b`$ is related to the temperature of the gas by
$$b=\left(\frac{2kT}{m_p}\right)^{1/2}=12.8T_4^{1/2}\mathrm{km}\mathrm{s}^1.$$
(9)
In reality, a number of mechanisms in addition to thermal broadening contribute to the widths of the lines, e.g. the differential Hubble flow across the absorbing region and the smoothing of small scale fluctuations by gas pressure (Bryan et al. 1999; Theuns, Schaye and Haenhelt 2000).
The first sets of simulations of the Ly$`\alpha `$ forest in CDM models appeared to show good agreement with the observed $`b`$-parameter distributions. However, subsequent simulations with higher spatial resolution produced a larger fraction of narrower lines than observed (Theuns et al. 1998a; Bryan et al. 1999). This is illustrated in Figure 5 which shows observations of the $`b`$ distribution at $`z3`$ compared to the results from numerical simulations. The figure to the left shows the distribution derived by fitting Voigt profiles to simulated spectra from reference model S (solid line). The lines in this model are clearly too narrow, suggesting that the temperature of the IGM is too low. From the assymptotic relation (2.2) one can see that the temperature of the IGM can be raised by increasing the baryon density and by lowering $`\mathrm{\Omega }_m`$ (the assymptotic temperature is extremely insentive to the amplitude of the photoionizing background long after reionization). The right hand panel shows the effect of increasing the baryon density to $`\omega _b=0.025`$. The dashed lines in each panel show the effect of lowering $`\mathrm{\Omega }_m`$ and introducing a cosmological constant so that the universe remains spatially flat. These variations in cosmological parameters can go some way to resolving the conflict with observations (Theuns et al. 1999), but cannot provide an exact match unless the baryon density is much higher than the value favoured from primordial nucleosythesis. This suggests that we are missing a significant heating source of the IGM. A number of mechanisms have been suggested that might boost the temperature of the IGM, e.g. photoelectric heating of dust grains (Nath, Sethi and Shchekinov 1999) and Compton heating by the hard X-ray background (Madau and Efstathiou 1999). However, the most plausible explanation (Abel and Haenhelt 1999) is that the simulations underestimate the temperature at $`z3`$ because they assume an optically thin IGM (and also a uniform photoionizing background that has already been reprocessed by Lyman $`\alpha `$ absorbing clouds, Haardt and Madau 1996). This is inconsistent and the simulations should properly include the effects of radiative transfer while the medium is still optically thick prior to complete reionization, because in this regime every photoionizing photon is absorbed and contributes to heating the IGM. Abel and Haehnelt estimate that correct inclusion of radiative transfer during the epoch of HeII reionization might increase the temperature of the IGM by a factor of $`2`$, perhaps enough to resolve the discrepancy with the observed $`b`$-parameter distributions illustrated in Figure 5. The idea that the temperature of the IGM at $`z3`$ is boosted by HeII reionization is supported by the analysis of the equation of state of the IGM described in the next subsection.
### 3.5 Constraining the equation of state of the IGM
As we have mentioned, a number of physical mechanisms contribute to the breadth of the $`b`$-distribution. However, the minimum line-width is set by the temperature of the gas, which in turn depends on the density (cf. Figure 1). By fitting the cut-off in the $`b`$-distribution as a function of column density (the ‘$`b(N)`$’ distribution) one can therefore reconstruct the effective equation of state of the IGM (Schaye, et al. 1999; Ricotti, Gnedin and Shull 2000; Bryan and Machacek 2000).
The method is illustrated by Figure 6 which shows the $`b(N)`$ distributions derived by applying the VPFIT line-fitting program to two cosmological simulations. Results for the standard reference model S are plotted in Figure (6b). Figure 6(a) shows a hotter model, which has the same parameters as model Lb plotted in Figure 5 but with the HeI and HeII photoheating rates doubled over those of the HM model (crudely representing ‘radiative transfer’ effects during the reionization of helium). The dashed line in the figure shows the best fit to the lower envelope of the $`b(N)`$ distribution of Figure (6a) determined by applying the iterative fitting algorithm of Schaye et al. (1999). The same line is plotted in Figure (6b) and passes almost through the middle of the $`b(N)`$ distribution of the colder model. The lower envelope of the $`b(N)`$ distribution is clearly a sensitive indicator of the characteristic temperature of the Ly$`\alpha `$ clouds and can be accurately determined from fits to high resolution quasar spectra.
What is more difficult is to convert the fit to the lower envelope of the $`b(N)`$ distribution, $`b=b_0(N/N_0)^{\mathrm{\Gamma }1}`$, into the parameters $`T_0`$ and $`\gamma `$ of the effective equation of state of the IGM (equation 2.1). This is done by calibrating $`b_0`$ and $`\mathrm{\Gamma }`$ against $`T_0`$ and $`\gamma `$ using numerical simulations with different equations of state (see Schaye et al. 1999 for details).
The results of applying this technique to nine high-quality quasar spectra spanning the redshift range $`2.0<z<4.5`$ are shown in Figure 7 (Schaye et al. 2000). Except for the two lowest redshift quasars, the Ly$`\alpha `$ forest spectra were divided in two to reduce the effects of redshift evolution and signal-to-noise variations across a single spectrum. From $`z3.5`$ to $`z3`$, the inferred temperature at the mean density of the IGM, $`T_0`$, increases and the gas is close to isothermal ($`\gamma 1`$). This behaviour differs markedly from that expected if helium were fully reionized at higher redshift. For example, the solid lines show the evolution of the equation of state in a simulation with the HM background. In this simulation, both hydrogen and helium are fully ionized by $`z4.5`$ and the temperature of the IGM declines slowly as the universe expands tending to the assymptotic form of equation (2.2). This model cannot account for the peak in the temperature at $`z3`$ inferred from the observations. Instead, we associate the peak in $`T_0`$, and the low value of $`\gamma `$, with reheating caused by the second reionization of helium (HeII $``$ HeIII). This interpretation is supported by the abrupt change in HeII opacity at $`z3`$ from the measurements of Heap et al. (2000) discussed in Section 3(a). It is also consistent with evidence from the ratio SiIV/CIV that the spectrum of the photoionizing background hardens abruptly at $`z3`$ (Songaila and Cowie 1996; Songaila 1998).<sup>2</sup><sup>2</sup>2Note, however, that Boksenberg, Sargent and Rauch (1998) find a more gradual change of the SiIV/CIV ratio with redshift.
The dashed lines in Figure 7 show a ‘designer model’ with parameters tuned to fit the observations. This simulation has a much softer UV background at high redshift than the model shown by the solid line, and consequently HeII reionizes at $`z3.2`$. In addition, the photoheating rates have been enhanced during reionization to boost the temperature of the IGM (again, crudely modelling the heating of optically thick gas). This model is suggestive that the jump in temperature at $`z3`$ may be associated with HeII reionization, but clearly more realistic simulations that include radiative transfer are required to determine the evolution of the equation of state during the reionization phase more accurately. Further observations would also be useful to assess how steeply the IGM temperature changes between $`z=4`$ and $`z=3`$.
## 4 Summary and Outlook
The work reviewed in this article provides a powerful case that the Ly$`\alpha `$ forest arises from a space-filling, highly photoionized diffuse IGM that contains most of the baryonic material in the Universe at high redshift. This model is a natural outcome of CDM theories of structure formation and can account for many observed properties of the Ly$`\alpha `$ forest in quantitative detail. The general features of the model thus seem to us to be reasonably secure.
However, a more detailed analysis of the thermal history of the IGM requires simulations that incorporate radiative transfer and a model for the spatial distribution of ionizing sources. Such calculations are now being done (Abel, Norman and Madau 1999; Gnedin 2000; Madau these proceedings) but the computational problems are formidable. Some outstanding problems that deserve further attention include:
$``$ detailed simulations of the inhomogeneous reionization of hydrogen and helium;
$``$ extending the analysis of Ly$`\alpha `$ line widths to redshifts $`\stackrel{>}{}4`$, perhaps leading to constraints on the epoch of reionization of hydrogen;
$``$ analysis of inhomogeneities in the temperature of the IGM. Are there, for example, regions in the spectra of quasars in which Ly$`\alpha `$ line-widths are systematically broader or narrower than in other regions?
$``$ further observations of absorption gaps in HeII Ly$`\alpha `$ absorption (reported by Heap et al. (2000) and others) and the development of a model to understand their sizes;
$``$ searching for signatures of outflows around protogalaxies in the Ly$`\alpha `$ forest;
$``$ determining the mean metallicity of the IGM and understanding how the metals were transported from protogalaxies.
###### Acknowledgements.
J. Schaye thanks the Isaac Newton Trust for financial support and PPARC for the award of a studentship. We also thank Anthony Leonard, Michael Rauch and Wal Sargent for their contributions to this work.
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# About the Dirac Equation with a 𝛿 potential
## Acknowledgments
R. B. acknowledges support from Fondecyt, under grant 1990427. M. L. acknowledges support from Fondecyt, under grant 1980577. H. C. acknowledges financial support from a PUC fellowship.
## References
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# 1 Introduction
## 1 Introduction
Semiclassical theories for spectral statistics have been developed to find an explanation for the observed universality in energy spectra of quantum systems with a chaotic classical limit, the agreement of correlations in energy spectra with those between eigenvalues of random matrices . They are based on semiclassical trace formulas that approximate the density of states in terms of classical trajectories . It has been shown by these theories that in the asymptotic limit of long-range correlations two-point correlation functions do coincide with those of random matrix theory . These results are based on mean properties of periodic orbits . To go beyond the leading asymptotic term requires information about correlations between periodic orbits which are presently not available .
One of the expectations, on basis of the random matrix hypothesis , is that a perturbation of a chaotic system should not change the statistical distribution of the energy levels of the quantum system, if it does not change the chaotic nature of the classical motion. In the present article we investigate, on the level of the semiclassical approximation, whether the perturbation by a point-like scatterer has this property. One argument in favour of this invariance is that the semiclassical approximation for the density of states is not changed in leading order of $`\mathrm{}`$ for this perturbation. The influence of the scatterer is described semiclassically by a certain class of trajectories, so-called diffractive orbits that start from the scatterer and return to it. They contribute to the density of states in higher order of $`\mathrm{}`$ than the leading order contribution from periodic orbits.
The present article is motivated by the observation in that a scatterer could nevertheless have an influence on spectral statistics. When spectral correlation functions are calculated by using mean properties of diffractive orbits, the so-called diagonal approximation, they show modifications which, in general, do not vanish in the semiclassical limit ($`\mathrm{}0`$). In order that this does not lead to deviations from random matrix statistics, these terms have to be cancelled by off-diagonal terms which contain information about correlations between different trajectories. As remarked above, the calculation of correlations between trajectories is an unsolved problem in general systems. For the diffractive orbits that describe the influence of a scatterer, however, off-diagonal terms can be calculated explicitly. This is done in the following sections. The results show that diagonal and off-diagonal terms indeed cancel each other. Furthermore, the results can be used to investigate parametric spectral correlations, i. e. correlations between spectra of the system for different parameter values, where the parameter is the strength of the scatterer. It is shown that the parametric spectral correlations are universal for small changes of the parameter.
## 2 The spectral form factor
The perturbation by a point-like scatterer is represented, formally, by a delta-potential $`\widehat{H}=\widehat{H}_0+\lambda \delta (𝒓𝒓_0)`$, where $`\lambda `$ and $`𝒓_0`$ are the strength and position of the scatterer, respectively. In more than one dimension such a delta-potential is, however, not well defined. For example, it leads to a divergent expression for the Green function. The problem can be regularised by the method of self-adjoint extensions, leading to a one-parameter family of Hamiltonians. A detailed monograph with references on the history and on applications of delta-like potentials is . We use in the following the property that the semiclassical approximation for the density of states has the same form as in the geometrical theory of diffraction , (see also for applications on spectral statistics).
We consider chaotic systems whose Hamiltonian is given in terms of a scalar and a vector potential. Billiard systems can be included in this description by letting the scalar potential be infinite outside the billiard region. The statistical distribution of the energy levels is investigated by semiclassically approximating the spectral form factor. We restrict to two-dimensional systems in order to keep the notation simple, but analogous calculations can be performed in higher dimensions.
The spectral form factor is defined as Fourier transform of the spectral two-point correlation function
$$K(\tau )=_{\mathrm{}}^{\mathrm{}}\frac{\text{d}\eta }{\overline{d}(E)}d_{\text{osc}}\left(E+\frac{\eta }{2}\right)d_{\text{osc}}\left(E\frac{\eta }{2}\right)_E\mathrm{exp}\left(2\pi i\eta \tau \overline{d}(E)\right).$$
(1)
The function $`d_{\text{osc}}(E)=d(E)\overline{d}(E)`$ is the oscillatory part of the density of states, and $`\overline{d}(E)`$ is the smooth part which is given in two dimensions by $`\overline{d}(E)\mathrm{\Sigma }(E)(2\pi \mathrm{})^2`$, $`E\mathrm{}`$, where $`\mathrm{\Sigma }(E)`$ is the volume of the surface of constant energy in phase space. The statistics is evaluated by averaging over an energy interval that is small in comparison to $`E`$ but contains a large number of energy levels.
The semiclassical approximation for $`K(\tau )`$ is obtained by inserting into (1) the approximation for the oscillatory part of the density of states
$$d_{\text{osc}}(E)\frac{1}{\pi \mathrm{}}\mathrm{Re}\underset{\gamma }{}A_\gamma \mathrm{exp}\left(\frac{i}{\mathrm{}}S_\gamma (E)\right).$$
(2)
In systems with a delta-like potential the sum in (2) runs over all periodic orbits , and further over all diffractive orbits that start from the scatterer and return to it an arbitrary number of times $`n`$ . An example for a double-diffractive orbit ($`n=2`$) in a billiard system is shown in figure 1. For $`n`$-fold diffractive orbits the amplitude $`A_\gamma `$ has an $`\mathrm{}`$-dependence of $`\mathrm{}^{n/2}`$, and $`S_\gamma `$ denotes the action of an orbit.
With (2) one obtains the following approximation for the spectral form factor
$$K(\tau )=\frac{1}{2\pi \mathrm{}\overline{d}(E)}\underset{\gamma ,\gamma ^{}}{}A_\gamma A_\gamma ^{}^{}\mathrm{exp}\left\{\frac{i}{\mathrm{}}\left(S_\gamma (E)S_\gamma ^{}(E)\right)\right\}\delta \left(T\frac{T_\gamma +T_\gamma ^{}}{2}\right)_E,$$
(3)
where $`T=2\pi \mathrm{}\overline{d}(E)\tau `$, and $`T_\gamma `$ is the period of an orbit. For small values of $`\tau `$ one can evaluate the double sum in (3) in the diagonal approximation . One obtains in this way from the periodic orbits the correct random matrix result $`K(\tau )\frac{2}{\beta }\tau `$, $`\tau 0`$, where $`\beta =1`$ or $`2`$ for systems with or without time-reversal symmetry, respectively.
The diagonal contributions from diffractive orbits to the form factor have been calculated in . The result for $`n`$-fold diffractive orbits is
$$K_d^{(n)}(\tau )=\frac{|𝒟|^{2n}}{(2\beta )^n}\frac{\tau ^{n+1}}{n},$$
(4)
where $`𝒟`$ is the diffraction coefficient for the diffraction on the singularity of the potential . It can be parameterised in the following form
$$𝒟=\frac{2\pi }{i\frac{\pi }{2}\gamma \mathrm{log}\left(\frac{ka}{2}\right)}.$$
(5)
Here $`k=\sqrt{2m(EV(𝒓_0))}/\mathrm{}`$, $`𝒓_0`$ is the position of the scatterer, $`a`$ is a parameter describing the strength of the potential, and $`\gamma `$ is Euler’s constant. In order that the terms (4) do not lead to a deviation from random matrix statistics they have to be cancelled by off-diagonal terms involving diffractive orbits. By calculating off-diagonal terms for order $`\tau ^2`$ and $`\tau ^3`$ we show in the following that such a cancellation does indeed occur.
We note that the diffraction coefficient satisfies the identities
$$|𝒟|^2=4\mathrm{Im}𝒟,|𝒟|^4=8(|𝒟|^2\mathrm{Re}𝒟^2),$$
(6)
that will be used in the following. The first of these relations expresses the conservation of probability, and the second is a consequence of the first one.
## 3 First-order correction
The first-order correction to the diagonal approximation for the form factor arises from off-diagonal terms in (3) between periodic orbits and single-diffractive orbits. In leading order, these orbits are only correlated if the diffractive orbit follows the periodic orbit very closely. This happens, if the diffractive orbit is almost periodic, i. e. if the final momentum is almost identical to the initial momentum. An example is shown in figure 2.
The periodic orbit can be described by linearising the motion around the diffractive orbit. The condition that a trajectory in the vicinity of the diffractive orbit is periodic leads to the following equation
$$\left(\begin{array}{c}\delta \\ p_v\gamma \end{array}\right)=M\left(\begin{array}{c}\delta \\ p_v(\gamma \epsilon )\end{array}\right).$$
(7)
Here $`M`$ is the stability matrix of the diffractive orbit, $`\epsilon `$ is the angle between the initial and final direction of the diffractive orbit, $`\gamma `$ is the angle between the direction along the periodic orbit and the final direction of the diffractive orbit, and $`\delta `$ is the spacial distance between periodic and diffractive orbit (see figure 2). The quantity $`p_v`$ is defined as mass times velocity, $`p_v=mv`$. The index $`v`$ is used in order to distinguish it from the canonical momentum $`p`$ in systems with magnetic field. The stability matrix for the motion in a magnetic field is discussed in the appendix.
In the linear approximation the difference in actions is obtained by expanding the action up to second order
$$\mathrm{\Delta }S(E)=S^{\text{po}}(E)S^{\text{1do}}(E)p_v\delta \epsilon +\frac{1}{2}(\mathrm{\Delta }p_f\mathrm{\Delta }p_i)\delta =\frac{1}{2}\delta \epsilon p_v,$$
(8)
where $`\mathrm{\Delta }p_f`$ and $`\mathrm{\Delta }p_i`$ are the differences between the initial and final momenta of the periodic orbit and the diffractive orbit, respectively. The solution of the linear equation (7) yields the following relation between $`\delta `$ and $`\epsilon `$
$$\delta =\frac{M_{12}}{\mathrm{Tr}M2}\epsilon p_v,$$
(9)
so that $`\mathrm{\Delta }S(E)`$ depends quadratically on $`\epsilon `$.
With this approximation the off-diagonal terms are calculated. The amplitude of the diffractive orbit is given by
$$A_\gamma ^{\text{1do}}=\frac{T_\gamma 𝒟}{4\pi p_v}\sqrt{\frac{2\pi \mathrm{}}{|(M_\gamma )_{12}|}}\mathrm{exp}\left\{i\frac{\pi }{2}\nu _\gamma i\frac{3\pi }{4}\right\},$$
(10)
where $`T_\gamma `$ is the time along the orbit, $`M_\gamma `$ is its stability matrix, and $`\nu _\gamma `$ is the number of conjugate points along the orbit. For the periodic orbit the corresponding amplitude is
$$A_\gamma ^{\text{po}}=\frac{T_\gamma }{\sqrt{|\mathrm{Tr}M_\gamma 2|}}\mathrm{exp}\left\{i\frac{\pi }{2}\mu _\gamma \right\}.$$
(11)
The stability matrix is the same for both orbits in leading order, but the Maslov index of the periodic orbit can differ from the number of conjugate points $`\nu _\gamma `$ by 1
$$\mu _\gamma =\nu _\gamma +\frac{1}{2}(1\kappa _\gamma ),\kappa _\gamma =\mathrm{sign}\left(\frac{(M_\gamma )_{12}}{\mathrm{Tr}M_\gamma 2}\right).$$
(12)
In the following we sum over all diffractive orbits that are almost periodic. This is done by applying first the sum rule for diffractive orbits for which the angle difference between initial and final direction has a fixed value $`\epsilon `$, and then integrating over the angle $`\epsilon `$. The sum rule is given by
$$\underset{\gamma }{}^{(\epsilon )}\frac{1}{|(M_\gamma )_{12}|}\delta (TT_\gamma )\frac{2\pi p_v^2}{\mathrm{\Sigma }(E)},$$
(13)
where $`\mathrm{\Sigma }(E)`$ is the volume of the energy shell. (It is implied in (13) that the left-hand side is smoothed over small intervals of $`T`$ and $`\epsilon `$ in order to obtain a non-singular expression.)
Finally, one has to determine the multiplicity factor of the contribution. First, each off-diagonal term in (3) has a corresponding complex conjugate partner. If the summation is carried out over only one of these terms one has to take twice the real part of the sum. Furthermore, the periodic orbit and the diffractive orbit both have multiplicities $`2\beta ^1`$, but a particular constellation occurs $`2\beta ^1`$ times in the sum over $`\epsilon `$ (for systems with time-reversal symmetry for $`\epsilon `$ and $`\epsilon `$), so the total multiplicity is $`g=2\beta ^1`$.
Inserting the amplitudes (10) and (11), and the action difference (8) with (9) into (3) we obtain
$`K_{\text{off}}^{(1)}(\tau )`$ $`={\displaystyle \frac{g}{2\pi \mathrm{}\overline{d}(E)}}\mathrm{\hspace{0.17em}2}\mathrm{Re}{\displaystyle _{\mathrm{}}^{\mathrm{}}}\text{d}\epsilon {\displaystyle \underset{\gamma }{}^{(\epsilon )}}A_\gamma ^{\text{1do}}(A_\gamma ^{\text{po}})^{}\mathrm{exp}\left({\displaystyle \frac{i}{\mathrm{}}}\mathrm{\Delta }S_\gamma (E)\right)\delta (TT_\gamma )`$
$`={\displaystyle \frac{4}{2\pi \mathrm{}\overline{d}(E)\beta }}\mathrm{Re}{\displaystyle _{\mathrm{}}^{\mathrm{}}}\text{d}\epsilon {\displaystyle \underset{\gamma }{}^{(\epsilon )}}{\displaystyle \frac{T_\gamma ^2𝒟\sqrt{2\pi \mathrm{}}\delta (TT_\gamma )}{4\pi p_v\sqrt{|(M_\gamma )_{12}(\mathrm{Tr}M_\gamma 2)|}}}e^{\frac{i\epsilon ^2p_v^2}{2\mathrm{}}\frac{(M_\gamma )_{12}}{\mathrm{Tr}M_\gamma 2}i\frac{\pi }{4}(2+\kappa _\gamma )}`$
$`={\displaystyle \frac{4}{2\pi \mathrm{}\overline{d}(E)\beta }}\mathrm{Re}{\displaystyle _{\mathrm{}}^{\mathrm{}}}\text{d}\epsilon ^{}{\displaystyle \underset{\gamma }{}^{(\epsilon ^{})}}{\displaystyle \frac{T_\gamma ^2𝒟\sqrt{2\pi \mathrm{}}\delta (TT_\gamma )}{4\pi p_v|(M_\gamma )_{12}|}}e^{\frac{\epsilon _{}^{}{}_{}{}^{2}p_v^2}{2\mathrm{}}i\frac{\pi }{2}}`$
$`={\displaystyle \frac{2T^2}{2\pi \mathrm{}\overline{d}(E)\beta }}\mathrm{Im}{\displaystyle _{\mathrm{}}^{\mathrm{}}}\text{d}\epsilon ^{}{\displaystyle \frac{𝒟p_v\sqrt{2\pi \mathrm{}}}{\mathrm{\Sigma }(E)}}e^{\frac{\epsilon _{}^{}{}_{}{}^{2}p_v^2}{2\mathrm{}}}`$
$`={\displaystyle \frac{2\tau ^2}{\beta }}\mathrm{Im}𝒟.`$ (14)
The integration over $`\epsilon `$ can be carried out from minus to plus infinity since the main contribution comes from the vicinity of $`\epsilon =0`$. Furthermore, the following steps have been carried out. First the integration variable has been changed to make the exponent independent of the stability matrix. Then the sum rule (13) has been applied, assuming that the distribution of angles between initial and final momenta of a diffractive orbit is independent of the distribution of the elements of the stability matrix
$$\underset{\gamma }{}^{(\epsilon )}\frac{1}{|(M_\gamma )_{12}|}\delta (TT_\gamma )\underset{\gamma }{}^{(\epsilon ^{})}\frac{1}{|(M_\gamma )_{12}|}\delta (TT_\gamma ),\epsilon ^{}=\epsilon \sqrt{\frac{i(M_\gamma )_{12}}{\mathrm{Tr}M_\gamma 2}},$$
(15)
and finally the integration has been carried out.
$`K_{\text{off}}^{(1)}(\tau )`$ is the leading order correction to the diagonal approximation for the form factor and it cancels exactly the diagonal contribution from single diffractive orbits ((4) with $`n=1`$). This can be seen by using (6)
$$K_\text{d}^{(1)}(\tau )+K_{\text{off}}^{(1)}(\tau )=\frac{|𝒟|^2}{2\beta }\tau ^2+\frac{2}{\beta }\tau ^2\mathrm{Im}𝒟=0.$$
(16)
It shows that the presence of a point-like scatterer does not modify the spectral form factor up to order $`\tau ^2`$ in systems with a chaotic classical limit.
In order to find the geometries of orbits which contribute to a given order in $`\tau `$ it is helpful to count the orders of $`\mathrm{}`$. The $`m`$-th order off-diagonal correction to the form factor is a $`\tau ^{m+1}`$-term with a coefficient that has to be $`\mathrm{}`$-independent. The prefactor of the double sum over orbits in (3) is of order $`\mathrm{}^1`$ and the product of the amplitudes of a $`n_1`$-fold and a $`n_2`$-fold diffractive orbit is of order $`\mathrm{}^{(n_1+n_2)/2}`$, where periodic orbits are denoted here as 0-fold diffractive orbits. The conversion of time $`T^{m+1}`$ into $`\tau ^{m+1}`$ gives an order $`\mathrm{}^{m+1}`$, which yields altogether an order of $`\mathrm{}^{(2mn_1n_2)/2}`$. Furthermore, every integration over a small parameter $`\epsilon `$ gives an additional order $`\mathrm{}^{1/2}`$, if the action difference is quadratic in this parameter. As a consequence, $`2mn_1n_2`$ small parameters are necessary in order that the prefactor of $`\tau ^{(m+1)}`$ is $`\mathrm{}`$-independent. For the first order correction in this section ($`m=n_1=1`$ and $`n_2=0`$) this estimate gives one small parameter $`\epsilon `$.
## 4 Second-order corrections
For the second-order corrections we consider orbits that return twice to the region in coordinate space from which they started. These orbits are close to double-diffractive orbits. Double-diffractive orbits have the semiclassical amplitude
$$A_\gamma ^{\text{2do}}=\frac{\mathrm{}T_\gamma 𝒟^2}{16\pi p_v^2}\frac{1}{\sqrt{|(R_\gamma )_{12}(L_\gamma )_{12}|}}\mathrm{exp}\left\{i\frac{\pi }{2}(\nu _{\gamma ,L}+\nu _{\gamma ,R})i\frac{3\pi }{2}\right\},$$
(17)
where $`T_\gamma `$ is the total time along the trajectory, $`L_\gamma `$ and $`R_\gamma `$ are the stability matrices for the two loops (’left’ and ’right’), and $`\nu _{\gamma ,L}`$ and $`\nu _{\gamma ,R}`$ are the number of conjugate points along the loops. (17) is the amplitude for one particular sequence in which the loops are traversed. In systems without time-reversal symmetry the degeneracy of the trajectory is thus two, meaning that there is another trajectory with exactly the same semiclassical amplitude and action. This trajectory traverses first the second loop of $`\gamma `$, and then the first loop of $`\gamma `$. In systems with time-reversal symmetry the degeneracy is eight.
The sum rule for double-diffractive orbits is given by
$$\underset{\gamma }{\overset{(\epsilon _1,\mathrm{},\epsilon _n)}{}}\frac{1}{|(R_\gamma )_{12}(L_\gamma )_{12}|}\delta (TT_\gamma )\frac{(2\pi p_v)^4}{\mathrm{\Sigma }(E)^2}\frac{T}{(2\pi )^n},$$
(18)
if there are $`n`$ restrictions to the four directions of the velocities at the point from which the trajectories start and to which they return. As will be seen in the following, it follows from this sum rule that the contributions are of order $`\tau ^3`$ (there is a factor $`T`$ from every semiclassical amplitude, and a factor $`T`$ from the sum rule).
There are several possibilities in which a double-diffractive orbit can have an action which is almost identical to the action of a single-diffractive or a periodic orbit. A necessary condition is that there is always at least one small relative angle between the different initial and final directions of the orbit at the scattering point. In order to find the relevant cases one has to consider all possibilities and take into account the $`\mathrm{}`$-argument that was given at the end of section 3. The result is that there are three relevant configurations for systems without time-reversal symmetry and five configurations for systems with time-reversal symmetry. They are discussed in the next sections.
### 4.1 Correlations between double-diffractive and single-diffractive orbits
Correlations between double-diffractive and single-diffractive orbits exist if the double-diffractive orbit is almost single-diffractive. This occurs if the final velocity of one loop deviates by a small angle $`\epsilon `$ from the initial velocity of the other loop. An example is shown in figure 3. There is only one small parameter here which agrees with the estimate $`2mn_1n_2`$ for $`m=2`$, $`n_1=2`$ and $`n_2=1`$.
The further calculations are done analogously to the last section. The motion in the vicinity of the double-diffractive orbit is linearised and one obtains in this approximation a condition for the neighbouring single-diffractive orbit
$$\left(\begin{array}{c}\delta \\ p_v\gamma _2\end{array}\right)=L\left(\begin{array}{c}0\\ p_v\gamma _1\end{array}\right),\left(\begin{array}{c}0\\ p_v\gamma _3\end{array}\right)=R\left(\begin{array}{c}\delta \\ p_v(\gamma _2\epsilon )\end{array}\right).$$
(19)
The angles $`\gamma _1`$, $`\gamma _2`$, $`\gamma _3`$ and the distance $`\delta `$ are shown in figure 3. The difference in action is obtained by expanding the action up to second order
$$\mathrm{\Delta }S(E)=S^{\text{1do}}(E)S^{\text{2do}}(E)=\frac{1}{2}\delta \epsilon p_v=\frac{\epsilon ^2p_v^2}{2}\frac{L_{12}R_{12}}{M_{12}},$$
(20)
where $`M=RL`$ is the stability matrix of the single-diffractive orbit. The last step in (20) follows from the solution of (19) for $`\delta `$.
The number of conjugate points $`\nu _\gamma `$ along the single-diffractive orbit can differ from the sum of the number of conjugate points along the two loops, $`\nu _{\gamma ,L}`$ and $`\nu _{\gamma ,R}`$, by 1. The general condition for this is
$$\nu _\gamma =\nu _{\gamma ,L}+\nu _{\gamma ,R}+\frac{1}{2}(1\sigma _\gamma ),\sigma _\gamma =\mathrm{sign}\left(\frac{(L_\gamma )_{12}(R_\gamma )_{12}}{(M_\gamma )_{12}}\right).$$
(21)
The single- and double-diffractive orbits have multiplicity $`2\beta ^1`$ and $`8\beta ^2`$, respectively, but each configuration occurs for $`2\beta ^1`$ different values of $`\epsilon `$. Therefore the total multiplicity is $`g=8\beta ^2`$. The contribution to the form factor from all pairs of orbits is thus given by
$`K_{\text{off}}^{(2a)}(\tau )`$ $`={\displaystyle \frac{g}{2\pi \mathrm{}\overline{d}(E)}}\mathrm{\hspace{0.17em}2}\mathrm{Re}{\displaystyle _{\mathrm{}}^{\mathrm{}}}\text{d}\epsilon {\displaystyle \underset{\gamma }{}^{(\epsilon )}}A_\gamma ^{\text{2do}}(A_\gamma ^{\text{1do}})^{}\mathrm{exp}\left({\displaystyle \frac{i}{\mathrm{}}}\mathrm{\Delta }S_\gamma (E)\right)\delta (TT_\gamma )`$
$`={\displaystyle \frac{16}{2\pi \mathrm{}\overline{d}(E)\beta ^2}}\mathrm{Re}{\displaystyle _{\mathrm{}}^{\mathrm{}}}\text{d}\epsilon {\displaystyle \underset{\gamma }{}^{(\epsilon )}}{\displaystyle \frac{\mathrm{}T_\gamma ^2𝒟^2𝒟^{}\sqrt{2\pi \mathrm{}}\delta (TT_\gamma )}{64\pi ^2p_v^3\sqrt{|(L_\gamma )_{12}(R_\gamma )_{12}(M_\gamma )_{12}|}}}e^{\frac{i\epsilon ^2p_v^2}{2\mathrm{}}\frac{(L_\gamma )_{12}(R_\gamma )_{12}}{(M_\gamma )_{12}}i\frac{\pi }{4}(2+\sigma _\gamma )}`$
$`={\displaystyle \frac{1}{2\pi \mathrm{}\overline{d}(E)\beta ^2}}\mathrm{Re}{\displaystyle \underset{\gamma }{}^{(\epsilon ^{})}}{\displaystyle \frac{\mathrm{}^2T^2𝒟^2𝒟^{}\delta (TT_\gamma )}{2\pi p_v^4|(L_\gamma )_{12}(R_\gamma )_{12}|}}e^{i\frac{\pi }{2}}`$
$`={\displaystyle \frac{\tau ^3}{\beta ^2}}\mathrm{Im}(𝒟^2𝒟^{}).`$ (22)
Here we slightly abbreviated the procedure of (3) and performed the integration directly.
### 4.2 Correlations between double-diffractive orbits and periodic orbits
In order that a double-diffractive orbit is close to a periodic orbit it has to be almost periodic. This means that the final direction of each loop has to be almost identical to the initial direction of the other loop, so there are two small relative angles as is shown in figure 4.
The linearisation of the motion in the vicinity of the double-diffractive orbit leads to the following condition for the periodic orbit
$$\left(\begin{array}{c}\delta _2\\ p_v\gamma _2\end{array}\right)=L\left(\begin{array}{c}\delta _1\\ p_v(\gamma _1\epsilon _1)\end{array}\right),\left(\begin{array}{c}\delta _1\\ p_v\gamma _1\end{array}\right)=R\left(\begin{array}{c}\delta _2\\ p_v(\gamma _2\epsilon _2)\end{array}\right).$$
(23)
The angles $`\gamma _i`$ and distances $`\delta _i`$ are defined analogously as before in terms of the local coordinate systems that are oriented along the two final directions of the diffractive orbit.
The difference in actions is given by
$`\mathrm{\Delta }S(E)`$ $`=S^{\text{po}}(E)S^{\text{2do}}(E)`$
$`={\displaystyle \frac{p_v}{2}}(\delta _1\epsilon _1+\delta _2\epsilon _2)`$
$`={\displaystyle \frac{p_v^2}{2}}{\displaystyle \frac{(RL)_{12}\epsilon _1^2+(LR)_{12}\epsilon _2^2+2(L_{12}+R_{12})\epsilon _1\epsilon _2}{\mathrm{Tr}M2}},`$ (24)
where $`M=RL`$ is the stability matrix for the periodic orbit. Equation (4.2) can be written in the terms of a symmetric matrix $`A`$ such that $`\mathrm{\Delta }S(E)=\frac{1}{2}p_v^2_{i,j}A_{ij}\epsilon _i\epsilon _j`$. The number of negative eigenvalues of $`A`$ is given by the number of sign changes in the sequence of sub-determinants ($`1,A_{11},detA`$). With $`A_{11}=M_{12}/(\mathrm{Tr}M2)`$ and $`detA=L_{12}R_{12}/(\mathrm{Tr}M2)`$ one finds that the two signs of the eigenvalues are given by
$$\kappa =\mathrm{sign}\left(\frac{M_{12}}{\mathrm{Tr}M2}\right),\sigma =\mathrm{sign}\left(\frac{L_{12}R_{12}}{M_{12}}\right).$$
(25)
The Maslov index of the periodic orbit can now differ by zero, one or two from the sum of conjugate points along the left and right loop. The criterion for this is the combination of (12) and (21) and has the form
$$\mu =\nu _L+\nu _R+\frac{1}{2}(2\kappa \sigma ),$$
(26)
where $`\kappa `$ and $`\sigma `$ are given in (25).
The multiplicities of double-diffractive and periodic orbits are $`8\beta ^2`$ and $`2\beta ^1`$, respectively, but a particular configuration of them occurs $`4\beta ^1`$ times in the integral over the angles (for example, for systems without time reversal symmetry the two angles can be interchanged), so the total multiplicity is $`g=4\beta ^2`$.
After inserting the amplitudes (17) and (11) and the action difference (4.2) into (3) we obtain
$`K_{\text{off}}^{(2b)}(\tau )`$ $`={\displaystyle \frac{g}{2\pi \mathrm{}\overline{d}(E)}}\mathrm{\hspace{0.17em}2}\mathrm{Re}{\displaystyle _{\mathrm{}}^{\mathrm{}}}\text{d}\epsilon _1\text{d}\epsilon _2{\displaystyle \underset{\gamma }{}^{(\epsilon _1,\epsilon _2)}}A_\gamma ^{\text{2do}}(A_\gamma ^{\text{po}})^{}\mathrm{exp}\left({\displaystyle \frac{i}{\mathrm{}}}\mathrm{\Delta }S_\gamma (E)\right)\delta (TT_\gamma )`$
$`={\displaystyle \frac{8}{2\pi \mathrm{}\overline{d}(E)\beta ^2}}\mathrm{Re}{\displaystyle _{\mathrm{}}^{\mathrm{}}}\text{d}\epsilon _1\text{d}\epsilon _2{\displaystyle \underset{\gamma }{}^{(\epsilon _1,\epsilon _2)}}{\displaystyle \frac{\mathrm{}T_\gamma ^2𝒟^2\delta (TT_\gamma )e^{\frac{ip_v^2}{2\mathrm{}}{\scriptscriptstyle A_{i,j}\epsilon _i\epsilon _j}i\frac{\pi }{4}(4+\sigma _\gamma +\kappa _\gamma )}}{16\pi p_v^2\sqrt{|(L_\gamma )_{12}(R_\gamma )_{12}(\mathrm{Tr}M_\gamma 2)|}}}`$
$`={\displaystyle \frac{1}{2\pi \mathrm{}\overline{d}(E)\beta ^2}}\mathrm{Re}{\displaystyle \underset{\gamma }{}^{(\epsilon _1^{},\epsilon _2^{})}}{\displaystyle \frac{\mathrm{}^2T^2𝒟^2\delta (TT_\gamma )}{p_v^4|(L_\gamma )_{12}(R_\gamma )_{12}|}}e^{i\pi }`$
$`={\displaystyle \frac{\tau ^3}{\beta ^2}}\mathrm{Re}(𝒟^2).`$ (27)
### 4.3 Correlations between pairs of single-diffractive orbits
For exactly the same kind of double-diffractive orbits as in the last section, there is one further type of correlation that has to be considered. It occurs because there are two possible ways in which the double-diffractive orbit can be deformed into a nearby single-diffractive orbit, and consequently, there are correlations between these single-diffractive orbits.
The action difference between the two orbits can be obtained from the action difference between each of these orbits and the double-diffractive orbit (20)
$$\mathrm{\Delta }S(E)=S^{\text{1do,1}}(E)S^{\text{1do,2}}(E)=\frac{1}{2}\delta _1\epsilon _1p_v+\frac{1}{2}\delta _2\epsilon _2p_v=\frac{p_v^2}{2}\left(\frac{L_{12}R_{12}}{N_{12}}\epsilon _1^2\frac{L_{12}R_{12}}{M_{12}}\epsilon _2^2\right),$$
(28)
where $`M=RL`$ and $`N=LR`$ are the stability matrices of the two single-diffractive orbits. The number of conjugate points along the orbits are given by
$`\nu _1=\nu _L+\nu _R+{\displaystyle \frac{1}{2}}(1\sigma _1),`$ $`\nu _2=\nu _L+\nu _R+{\displaystyle \frac{1}{2}}(1\sigma _2),`$
$`\sigma _1=\mathrm{sign}\left({\displaystyle \frac{L_{12}R_{12}}{N_{12}}}\right),`$ $`\sigma _2=\mathrm{sign}\left({\displaystyle \frac{L_{12}R_{12}}{M_{12}}}\right).`$ (29)
The multiplicities of the two orbits is both $`2\beta ^1`$ and, as before, the configuration occurs $`4\beta ^1`$ times in the double integral over the angles, so the total multiplicity is $`g=\beta ^1`$.
Inserting the amplitude (10) with stability matrix $`M`$ and $`N`$, respectively, and the action difference (28) into (3) one obtains
$`K_{\text{off}}^{(2c)}(\tau )`$ $`={\displaystyle \frac{g}{2\pi \mathrm{}\overline{d}(E)}}\mathrm{\hspace{0.17em}2}\mathrm{Re}{\displaystyle _{\mathrm{}}^{\mathrm{}}}\text{d}\epsilon _1\text{d}\epsilon _2{\displaystyle \underset{\gamma }{}^{(\epsilon _1,\epsilon _2)}}A_\gamma ^{\text{1do,2}}(A_\gamma ^{\text{1do,1}})^{}\mathrm{exp}\left({\displaystyle \frac{i}{\mathrm{}}}\mathrm{\Delta }S_\gamma (E)\right)\delta (TT_\gamma )`$
$`={\displaystyle \frac{2}{2\pi \mathrm{}\overline{d}(E)\beta }}\mathrm{Re}{\displaystyle _{\mathrm{}}^{\mathrm{}}}\text{d}\epsilon _1\text{d}\epsilon _2{\displaystyle \underset{\gamma }{}^{(\epsilon _1,\epsilon _2)}}{\displaystyle \frac{2\pi \mathrm{}T_\gamma ^2|𝒟|^2\delta (TT_\gamma )e^{\frac{i}{\mathrm{}}\mathrm{\Delta }S_\gamma (E)i\frac{\pi }{4}(\sigma _{\gamma ,1}\sigma _{\gamma ,2})}}{16\pi ^2p_v^2\sqrt{|(M_\gamma )_{12}(N_\gamma )_{12}|}}}`$
$`={\displaystyle \frac{1}{2\pi \mathrm{}\overline{d}(E)\beta }}\mathrm{Re}{\displaystyle \underset{\gamma }{}^{(\epsilon _1^{},\epsilon _2^{})}}{\displaystyle \frac{\mathrm{}^2T^2|𝒟|^2\delta (TT_\gamma )}{2p_v^4|(L_\gamma )_{12}(R_\gamma )_{12}|}}`$
$`={\displaystyle \frac{\tau ^3}{2\beta }}|𝒟|^2.`$ (30)
The contributions $`K_{\text{off}}^{(2a)}(\tau )`$ $`K_{\text{off}}^{(2b)}(\tau )`$ and $`K_{\text{off}}^{(2c)}(\tau )`$ are the only second-order off-diagonal corrections in systems without time-reversal symmetry. As will be shown in the following, these contributions cancel exactly the diagonal term $`K_\text{d}^{(2)}(\tau )`$ for $`\beta =2`$. For systems with time-reversal symmetry there are further contributions. They arise from the possibility that one trajectory can follow one loop of another trajectory in the same direction, but the other loop in the time-reversed direction. We assume in the following that the relevant contributions come from orbits which are close to double diffractive orbits in coordinate space and we evaluate their contributions in the next two subsections.
### 4.4 Correlations between pairs of single-diffractive orbits involving time-reversed loops
The first possibility involves two single-diffractive orbits. These orbits follow closely a double-diffractive orbit and traverse one loop in the same direction and the other loop in the opposite direction. In order that this can occur there must be one loop which has a very small opening angle, and it must be almost aligned to the final direction of the other loop as shown in figure 6.
The action difference and the indices are given by the equations (28) and (4.3) but now with $`M=RL`$ and $`N=R^iL`$ where $`R^i`$ is the stability matrix for the time-reversed second loop. In term of the elements of $`R`$ the matrix $`R^i`$ is given by
$$R^i=\left(\begin{array}{cc}R_{22}& R_{12}\\ R_{21}& R_{11}\end{array}\right).$$
(31)
The multiplicity of each orbit is two, and the configuration occurs two times in the integral over the angles, so the total multiplicity is $`g=2`$. Inserting the amplitude (10) with stability matrix $`M`$ and $`N`$, respectively, and the action difference (28) into (3) results in
$`K_{\text{off}}^{(2d)}(\tau )`$ $`={\displaystyle \frac{g}{2\pi \mathrm{}\overline{d}(E)}}\mathrm{\hspace{0.17em}2}\mathrm{Re}{\displaystyle _{\mathrm{}}^{\mathrm{}}}\text{d}\epsilon _1\text{d}\epsilon _2{\displaystyle \underset{\gamma }{}^{(\epsilon _1,\epsilon _2)}}A_\gamma ^{\text{1do,2}}(A_\gamma ^{\text{1do,1}})^{}\mathrm{exp}\left({\displaystyle \frac{i}{\mathrm{}}}\mathrm{\Delta }S_\gamma (E)\right)\delta (TT_\gamma )`$
$`={\displaystyle \frac{4}{2\pi \mathrm{}\overline{d}(E)}}\mathrm{Re}{\displaystyle _{\mathrm{}}^{\mathrm{}}}\text{d}\epsilon _1\text{d}\epsilon _2{\displaystyle \underset{\gamma }{}^{(\epsilon _1,\epsilon _2)}}{\displaystyle \frac{2\pi \mathrm{}T_\gamma ^2|𝒟|^2\delta (TT_\gamma )e^{\frac{i}{\mathrm{}}\mathrm{\Delta }S_\gamma (E)i\frac{\pi }{4}(\sigma _{\gamma ,1}\sigma _{\gamma ,2})}}{16\pi ^2p_v^2\sqrt{|(M_\gamma )_{12}(N_\gamma )_{12}|}}}`$
$`={\displaystyle \frac{1}{2\pi \mathrm{}\overline{d}(E)}}\mathrm{Re}{\displaystyle \underset{\gamma }{}^{(\epsilon _1^{},\epsilon _2^{})}}{\displaystyle \frac{\mathrm{}^2T^2|𝒟|^2\delta (TT_\gamma )}{p_v^4|(L_\gamma )_{12}(R_\gamma )_{12}|}}`$
$`=\tau ^3|𝒟|^2.`$ (32)
### 4.5 Correlations between single-diffractive orbits and periodic orbits involving time-reversed loops
The last relevant configuration occurs if all initial and final velocities of the double-diffractive orbit lie almost in one line as in figure 7. Then there exist neighbouring single-diffractive and periodic orbits which follow one loop in the same direction and the other loop in the opposite direction.
The action difference between the periodic orbit and the diffractive orbit is obtained from (20) and (4.2)
$`\mathrm{\Delta }S(E)`$ $`=S^{\text{po}}(E)S^{\text{1do}}(E)`$
$`={\displaystyle \frac{p_v}{2}}(\delta _1\epsilon _1+\delta _2\epsilon _2)+{\displaystyle \frac{p_v}{2}}\delta _3\epsilon _3`$
$`={\displaystyle \frac{p_v^2}{2}}{\displaystyle \frac{(RL)_{12}\epsilon _1^2+(LR)_{12}\epsilon _2^2+2(L_{12}+R_{12})\epsilon _1\epsilon _2}{\mathrm{Tr}M2}}+{\displaystyle \frac{1}{2}}\epsilon _3^2p_v^2{\displaystyle \frac{L_{12}R_{12}}{N_{12}}},`$ (33)
where $`M=RL`$ and $`N=R^iL`$ are the stability matrices of the periodic and diffractive orbit, respectively, and the indices $`\nu `$ and $`\mu `$ of the diffractive and the periodic orbit are
$$\nu =\nu _L+\nu _R+\frac{1}{2}(1\sigma _2),\mu =\nu _L+\nu _R+\frac{1}{2}(2\kappa \sigma _1),$$
$$\sigma _2=\mathrm{sign}\left(\frac{L_{12}R_{12}}{N_{12}}\right),\kappa =\mathrm{sign}\left(\frac{M_{12}}{\mathrm{Tr}M2}\right),\sigma _1=\mathrm{sign}\left(\frac{L_{12}R_{12}}{M_{12}}\right).$$
(34)
For each double-diffractive orbit like the one in figure 7 there are two periodic orbits of multiplicity two and four single-diffractive orbits of multiplicity two in the vicinity, but only half of the possible pairs involve exactly one time reversed loop which makes a total of 16. Furthermore, the double-diffractive orbit occurs 8 times in the integral over the angles and the total multiplicity is thus $`g=2`$.
Inserting the amplitudes (10) and (11) and the action difference (4.5) into (3) yields
$`K_{\text{off}}^{(2e)}(\tau )`$ $`={\displaystyle \frac{g}{2\pi \mathrm{}\overline{d}(E)}}\mathrm{\hspace{0.17em}2}\mathrm{Re}{\displaystyle _{\mathrm{}}^{\mathrm{}}}\text{d}\epsilon _1\text{d}\epsilon _2\text{d}\epsilon _3{\displaystyle \underset{\gamma }{}^{(\epsilon _1,\epsilon _2,\epsilon _3)}}A_\gamma ^{\text{1do}}(A_\gamma ^{\text{po}})^{}\mathrm{exp}\left({\displaystyle \frac{i}{\mathrm{}}}\mathrm{\Delta }S_\gamma (E)\right)\delta (TT_\gamma )`$
$`={\displaystyle \frac{4}{2\pi \mathrm{}\overline{d}(E)}}\mathrm{Re}{\displaystyle _{\mathrm{}}^{\mathrm{}}}\text{d}\epsilon _1\text{d}\epsilon _2\text{d}\epsilon _3{\displaystyle \underset{\gamma }{}^{(\epsilon _1,\epsilon _2,\epsilon _3)}}{\displaystyle \frac{T_\gamma ^2𝒟\sqrt{2\pi \mathrm{}}\delta (TT_\gamma )e^{\frac{i}{\mathrm{}}\mathrm{\Delta }S_\gamma (E)i\frac{\pi }{4}(2\sigma _{\gamma ,2}+\kappa _\gamma +\sigma _{\gamma ,1})}}{4\pi p_v\sqrt{|(M_\gamma )_{12}(\mathrm{Tr}M_\gamma 2)|}}}`$
$`={\displaystyle \frac{1}{2\pi \mathrm{}\overline{d}(E)}}\mathrm{Re}{\displaystyle \underset{\gamma }{}^{(\epsilon _1^{},\epsilon _2^{},\epsilon _3^{})}}{\displaystyle \frac{4\pi \mathrm{}^2T^2𝒟\delta (TT_\gamma )}{p_v^4|(L_\gamma )_{12}(R_\gamma )_{12}|}}e^{i\frac{\pi }{2}}`$
$`=2\tau ^3\mathrm{Im}𝒟.`$ (35)
The sum of all contributions can be written in the form
$`K^{(2)}(\tau )`$ $`=K_\text{d}^{(2)}(\tau )+K_{\text{off}}^{(2a)}(\tau )+K_{\text{off}}^{(2b)}(\tau )+K_{\text{off}}^{(2c)}(\tau )+K_{\text{off}}^{(2d)}(\tau )+K_{\text{off}}^{(2e)}(\tau )`$
$`={\displaystyle \frac{\tau ^3}{\beta ^2}}\left({\displaystyle \frac{1}{8}}|𝒟|^4+|𝒟|^2\mathrm{Im}𝒟\mathrm{Re}𝒟^2+{\displaystyle \frac{\beta }{2}}|𝒟|^2+(2\beta )|𝒟|^2+2(2\beta )\mathrm{Im}𝒟\right)`$
$`={\displaystyle \frac{\tau ^3}{\beta ^2}}\left({\displaystyle \frac{1}{8}}|𝒟|^4+|𝒟|^2\mathrm{Im}𝒟\mathrm{Re}𝒟^2+|𝒟|^2\right)`$
$`=0,`$ (36)
which can be seen by using (6). This shows that off-diagonal terms cancel the diagonal term also in this order. It implies that the form factor is determined by periodic orbits alone, because the different terms which involve diffractive orbits cancel each other. This might be true also for other point-like sources of diffraction like e. g. Aharonov-Bohm flux lines (see ), although a quantitative analysis would require here the use of uniform approximations.
## 5 Universality in parametric correlations
The cancellation of off-diagonal and diagonal terms is conform with the expected universality of spectral statistics in chaotic systems. Universality is, however, not only expected in the properties of single systems, but also in the way in which system properties vary when a parameter of the system is changed. For example, random matrix theory makes predictions about correlations between densities of states for different parameter values. The semiclassical calculation of diagonal and off-diagonal terms allows to test this prediction for systems with a point-like scatterer where the parameter is the strength of the scatterer.
In analogy to the spectral form factor, a parametric form factor can be defined as Fourier transform of the parametric two-point correlation function
$$K(\tau ,x)=_{\mathrm{}}^{\mathrm{}}\frac{\text{d}\eta }{\overline{d}(E)}d_{\text{osc}}(E+\frac{\eta }{2},X+\frac{x}{2})d_{\text{osc}}(E\frac{\eta }{2},X\frac{x}{2})_E\mathrm{exp}\left(2\pi i\eta \tau \overline{d}(E)\right).$$
(37)
Here $`x`$ is the parameter difference between two systems. In order for this statistics to be universal the parameter has to be chosen in a particular way. The requirement is that the variance of the velocities, the derivatives of the unfolded energies with respect to the parameter ($`\epsilon _n/x`$), is equal to unity. The unfolded energies are obtained from the quantum energies of the system by a scaling that leads to mean level distance of one.
For random matrix ensembles, the parametric two-point correlation function was derived in in the context of disordered metallic systems. For the GUE-result, which we discuss first, the Fourier transform in (37) can be evaluated in a closed form. It results in
$$K^{\text{GUE}}(\tau ,x)=\{\begin{array}{cc}\frac{\mathrm{sinh}(2\pi ^2x^2\tau ^2)}{2\pi ^2x^2\tau }\mathrm{exp}(2\pi ^2x^2\tau )\hfill & \text{if }\tau <1\text{,}\hfill \\ \frac{\mathrm{sinh}(2\pi ^2x^2\tau )}{2\pi ^2x^2\tau }\mathrm{exp}(2\pi ^2x^2\tau ^2)\hfill & \text{if }\tau >1\text{,}\hfill \end{array}$$
(38)
and has for small values of $`\tau `$ the expansion
$$K^{\text{GUE}}(\tau ,x)=\tau 2\pi ^2x^2\tau ^2+2\pi ^4x^4\tau ^3+\mathrm{}.$$
(39)
We examine in the following whether the perturbation by a point-like scatterer leads to universal correlations. For large parameter differences $`x`$ the parametric correlations for these systems cannot be expected to be universal. The treatment of a delta-scatterer by the method of self-adjoint extensions leads to a quantisation condition with the property that there is exactly one eigenvalue of the perturbed system within each pair of neighbouring eigenvalues of the unperturbed system. This puts a restriction to the movement of eigenvalues when the parameter is changed. For this reason, one can expect universal properties only for small parameter differences.
We choose first $`a`$ in (5) as parameter of the system. The two densities in (37) differ then only in the diffraction coefficient $`𝒟`$. As a consequence, the results for the parametric form factor can be obtained directly from the spectral form factor without further calculations. One has to express the contributions to the form factor, (16) and (4.5), in terms of $`𝒟_1`$ and $`𝒟_1^{}`$ and replace $`𝒟_1^{}`$ by $`𝒟_2^{}`$. For $`\beta =2`$ one obtains in this way
$`\stackrel{~}{K}_{\text{sc}}(\tau ,x)\stackrel{~}{K}_{\text{sc}}(\tau ,0)`$ $`{\displaystyle \frac{\tau ^2}{4}}[𝒟_1𝒟_2^{}2i𝒟_1+2i𝒟_2^{}]+{\displaystyle \frac{\tau ^3}{32}}[(𝒟_1𝒟_2^{})^2`$
$`4i𝒟_1𝒟_2^{}(𝒟_1𝒟_2^{})4𝒟_1𝒟_14𝒟_2^{}𝒟_2^{}+8𝒟_1𝒟_2^{}]`$
$`={\displaystyle \frac{\tau ^2}{4}}\left[𝒟_1𝒟_2^{}2i𝒟_1+2i𝒟_2^{}\right]+{\displaystyle \frac{\tau ^3}{32}}\left[𝒟_1𝒟_2^{}2i𝒟_1+2i𝒟_2^{}\right]^2.`$ (40)
A first point to notice is that the parametric form factor in (5) is, in general, not real. This seems to be in contrast to the random matrix result (38) which is real. The reason for this lies, however, in the correct choice of the unfolding procedure. The definition (37) yields only the universal form factor in case that the mean density of states $`\overline{d}(E)`$ does not depend on the parameter of the system. However, in the present case the mean density changes slightly with the parameter $`x`$ of the system, and this leads to a slight shift of the spectrum with $`x`$ . As a consequence, the argument $`\eta `$ of the two-level correlation function is shifted, and its Fourier transform, the form factor, is multiplied by a term of the form $`e^{ic\tau }`$, where $`c`$ is determined by the shift of the levels. By rewriting (5) up to the considered order in $`\tau `$ in the form
$`\stackrel{~}{K}_{\text{sc}}(\tau ,x)\stackrel{~}{K}_{\text{sc}}(\tau ,0)`$ $`\left(\tau +{\displaystyle \frac{\tau ^2}{4}}\mathrm{Re}\left[𝒟_1𝒟_2^{}2i𝒟_1+2i𝒟_2^{}\right]+{\displaystyle \frac{\tau ^3}{32}}\left(\mathrm{Re}\left[𝒟_1𝒟_2^{}2i𝒟_1+2i𝒟_2^{}\right]\right)^2\right)`$
$`\times \mathrm{exp}\left({\displaystyle \frac{i\tau }{4}}\mathrm{Im}\left[𝒟_1𝒟_2^{}2i𝒟_1+2i𝒟_2^{}\right]\right)\tau ,`$ (41)
one can extract the result that corresponds to a proper unfolding by dropping the exponential, and one obtains
$`K_{\text{sc}}(\tau ,x)K_{\text{sc}}(\tau ,0)`$ $`{\displaystyle \frac{\tau ^2}{4}}\mathrm{Re}\left[𝒟_1𝒟_2^{}2i𝒟_1+2i𝒟_2^{}\right]+{\displaystyle \frac{\tau ^3}{32}}\left(\mathrm{Re}\left[𝒟_1𝒟_2^{}2i𝒟_1+2i𝒟_2^{}\right]\right)^2`$
$`={\displaystyle \frac{\tau ^2}{8}}|𝒟_1𝒟_2|^2+{\displaystyle \frac{\tau ^3}{128}}|𝒟_1𝒟_2|^4,`$ (42)
where (6) has been used.
A comparison with (39) shows that it has now the same form as the random matrix result. The remaining step is to connect the parameter $`a`$ with the universal parameter $`x`$. Since we consider small parameter differences, $`x`$ is given by $`x=\mathrm{\Delta }a\sigma _v`$, where $`\sigma _v`$ is the square root of the variance of the velocities with respect to the parameter $`a`$. In order to evaluate $`\sigma _v`$ we employ a semiclassical method which expresses it in the form
$$\sigma _v^2=\left(\frac{\epsilon _n}{a}\right)^2_E=\underset{\eta 0}{lim}2\pi \eta \left[d_v^\eta (\epsilon )\right]^2_E,$$
(43)
where $`d_v^\eta `$ is a Lorentzian smoothed density of states, that is weighted by the velocities and expressed semiclassically by
$$d_v^\eta (\epsilon )=\underset{n}{}\frac{\epsilon _n}{a}\frac{1}{\pi }\frac{\eta }{(\epsilon \epsilon _n)^2+\eta ^2}\mathrm{Im}\frac{}{a}\underset{\gamma }{}\frac{A_\gamma }{\pi T_\gamma }\mathrm{exp}\left\{\frac{i}{\mathrm{}}S_\gamma \frac{\eta T_\gamma }{\mathrm{}\overline{d}(E)}\right\}.$$
(44)
Here $`\eta `$ is the width of the Lorentzian and all other quantities are defined as before.
For small parameter differences we can express the derivative of the diffraction coefficient with respect to $`a`$ by $`(𝒟_2𝒟_1)/\mathrm{\Delta }a`$. After inserting (44) into (43) the leading order contribution comes from single-diffractive orbits, and we evaluate the double sum in the diagonal approximation with the sum rule (13) and amplitudes (10)
$`\sigma _v^2`$ $`=\underset{\eta 0}{lim}{\displaystyle _0^{\mathrm{}}}\text{d}T{\displaystyle \underset{\gamma }{}}{\displaystyle \frac{\eta \mathrm{}\delta (TT_\gamma )}{4\pi ^2p_v^2\beta |M_{12}|}}{\displaystyle \frac{|𝒟_2𝒟_1|^2}{(\mathrm{\Delta }a)^2}}\mathrm{exp}\left\{{\displaystyle \frac{2\eta T_\gamma }{\mathrm{}\overline{d}(E)}}\right\}`$
$`=\underset{\eta 0}{lim}{\displaystyle _0^{\mathrm{}}}\text{d}T{\displaystyle \frac{\eta \mathrm{}}{\beta \mathrm{\Sigma }(E)}}{\displaystyle \frac{|𝒟_2𝒟_1|^2}{(\mathrm{\Delta }a)^2}}\mathrm{exp}\left\{{\displaystyle \frac{2\eta T}{\mathrm{}\overline{d}(E)}}\right\}`$
$`={\displaystyle \frac{2}{\beta }}{\displaystyle \frac{|𝒟_2𝒟_1|^2}{(4\pi \mathrm{\Delta }a)^2}},`$ (45)
For $`\beta =2`$, the universal parameter is given by $`x=\mathrm{\Delta }a\sigma _v=|𝒟_2𝒟_1|/(4\pi )`$, and after insertion into (5) we reproduce the random matrix result (39)
$$K_{sc}(\tau ,x)K_{sc}(\tau ,0)2\pi ^2x^2\tau ^2+2\pi ^4x^4\tau ^3.$$
(46)
For the GOE ensemble, the parametric density correlation function is given by a triple-integral which cannot be expressed in closed form . We consider here only the first-order correction for which the GOE-result can be obtained from the asymptotic form of the correlation function for long-range correlations in : $`K^{\text{GOE}}(\tau ,x)K^{\text{GOE}}(\tau ,0)2\pi ^2x^2\tau ^2`$. Again we find an agreement with the semiclassical result $`K_{sc}(\tau ,x)K_{sc}(\tau ,0)\tau ^2|𝒟_1𝒟_2|^2/4=2\pi ^2x^2\tau ^2`$.
## 6 Discussion
We have investigated in this article the influence of a point-like scatterer on the spectral statistics of quantum systems with chaotic classical limit. It has been shown that the modification of the form factor $`K(\tau )`$ due to the scatterer can be evaluated systematically in a semiclassical approximation. The expansion of the form factor in powers of $`\tau `$ corresponds on the semiclassical side to an expansion in the number of loops of the diffractive orbits. We have calculated off-diagonal contributions to the $`\tau ^2`$\- and $`\tau ^3`$-term, but the method can be extended to higher order terms.
The results lead to the conclusion that the delta-perturbation does not modify the form factor. Off-diagonal terms from pairs of different diffractive orbits and from pairs of diffractive and periodic orbits cancel exactly the diagonal terms from diffractive orbits. This requires the existence of correlations between different orbits. These correlations arise from pairs of orbits which are very close in coordinate space.
The results provide a support for the random-matrix conjecture. They imply, up to the considered order, that the statistics of chaotic systems are invariant under the perturbation by a point-like scatterer as is expected from universality. They show also that correlations between two energy spectra for different parameter values are universal, provided that the parameter difference is small. Furthermore, they indicate indirectly that the spectral statistics of the unperturbed system (and thus also of the perturbed system) are identical with those of random matrix theory. The reason for this is that independent results on the invariance of spectral statistics under a delta-perturbation are based on the assumption that the unperturbed energy levels and wave functions have random matrix distributions . Since the semiclassical results show this invariance for chaotic systems, the combination of both results provides a theoretical indication that chaotic systems follow the random matrix hypothesis.
Finally, the results are a support for the semiclassical method. They show that semiclassical approximations are capable to go beyond the leading term in $`\tau `$ and are an appropriate tool for investigating spectral statistics in the semiclassical limit.
I would like to thank K. Richter and P. Šeba for helpful discussions. After completion of this article I learned about work by E. Bogomolny, P. Leboeuf and C. Schmit with related semiclassical results for the first-order correction.
## Appendix A Stability matrix for the motion in a magnetic field
The stability matrix $`\stackrel{~}{M}`$ of a trajectory determines infinitesimal orthogonal deviations from the final point of the trajectory in terms of the deviations from the initial point of the trajectory
$$\left(\begin{array}{c}\text{d}r_f\\ \text{d}p_f\end{array}\right)=\stackrel{~}{M}\left(\begin{array}{c}\text{d}r_i\\ \text{d}p_i\end{array}\right).$$
(47)
In systems with a magnetic field the momentum has the form $`𝒑=m𝒗+\frac{q}{c}𝑨(𝒓)`$, where $`q`$ is the charge of the particle and $`𝑨`$ is the vector potential. In this case it is often more convenient to consider a matrix $`M`$ that describes deviations of the velocities instead of those of the momenta. The relation between both matrices is given by
$$\left(\begin{array}{c}\text{d}r_f\\ m\text{d}v_f\end{array}\right)=M\left(\begin{array}{c}\text{d}r_i\\ m\text{d}v_i\end{array}\right),M=A_f^1\stackrel{~}{M}A_i,A_{i,f}=\left(\begin{array}{cc}1& 0\\ a_{i,f}& 1\end{array}\right),$$
(48)
where $`a_{i,f}=\frac{q}{c}(\widehat{n}_{i,f})(\widehat{n}_{i,f}𝑨(𝒓_{i,f}))`$, and $`\widehat{n}_{i,f}`$ is the direction orthogonal to the trajectory at the initial and final point of the trajectory, respectively. The matrix $`M`$ has unit determinant and satisfies $`M_{12}=\stackrel{~}{M}_{12}`$ and $`\mathrm{Tr}M=\mathrm{Tr}\stackrel{~}{M}+\stackrel{~}{M}_{12}(a_ia_f)`$. In cases where the initial and final points are identical and the initial and final velocities differ by a small angle $`\epsilon `$, the traces are identical in leading order of $`\epsilon `$. Since the semiclassical approximations in this article involve only $`\stackrel{~}{M}_{12}`$ and $`\mathrm{Tr}\stackrel{~}{M}2`$, we express all quantities in terms of $`M`$ instead of $`\stackrel{~}{M}`$, and we use also the term stability matrix for it.
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# No Ghost State in the Brane World
## I Introduction
Recently, there have been lots of interest in the phenomenon of localization of gravity proposed by Randall and Sundrum (RS) (for previous relevant work see references therein). RS assume a single positive tension 3-brane and a negative bulk cosmological constant in the five dimensional spacetime. By considering metric fluctuations from a background which is isomorphic to sections of $`AdS_5`$, they have shown that it reproduces the effect of four dimensional gravity localized on the brane without the need to compactify the extra dimension due to the “warping” in the fifth dimensional space. In more detail, the solution to linearized equations in the five dimensions results in a zero mode, which can be identified with the four dimensional massless graviton, and the massive continuum Kaluza-Klein (KK) modes. Surprisingly, the wavefunctions of the massive continuum KK modes are suppressed at the brane for small energies, and thus ordinary gravity localized on the brane is reproduced at large distances.
On the other hand, Gregory, Rubakov and Sibiryakov (GRS) have recently considered a brane model which is not asymptotically $`AdS_5`$, but Minkowski flat. In the GRS model, however, the ordinary 4D Newton potential is reproduced at intermediate scales only not because of the massless zero mode, but for the resonance of zero mass in the continuum KK spectrum . In Ref. , however, it is pointed out that the $`m_h0`$ limit of a massive graviton propagator does not reproduce the massless graviton propagator due to the missmatch of the number of polarizations. In this sense the GRS model of “quasi-localization” of gravity would differ from the RS model. Contrary to it, Csáki, Erlich and Hollowood recently have argued that in the presence of localized source at $`z=0`$ the bending of the brane exactly compensates for the effects of the extra polarization in the massive graviton propagator. Thus the graviton propagator at intermediate scales is equivalent to the massless propagator of the Einstein theory just as in the RS scenario. However, at ultra large scales this effective theory includes scalar anti-gravity . This problem may be cured by the RG analysis . Also the authors in Ref. point out that the mechanism to cancel the unwanted extra polarization leads to the presence of ghost. In order to have a well-defined theory, the ghost should disappear.
In this paper, we investigate the non-traceless metric fluctuations in the presence of uniform source along $`z`$-axis in the RS model. We introduce the trace field ($`h`$) here instead of $`\xi ^5`$ in Ref. . It shows that massive graviton modes contain ghost states which can be removed by assuming a further condition on the matter source. Our work corresponds to an alternative realization of the results in Ref. .
## II Linearized perturbations
The Randall-Sundrum model with a single domain wall (or brane) perpendicular to the infinite fifth direction can be described by the following action:
$$I=d^4x_{\mathrm{}}^{\mathrm{}}𝑑z\left[\frac{1}{16\pi G_5}\sqrt{\widehat{g}}(\widehat{R}2\mathrm{\Lambda })\sqrt{\widehat{g}_B}\sigma (z)+_M\right].$$
(1)
Here $`G_5`$ is the five dimensional Newton’s constant, $`\mathrm{\Lambda }`$ the bulk cosmological constant of five dimensinal spacetime, $`\widehat{g}_B`$ the determinant of the metric describing the brane, and $`\sigma (z)=\sigma \delta (z)`$, $`\sigma `$ the tension of the brane. $`I_M=d^4x𝑑z_M`$ denotes the matter action, and it contributes only in the linearized level. In this paper, we use the signature $`(,+,+,+,+)`$.
If we introduce a conformal factor as follows
$$ds^2=\widehat{g}_{MN}dx^Mdx^N=H^2g_{MN}dx^Mdx^N,$$
(2)
the field equation becomes
$`G_{MN}+3{\displaystyle \frac{_M_NH}{H}}3g_{MN}\left[{\displaystyle \frac{_P^PH}{H}}2{\displaystyle \frac{_PH^PH}{H^2}}\right]`$ (3)
$`=8\pi G_5\left[{\displaystyle \frac{\mathrm{\Lambda }}{8\pi G_5H^2}}g_{MN}{\displaystyle \frac{\sqrt{g_B}}{\sqrt{g}}}{\displaystyle \frac{|H|}{H^2}}\sigma (z)g_{\mu \nu }\delta _M^\mu \delta _N^\nu {\displaystyle \frac{2}{\sqrt{\widehat{g}}}}{\displaystyle \frac{\delta I_M}{\delta \widehat{g}^{MN}}}\right]`$ (4)
with the Einstein tensor $`G_{MN}`$ constructed from the metric $`g_{MN}`$. Now it is straightforward to see that, in the absence of matter source except for the domain wall itself (i.e., $`\delta I_M/\delta \widehat{g}^{MN}=0`$), the most general solution having a four dimensional Poincaré symmetry is
$$ds^2=H^2(z)(\eta _{\mu \nu }dx^\mu dx^\nu +dz^2),$$
(5)
where $`H=k|z|+1`$, $`\mathrm{\Lambda }=6k^2(<0)`$, and $`\sigma =3k/4\pi G_5`$.
Let us consider metric fluctuations around this background spacetime as follows :
$$g_{MN}=\eta _{MN}+h_{MN}.$$
(6)
Defining $`\overline{h}_{MN}=h_{MN}\frac{1}{2}\eta _{MN}h`$ where $`h=\eta ^{MN}h_{MN}`$, the linearized perturbation equation of Eq. (4) is
$`{\displaystyle \frac{1}{2}}\mathrm{}\overline{h}_{MN}+_{(M}^P\overline{h}_{N)P}{\displaystyle \frac{1}{2}}\eta _{MN}^P^Q\overline{h}_{PQ}{\displaystyle \frac{3^PH}{2H}}(_Mh_{NP}+_Nh_{MP}_Ph_{MN})`$ (9)
$`3\eta _{MN}\left[\left({\displaystyle \frac{^P^QH}{H}}+2{\displaystyle \frac{^PH^QH}{H^2}}\right)h_{PQ}{\displaystyle \frac{^QH}{H}}^P\overline{h}_{PQ}\right]3\left({\displaystyle \frac{\mathrm{}H}{H}}2{\displaystyle \frac{_PH^PH}{H^2}}\right)h_{MN}`$
$`+8\pi G_5H^2\left\{{\displaystyle \frac{\mathrm{\Lambda }}{8\pi G_5}}h_{MN}+\right|H|\sigma \delta (z)\left[{\displaystyle \frac{1}{2}}(\eta ^{\alpha \beta }h_{\alpha \beta }\eta ^{PQ}h_{PQ})\eta _{\mu \nu }\delta _M^\mu \delta _N^\nu +\delta _M^\mu \delta _N^\nu h_{\mu \nu }\right]\}`$
$`=`$ $`8\pi G_5\stackrel{~}{T}_{MN},`$ (10)
where the linearized five dimensional matter source $`\stackrel{~}{T}_{MN}=\delta (2\delta I_M/\sqrt{\widehat{g}}\delta \widehat{g}^{MN})`$ is included, and $`\mathrm{}=\eta ^{MN}_M_N`$.
Taking the harmonic gauge condition,
$$^M\overline{h}_{MN}=0\mathrm{or}^Mh_{MN}=\frac{1}{2}_Nh,$$
(11)
the linearized field equation becomes
$`\mathrm{}\overline{h}_{MN}+3{\displaystyle \frac{_5H}{H}}(_Mh_{5N}+_Nh_{5M}_5h_{MN})+\eta _{MN}{\displaystyle \frac{12k^2}{H^2}}h_{55}`$ (12)
$`+{\displaystyle \frac{12k}{H}}\delta (z)\left[(h_{MN}\eta _{MN}h_{55})(h_{\mu \nu }{\displaystyle \frac{1}{2}}\eta _{\mu \nu }h_{55})\delta _M^\mu \delta _N^\nu \right]=16\pi G_5\stackrel{~}{T}_{MN}.`$ (13)
In the case of $`h=0`$ and $`\stackrel{~}{T}_{MN}=0`$, this leads to the equation (5) in Ref. . In components, it becomes
$`(\mathrm{}+3{\displaystyle \frac{_5H}{H}}_5)\overline{h}_{55}+{\displaystyle \frac{12k^2}{H^2}}\overline{h}_{55}({\displaystyle \frac{_5H}{H}}_5+{\displaystyle \frac{4k^2}{H^2}})\overline{h}=16\pi G_5\stackrel{~}{T}_{55},`$ (14)
$`(\mathrm{}+{\displaystyle \frac{12k}{H}}\delta (z))\overline{h}_{5\mu }+3{\displaystyle \frac{_5H}{H}}_\mu \overline{h}_{55}{\displaystyle \frac{_5H}{H}}_\mu \overline{h}=16\pi G_5\stackrel{~}{T}_{5\mu },`$ (15)
$`(\mathrm{}3{\displaystyle \frac{_5H}{H}}_5)\overline{h}_{\mu \nu }+3{\displaystyle \frac{_5H}{H}}(_\mu \overline{h}_{5\nu }+_\nu \overline{h}_{5\mu })+\eta _{\mu \nu }({\displaystyle \frac{12k^2}{H^2}}{\displaystyle \frac{6k}{H}}\delta (z))\overline{h}_{55}`$ (16)
$`+\eta _{\mu \nu }({\displaystyle \frac{_5H}{H}}_5{\displaystyle \frac{4k^2}{H^2}}+{\displaystyle \frac{2k}{H}}\delta (z))\overline{h}=16\pi G_5\stackrel{~}{T}_{\mu \nu },`$ (17)
where $`\overline{h}=\eta ^{MN}\overline{h}_{MN}=\frac{3}{2}h`$.
For longitudinal metric fluctuations (i.e., $`h_{55}=h_{5\mu }=0`$), one has
$$h=\eta ^{\mu \nu }h_{\mu \nu }+h_{55}=h_\mu ^\mu =\overline{h}_\mu ^\mu ,\overline{h}_{55}=\frac{1}{2}h_\mu ^\mu ,\overline{h}_{5\mu }=0,\overline{h}=3\overline{h}_{55}.$$
(18)
The harmonic gauge condition in Eq. (11) gives
$$^\mu \overline{h}_{\mu \nu }=0,_5\overline{h}_{55}=0.$$
(19)
Consequently, we have $`_5\overline{h}=_5h=_5h_\mu ^\mu =_5\overline{h}_\mu ^\mu =0`$. Using these results, we find from Eq. (15) that $`\stackrel{~}{T}_{5\mu }=0`$. And other two linearized equations become
$`\mathrm{}\overline{h}_{55}=16\pi G_5\stackrel{~}{T}_{55},`$ (20)
$`(\mathrm{}3{\displaystyle \frac{_5H}{H}}_5)\overline{h}_{\mu \nu }=16\pi G_5\stackrel{~}{T}_{\mu \nu }.`$ (21)
Acting $`^\mu `$ on Eq. (21) and using the gauge condition Eq. (19) lead to the source conservation law
$$^\mu \stackrel{~}{T}_{\mu \nu }=0.$$
(22)
This is a relic of the 4D general covariance on the brane. By taking the trace of Eq. (21), we also find
$$\mathrm{}_4h_\mu ^\mu =16\pi G_5\stackrel{~}{T}_\mu ^\mu \mathrm{with}\mathrm{}_4=\eta ^{\mu \nu }_\mu _\nu .$$
(23)
This means that the trace can propagate on the brane if one includes the matter source. Note, however, this corresponds to a massless scalar propagation. Considering $`h_\mu ^\mu =2\overline{h}_{55}`$, the consistency between Eq. (20) and Eq. (23) requires the following relation
$$\stackrel{~}{T}_{55}=\frac{1}{2}\stackrel{~}{T}_\mu ^\mu .$$
(24)
This is exactly the stabilization condition implemented in Refs. . From Eqs. (19) and (23) one obtains additional constraints as
$$_5\stackrel{~}{T}_\mu ^\mu =_5\stackrel{~}{T}_{55}=0.$$
(25)
For our purpose, let us choose here the uniform source along $`z`$-axis
$$\stackrel{~}{T}_{MN}=\left(\begin{array}{cc}\frac{T_{\mu \nu }(x)}{L}& 0\\ 0& \frac{T_{55}(x)}{L}\end{array}\right),$$
(26)
which satisfies Eq. (25). Here the size $`L`$ of the extra space is still finite as is shown by
$$L=2_0^{\mathrm{}}\sqrt{g_{55}}𝑑z=2_0^{\mathrm{}}\frac{dz}{kz+1}=\frac{2}{k}\mathrm{ln}[kz+1]|_0^{\mathrm{}}\frac{1}{k}.$$
(27)
We note that $`\mathrm{ln}[k\mathrm{}+1]`$ is still finite but it is very small compared with $`1/k`$. This is because $`k`$ is allowed for up to very small quantity as the Plank scale ($`10^{34}`$ cm). Then we find
$$8\pi G_5\stackrel{~}{T}_{MN}=8\pi G\left(\begin{array}{cc}T_{\mu \nu }(x)& 0\\ 0& T_{55}(x)\end{array}\right),$$
(28)
where $`G=G_5/L`$ is the four dimensional Newton’s constant.
Using Eq. (23), Eq. (21) takes the form
$$(\mathrm{}_4m_h^2)h_{\mu \nu }=16\pi G_5(\stackrel{~}{T}_{\mu \nu }\frac{1}{2}\eta _{\mu \nu }\stackrel{~}{T}),$$
(29)
where $`\stackrel{~}{T}=\eta ^{\rho \sigma }\stackrel{~}{T}_{\rho \sigma }=\stackrel{~}{T}_\rho ^\rho `$. Here the mass squared $`m_h^2`$ is defined by the Schrödinger-like equation
$$\left[\frac{1}{2}_5^2+\frac{15k^2}{8H^2}\frac{3k}{2H}\delta (z)\right]\psi (z)=\frac{1}{2}m_h^2\psi (z)$$
(30)
with $`h_{\mu \nu }(x,z)=H^{3/2}\psi (z)\widehat{h}_{\mu \nu }(x)`$.
Now we examine the graviton propagator on the brane at $`z=0`$ by considering only $`h_{\mu \nu }(x,0)\widehat{h}_{\mu \nu }(x)`$, which satisfies
$$(\mathrm{}_4m_h^2)\widehat{h}_{\mu \nu }=16\pi G(T_{\mu \nu }\frac{1}{2}\eta _{\mu \nu }T).$$
(31)
It requires the bilinear forms of the source with the inverse propagator to isolate the physical modes. As the present analysis is on the classical level, we express $`\widehat{h}_{\mu \nu }`$ in terms of source. Taking Fourier transformation to momentum space results in
$$\widehat{h}_{\mu \nu }(p)=\frac{16\pi G}{p^2+m_h^2}\left[T_{\mu \nu }(p)\frac{1}{2}\eta _{\mu \nu }T(p)\right].$$
(32)
Then the one graviton exchange amplitude for the source $`T_{\mu \nu }`$ is given by
$$A^{\mathrm{class}}=\frac{1}{4}\widehat{h}_{\mu \nu }(p)T^{\mu \nu }(p)=\frac{4\pi G}{p^2+m_h^2}(T^{\mu \nu }T_{\mu \nu }\frac{1}{2}T^2).$$
(33)
In order to study the massive states, it is best to use the rest frame in which
$$p_10,p_2=p_3=p_4=0.$$
(34)
Considering Eqs. (22) and (34) leads to the following source relations
$$T_{11}=T_{12}=T_{13}=T_{14}=0.$$
(35)
Thus, one obtains
$$T^{\mu \nu }T_{\mu \nu }\frac{1}{2}T^2=|T_{+2}|^2+|T_2|^2+|T_{+1}|^2+|T_1|^2+T_{44}\left[\frac{1}{2}T_{44}(T_{22}+T_{33})\right],$$
(36)
where the first two terms correspond to the exchange of graviton with $`T_{\pm 2}=\frac{1}{2}(T_{22}T_{33})\pm iT_{23}`$, and the third and fourth terms are the exchange of the graviphoton with $`T_{\pm 1}=T_{24}\pm iT_{34}`$. We note here that the last term in the above equation is not positive definite. This means that there exist ghost states (negative norm states). However, if one requires
$$T_{44}=2(T_{22}+T_{33}),$$
(37)
one immediately finds that
$$T^{\mu \nu }T_{\mu \nu }\frac{1}{2}T^2=|T_{+2}|^2+|T_2|^2+|T_{+1}|^2+|T_1|^2$$
(38)
with all positive norm states.
In the limit of $`m_h^20`$, the graviphoton propagation can be decoupled from the brane . Hence we can neglect $`|T_{\pm 1}|^2`$-terms. Finally the amplitude takes the form
$$A_{m_h^20}^{\mathrm{class}}=\underset{m_h^20}{lim}\frac{4\pi G}{p_1^2+m_h^2}\left[|T_{+2}|^2+|T_2|^2\right],$$
(39)
which corresponds to the massless spin-2 amplitude. This is our key result. Although this is based on the second RS model, the results obtained above are directly applicable to the situation in the intermediate scales of the GRS model .
## III Discussion
We resolve the problem raised in the mechanism to cancel the unwanted extra polarization in the quasi-localization of gravity. This is done with introducing both the trace ($`h`$) and the uniform source ($`\stackrel{~}{T}_{MN}`$) at the linearized level. In the conventional RS approach, the trace ($`h`$) is just a gauge-dependent scalar and hence it can be gauged away. However, including the uniform matter source, this plays the role of $`\xi ^5`$ in the brane bending model . This is because $`h`$ ($`\xi ^5`$) satisfy the nearly same massless equations of $`\mathrm{}_4h=16\pi G_5\stackrel{~}{T}_\mu ^\mu `$ ($`\mathrm{}_4\xi ^5=\frac{8\pi G_5}{6}S_\mu ^\mu `$ in Ref. ). And the comparison of the equation $`\overline{h}_{\mu \nu }=h_{\mu \nu }\frac{1}{2}\eta _{\mu \nu }h`$ with $`\overline{h}_{\mu \nu }=h_{\mu \nu }^{(m)}+2k\eta _{\mu \nu }\xi ^5`$ in Ref. confirms the close relationship between $`h`$ and $`\xi ^5`$.
If $`T_\mu ^\mu =0`$, one finds from Eq. (33) that the massive spin-2 states have 5 polarizations with all positive norm states . In the case of $`h0`$, $`T_\mu ^\mu 0`$, requiring the additional condition $`T_{44}=2(T_{22}+T_{33})`$ in Eq. (37), we find the massless spin-2 state with 2 polarizations in the limit of $`m_h^20`$. In this case the ghost states disappear.
Now we wish to comment a recent paper by Kogan and Ross . They require only $`\stackrel{~}{T}_{55}=\frac{1}{2}\stackrel{~}{T}_\mu ^\mu `$ in Eq. (24). This condition can be interpreted as follows: $`\stackrel{~}{T}_\mu ^\mu 2\stackrel{~}{T}_{55}`$ is the source for the scalar radion. In the case of a mechanism that stabilizes the extra dimension, the source for the constant mode of this scalar is identically zero. In our case this comes from the consistency between Eqs. (20) and (23) with $`h_\mu ^\mu =2\overline{h}_{55}`$. In Ref. , Kogan and Ross neither choose the massive frame of Eq. (34) nor use the source conservation law in Eq. (22). These two are essential steps for obtaining the massive spin-2 amplitude of Eq. (39). In the massless frame of $`p_1=p_4`$, $`p_2=p_3=0`$ with $`p^\mu T_{\mu \nu }=0`$ , $`T^{\mu \nu }T_{\mu \nu }\frac{1}{2}T^2`$ exactly reduces to $`|T_{+2}|^2+|T_2|^2`$ without the ghost states. But in the massive frame it leads to Eq. (36). Hence one finds the ghost states. To eliminate these, we require a further condition as Eq. (37).
We are still under the harmonic gauge in Eq. (11). We note that the expression for $`A^{\mathrm{class}}`$ is simplified by refering it to an appropriate Lorentz frame. Usually one chooses the rest frame $`p_\mu =(p,0,0,0)`$ to see the massive states . On the other hand, we use the light-cone frame of $`p_\mu =(p,0,0,p)`$ for the massless states. The ghost-free condition requires a further relation among diagonal elements of $`T_{MN}`$. We are free from the ghost provided that $`T_\mu ^\mu =3(T_{22}+T_{33})=2T_{55}`$.
Finally, we comment on our source $`T_{MN}`$ in Eq. (26). In the brane-bending approach , the authors choose a localized source on the brane as $`T_{\mu \nu }(x,z)=T_{\mu \nu }(x)\delta (z)`$. Here we choose $`\stackrel{~}{T}_{\mu \nu }(x,z)=T_{\mu \nu }(x)/L`$. But, up to the integration over $`z`$, these two expressions lead to the same one. In the case of “brane-bending,” the source is located on the brane at $`z=0`$ whereas our source is uniformly distributed along the extra dimension $`z`$.
## Acknowledgments
The authors thank Hyungwon Lee for helpful discussions. This work was supported by the Brain Korea 21 Programme, Ministry of Education, Project No. D-0025.
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# Electronic structure of polychiral carbon nanotubes
## I Introduction
Multiwall nanotubes (MWNT) have not attracted as much attention from the theoreticians as the single-wall carbon tubes did, although they may be useful in many applications. Multiwall nanotubes are obviously more complex than the one-layer tubules, which explains why the former are much less documented. As far as the electronic structure is concerned, calculations have already been performed for the simplest MWNTs, those made of non-chiral layers . The systems that were investigated in these works were all made of either zig-zag or armchair nanotubes, without mixing. The reason was simply that mixing nanotubes with incommensurate periods along their axis leads to an non-periodic system, which therefore precludes the use of Bloch theorem and makes the calculations more difficult. However, there are many indications from electron diffraction and STM that the layers in a multiwall nanotube often have different helicities with nearly random distribution.
The aim of the present paper is to investigate the electronic properties of polychiral nanotubes, namely multiwall structures mixing layers with different helicities. The motivation of this work is twofold. First, to find out in which way electronic states can be induced in the band gap of a semiconducting layer by its coupling to a metallic tube of different chirality. The resulting interlayer coupling varies from site to site in a pseudo-random manner that could be responsible for localization of the electronic wave functions. No such effects were found in the present study. The second motivation was to see whether the electronic local density of states of the external layer of a MWNT can reflect the atomic structure of the underlying layers, leading to a pattern that could be observed with an STM. In some cases, MWNT analyzed with the STM appear like graphite, where only every other two atoms of the external layer are clearly imaged. In other cases, there is a spatial modulation of the image intensity, most obviously because a perfect lattice coherence cannot be realized between two cylindrical graphitic sheets when the layers have different helicities. Our calculations show that, indeed, a site asymmetry of the STM image of a MWNT similar to that of graphite may appear, but this requires a special monochiral geometry like (5,5)@(10,10). In polychiral nanotubes, by contrast, there is no site asymmetry and no Moiré pattern in the STM images computed for bilayer systems. This conclusion is in agreement with recent STM images with atomic resolution obtained on MWNTs, although Moiré patterns have frequently been observed in other experiments as recalled here above.
All these effects were explored within a tight-binding description of the $`\pi `$ electronic states, using the methodology presented in section 2. The results on the local DOS calculations for nanotubes mixing semiconducting and metallic layers are given in sections 3-4, and the STM image simulations are reported to section 5. All the calculations were restricted to bilayer nanotubes to keep the computing load reasonable.
## II Methodology
Several two-wall nanotubes were generated on the computer, with the requirement that the layer radii differ by approximately 0.34 nm which corresponds to the observed interlayer distance in MWNTs. In the $`\pi `$-electron tight-binding Hamiltonian used, the first-neighbor C pairs within a same shell received a hopping interaction $`\gamma _0`$ = -2.75 eV. This value was used for consistency with previous calculations of ours. That value slightly underestimates a recent experimental estimation of $`\gamma _0`$ (-2.9 eV), which simply means that our energies should be scaled by a factor 1.05. This scaling does not alter any of the conclusions of the present work related to the interlayer interaction and STM imaging. The interlayer interactions were written as $`W\mathrm{cos}\varphi \mathrm{exp}[(d\delta )/L]`$ with $`d`$ the distance between the coupled atoms, $`\varphi `$ the angle between the $`\pi `$ orbitals on these two atoms, $`\delta `$ = 0.334 nm, and $`L`$ = 0.045 nm. Two values of $`W`$ where used to describe the graphite AA, BB or AB-like interactions: $`W`$ = 0.36 eV for the first two and $`W`$ = 0.16 eV for the latter. The range of the interlayer interactions was limited to a maximum distance $`d`$ = 0.39 nm. This parameters reproduce well first-principles calculations for MWNTs.
Local densities of states in the multiwall nanotubes were computed by the recursion method. This technique does not rely on the Bloch theorem. It gives rise to a continued-fraction development of the Green’s function diagonal elements in the complex energy plane, $`G_{ii}(z)=i|(zH)^1|i`$. For each atomic site $`i`$ of interest, $`n`$ = 500 levels of continued fraction were computed. When the continued fraction is truncated after $`n`$ levels, the resulting density of states is composed of $`n`$ Dirac delta peaks. With $`n1`$ levels, another set of $`n1`$ peaks is obtained. These two, interlaced sets were mixed with equal weights, and each delta peak was represented by a Gaussian function with standard deviation 0.023 eV (the band width, $`6|\gamma _0|`$, divided by $`\sqrt{2}n`$). Due to this broadening, all the singularities of the densities of states, including the band edges, are slightly smoothed out. Due to its smallness, this broadening should not alter a main conclusion of the work, namely the absence of localized states in the gap of the semiconducting layer.
The change of density of states brought about by the interlayer coupling is expected to be small. The first-order perturbation expression of the Green function is indeed $`G=G_0+G_0WG_0`$. The unperturbed Green function $`G_0`$ is made of blocks corresponding to the individual layers that $`W`$ couples together. Due to that structure, all the diagonal elements of $`G_0WG_0`$ are zero, which means that the density of states is not perturbed at first order in the interlayer interaction.
STM image simulations of multiwall nanotubes were performed for comparison with experiment. These calculations are based on a simple tight-binding theory of the STM current
$$I=(2\pi )^2\frac{e}{h}_{E_F^seV}^{E_F^s}𝑑En_t(E_F^tE_F^s+eV+E)\underset{i,i^{}s}{}v_{ti}v_{ti^{}}^{}n_{ii^{}}^s(E)$$
where the $`E_F`$’s are the Fermi levels of the tip (t) and sample (s) and $`V`$ is the tip–sample bias potential. The tip is treated as a single atom with an s orbital and a Gaussian density of states $`n_t(E)`$. $`v_{ti}`$ is the tight-binding hopping interaction between the tip atom and $`\pi `$ orbital located on site $`i`$ of the nanotube sample, and $`n_{ii^{}}^s(E)=(1/\pi )\text{ Im }G_{ii^{}}^s(E)`$. The Green function elements $`G_{ii^{}}`$ of the nanotube were computed by recursion with 200 continued-fraction levels that give converged results for the present imaging studies. Similarly to the density of states calculations, a small imaginary part was added to the energy to force the convergence.
## III Metal-semiconductor nanotubes
We first consider two-wall nanotubes having a metal at the inner shell and a semiconductor at the outer shell. Table I gives a few such metal-semiconductor nanotubes. The first system mixes an armchair and a zig-zag nanotube. Although these two nanotubes are non-chiral, their chiral angles differ by 30, their translation periods differ by a factor of $`\sqrt{3}`$, and the combined system may be described as polychiral. The next three nanotubes of table I are real polychiral systems. By contrast, the fifth nanotube is monochiral since its layers have the same helicity. Hence, the two layers in the (15,-6)@(15,10) nanotube have the same Bravais period but have opposite chiral angles. This last nanotube differs from the fourth one by the fact the inner layers (15,-6) and (9,6) are enantiomers.
Local density of states in the external layer were computed in a slice of 0.2 nm height, which contained between 32 and 40 atoms, depending on the nanotube. Although the coupling to the inner layer varies from site to site, all the atoms of the external layer were found to have pretty much the same density of states.
For all the polychiral nanotubes investigated, the density of states of the semiconducting layer was found to be weakly affected by its coupling to the inner metallic layer, at least in an interval between -1 and +1 eV (the zero of energy is always considered to be at the Fermi level). In that interval of energy, the metallic tube presents a constant density of states – hereafter called the metallic plateau – with no van Hove singularities. Fig. 1(a), which concerns with (15,-6)@(15,10), is a typical example of this effect. As compared to the single-wall (15,10) nanotube (dashed curve), there is a minute downshift of the bottom of the conduction band of the semiconducting layer (full curve), whereas the top of the valence band does not move. The shapes of the gap edges are not modified. If the hopping interactions between the layers introduce a tailing of the valence and conduction states inside the gap, the decay will take place in an energy range shorter than the peak broadening used (0.023 eV). The metallic layer induces a few states in the band gap of the semiconductor (the total number of states in the band gap is 0.42$`\times 10^4`$ per atom of the semiconducting layer, see table I). The density of states in the gap region is approximately uniform and very small.
The interlayer coupling is much more efficient in the monochiral (9,6)@(15,10) nanotube, as shown in fig. 1(b). The density of states in the band gap region is approximately five times larger than in fig. 1(a) and reaches 2.5% of that of the metallic layer. It must be stressed out that both nanotubes in fig. 1 exhibit exactly the same electronic structure when the interlayer coupling is switched off. In other words, all the differences between the full curves in figs. 1(a) and (b) come out from the different environments the (15,10) layer feel in both systems. The monochiral tubes constitute a particular case in which the electronic properties are much affected by specific symmetries in MWNTs, both as pseudogaps in the local density of states or as a change in the intensity of every two atom of the STM image (see below). A change of the width of the metallic gap was also found in the case of commensurate three-wall armchair tubes.
## IV Semiconductor-metal nanotubes
We now consider two-wall nanotubes having the semiconducting layer at the interior and the metallic layer outside. A list of such nanotubes which mix different helicities is given in table II. Local density of states in the inner layer were computed in a slice of 0.3 nm height, which contained between 24 and 32 atoms depending on the nanotube. Here again, the fluctuations of densities of states in the semiconducting layer, although two times as large as in Sect. III, remained small (see table II). The local densities of states were then averaged as before.
The ratio between the band gap of the semiconductor and the width of the metallic plateau of the metal is now approximately 2:3, instead of 1:6 as for the previous configurations (Sect. III). For instance, the metallic plateau of the (10,10) nanotube is bounded by two Van Hove singularities at $`E=\pm `$0.9 eV. These singularities can be seen in the density of states of the (6,4) layer in the (6,4)@(10,10) bilayer (fig. 2(a)). In the conduction band for instance, the singularity leads to a resonance followed by an anti-resonance. This kind of structure was frequently observed among the nanotubes of table II. In a systematic way, also, the band gap of the semiconductor is reduced by the coupling to the outer layer: the top of the valence band has moved upwards by approximately 0.03 eV.
As shown in fig. 2(a), the density of states in the band gap of the inner semiconducting tube is small, but still significantly higher than in fig. 1(a) for instance. The difference between these two situations is that the semiconducting layer is now at the interior rather than at the exterior, and the inner layer is more perturbed than the outer one. Indeed the average numbers of intersheet bonds per atom in layers 1 and 2 are inversely proportional to the number of atoms in these layers. Since there are approximately two times less atoms in the inner shell than in the outer (for those systems we are investigating), the interlayer coupling is two times more efficient on layer 1 than on layer 2. Since, in addition, the (6,4) semiconductor has a larger band gap than those of the semiconducting layers of table I (due to its smaller diameter), the number of states in the gap has increased. As revealed by table II, the number of states in the gap looks remarkably constant for all the nanotubes examined, around 2.3$`\times 10^4`$ per atom.
As for the metallic layer, one can hardly see any change in its density of states around the Fermi level (fig. 2(b)). The site-to-site fluctuations of densities of states near the Fermi level are less than 1% of the (10,10) density of states at $`E_F`$. This kind of weak disorder, here due to the coupling with a chiral nanotube, is small but is perhaps sufficient to affect the transport properties of a MWNT in the weak localization regime. Clear effects of the coupling to the inner layer appear below -0.5 eV and above +0.5 eV, where the (6,4) nanotube has its band edges. It is clear from this example that the amplitudes of the Van Hove singularities at $`E=\pm `$0.9 eV are reduced as compared with the single-wall nanotube. This is due to the breaking of the translational symmetry brought about by the coupling between layers of different chiralities.
## V Simulation of the STM image of two-wall nanotubes
Several constant-current STM images of multiwall nanotubes show intensity or contrast modulations. Such modulations have been interpreted as being a Moiré pattern formed by the atomic structure of the last two layers. In graphite, Moiré pattern effects have clearly been identified with an STM in regions where the last layer was folded back on the surface with a misorientation of its crystallographic directions.
As pointed out in the previous sections, the local DOS show only little variations on going from one site to another in a multiwall nanotube. To quantify that property, the fluctuations (rms) of the local DOS were computed on a chain of 25 first-neighbor atoms located as close as possible to a generator of the external layer. These atoms were selected because they would be probed in a scan of the topmost part of the nanotube by a STM tip. The DOS fluctuations, averaged over the energy interval (-1,+1) eV, are listed in the last row of tables I and II. They are small, less than 1% of the mean density of states. According to these data, the spatial variations of the density of states in polychiral nanotubes cannot explain the modulations of the STM image intensity.
Fig. 3 is a simulation of the STM image of (6,6)@(19,0) computed with the methodology described in section II. The tip is at a potential of 0.5 V with respect to the sample. For that polarity, the most prominent features in the STM image of the semiconducting (19,0) layer are the CC bonds not parallel to the axis, which appear as bright stripes at 60 to the axis. All along the portion of nanotube displayed, the periodicity of the image is that ($`\sqrt{3}a`$, with $`a`$ the lattice parameter of graphene) of the external zig-zag layer. There is no visible sign of the underlying armchair nanotube with its shorter period $`a`$. The topographic line cut shown at the bottom of fig. 3 clearly proves that statement. The sharp minima of the curve correspond to the centers of the hexagons. The apparent variations of their depth are due to the pixel discretization. The maxima correspond to the atoms, the secondary minima are at the center of the CC bonds parallel to the axis.
Nothing similar to a Moiré pattern appears in the computed STM image of fig. 3 nor in the simulations we carried out for other polychiral nanotubes. However, these patterns occur occasionally in the experimental images of multiwall nanotubes, as reported above. It is not impossible that a mechanical deformation of outer layer of the tube caused by the STM tip induces metallic islands (with a much larger density of states) at special places where the layers are in suitable registry. The pressure of the tip may also induce better electric contact with the substrate at some places, leading to a larger tunneling current.
The atoms in fig. 3 all look the same, unlike the case of multilayer graphite where the STM current at low bias ($``$ 0.1 V) shows a strong site asymmetry. This asymmetry is already present with two layers only, since the coupling makes the atoms having a neighbor underneath different from those that have not (A and B atoms, respectively).
In a multiwall nanotube, it is impossible to realize the same stacking as in natural graphite all around the cross section. However, at least one multiwall nanotube exists where the topographic STM image is predicted to look like that of graphite. This case is (5,5)@(10,10). This system is known to exhibit small pseudo-gaps near $`E_F`$ as the consequence of avoided band crossings for relative tube orientations such that the mirror planes of (5,5) do not coincide with those of (10,10). Local DOS calculations then show that the atoms (of the external layer) are not equivalent, at least in a small interval around the Fermi level. Interestingly, first-neighbor atoms have a peak or a deep at $`E_F`$, alternatively, very much like in graphite where the A and B atoms alternate. This effect is shown in fig. 4. The explanation of this bi-partition of the honeycomb lattice is presented in the Appendix. A consequence of it is that the STM image of the (5,5)@(10,10) nanotube at low bias resembles that of graphite, with maxima of protrusion on every other two atoms (those with the largest DOS at $`E_F`$), see fig. 5. However, for other relative orientations of the tubes, no bipartition effect is observed (in agreement with first-principle calculations ). What is special about (5,5)@(10,10) is that this system has at least a five-fold common symmetry. In (6,6)@(11,11) for instance, there is no axial symmetry, which destroys the effects of the interlayer coupling on the DOS around the Fermi level, very much like in polychiral nanotubes. The resulting intertube interaction averaging reduces any symmetry related feature such as the opening of pseudo-gaps and the bipartition of the honeycomb lattice. The STM image of (6,6)@(11,11) is then similar to that of the isolated single-wall (11,11) nanotube (see fig. 5). With other metallic nanotubes such as (7,4)@(12,9) which we also have examined, some variations of the local DOS from one atom to the other were detected in the metallic plateau, but these were two small to lead to a clear site asymmetry in the STM image.
The property that (5,5)@(10,10) presents a site asymmetry that depends on the relative orientation of the layers can be illustrated by giving a uniform torsion to the (10,10) nanotube. At regular intervals along the axis, the planes that bisect the CC bonds perpendicular to the axis of (10,10) coincide with the mirror planes of (5,5) ($`C_{5v}`$ symmetry, no site asymmetry). Away from these positions, the local symmetry of the atomic structure is lower and the two-site asymmetry of the DOS should come out and reach a maximum in between.
The twist is equivalent to applying a shear of the honeycomb network, which affects the bond lengths and opens a small gap at the Fermi level. The calculations were performed for a twist of 1.5 (shear strain) corresponding to a torsion angle of 2.2/nm, which leads to a band gap of 0.25 eV (the $`\gamma _0`$ parameter was scaled according to a $`d^2`$ law, with $`d`$ the bond length). In the (5,5)@(10,10) distorted nanotube, the local DOS of the twisted (10,10) layer has a peak near the Fermi energy induced by the interactions with the inner (5,5) nanotube (fig. 6(b)). The shear also affects the Brillouin zone of the rolled-up graphene sheet, which moves the Fermi points of the twisted (10,10) nanotube away from those of the (5,5) layer. As a consequence, the minigaps of the bilayer are no longer located at the center of the metallic plateau (where the semiconducting gap of the twisted nanotube has opened) but are shifted 0.25 eV on both sides of the Fermi level. As can be seen in fig. 6(a), the DOS features in the mini-gaps at $`\pm 0.25`$ eV resemble the two-site asymmetry observed near $`E_F`$ in fig. 4 for the perfect (5,5)@(10,10), except that the magnitude of the asymmetry now varies along the tube. The DOS curves in fig. 6(a) correspond to 25 successive atoms along a longitudinal zig-zag chain on the external tube. The curves at the center look all the same, where the local symmetry of the (5,5)@(10,10) distorted nanotube is close to $`C_{5v}`$. By contrast, the curves at the bottom and at the top have peaks and deeps at $`E=\pm `$0.25 eV that alternate from one site to the next. The local symmetry has been changed to $`C_5`$ in these regions. The modulation of the degree of unequivalence between the atoms is due to the continuous change of the local symmetry along the nanotube axis. Unfortunately, this modulation did not appear clearly in the STM images that we computed for bias potentials of $`0.3`$ and $`+0.3`$ V. The reason is that the expression of the STM current (sect. II) is an integral of the Green’s function elements over the bias window, to which the site-dependent features in the density of states in fig. 6(a) contribute little. However, these features might be observed by current imaging tunneling spectroscopy.
## VI Conclusion
Two-wall nanotubes mixing metallic and semiconducting layers retain the basic properties of the uncoupled constituents, as shown previously for monochiral nanotubes. The intertube interactions induce a small continuous distribution of states in the band gap of the semiconducting layer. The electronic states near the Fermi level come from the metallic tube and they will dominate the transport properties in the weak localization regime. Constant-current STM images computed for polychiral nanotubes are pretty much the same as the ones obtained on the isolated external layer. It is only in the case of monochiral and commensurate structures like (5,5)@(10,10) that interlayer effects can be seen in the STM topography. The interlayer coupling gives rise to a site asymmetry in the STM image at low voltage ($``$0.1 V), similar to that obtained on multilayer graphite. This site bipartition is maximum when the symmetry of the two-wall nanotube is reduced to $`C_5`$, and it disappears when the symmetry is higher. The site asymmetry also disappears in relationally-incommensurate nanotubes like (6,6)@(11,11). In polychiral nanotubes, there is no site asymmetry and no Moiré patterns appear in the computed STM images. From the few cases we have investigated, it can be concluded that the superstructures often observed in the STM images of MWNTs cannot be ascribed to pure electronic effects.
###### Acknowledgements.
This work has been partly funded by the interuniversity research project on reduced dimensionality systems (PAI P4/10) of the Belgian Office for Scientific, Cultural and Technical affairs, the EU NAMITECH contract: ERBFMRX-CT96-0067 (DG12-MITH) and by a JCyL (Grant: VA28/99). The authors acknowledge Laszlo P. Biró for helpful discussions.
## Appendix: Multi-wall armchair nanotubes
In a single-wall armchair nanotube, the two bands that cross each-other at the Fermi level correspond to the irreducible representation $`A_1`$ and $`A_2`$ of the symmetry group $`C_{5v}`$ of the wave function for a general Bloch wave vector $`k`$. These wave functions are respectively symmetric and antisymmetric upon a reflection on the “vertical” mirror planes that bisect the CC bonds perpendicular to the nanotube axis. When all curvature effects are neglected in the Hamiltonian as here, the Fermi wave vector is independent on the tube diameter. This means that in a two-wall nanotube such as (5,5)@(10,10), the states at the Fermi points have a fourfold degeneracy when the interlayer coupling is ignored. In the presence of the coupling, the wave functions adapted to the perturbation $`W`$ are linear combinations of the four Fermi states $`\psi _1^i`$, $`\psi _2^i`$, $`\psi _1^e`$, and $`\psi _2^e`$ of the internal and external layers (upper indices $`i`$ and $`e`$) corresponding to $`A_1`$ and $`A_2`$ symmetries (lower indices 1 and 2). These combinations of states diagonalize the perturbation matrix
$$\left(\begin{array}{cccc}0\hfill & 0\hfill & w_{11}^{ie}\hfill & w_{12}^{ie}\hfill \\ 0\hfill & 0\hfill & w_{21}^{ie}\hfill & w_{22}^{ie}\hfill \\ w_{11}^{ei}\hfill & w_{12}^{ei}\hfill & 0\hfill & 0\hfill \\ w_{21}^{ei}\hfill & w_{22}^{ei}\hfill & 0\hfill & 0\hfill \end{array}\right)$$
where $`w_{12}^{ie}=\psi _1^i|W|\psi _2^e`$, etc. In any case, the perturbation is sufficient to split off the degeneracy of the Fermi states. In general, also, the eigenvector of the perturbation matrix will mix the four Fermi states. This means in particular that the perturbed wave functions at the Fermi level mix the $`\psi _1^e`$ and $`\psi _2^e`$ states of the external layer. By mixing these states, which are respectively even and odd with respect to the center of each CC bond, one forms a disymmetric combination. As a result, there are now two kinds of unequivalent atoms in each layer, as shown by the density of states in fig. 4, which explains the site asymmetry of the computed STM images.
An exception to this explanation arises when the four matrix elements of the kind $`w_{12}^{ie}`$ along the ascending diagonal vanish for symmetry reason. This takes place when the symmetry of the two-wall nanotube preserves the mirror planes of the inner layer. Then, $`A_1`$ and $`A_2`$ remain valid irreducible representations of the symmetry group of the coupled system, band crossings remain allowed, and there is no pseudo-gap formation in the density of states. This also means that the site asymmetry of the STM image disappears. For the (5,5)@(10,10) nanotube, this happens with $`C_{5v}`$, $`D_{5h}`$ and $`D_{5d}`$ configurations of the layers. The latter configuration corresponds to the minimum of the total energy of the nanotube.
In the case where the two nanotubes have no axial symmetry in common, such as for instance with (6,6)@(11,11), all band crossings are avoided and two kinds of unequivalent atoms are formed in each layer, as above. However, due to a cancellation effect, all the elements in the perturbation matrix are found to be small except the ones derived from the totally-symmetric states, $`w_{11}^{ie}`$ and $`w_{11}^{ei}`$. In practice, then, all the atoms look equivalent. Also, the pseudo-gaps near the Fermi level are much weaker than for the (5,5)@(10,10) case.
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# Private Quantum Channels and the Cost of Randomizing Quantum Information
## 1 Introduction
Secure transmission of classical information is a well studied topic. Suppose Alice wants to send an $`n`$-bit message $`M`$ to Bob over an insecure (i.e. spied-on) channel, in such a way that the eavesdropper Eve cannot obtain any information about $`M`$ from tapping the channel. If Alice and Bob share some secret $`n`$-bit key $`K`$, then here is a simple way for them to achieve their goal: Alice exclusive-ors $`M`$ with $`K`$ and sends the result $`M^{}=MK`$ over the channel, Bob then xors $`M^{}`$ again with $`K`$ and obtains the original message $`M^{}K=M`$. Eve may see the encoded message $`M^{}`$, but if she does not know $`K`$ then this will give her no information about the real message $`M`$, since for any $`M`$ there is a key $`K^{}`$ giving rise to the same encoding $`M^{}`$. This scheme is known as the Vernam cipher or one-time pad (“one-time” because $`K`$ can be used only once if we want information-theoretic security). It shows that $`n`$ bits of shared secret key are sufficient to securely transmit $`n`$ bits of information. Shannon \[Sha48, Sha49\] has shown that this scheme is optimal: $`n`$ bits of shared key are also necessary in order to transmit an $`n`$-bit message in an information-theoretically secure way.
Now let us consider the analogous situation in the quantum world. Alice and Bob are connected by a one-way quantum channel, to which an eavesdropper Eve has complete access. Alice wants to transmit to Bob some $`n`$-qubit state $`\rho `$ taken from some set $`𝒮`$, without allowing Eve to obtain any information about $`\rho `$. Alice and Bob could easily achieve such security if they share $`n`$ EPR-pairs (or if they were able to establish EPR-pairs over a secure quantum channel), for then they can apply teleportation \[BBC<sup>+</sup>93\] and transmit every qubit via 2 random classical bits, which will give Eve no information whatsoever. But now suppose Alice and Bob do not share EPR-pairs, but instead they only have the resource of shared randomness, which is weaker but easier to maintain.
A first question is: is it at all possible to send quantum information fully securely using only a finite amount of randomness? At first sight this may seem hard: Alice and Bob have to “hide” the amplitudes of a quantum state, which are infinitely precise complex numbers. Nevertheless, the question has a positive answer. More precisely, to privately send $`n`$ qubits, a $`2n`$-bit classical key is sufficient. The encryption technique is fairly natural. Alice applies to the state $`\rho `$ she wants to transmit a reversible quantum operation specified by the shared key $`K`$ (basically, she applies a random Pauli matrix to each qubit), and she sends the result $`\rho ^{}`$ to Bob. In the most general setting this reversible operation can be represented as doing a unitary operation on the state $`\rho `$ augmented with a known fixed ancilla state $`\rho _a`$. Knowing the key $`K`$ that Alice used, Bob knows which operation Alice applied and he can reverse this, remove the ancilla, and retrieve $`\rho `$. In order for this scheme to be information-theoretically secure against the eavesdropper, we have to require that Eve always “sees” the same density matrix $`\rho _0`$ on the channel, no matter what $`\rho `$ was. Because Eve does not know $`K`$, this condition can indeed be satisfied. Accordingly, an insecure quantum channel can be made secure (private) by means of shared classical randomness.
A second question is, then, how much key Alice and Bob need to share in order to be able to privately transmit any $`n`$-qubit state. A good way to measure key size is by the amount of entropy required to create it. As one might imagine, showing that $`2n`$ bits of key are also necessary is the most challenging part of the article. We prove this in Section 5.<sup>1</sup><sup>1</sup>1Note that if Alice and Bob share an insecure two-way channel, then they can do quantum key exchange \[BB84\] in order to establish a shared random key, so in this case no prior shared key (or only a very small one) is required. Accordingly, in analogy with the classical one-time pad, we have an optimal quantum one-time pad which uses $`2n`$ classical bits to completely “hide” $`n`$ qubits from Eve. In particular, hiding a qubit is only twice as hard as hiding a classical bit, despite the fact that in the qubit we now have to hide amplitudes coming from a continuous set.
Now imagine an alternative scenario. Alice has a state $`\rho `$ from some specific set and she wants to randomize it completely. How much entropy does she need for this? That is, what is the thermodynamical cost of forgetting quantum information? A natural and general way to do that is for Alice to perform a unitary transformation to $`\rho `$ augmented with an ancilla and then to forget which one. The thermodynamical price of this operation is now the entropy of the probability distribution over the set of unitary transformations. The parallel between these two scenarios should be clear. If one has a private quantum channel, one automatically has a related randomization procedure. Consequently, we obtain the result that $`2n`$ bits are necessary and sufficient to randomize an $`n`$-qubit quantum register. For the case $`n=1`$, this result has also been obtained by Braunstein, Lo, and Spiller \[BLS99, Lo99\].
The article is organized as follows. Section 2 introduces some notation and some properties of Von Neumann entropy. In Section 3 we give a formal definition of a private quantum channel (PQC). In Section 4 we give some examples of PQCs. In particular we show that there is a PQC that privately sends any $`n`$-qubit state using $`2n`$ bits of randomness (shared key). We also exhibit a non-trivial set of $`n`$-qubit states (namely the tensor products of qubits with real amplitudes) for which there is PQC requiring only $`n`$ bits of randomness. The latter result includes the classical one-time pad. In Section 5 we show that $`2n`$ bits of randomness are necessary if we want to be able to send any $`n`$-qubit state privately. Finally, in Section 6 we restate the previous results in terms of the thermodynamical cost of randomization of quantum information.
Remark about related work. A number of recent papers independently discussed issues similar to our work. We already mentioned the result of Braunstein, Lo, and Spiller \[BLS99, Lo99\] for state randomization. Very recently, Boykin and Roychowdhury \[BR00\] exhibited the $`2n`$-bit Pauli-matrix one-time pad and proved a $`2n`$-bit lower bound for the case where the encryption scheme does not allow the use of an ancilla state (they also give a general characterization of all possible encryption schemes without ancilla). In Section 5 we give a simpler proof of this lower bound for the no-ancilla case and give a different and more complicated proof for the lower bound in the case where we do allow an ancilla.
## 2 Preliminaries
### 2.1 States and operators
We use $`v`$ for the Euclidean norm of vector $`v`$. If $`A`$ is a matrix, then we use $`A^{}`$ for its conjugate transpose and $`\mathrm{Tr}(A)`$ for its trace (the sum of its diagonal entries). A square matrix $`A`$ is Hermitian if $`A=A^{}`$, and unitary if $`A^1=A^{}`$. Important examples of unitary transformations are the 4 Pauli matrices:
$$\sigma _0=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),\sigma _1=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),\sigma _2=\left(\begin{array}{cc}0& i\\ i& 0\end{array}\right),\sigma _3=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right).$$
Let $`|0,\mathrm{},|M1`$ denote the basis states of some $`M`$-dimensional Hilbert space $`_M`$. We use $`_{2^n}`$ for the Hilbert space whose basis states are the $`2^n`$ classical $`n`$-bit strings. A pure quantum state $`|\varphi `$ is a norm-1 vector in $`_M`$. We treat $`|\varphi `$ as an $`M`$-dimensional column vector and use $`\varphi |`$ for the row vector that is its conjugate transpose. The inner product between pure states $`|\varphi `$ and $`|\psi `$ is $`\varphi |\psi `$. A mixed quantum state or density matrix $`\rho `$ is a non-negative Hermitian matrix that has trace $`\mathrm{Tr}(\rho )=1`$. The density matrix corresponding to a pure state $`|\varphi `$ is $`|\varphi \varphi |`$. Because a density matrix $`\rho `$ is Hermitian, it has a diagonalization $`\rho =_{i=1}^Np_i|\varphi _i\varphi _i|`$, where the $`p_i`$ are its eigenvalues, $`p_i0`$, $`_ip_i=1`$, and the $`|\varphi _i`$ form an orthonormal set. Thus $`\rho `$ can be viewed as describing a probability distribution over pure states. We use $`\stackrel{~}{I}_M=\frac{1}{M}I_M=\frac{1}{M}_{i=1}^M|ii|`$ to denote the totally mixed state, which represents the uniform distribution on all basis states. If two systems are in pure states $`|\varphi `$ and $`|\psi `$, respectively, then their joint state is the tensor product pure state $`|\varphi |\psi =|\varphi |\psi `$. If two systems are in mixed states $`\rho _1`$ and $`\rho _2`$, respectively, then their joint state is the tensor product $`\rho _1\rho _2`$. Note that $`(|\varphi |\psi )(\varphi |\psi |)`$ is the same as $`|\varphi \varphi ||\psi \psi |`$.
Applying a unitary transformation $`U`$ to a pure state $`|\varphi `$ gives pure state $`U|\varphi `$, applying $`U`$ to a mixed state $`\rho `$ gives mixed state $`U\rho U^{}`$. We will use $`=\{\sqrt{p_i}U_i1iN\}`$ to denote the superoperator which applies $`U_i`$ with probability $`p_i`$ to its argument (we assume $`_ip_i=1`$). Thus $`(\rho )=_ip_iU_i\rho U_i^{}`$. Quantum mechanics allows for more general superoperators, but this type suffices for our purposes. If two superoperators $`=\{\sqrt{p_i}U_i1iN\}`$ and $`^{}=\{\sqrt{p_i^{}}U_i^{}1iN^{}\}`$ are identical ($`(\rho )=^{}(\rho )`$ for all $`\rho `$), then they are unitarily related in the following way \[Nie98, Section 3.2\] (where we assume $`NN^{}`$ and if $`N>N^{}`$ we pad $`^{}`$ with zero operators to make $``$ and $`^{}`$ of equal size): there exists a unitary $`N\times N`$ matrix $`A`$ such that for all $`i`$
$$\sqrt{p_i}U_i=\underset{j=1}{\overset{N}{}}A_{ij}\sqrt{p_j^{}}U_j^{}.$$
### 2.2 Von Neumann entropy
Let density matrix $`\rho `$ have the diagonalization $`_{i=1}^Np_i|\varphi _i\varphi _i|`$. The Von Neumann entropy of $`\rho `$ is $`S(\rho )=H(p_1,\mathrm{},p_N)=_{i=1}^Np_i\mathrm{log}p_i`$, where $`H`$ is the classical entropy function. This $`S(\rho )`$ can be interpreted as the minimal Shannon entropy of the measurement outcome, minimized over all possible complete measurements. Note that $`S(\rho )`$ only depends on the eigenvalues of $`\rho `$. The following properties of Von Neumann entropy will be useful later (for proofs see for instance \[Weh78\]).
1. $`S(|\varphi \varphi |)=0`$, for every pure state $`|\varphi `$.
2. $`S(\rho _1\rho _2)=S(\rho _1)+S(\rho _2)`$.
3. $`S(U\rho U^{})=S(\rho )`$.
4. $`S(\lambda _1\rho _1+\lambda _2\rho _2+\mathrm{}+\lambda _n\rho _n)\lambda _1S(\rho _1)+\lambda _2S(\rho _2)+\mathrm{}+\lambda _nS(\rho _n)`$ if $`\lambda _i0`$ and $`_i\lambda _i=1`$.
5. If $`\rho =_{i=1}^Np_i|\varphi _i\varphi _i|`$ with the $`|\varphi _i`$ not necessarily orthogonal, then $`S(\rho )H(p_1,\mathrm{},p_N)`$.
## 3 Private Quantum Channel
Let us sketch the scenario for a private quantum channel. There are $`N`$ possible keys, which we identify for convenience with the numbers $`1,\mathrm{},N`$. The $`i`$th key has probability $`p_i`$, so the key has entropy $`H(p_1,\mathrm{},p_N)`$ when viewed as a random variable. Each key $`i`$ corresponds to a unitary transformation $`U_i`$. Suppose Alice wants to send a pure state $`|\varphi `$ from some set $`𝒮`$ to Bob. She appends some fixed ancilla qubits in state $`\rho _a`$ to $`|\varphi \varphi |`$ and then applies $`U_i`$ to $`|\varphi \varphi |\rho _a`$, where $`i`$ is her key. She sends the resulting state to Bob. Bob, who shares the key $`i`$ with Alice, applies $`U_i^1`$ to obtain $`|\varphi \varphi |\rho _a`$, removes the ancilla $`\rho _a`$, and is left with Alice’s message $`|\varphi \varphi |`$. Now in order for this to be secure against an eavesdropper Eve, we have to require that if Eve does not know $`i`$, then the density matrix $`\rho _0`$ that she gets from monitoring the channel is independent of $`|\varphi `$. This implies that she gets no information at all about $`|\varphi `$. Of course, Eve’s measuring the channel might destroy the encoded message, but this is like classically jamming the channel and cannot be avoided. The point is that if Eve measures, then she receives no information about $`|\varphi `$. It is not hard to see that this is the most general quantum mechanical scenario which allows Bob to recover the message perfectly and at the same time gives Eve zero information.
We formalize this scenario as follows.
###### Definition 3.1
Let $`𝒮_{2^n}`$ be a set of pure $`n`$-qubit states, $`=\{\sqrt{p_i}U_i1iN\}`$ be a superoperator where each $`U_i`$ is a unitary mapping on $`_{2^m}`$, $`_{i=1}^Np_i=1`$, $`\rho _a`$ be an $`(mn)`$-qubit density matrix, and $`\rho _0`$ be an $`m`$-qubit density matrix. Then $`[𝒮,,\rho _a,\rho _0]`$ is called a Private Quantum Channel (PQC) if and only if for all $`|\varphi 𝒮`$ we have
$$(|\varphi \varphi |\rho _a)=\underset{i=1}{\overset{N}{}}p_iU_i\left(|\varphi \varphi |\rho _a\right)U_i^{}=\rho _0.$$
If $`n=m`$ (i.e. no ancilla), then we omit $`\rho _a`$.
Note that by linearity, if the PQC works for all pure states in $`𝒮`$, then it also works for density matrices over $`𝒮`$: applying the PQC to a mixture of states from $`𝒮`$ gives the same $`\rho _0`$ as when we apply it to a pure state. Accordingly, if $`[𝒮,\{\sqrt{p_i}U_i1iN\},\rho _a,\rho _0]`$ is a PQC, then $`H(p_1,\mathrm{},p_N)`$ bits of shared randomness are sufficient for Alice to send any mixture $`\rho `$ of $`𝒮`$-states to Bob in a secure way. Alice encodes $`\rho `$ in a reversible way depending on her key $`i`$ and Bob can decode because he knows the same $`i`$ and hence can reverse Alice’s operation $`U_i`$. On the other hand, Eve has no information about the key $`i`$ apart from the distribution $`p_i`$, so from her point of view the channel is in state $`\rho _{Eve}=\rho _0`$. This is independent of the $`\rho `$ that Alice wants to send, and hence gives Eve no information about $`\rho `$.
## 4 Examples of Private Quantum Channels
In this section we exhibit some private quantum channels. The first uses $`2n`$ bits key to send privately any $`n`$-qubit state. The idea is simply to apply a random Pauli matrix to each bit individually. This takes 2 random bits per qubit and it is well known that the resulting qubit is in the completely mixed state. For notational convenience we identity the numbers $`\{0,\mathrm{},2^{2n}1\}`$ with the set $`\{0,1,2,3\}^n`$. For $`x\{0,1,2,3\}^n`$ we use $`x_i\{0,1,2,3\}`$ for its $`i`$th entry, and we use $`\overline{\sigma _x}`$ to denote the $`n`$-qubit unitary transformation $`\sigma _{x_1}\mathrm{}\sigma _{x_n}`$.
###### Theorem 4.1
If $`=\{\frac{1}{\sqrt{2^{2n}}}\overline{\sigma _x}x\{0,1,2,3\}^n\}`$, then $`[_{2^n},,\stackrel{~}{I}_{2^n}]`$ is a PQC.
Proof It is easily verified that applying each $`\sigma _i`$ with probability $`1/4`$ to a qubit puts that qubit in the totally mixed state $`\stackrel{~}{I}_2`$ (no matter if it is entangled with other qubits). Operator $``$ just applies this treatment to each of the $`n`$ qubits, hence $`(|\varphi \varphi |)=\stackrel{~}{I}_{2^n}`$ for every $`|\varphi _{2^n}`$. $`\mathrm{}`$
Since the above $``$ contains $`2^{2n}`$ operations and they have uniform probability, it follows that $`2n`$ bits of private key suffice to privately send any state from $`_{2^n}`$.
The next theorem shows that there is some nontrivial subspace of $`_{2^n}`$ where $`n`$ bits of private key suffice, namely the set of all tensor products of real-amplitude qubits.
###### Theorem 4.2
If $`B=\{\mathrm{cos}(\theta )|0+\mathrm{sin}(\theta )|10\theta <2\pi \}`$, $`𝒮=B^n`$, and $`=\{\frac{1}{\sqrt{2^n}}\overline{\sigma _x}x\{0,2\}^n\}`$, then $`[𝒮,,\stackrel{~}{I}_{2^n}]`$ is a PQC.
Proof This is easily verified: applying $`\sigma _0`$ and $`\sigma _2`$, each with probability 1/2, puts any qubit from $`B`$ in the totally mixed state. Operator $``$ does this to each of the $`n`$ qubits individually. $`\mathrm{}`$
Note that if we restrict $`B`$ to classical bits (i.e. $`\theta \{0,\pi /2\}`$) then the above PQC reduces to the classical one-time pad: flipping each bit with probability 1/2 gives information-theoretical security.
In the previous PQCs, $`\rho _0`$ was the completely mixed state $`\stackrel{~}{I}_{2^n}`$. This is no accident, and holds whenever $`n=m`$ and $`\stackrel{~}{I}_{2^n}`$ is one of the states that the PQC can send:
###### Theorem 4.3
If $`[𝒮,,\rho _0]`$ is a PQC without ancilla and $`\stackrel{~}{I}_{2^n}`$ can be written as a mixture of $`S`$-states, then $`\rho _0=\stackrel{~}{I}_{2^n}`$.
Proof If $`\stackrel{~}{I}_{2^n}`$ can be written as a mixture of $`S`$-states, then
$`\rho _0=(\stackrel{~}{I}_{2^n})={\displaystyle \underset{i=1}{\overset{N}{}}}p_iU_i\stackrel{~}{I}_{2^n}U_i^{}={\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{p_i}{2^n}}U_iU_i^{}={\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{p_i}{2^n}}I_{2^n}=\stackrel{~}{I}_{2^n}.`$ $`\mathrm{}`$
In general $`\rho _0`$ need not be $`\stackrel{~}{I}_{2^n}`$. For instance, let $`𝒮=\{|0,\frac{1}{\sqrt{2}}(|0+|1)\}`$, $`=\{\sqrt{p_1}I_2,\frac{\sqrt{p_2}}{\sqrt{2}}\left(\begin{array}{cc}1& 1\\ 1& 1\end{array}\right)\}`$ with $`p_1=p_2=1/2`$, and $`\rho _0=\left(\begin{array}{cc}\frac{3}{4}& \frac{1}{4}\\ \frac{1}{4}& \frac{1}{4}\end{array}\right)`$. Then it is easily verified that $`[𝒮,,\rho _0]`$ is a PQC.
## 5 Lower Bound on the Entropy of PQCs
In the previous section we showed that $`2n`$ bits of entropy suffice for a PQC that can send arbitrary $`n`$-qubit states. In this section we will show that $`2n`$ bits are also necessary for this. Very recently and independently of our work, this $`2n`$-bit lower bound was also proven by Boykin and Roychowdhury \[BR00\] for the special case where the PQC is not allowed to use any ancilla qubits. We will first give a shorter version of their proof, basically by observing that a large part of it can be replaced by a reference to the unitary equivalence of identical superoperators stated at the end of Section 2.1.
###### Theorem 5.1
If $`[_{2^n},\{\sqrt{p_i}U_i1iN\},\stackrel{~}{I}_{2^n}]`$ is a PQC, then $`H(p_1,\mathrm{},p_N)2n`$.
Proof Let $`=\{\sqrt{p_i}U_i\}`$, $`^{}=\{\frac{1}{\sqrt{2^{2n}}}\overline{\sigma _x}x\{0,1,2,3\}^n\}`$ be the superoperator of Theorem 4.1, and let $`K=\mathrm{max}(2^{2n},N)`$. Since $`(\rho )=^{}(\rho )=\stackrel{~}{I}_{2^n}`$ for all $`n`$-qubit states $`\rho `$, we have that $``$ and $`^{}`$ are unitarily related in the way mentioned in Section 2.1: there exists a unitary $`K\times K`$ matrix $`A`$ such that for all $`1iN`$ we have
$$\sqrt{p_i}U_i=\underset{x\{0,1,2,3\}^n}{}A_{ix}\frac{1}{\sqrt{2^{2n}}}\overline{\sigma _x}.$$
We can view the set of all $`2^n\times 2^n`$ matrices as a $`2^{2n}`$-dimensional vector space, with inner product $`M,M^{}=\mathrm{Tr}(M^{}M^{})/2^n`$ and induced matrix norm $`M=\sqrt{M,M}`$ (as done in \[BR00\]). Note that $`M=1`$ if $`M`$ is unitary. The set of all $`\overline{\sigma _x}`$ forms an orthonormal basis for this vector space, so we get:
$$p_i=\sqrt{p_i}U_i^2=\underset{x}{}A_{ix}\frac{1}{\sqrt{2^{2n}}}\overline{\sigma _x}^2=\frac{1}{2^{2n}}\underset{x}{}|A_{ix}|^2\frac{1}{2^{2n}}.$$
Hence $`N2^{2n}`$ and $`H(p_1,\mathrm{},p_N)2n`$. $`\mathrm{}`$
However, even granted this result it is still conceivable that a PQC might require less randomness if it can “spread out” its encoding over many ancilla qubits — it is even conceivable that those ancilla qubits can be used to establish privately shared randomness using some variant of quantum key distribution. The general case with ancilla is not addressed in \[BR00\], and proving that the $`2n`$-bit lower bound extends to this case requires more work. The next few theorems will do this. These show that a PQC that can transmit any $`n`$-qubit state requires $`2n`$ bits of randomness, no matter how many ancilla qubits it uses. Thus Theorem 4.1 exhibits an optimal quantum one-time pad, analogous to the optimal classical one-time pad mentioned in the introduction.
We will use the notation $`𝒞_k=\{|i0ik1\}`$ for the set of the first $`k`$ classical states. The next theorem states that a PQC that privately conveys $`n`$ qubits using $`m`$ bits of key, can be transformed into a PQC that privately conveys any state from $`𝒞_{2^{2n}}`$, still using only $`m`$ bits of key.
###### Theorem 5.2
If there exists a PQC $`[_{2^n},=\{\sqrt{p_i}U_i1iN\},\rho _a,\rho _0]`$, then there exists a PQC $`[𝒞_{2^{2n}},^{}=\{\sqrt{p_i}U_i^{}1iN\},\rho _a,\stackrel{~}{I}_{2^n}\rho _0]`$.
Proof For ease of notation we assume without loss of generality that $``$ uses no ancilla, so we assume $`\rho _0`$ is an $`n`$-qubit state and omit $`\rho _a`$ (this does not affect the proof in any way). We first show that $`(|xy|)=0`$ whenever $`x,y𝒞_{2^{2n}}`$ and $`xy`$ ($`(|xy|)`$ is well-defined but somewhat of an abuse of notation, since the matrix $`|xy|`$ is not a density matrix). This is implied by the following 3 equalities:
$`\rho _0=\left({\displaystyle \frac{1}{2}}(|xx|+|yy|)\right)={\displaystyle \frac{1}{2}}\left((|xx|)+(|yy|)\right).\rho _0=\left(({\displaystyle \frac{1}{\sqrt{2}}}(|x+|y))({\displaystyle \frac{1}{\sqrt{2}}}(x|+y|))\right)={\displaystyle \frac{1}{2}}\left((|xx|)+(|yy|)+(|xy|)+(|yx|)\right).\rho _0=\left(({\displaystyle \frac{1}{\sqrt{2}}}(|x+i|y))({\displaystyle \frac{1}{\sqrt{2}}}(x|iy|))\right)={\displaystyle \frac{1}{2}}\left((|xx|)+(|yy|)i(|xy|)+i(|yx|)\right).`$ The first and second equality imply $`(|xy|)+(|yx|)=0`$, the first and third equality imply $`(|xy|)(|yx|)=0`$. Hence $`(|xy|)=(|yx|)=0`$.
We now define $`^{}`$ and show that it is a PQC. Intuitively, $`^{}`$ will map every state from $`𝒞_{2^{2n}}`$ to a tensor product of $`n`$ Bell states by mapping pairs of bits to one of the four Bell states.<sup>2</sup><sup>2</sup>2The 4 Bells states are $`\frac{1}{\sqrt{2}}(|00\pm |11)`$ and $`\frac{1}{\sqrt{2}}(|01\pm |10)`$. The second bits of the pairs are then moved to the second half of the state and randomized by applying $``$ to them. Because of the entanglement between the two halves of each Bell state, the resulting $`2n`$-qubit density matrix will be $`\stackrel{~}{I}_{2^n}\rho _0`$. More specifically, define
$$U|x=\left(\overline{\sigma _x}I_{2^n}\right)\frac{1}{\sqrt{2^n}}\underset{i=0}{\overset{2^n1}{}}|i|i,$$
with $`\overline{\sigma _x}=\sigma _{x_1}\mathrm{}\sigma _{x_n}`$ as in Theorem 4.1. Also define $`U_i^{}=(I_{2^n}U_i)U`$. It remains to show that $`^{}(|xx|)=\stackrel{~}{I}_{2^n}\rho _0`$ for all $`|x𝒞_{2^{2n}}`$:
$`^{}(|xx|)`$
$`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}p_i(I_{2^n}U_i)\left[(\overline{\sigma _x}I_{2^n})\left({\displaystyle \frac{1}{\sqrt{2^n}}}{\displaystyle \underset{y=0}{\overset{2^n1}{}}}|y|y\right)\left({\displaystyle \frac{1}{\sqrt{2^n}}}{\displaystyle \underset{z=0}{\overset{2^n1}{}}}z|z|\right)(\overline{\sigma _x}I_{2^n})^{}\right](I_{2^n}U_i)^{}`$
$`=`$ $`(\overline{\sigma _x}I_{2^n})\left[{\displaystyle \frac{1}{2^n}}{\displaystyle \underset{i=1}{\overset{N}{}}}p_i(I_{2^n}U_i)\left({\displaystyle \underset{y,z\{0,2^n1\}}{}}|yz||yz|\right)(I_{2^n}U_i)^{}\right](\overline{\sigma _x}I_{2^n})^{}`$
$`=`$ $`(\overline{\sigma _x}I_{2^n})\left[{\displaystyle \frac{1}{2^n}}{\displaystyle \underset{y,z\{0,2^n1\}}{}}|yz|\left({\displaystyle \underset{i=1}{\overset{N}{}}}p_iU_i|yz|U_i^{}\right)\right](\overline{\sigma _x}I_{2^n})^{}`$
$`=`$ $`(\overline{\sigma _x}I_{2^n})\left[{\displaystyle \frac{1}{2^n}}{\displaystyle \underset{y,z\{0,2^n1\}}{}}|yz|(|yz|)\right](\overline{\sigma _x}I_{2^n})^{}`$
$`\stackrel{()}{=}`$ $`(\overline{\sigma _x}I_{2^n})\left[{\displaystyle \frac{1}{2^n}}{\displaystyle \underset{y=0}{\overset{2^n1}{}}}|yy|(|yy|)\right](\overline{\sigma _x}I_{2^n})^{}`$
$`=`$ $`(\overline{\sigma _x}I_{2^n})\left[\stackrel{~}{I}_{2^n}\rho _0\right](\overline{\sigma _x}I_{2^n})^{}`$
$`=`$ $`\stackrel{~}{I}_{2^n}\rho _0.`$
In the step marked by $`()`$ we used that $`(|yz|)=0`$ if $`yz`$. $`\mathrm{}`$
Before proving a lower bound on the entropy required for sending arbitrary $`n`$-qubit states, we first prove a lower bound on the entropy required for sending states from $`𝒞_{2^m}`$. Privately sending any state from $`𝒞_{2^m}`$ corresponds to privately sending any classical $`m`$-bit string. If communication takes place through classical channels, then Shannon’s theorem implies that $`m`$ bits of shared key are required to achieve such security. Shannon’s classical lower bound does not translate automatically to the quantum world (it is in fact violated if a two-way quantum channel is available, see Footnote 1). Nevertheless, if Alice and Bob communicate via a one-way quantum channel, then Shannon’s theorem does generalize to the quantum world:
###### Theorem 5.3
If $`[𝒞_{2^m},\{\sqrt{p_i}U_i1iN\},\rho _a,\rho _0]`$ is a PQC, then $`H(p_1,\mathrm{},p_N)m`$.
Proof Diagonalize the ancilla as $`\rho _a=_{j=1}^rq_j|\psi _j\psi _j|`$, so $`S(\rho _a)=H(q_1,\mathrm{},q_r)`$. First note that the properties of Von Neumann entropy (Section 2) imply:
$`S(\rho _0)`$ $`=`$ $`S\left({\displaystyle \underset{i=1}{\overset{N}{}}}p_iU_i(|00|\rho _a)U_i^{}\right)=S\left({\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \underset{j=1}{\overset{r}{}}}p_iq_jU_i(|00||\psi _j\psi _j|)U_i^{}\right)`$
$``$ $`H(p_1q_1,p_1q_2,\mathrm{},p_Nq_{r1},p_Nq_r)=H(p_1,\mathrm{},p_N)+H(q_1,\mathrm{},q_r).`$
Secondly, note that
$$S(\rho _0)=S\left(\underset{i=1}{\overset{N}{}}p_iU_i(\stackrel{~}{I}_{2^m}\rho _a)U_i^{}\right)\underset{i=1}{\overset{N}{}}p_iS\left(\stackrel{~}{I}_{2^m}\rho _a\right)=\underset{i=1}{\overset{N}{}}p_i(m+S(\rho _a))=m+S(\rho _a).$$
Combining these two inequalities gives the theorem. $`\mathrm{}`$
In particular, for sending arbitrary states from $`𝒞_{2^{2n}}`$ we need entropy at least $`2n`$. Combining Theorems 5.2 and 5.3 we thus obtain:
###### Corollary 5.4
If $`[_{2^n},\{\sqrt{p_i}U_i1iN\},\rho _a,\rho _0]`$ is a PQC, then $`H(p_1,\mathrm{},p_N)2n`$ (and hence in particular $`N2^{2n}`$).
In relation to Theorem 4.2, note that $`𝒞_{2^n}B^n`$. Hence another corollary of Theorem 5.3 is the optimality of the PQC of Theorem 4.2:
###### Corollary 5.5
If $`[B^n,\{\sqrt{p_i}U_i1iN\},\rho _a,\rho _0]`$ is a PQC, then $`H(p_1,\mathrm{},p_N)n`$ (and hence in particular $`N2^n`$).
## 6 Randomization of Quantum States
The above concepts and results were motivated by cryptographic goals, namely to enable private transmission of quantum information using a shared classical key. However, our results can also be stated in terms of the problem of “forgetting” or “randomizing” quantum information, as discussed recently by Braunstein, Lo, and Spiller \[BLS99\].
The randomization of a quantum source $`𝒮`$ is a procedure that maps any state $`\rho `$ coming from $`𝒮`$ to some fixed constant state $`\rho _0`$ (for instance the completely mixed state). The process thus “forgets” what was specific to $`\rho `$. To help in this process, we allow the randomizing process to make use of a piece of the environment which is in some fixed state $`\rho _a`$ (ancilla qubits). We also give it access to some source of classical randomness. Because every quantum operation can be viewed as a unitary transformation on a larger space, we can assume without loss of generality that the randomization process has the following form: it uses the source of randomness to pick some $`i`$ with probability $`p_i`$, then it applies some unitary transformation $`U_i`$ to $`\rho `$ and the ancillary environment, and then it forgets $`i`$. The resulting mixed state should be $`\rho _0`$. At this point it should be clear to the reader that if $`[𝒮,\{\sqrt{p_i}U_i1iN\},\rho _a,\rho _0]`$ is a PQC, then it also constitutes a randomization procedure, and vice versa.
We are interested in the amount of entropy that such a randomization procedure needs to generate. This is the entropy of forgetting the random classical input $`i`$. It quantifies the thermodynamic cost of the process. Braunstein, Lo, and Spiller \[BLS99\] have shown that 2 bits of entropy are necessary and sufficient for the randomization of 1 qubit. By translating our PQC-results to the randomization-context, we can generalize their result to:
###### Corollary 6.1
The generation of $`2n`$ bits of entropy is sufficient and necessary in order to randomize arbitrary $`n`$-qubit states.
Proof Sufficiency follows from Theorem 4.1 and necessity from Corollary 5.4. $`\mathrm{}`$
For the more limited set of states $`𝒮=B^n`$ we have:
###### Corollary 6.2
The generation of $`n`$ bits of entropy is sufficient and necessary in order to randomize arbitrary tensor products of $`n`$ real-amplitude qubits.
Proof Sufficiency follows from Theorem 4.2 and necessity from Corollary 5.5. $`\mathrm{}`$
## 7 Summary
The main result of this paper is an optimal quantum version of the classical one-time pad. On the one hand, if Alice and Bob share $`2n`$ bits of key, Alice can send Bob any $`n`$-qubit state $`\rho `$, encoded in another $`n`$-qubit state in a way which conveys no information about $`\rho `$ to the eavesdropper. This is a simple scheme which works locally (i.e. deals with each qubit separately) and uses no ancillary qubits. On the other hand, we showed that even if Alice and Bob are allowed to use any number of ancilla qubits, then they still require $`2n`$ bits of entropy. In the context of state randomization, it follows that $`2n`$ bits of entropy are necessary and sufficient for randomization of $`n`$-qubit states.
### Acknowledgment
We thank Richard Cleve, Hoi-Kwong Lo, Michael Nielsen, Harry Buhrman, and P. Oscar Boykin for useful discussions and comments.
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# Multibaryons in the collective coordinate approach to the SU(3) Skyrme model
## I Introduction
In the Skyrme model and its generalizations, baryons arise as topological excitations of a non-linear chiral Lagrangian written in terms of meson fields. These type of models have been quite successful in describing the properties of single baryons such as the nucleon and the strange hyperons (see e.g., Refs.). This has lead people to investigate the lowest energy Skyrmion configurations with topological number greater than one, which are of inherent interest as examples of three dimensional solitonic structures and may also be relevant for nuclear physics. These studies were already started by Skyrme in his pioneer papers at the beginning of the sixties. However, it was only in 1987 that the minimum energy $`B=2`$ Skyrmion was correctly identified. Some time later the authors of Ref. found the solutions with $`B=3,4`$ and $`5`$ by numerical relaxation calculations. Finally, a few years ago, after some demanding numerical work, the global minimum energy configurations with topological number up to $`B=9`$ were constructed. One particularly interesting aspect of all these multi-Skyrmion fields is that they are very symmetric. While for $`B=2`$ the symmetry group corresponds to that of a torus, for $`B=3,4,7,9`$ they possess the symmetries of the platonic polihedra $`T_d,O_h,I_h`$ and $`T_d`$, respectively, and for $`B=5,6,8`$ the dihedral symmetries $`D_{2d},D_{4d},D_{6d}`$, respectively. It should be stressed that, in spite of this, the multi-Skyrmion fields are very complicated functions of the space coordinates which are only known numerically. Fortunately, rather simple and accurate approximations to these configurations have been found. They are based on some Ansätze which are written in terms of rational maps and take advantage of the similarities between multi-Skyrmion fields and Bogomol’nyi-Prasad-Sommerfield (BPS) monopoles. These developments triggered several investigations concerning the properties of the multi-Skyrmions (such as e.g., vibrational excitations) as well as their application to baryonic systems containing strangeness and heavier flavors. The extension to flavored multibaryons is also motivated by the advent of heavy ion colliders with the possibility of producing strange and even charmed multibaryonic states with rather low baryon number in the laboratory. To describe the strange multibaryons one has to extend the model to SU(3) flavor space. The classical background configurations are simply obtained by embedding the SU(2) static multi-Skyrmions in the isospin subgroup of SU(3). In order to obtain the spectrum with states of well defined spin and isospin quantum numbers, as well as their splittings, we have to perform the quantization of this system. However, the presence of the rather important symmetry breaking terms associated with the mass of the strange quark makes the quantization process not completely trivial. In fact, two alternative methods have been suggested in the literature. One is known as the bound state approach (BSA) in which strange baryons are described as SU(2) rotating soliton-kaon bound systems. The other scheme assumes that the strange degrees of freedom can still be treated as rotational modes but the corresponding collective Hamiltonian is to be diagonalized exactly. This method is usually called the rigid rotator approach (RRA). In two recent articles SU(3) multi-Skyrmions have been investigated following the BSA. In this work we complement such investigations by considering these configurations within the framework of the RRA.
This paper is organized as follows. In Sec. II we provide a brief description of the model and obtain the collective Hamiltonian for the different baryon numbers. In Sec. III we focus on the determination of the multibaryon quantum numbers and wavefunctions. In Sec. IV we present the numerical results and in Sec. V our conclusions. Finally, in the Appendix we give the explicit form of the collective Hamiltonians for $`3B9`$.
## II The Model
We start with the effective action of the SU(3) Skyrme model supplemented with an appropriate symmetry breaking term . Expressed in terms of the SU(3)-valued chiral field $`U(x)`$ it reads
$`\mathrm{\Gamma }={\displaystyle d^4x\left\{\frac{f_\pi ^2}{4}\mathrm{Tr}\left[_\mu U^\mu U^{}\right]+\frac{1}{32e^2}\mathrm{Tr}\left[[U^{}_\mu U,U^{}_\nu U]^2\right]\right\}}+\mathrm{\Gamma }_{WZ}+\mathrm{\Gamma }_{SB},`$ (1)
where $`f_\pi `$ is the pion decay constant ( $`=93\mathrm{MeV}`$ empirically) and $`e`$ is the so-called Skyrme parameter. In Eq. (1), the symmetry breaking term $`\mathrm{\Gamma }_{SB}`$ accounts for the different masses and decay constants of the pion and kaon fields while $`\mathrm{\Gamma }_{WZ}`$ is the usual Wess-Zumino action. Their explicit forms are
$`\mathrm{\Gamma }_{SB}`$ $`=`$ $`{\displaystyle }d^4x\{{\displaystyle \frac{f_\pi ^2m_\pi ^2+2f_K^2m_K^2}{12}}\mathrm{Tr}[U+U^{}2]+{\displaystyle \frac{f_\pi ^2m_\pi ^2f_K^2m_K^2}{6}}\mathrm{Tr}\left[\sqrt{3}\lambda ^8(U+U^{})\right]`$ (3)
$`+{\displaystyle \frac{f_K^2f_\pi ^2}{12}}\mathrm{Tr}\left[(1\sqrt{3}\lambda ^8)(U_\mu U^{}^\mu U+U^{}_\mu U^\mu U^{})\right]\},`$
$`\mathrm{\Gamma }_{WZ}`$ $`=`$ $`i{\displaystyle \frac{N_c}{240\pi ^2}}{\displaystyle d^5x\epsilon ^{\mu \nu \alpha \beta \gamma }\mathrm{Tr}\left[U^{}(_\mu U)U^{}(_\nu U)U^{}(_\alpha U)U^{}(_\beta U)U^{}(_\gamma U)\right]},`$ (4)
where $`\lambda ^8`$ is the eighth Gell-Mann matrix, $`N_c`$ the number of colors, $`m_\pi `$ and $`m_K`$ are the pion and kaon masses, respectively, and $`f_K`$ is the kaon decay constant.
We proceed by introducing the following Ansatz for the time dependent chiral field
$$U(\stackrel{\mathbf{}}{𝒓},t)=𝒜(t)\left(\begin{array}{cc}\mathrm{exp}\left[i\stackrel{\mathbf{}}{𝝉}\stackrel{\mathbf{}}{𝝅}(^1(t)\stackrel{\mathbf{}}{𝒓})\right]& 0\\ 0& 1\end{array}\right)𝒜^{}(t),$$
(5)
where the embedded SU(2) background configuration is rigidly rotated both in SU(3) flavor space and real space, the collective coordinates given by $`𝒜(t)`$ SU(3) and $`(t)`$ SO(3), respectively. Substituting $`U(\stackrel{\mathbf{}}{𝒓},t)`$ given by Eq. (5) into the effective action yields a Lagrangian of the general form
$$L=M_{sol}+L_{coll},$$
(6)
where $`M_{sol}`$ is the static SU(2) soliton mass and $`L_{coll}`$ is the collective Lagrangian, whose general expression will be given below. Following the usual steps in the RRA, we first find the soliton background configuration by minimizing $`M_{sol}`$. For this purpose we introduce the rational map Ansätze for the pion field
$$\stackrel{\mathbf{}}{𝝅}(\stackrel{\mathbf{}}{𝒓})=F(r)\widehat{𝒏}.$$
(7)
Here, $`F(r)`$ is the multi-Skyrmion profile which depends on the radial coordinate only and $`\widehat{𝒏}`$ is a unit vector given by
$$\widehat{𝒏}=\frac{1}{1+|R|^2}\left[2\mathrm{}(R)\widehat{\mathit{ı}}+2\mathrm{}(R)\widehat{\mathit{ȷ}}+(1|R|^2)\widehat{𝒌}\right],$$
(8)
with $`R=R(z)`$ the rational map corresponding to a certain winding number $`B`$ which is identified with the baryon number. The complex variable $`z`$ is related to the usual two spherical coordinates $`(\theta ,\varphi )`$ via stereographic projection, namely, $`z=\mathrm{tan}(\theta /2)\mathrm{exp}(i\varphi )`$. For example, the map corresponding to the $`B=1`$ hedgehog Ansatz is the identity map $`R=z`$. The explicit form of the rational maps corresponding to the other baryon numbers $`B9`$ and the resulting expression for the soliton mass $`M_{sol}`$ can be found in Refs.. The radial profile function $`F(r)`$ is determined by minimizing the classical soliton energy $`M_{sol}`$. Details of this procedure as well as plots of these profiles for different baryon numbers are given in Ref..
The collective Lagrangian written in terms of the collective degrees of freedom and the corresponding angular velocities $`\mathrm{\Omega }_a,\omega _\alpha `$ defined by <sup>*</sup><sup>*</sup>*Here and in the following the spin/isospin indices $`a,b,c`$ run over $`\{1,\mathrm{},3\}`$, the flavor index $`\alpha `$ over $`\{1,\mathrm{},8\}`$ and the $`k\{4,\mathrm{},7\}`$ index corresponds to excitations into strangeness directions.
$`\left(^1\dot{}\right)_{ab}`$ $`=`$ $`ϵ_{abc}\mathrm{\Omega }_c,`$ (9)
$`𝒜^1\dot{𝒜}`$ $`=`$ $`{\displaystyle \frac{i}{2}}\lambda _\alpha \omega _\alpha ,`$ (10)
takes the general form
$`L_{coll}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{a,b}{}}\left(\mathrm{\Theta }_{ab}^J\mathrm{\Omega }_a\mathrm{\Omega }_b+\mathrm{\Theta }_{ab}^N\omega _a\omega _b+2\mathrm{\Theta }_{ab}^M\mathrm{\Omega }_a\omega _b\right)+{\displaystyle \frac{1}{2}}\mathrm{\Theta }^S{\displaystyle \underset{k}{}}\omega _k^2{\displaystyle \frac{N_cB}{2\sqrt{3}}}\omega _8`$ (12)
$`{\displaystyle \frac{1}{2}}G_{SB}\left(1D_{88}\right),`$
with $`D_{88}=\frac{1}{2}\mathrm{Tr}\left[\lambda _8𝒜\lambda _8𝒜^{}\right]`$. The moment of inertia in the strangeness direction $`\mathrm{\Theta }^S`$ is
$$\mathrm{\Theta }^S=d^3r\frac{1c}{2}\left[f_K^2+\frac{1}{4e^2}\left(F_{}^{}{}_{}{}^{2}+2B\frac{s^2}{r^2}\right)\right],$$
(13)
where we have introduced the short hand notation $`s=\mathrm{sin}F`$, $`c=\mathrm{cos}F`$. The spin $`\mathrm{\Theta }_{ab}^J`$, isospin $`\mathrm{\Theta }_{ab}^N`$ and mixed moments of inertia $`\mathrm{\Theta }_{ab}^M`$ are
$`\mathrm{\Theta }_{ab}^J`$ $`=`$ $`{\displaystyle d^3rs^2r^2\left[\left(f_\pi ^2+\frac{F^2}{e^2}\right)+\frac{1}{2}\frac{s^2}{e^2}_c\widehat{𝒏}_c\widehat{𝒏}\right]_a\widehat{𝒏}_b\widehat{𝒏}},`$ (14)
$`\mathrm{\Theta }_{ab}^N`$ $`=`$ $`{\displaystyle d^3rs^2\left[\left(f_\pi ^2+\frac{F^2}{e^2}\right)(\delta _{ab}\widehat{n}_a\widehat{n}_b)+\frac{s^2}{e^2}(\delta _{ab}2\widehat{n}_a\widehat{n}_b)_c\widehat{𝒏}_c\widehat{𝒏}\right]},`$ (15)
$`\mathrm{\Theta }_{ab}^M`$ $`=`$ $`{\displaystyle d^3rs^2r\left[\left(f_\pi ^2+\frac{F^2}{e^2}\right)+\frac{1}{2}\frac{s^2}{e^2}_c\widehat{𝒏}_c\widehat{𝒏}\right]_a\widehat{n}_b}.`$ (16)
For the rational maps we are interested in all these moments of inertia are diagonal . This is a direct consequence of the symmetries of these Ansätze. Finally, the symmetry breaking parameter $`G_{SB}`$ is
$$G_{SB}=\frac{2}{3}(f_K^2f_\pi ^2)d^3r\left(F_{}^{}{}_{}{}^{2}+2B\frac{s^2}{r^2}\right)c+\frac{4}{3}(f_K^2m_K^2f_\pi ^2m_\pi ^2)d^3r(1c).$$
(17)
Given $`L_{coll}`$, the spin and flavor canonical momentum operators $`\widehat{J}_a`$ and $`\widehat{F}_\alpha `$ are defined in the usual way
$`\widehat{J}_a={\displaystyle \frac{L_{coll}}{\mathrm{\Omega }_a}};\widehat{F}_\alpha ={\displaystyle \frac{L_{coll}}{\omega _\alpha }}.`$ (18)
The collective Hamiltonian is conventionally obtained as the Legendre transformation $`H_{coll}=J_a\mathrm{\Omega }_a+F_\alpha \omega _\alpha L_{coll}`$, resulting in
$$H_{coll}=K^S\left[C_2(SU(3))\frac{3}{4}B^2\widehat{N}^2+\gamma (1D_{88})\right]+H_B^{JN}.$$
(19)
Here, $`C_2(SU(3))=_\alpha \widehat{F}_\alpha ^2`$ stands for the quadratic SU(3) Casimir operator, $`\widehat{N}_a\widehat{F}_a`$ is the isospin operator in the soliton frame, $`\gamma =\mathrm{\Theta }^SG_{SB}`$ is the dimensionless flavor symmetry breaking parameter and $`K^S=1/(2\mathrm{\Theta }^S)`$. In order to obtain Eq. (19) we have used $`N_c=3`$ and the constraint $`F_8=\frac{\sqrt{3}}{2}B`$. Finally, the detailed form of the spin-isospin collective Hamiltonians $`H_B^{JN}`$ depends on the soliton symmetry group. The method to derive them is very similar to the one described in Sec. III of Ref.. For $`B=1,2`$ the corresponding groups are the continuous groups $`O(3)`$ and $`D_\mathrm{}h`$, respectively. In those cases there are some relations between the spin and isospin operators which lead to the well known expressions for the spin-isospin collective Hamiltonians
$`H_{B=1}^{JN}`$ $`=`$ $`{\displaystyle \frac{1}{2\mathrm{\Theta }^J}}\widehat{J}^2,`$ (20)
$`H_{B=2}^{JN}`$ $`=`$ $`{\displaystyle \frac{1}{2\mathrm{\Theta }_1^J}}\left(\widehat{J}^2\widehat{J}_3^2\right)+{\displaystyle \frac{1}{2\mathrm{\Theta }_1^N}}\left(\widehat{N}^2\widehat{N}_3^2\right)+{\displaystyle \frac{1}{2\mathrm{\Theta }_3^N}}\widehat{N}_3^2.`$ (21)
On the other hand, for $`B3`$ the symmetry groups are finite. Thus, the general form of the spin-isospin collective Hamiltonian is
$`H_{B3}^{JN}`$ $`=`$ $`{\displaystyle \underset{a}{}}\left(K_a^J\widehat{J}_a^2+K_a^N\widehat{N}_a^22K_a^M\widehat{J}_a\widehat{N}_a\right),`$ (22)
where
$$K_a^J=\frac{1}{2}\frac{\mathrm{\Theta }_a^N}{\mathrm{\Delta }_a},K_a^N=\frac{1}{2}\frac{\mathrm{\Theta }_a^J}{\mathrm{\Delta }_a},K_a^M=\frac{1}{2}\frac{\mathrm{\Theta }_a^M}{\mathrm{\Delta }_a}.$$
(23)
and $`\mathrm{\Delta }_a\mathrm{\Theta }_a^J\mathrm{\Theta }_a^N(\mathrm{\Theta }_a^M)^2`$. The explicit form for each baryon number can be found in the Appendix.
## III Quantum numbers and collective wave functions
In order to calculate the rotational corrections to the multi-Skyrmion masses we have to find the corresponding wave functions. Following the Yabu-Ando procedure , they should diagonalize the flavor symmetry breaking term in the collective Hamiltonian. At the same time they should satisfy the constraints imposed by the symmetries of the classical soliton configuration. Thus, the general form of such eigenfunctions will be
$$|BJJ_z,YII_z,N=\underset{J_3N_3}{}\alpha _{J_3N_3}^{JN}D_{J_zJ_3}^J\mathrm{\Psi }_{(Y,I,I_z),(B,N,N_3)},$$
(24)
Here, $`D_{J_zJ_3}^J`$ is the usual SU(2) Wigner function and $`\mathrm{\Psi }_{(Y,I,I_z),(B,N,N_3)}`$ is a function depending on the 8 Euler angles that parametrize the SU(3) manifold. To obtain $`\mathrm{\Psi }_{(Y,I,I_z),(B,N,N_3)}`$ we should solve the eigenvalue equation
$$K^S\left[h+\gamma (1D_{88})\right]\mathrm{\Psi }=ϵ\mathrm{\Psi },$$
(25)
where $`h=C_2(SU(3))\frac{3}{4}B^2N(N+1)`$. The coefficients $`\alpha _{J_3N_3}^{JN}`$ are determined in such a way that the full wavefunction transforms as some particular one-dimensional irreducible representation (irrep) of the soliton symmetry group $`G`$. This will be discussed in some detail below.
To solve Eq. (25) we expand $`\mathrm{\Psi }_{(Y,I,I_z),(B,N,N_3)}`$ in a basis of SU(3) Wigner functions $`D_{(Y,I,I_z),(B,N,N_3)}^{(p,q)}`$, where $`(p,q)`$ are the labels used to identify the SU(3) irrep. Namely,
$$\mathrm{\Psi }_{(Y,I,I_z),(B,N,N_3)}=\underset{(p,q)}{}\beta _{(p,q)}\sqrt{d_{(p,q)}}D_{(Y,I,I_z),(B,N,N_3)}^{(p,q)},$$
(26)
where $`d_{(p,q)}=(p+1)(q+1)(p+q+2)/2`$ is the dimension of the irrep. In such basis $`h`$ is diagonal and the matrix elements of the symmetry breaking term can be expressed as a product of two SU(3) Clebsch-Gordan coefficients. To determine, for a given value of $`B`$, the allowed values of the $`Y`$, $`I`$ and $`N`$ quantum numbers as well as which SU(3) irrep should be included in the basis we proceed as follows. As already seen, the value of the right hypercharge $`Y_R2F_8/\sqrt{3}`$ is fixed by the constraint $`Y_R=B`$. Thus, any SU(3) irrep that appears in the expansion, Eq.(26), should have a maximum value of hypercharge equal or larger than $`B`$. Thus, the possible values of $`(p,q)`$ should satisfy
$$\frac{p+2q}{3}=B+m,$$
(27)
with $`p`$ and $`q`$ non-negative integer numbers and $`m=0,1,2,\mathrm{}`$. The irreps corresponding to $`m=0`$ are the so-called minimal irreps that we will denote $`(p_0,q_0)`$. It is possible to show that the matrix element of $`h`$ in any state that belongs to a minimal irrep is $`<h>_0=3B/2`$. To determine the relevant values of $`Y`$, $`I`$ and $`N`$ it is enough to consider such irreps. Although for non-vanishing strangeness $`S`$ other values of the quantum numbers could be allowed, they will be of no interest to us. In fact, it is not difficult to show that for the first state with “non-minimal quantum numbers” the matrix element of $`h`$ is more than twice $`<h>_0`$. Therefore, such state is expected to appear as a highly excited state in the spectrum. Since the minimal irreps have maximum right hypercharge $`Y_R=B`$ it is clear that corresponding possible values of the body-fixed isospin $`N`$ are $`N=p_0/2`$. On the other hand, those of the hypercharge $`Y`$ are
$$\frac{2p_0+q_0}{3}Y\frac{p_0+2q_0}{3}.$$
(28)
Finally, given a value of $`Y`$ that satisfies this relation, the allowed values of the isospin $`I`$ are
$$\left|\frac{Y}{2}+\frac{p_0q_0}{3}\right|I\frac{p_0+q_0}{2}\frac{1}{2}\left|Y\frac{p_0q_0}{3}\right|.$$
(29)
In Table I we list, for each baryon number $`3B9`$, the minimal SU(3) irrep which lead to states with $`N<3`$ together with the allowed values of isospin for some values of strangeness. Given a set of possible $`(B,I,Y,N)`$ quantum numbers one should find all the SU(3) irreps with $`m>0`$ that enter in the expansion, Eq. (26). This is done by selecting from all the irreps which satisfy Eq. (27) those that contain a state with this same set of quantum numbers. This leads to different towers of SU(3) irreps for each set of quantum numbers. Once this is done, it is a simple task to transform Eq. (25) into an ordinary linear eigenvalue problem whose solution provides the energy eigenvalues $`ϵ`$ and the coefficients $`\beta _{(p,q)}`$. Of course, to do that one should work with a basis of finite size. Since we are interested only in the few lowest eigenvalues the minimum size is fixed by the condition that those eigenvalues remain unchanged under a further increase of such size.
Having determined $`\mathrm{\Psi }_{(Y,I,I_z),(B,N,N_3)}`$ and the corresponding possible quantum numbers we have still to obtain the coefficients $`\alpha _{J_3N_3}^{JN}`$ of Eq. (24) and the allowed values of $`J`$. For this purpose, only the spin $`J`$ and isospin $`N`$ are relevant. Thus, the situation is very similar to that of the $`S=0`$ case discussed in Sec. IV of Ref.. As already mentioned the full wave function should transform as a one dimensional irrep of the multisoliton symmetry group $`G`$. For the configurations we are dealing with we have that, except for the $`B=5`$ and $`B=6`$ cases, such one dimensional irrep is the trivial irrep of the corresponding symmetry groups. For $`B=5`$, $`\mathrm{\Gamma }`$ is the $`A_2`$ irrep of $`D_{2d}`$, while for $`B=6`$ the wave functions should transform as the $`A_2`$ irrep of $`D_{4d}`$. Using standard group theoretical arguments we know that the product representation $`J\times N`$ of SU(2) is in general a reducible representation of $`G`$. The projector operator into the one dimensional irrep $`\mathrm{\Gamma }`$ is
$$P_\mathrm{\Gamma }=\frac{1}{|G|}\underset{gG}{}\chi _\mathrm{\Gamma }^{}(g)\rho (g),$$
(30)
where $`|G|`$ is the rank of the group, $`\chi _\mathrm{\Gamma }(g)`$ the character of operation $`g`$, and $`\rho (g)`$ the representation of $`g`$ in $`J\times N`$
$$\rho (g)=D^J(g)\times D^N(D_g).$$
(31)
where $`D_g`$ is the isospin operation associated with the space operation $`g`$. The eigenvalues of $`P_\mathrm{\Gamma }`$ can either vanish or be equal to one. The eigenvectors corresponding to each non-vanishing eigenvalue provide precisely the coefficients $`\alpha _{J_3N_3}^{JN}`$ of Eq. (24), and there are as many wave functions as non-zero eigenvalues. If all eigenvalues vanish there is no collective state with the given $`J,N`$. If there is only one, the wavefunction is an eigenfunction of the collective Hamiltonian. In case there would be more than one, we choose those combinations which diagonalize the parity operator.
## IV Numerical results
To calculate the multibaryon spectra we use the following set of values for the parameters appearing in the effective action, Eqs. (1-4). We fix $`f_\pi `$, $`m_\pi `$ and $`m_K`$ to their empirical values and take $`e=4.1`$ and $`f_K/f_\pi =1.29`$. This set of parameters leads to a single baryon excitation spectrum which is in very good agreement with the one observed for the octet and decuplet baryons. As well known, however, the use of the empirical value for $`f_\pi `$ implies a $`B=1`$ Skyrmion mass of around $`1.7\mathrm{GeV}`$. Consequently, the absolute values of the calculated masses come out to be too large. This problem is nowadays known to be solved by the inclusion of Casimir effects. We will come back to this issue below. With these values we can calculate $`M_{sol}`$ and the different quantities that appear in the expression of $`L_{coll}`$ given by Eq. (12). The results for the different baryon numbers up to $`B=9`$ are tabulated in Tables II and III. From Table II we observe that although $`M_{sol}/B`$ tends, on average, to decrease as a function of $`B`$ it always lies above $`1.5\mathrm{GeV}`$. This clearly indicates that Casimir effects will be also important to determine the absolute masses of the configurations with $`B>1`$. In any case, as in previous works where $`f_\pi `$ was adjusted to reproduce the empirical nucleon mass, we observe some deviation from a smooth behaviour. Also listed in Table II are the strange inertia parameter $`K^S`$ and the symmetry breaking parameter $`\gamma `$. We see that, roughly, $`K^S`$ decreases as $`1/B`$ while $`\gamma `$ increases as $`B^2`$. As we will see this has important consequences on the amount of configuration mixing as a function of the baryon number. In Table III we list the spin, isospin and mixing inertia parameters for the different values of $`B`$. We find that the values we obtained behave, as a function of $`B`$, as those of Ref.. In fact this is to be expected since, as explained in that reference, such behaviour as well as the number of independent components depends only on general properties of the Ansätze.
Given the values of the inertia and symmetry breaking parameters we can proceed to calculate the matrix elements of the rotational Hamiltonian. For this purpose we have to find the solutions of the eigenvalue equation Eq. (25). As explained in the previous section this amounts to determine the coefficients $`\beta _{(p,q)}`$ appearing in Eq. (26). We have done this calculation for the different sets of allowed quantum numbers. It interesting to note that the amount of configuration mixing increases with $`B`$. This can be clearly observed in Fig. 1 where we display the decomposition of the lowest energy states with strangeness $`S=0`$ (full line) and $`S=B`$ (dashed line) for $`B=3`$ and $`B=9`$. In this figure the symbol $`i`$ labels the different $`(p,q)`$ irrep that appear in each decomposition. Of course, $`i=0`$ indicates the corresponding minimal irrep. We see that while for $`B=3`$ about $`80\%`$ of the wavefunction corresponds to the minimal irrep, for $`B=9`$ such irrep represents less that $`30\%`$ with the rest of strength distributed in almost 10 irreps. This kind of behaviour can be simply understood using second order perturbation theory. Within that approximation $`\beta _{i=1}`$, that is the coefficient of the first non-minimal irrep, will be proportional to $`\gamma /(<h>_1<h>_0)`$. It is not hard to show that, for the ground states with $`S=0`$, one has $`(<h>_1<h>_0)B`$. Since we have seen that $`\gamma B^2`$ we obtain that $`\beta _{i=1}`$ should increase roughly as $`B`$. Similar arguments can be used for the case $`S=B`$. This explains why the configuration mixing is quite independent of the value of strangeness as it can be seen in Fig. 1 by comparing the solid lines with the dashed ones. From the numerical point of view the increase of configuration mixing implies that as larger values of $`B`$ are considered one has to increase the size of the basis in which the eigenfunction is expanded in order to obtain convergence. In all the cases of interest we found that no more than $`15`$ to $`20`$ configurations were needed.
The resulting multibaryon spectra are summarized in Tables IV and V. In Table IV we report the rotational corrections to the masses of the $`S=0`$ states. They are given as excitation energies taken with respect to the corresponding lowest energy state whose absolute rotational energy is indicated in brackets. It is important to mention that for the $`B=1`$ systems it was shown that this rather large absolute value is almost completely cancelled by the Casimir corrections due to kaon loops. Since similar cancellations are expected to happen for $`B>1`$, the excitation energies result to be the most meaningful quantities to look at. We observe that the predicted spectra are in agreement with the ones obtained in the alternative bound state approach , except for a few changes in the ordering of the states in the case of $`S=B`$ and $`B=5,8,9`$. From the numbers presented in Tables II, IV and V it is apparent that there is a clear separation of three different energy scales. There is a 1 $`\mathrm{GeV}`$ scale related to the classical masses (per baryon number) and the eigenvalues of Eq. (25) for $`S=0`$ states, there is another scale of about 300 $`\mathrm{MeV}`$ for the excitation of one unit of strangeness, and finally a 10-100 MeV scale related to spin-isospin excitations. This last energy scale is evident in Table IV while it appears as a small correction in Table V. In this way we recover the three leading order contributions in the $`N_c`$ expansion $`N_c,N_c^0`$ and $`N_c^1`$, which are more explicitly separated in the BSA.
## V Conclusions
In this work we have studied the multibaryon spectra for baryon number $`3B9`$ and strangeness values $`S=0,1,B`$ within the SU(3) collective coordinate approach to the three flavor Skyrme model. To describe the classical background solutions we have used Ansätze based on rational maps , which provide very good approximations and also share the same symmetries as the exact solutions. The symmetry structure is responsible for the spin and isospin assignments to the spectrum states. Therefore, the collective Hamiltonians and wave functions we obtain are of general validity, while the mass splittings depend on the particular values of the moments of inertia and of the symmetry breaking parameter.
We have found that, in general, the ordering of the different spin/isospin states corresponding to a given baryon number as well as the energy separation between those states obtained by using the present approach are very similar to the results of the alternative bound state treatment of the SU(3) Skyrme model. This fact together with the observation that in the collective approach the relative strength of the flavor symmetry breaking term increases with increasing baryon number (cf., Fig. 1) seems to indicate that both approaches tend to coincide as $`B`$ grows. In this sense we can conclude that our finding that the increase of one unit of strangeness implies a cost in energy of about 300 MeV rather independently of $`B3`$ appears to be a rather general prediction of the SU(3) Skyrme model.
Finally, note that in the present calculation we have set the meson decay constants to their empirical values. Consequently, all the resulting absolute masses are too large. For example, we find values of $`M_{sol}/B`$ of about 1.60 GeV and $`S=0`$ ground state rotational corrections of about 0.8 GeV. These values are expected to be largely compensated by the pion and kaon contributions to the Casimir energies, respectively. In fact, this has been recently shown to happen in the $`B=1`$ sector of the model. Unfortunately, for $`B>1`$ the difficulties associated with the treatment of the fluctuations around non-spherically symmetric soliton backgrounds have prevented so far the explicit evaluation of the Casimir effect even in the SU(2) case.
###### Acknowledgements.
This work was supported in part by the grant PICT 03-00000-00133 from ANPCYT, Argentina. C.L.S. thanks FAPERJ for financial support and the kind hospitality of Laboratorio TANDAR, CNEA, where part of this work was done.
##
In this Appendix we give the explicit expressions of the spin-isospin collective Hamiltonian for $`B3`$. The form of these expressions depends only on the symmetries of the soliton configurations. The method to derive them is very similar to the one described in Sec. III of Ref. . In fact, the following expressions can be easily obtained from the ones given in that reference by setting the corresponding hyperfine splitting constants to zero.
$`H_{B=3}^{JN}`$ $`=`$ $`H_{B=9}^{JN}=K^J\widehat{J}^2+K^N\widehat{N}^22K^M\stackrel{}{\widehat{N}}\stackrel{}{\widehat{J}},`$ (32)
$`H_{B=4}^{JN}`$ $`=`$ $`K^J\widehat{J}^2+K_1^N\widehat{N}^2+(K_3^NK_1^N)\widehat{N}_{3}^{}{}_{}{}^{2},`$ (34)
$`H_{B=5}^{JN}`$ $`=`$ $`K_1^J(\widehat{J}^2\widehat{J}_{3}^{}{}_{}{}^{2})+K_1^N(\widehat{N}^2\widehat{N}_{3}^{}{}_{}{}^{2})2K_1^M(\stackrel{}{\widehat{N}}\stackrel{}{\widehat{J}}\widehat{N}_3\widehat{J}_3)`$ (37)
$`+K_3^J\widehat{J}_3^2+K_3^N\widehat{N}_3^22K_3^M\widehat{N}_3\widehat{J}_3,`$
$`H_{B=6}^{JN}`$ $`=`$ $`H_{B=8}^{JN}=K_1^J\widehat{J}^2+K_1^N\widehat{N}^2+(K_3^JK_1^J)\widehat{J}_3^2+(K_3^NK_1^N)\widehat{N}_3^22K_3^M\widehat{N}_3\widehat{J}_3,`$ (39)
$`H_{B=7}^{JN}`$ $`=`$ $`K^J\widehat{J}^2+K^N\widehat{N}^2.`$ (41)
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# Asymptotic behaviour of a solution for Kadomtsev-Petviashvili-2 equation This work was partially supported by RFBR (97-01-459), grant for scientific schools (001596038) and INTAS (No1068).
## 1 Introduction
Properties of solutions and process of integration for the Kadomtsev-Petviashvili equation (KP) :
$$_x(_tu+6u_xu+_x^3u)=3\sigma ^2_y^2u$$
(1)
depend on sign of $`\sigma ^2`$. This equation is called KP-2 if $`\sigma ^2=1`$.
In this work the asymptotic behaviour of decreasing solution of equation KP-2 is obtained as $`t\mathrm{}`$. The main term of the asymptotics has an order by $`O(t^1)`$ and fast oscillates. An envelope of these oscillations depends on $`\xi =x/t`$ and $`\eta =y/t`$. Results of the work are formulated in terms of scattering data for auxiliary linear problem, which is associated with equation KP-2 in the inverse scattering transform (IST) .
The equation KP plays important role in modern mathematical physics. Many applications of these equation are known in plasma physics, water waves and other fields of waves process . Therefore questions about solvability of this equation in different functional classes were studied in detail. In particular, the existence of global solution of Cauchy problem in class of distribution functions was proved in . A norm of solution in Sobolev space was estimated and an asymptotics of a solution were obtained as $`t\mathrm{}`$ for KP-like equation but with more high nonlinearity, which called generalized equation KP, in .
The IST formalism allows to reduce a constructing of the solution of nonlinear integrable equation into solving of linear problems. One of most important achievements of IST is the constructing of asymptotic behaviours of solutions as $`t\mathrm{}`$. These results are well-known for 1+1-dimensional equations (one spatial and one temporal variable) -. Asymptotic behaviours of solutions for (2+1)-dimensional equations were studied not so in detail. Rigorous results about temporal asymptotics of solutions, which is nonuniform with respect to spatial variables, for a special class of nondecreasing solutions were obtained in , . In work an formal asymptotics of decreasing solution of equation KP-1, which is nonuniform with respect to spatial variables, was constructed. Results about asymptotic behaviour of decreasing solutions of equation KP-2 as $`t\mathrm{}`$ and about a remainder of this asymptotics, which are uniform with respect to spatial variables, are obtained in present work.
The formalism IST for solving of Cauchy problem for equation KP-2 was presented in . This formalism will be used in this work. Below we remind basic steps of solving of Cauchy problem for equation KP-2 by IST.
Let us denote an initial condition for equation KP-2 as:
$$u|_{t=0}=u_0(x,y).$$
(2)
The process of solving of the Cauchy problem (1), (2) consists of several steps.
First step is solving direct scattering problem. On this step a boundary problem is solved for function $`\phi `$:
$$_y\phi +_x^2\phi +2ik_x\phi +u\phi =0,\phi |_{|k|\mathrm{}}=1,$$
(3)
and scattering data are constructed by a formula:
$`F(k)=`$ $`{\displaystyle \frac{\text{sgn}(\text{Re}(k))}{2\pi }}{\displaystyle }{\displaystyle _\text{}^2}dxdyu_0(x,y)\phi (x,y,k,0)\times `$ (4)
$`\times \mathrm{exp}(i(k+\overline{k})x(k^2\overline{k}^2)y).`$
It is useful to note, if $`u`$ is real, then the function $`F(k)`$ has property $`F(\overline{k})=\overline{F}(k)`$. This follows from formulas (3) and (4).
On the next step an evolution of scattering data is determined. This evolution is very simple:
$$(k;t)=F(k)\mathrm{exp}(4it(k^3+\overline{k}^3)).$$
Third step is solving of inverse scattering problem. This is reduced to so-called $`\overline{D}`$-problem:
$$\begin{array}{c}_{\overline{k}}\phi =\psi F(\overline{k})\mathrm{exp}(itS),\\ _k\psi =\phi F(k)\mathrm{exp}(itS);\end{array}\left(\begin{array}{c}\phi \\ \psi \end{array}\right)|_{|k|\mathrm{}}=\left(\begin{array}{c}1\\ 1\end{array}\right),$$
(5)
where $`S=4(k^3+\overline{k}^3)+(k+\overline{k})\xi i(k^2\overline{k}^2)\eta ,\xi =x/t,\eta =y/t`$, the function $`F(k)`$ is nonanalytic with respect to complex variable $`k\text{}`$. Solving of this problem allows to obtain the functions $`\phi `$ and $`\psi `$ at any time, if we know the evolution of scattering data.
At last we can obtain the solution of the Cauchy problem by using the formula :
$$u(x,y,t)=_x_{\text{}}𝑑kd\overline{k}F(k)\psi (k,x,y,t)\mathrm{exp}(itS).$$
(6)
However, this successive method of solving for Cauchy problems for general initial conditions allows to obtain very implicit answer with respect to initial data. Therefore the problem (5) and integral (6) are used usually for solving of integrable nonlinear equations and the scattering data is taking from some suitable functional class (see, for example, -). On one-hand side this gives for studying a functional class of solutions, but on the other hand-side this leads to implicit restrictions on class of initial conditions.
The auxiliary linear problem (3) and some related problems were studied in -. In the work it was proved that if the function $`u_0(x,y)`$ decreases exponentially, then the direct scattering problem (3) is solvable for $`k\text{}`$.
To construct the asymptotics of $`u`$ as $`t\mathrm{}`$ we use the asymptotics of $`\psi `$ and evaluate an asymptotics of the integral (6) by stationary phase method . As it turned out, the main term of asymptotics of the function $`u`$ may be defined by using only the main term of an asymptotic expansion of $`\psi `$ as $`t\mathrm{}`$. Taking into account the fast oscillating coefficients of the system from (5) one can guess, that $`\psi =1+o(1)`$ as $`t\mathrm{}`$. But we cannot say something definitely about the asymptotic behaviour of the function $`u`$, until we do not know analytic properties and an order of the remainder of asymptotics of $`\psi `$ more precisely. Thus we come to studying an asymptotic behaviour of a solution of the $`\overline{D}`$-problem (5).
The asymptotic behaviour of the solution of the $`\overline{D}`$-problem as $`t\mathrm{}`$ with continuous and fast oscillated coefficients was obtained in . Here we study the $`\overline{D}`$-problem with discontinuous coefficients on imaginary axis of complex parameter $`k`$. Constructing of the asymptotics of solution for such problem is more complicated. Not only stationary points of phase function of an oscillated exponent define a structure of the asymptotic solution for (5) as in (see also ), but the location of these stationary points with respect to line of discontinuity of coefficients of the equations (5) as well. It contributes additional difficulties into evaluating and leads to changes in results. The uniform asymptotics of the solution of the problem (5) is constructed by matching method .
## 2 Main result
###### Theorem 1
Let $`(1+|k|)FL_1C`$, $`^\alpha FL_1C`$ as $`\text{Re}(k)0`$, $`|a|2`$ and:
$$\underset{z\text{}}{sup}_\text{}^2𝑑\kappa 𝑑\lambda \left|\frac{F(\kappa +i\lambda )}{\kappa +i\lambda z}\right|<\mathrm{\hspace{0.17em}2}\pi ,$$
(7)
then the solution of the Cauchy problem of equation KP-2 for corresponding initial condition exists as $`t>0`$. The asymptotic behaviour of the solution as $`t\mathrm{}`$ differs in different domains of variables $`(x,y,t)`$:
as $`(12\xi +\eta ^2)t^{1/3}1`$:
$`u(x,y,t)=4t^1{\displaystyle \frac{\pi }{12i\sqrt{\eta ^212\xi }}}f({\displaystyle \frac{1}{2}}\sqrt{\eta ^212\xi }+{\displaystyle \frac{i\eta }{12}})\times `$
$`\times \mathrm{exp}(11it\sqrt{{\displaystyle \frac{y^2}{t^2}}12{\displaystyle \frac{x}{t}}})+c.c.+o(1).`$
as $`(12\xi +\eta ^2)t^{1/3}1`$:
$$u=o(t^1);$$
as $`|12\xi +12\eta ^2|1`$:
$`u(x,y,t)=8it^1\sqrt{\pi }f(i\eta /12)({\displaystyle _0^{\mathrm{}}}dp_1\sqrt{p_1}\mathrm{cos}(8p_1^3zp_1)+`$
$`+{\displaystyle _0^{\mathrm{}}}dp_1\sqrt{p_1}\mathrm{sin}(8p_1^3zp_1))+o(t^1).`$
Here $`\xi =x/t`$, $`\eta =y/t`$,
$$z=8\left(\frac{y^2}{12t^{4/3}}+\frac{x}{t^{1/3}}\right);$$
$`f(k)={\displaystyle \frac{1}{2\pi }}{\displaystyle }{\displaystyle _\text{}^2}dxdyu_0(x,y)\phi (x,y,k,0)\times `$
$`\times \mathrm{exp}(i(k+\overline{k})x(k^2\overline{k}^2)y).`$
The domains of validity for the asymptotics of the solution of equation KP-2 are intersected and therefore they give combined asymptotics of the solution uniformly on plane of $`x,y`$.
## 3 An analytic behaviours of the scattering data
In this section we demonstrate analytic behaviour of the scattering data which corresponds to sufficiently smooth and decreasing initial condition (2).
First of all we show, that the solution of (5) exists.
###### Theorem 2
Let $`F(k)`$ be such, that the condition (7) is fulfilled, then the solution of (5) exists in a space of continuous vector-functions bounded when $`k\text{}`$.
Proof. Let us consider a system of integral equations which is equivalent to (5):
$$\left(\begin{array}{c}\phi \\ \psi \end{array}\right)=\left(\begin{array}{c}1\\ 1\end{array}\right)+G[F]\left(\begin{array}{c}\phi \\ \psi \end{array}\right),$$
where
$`G[F]V={\displaystyle }{\displaystyle _m\text{}}dmd\overline{m}\times `$
$`\left(\begin{array}{cc}0& \frac{F(\overline{m})}{km}\mathrm{exp}(itS)\\ \frac{F(m)}{\overline{km}}\mathrm{exp}(itS)& 0\end{array}\right)V(m,\xi ,\eta ,t).`$
Using well-known results about integral operators (see, for example ) one can show that the operator $`G[F]`$ transforms the space of continuous vector-functions into itself.
The operator $`G[F]`$ is contracting operator. It is follows from inequality (7). Hence we obtain the theorem statement.
To construct the asymptotics we must study smoothness of scattering data $`F(k)`$ in neighborhood of line of discontinuity $`\text{Re}(k)=0`$.
###### Lemma 1
Let $`u_0(x,y)`$ be a finite function, then in neighborhood of some point $`k^{}\text{}`$ the scattering data have to be represented in the form:
$$F(k)=\text{sgn}[\text{Re}(k)]\underset{|\alpha |=0}{\overset{2}{}}f_{\alpha _1\alpha _2}(k^{})(kk^{})^{\alpha _1}\overline{(kk^{})^{\alpha _2}}+O(|kk^{}|^3),$$
where $`f_{\alpha _1,\alpha _2}(k^{})`$ are continuous functions, when $`\text{Re}(k^{})0`$ and
$$f_{\alpha _1,\alpha _2}(k^{})=f_{\alpha _1\alpha _2}^{(1)}(k^{})+\text{sgn}[\text{Re}(k)]f_{\alpha _1\alpha _2}^{(2)}(k^{}),$$
when $`\text{Re}(k^{})=0`$.
Proof of this lemma consists of successive evaluating of partial derivatives of scattering data with respect to $`k`$ and $`\overline{k}`$. For example let us evaluate $`f_{10}^{(j)}(k^{})`$ and $`f_{10}^{(j)}(k^{})`$ as $`\text{Re}(k^{})=0`$.
Denote
$$f(k)=\frac{1}{2\pi }_\text{}^2𝑑x𝑑yu_0(x,y)\phi (x,y,k,0)\mathrm{exp}(i(k+\overline{k})x(k^2\overline{k}^2)y).$$
Then $`f_{00}^{(1)}(k^{})=f(k^{})`$, $`f_{00}^{(0)}(k^{})0`$.
Evaluate $`f_{10}(k^{})=f_{10}^{(1)}(k^{})+\text{sgn}[\text{Re}(k)]f_{10}^{(0)}(k^{})`$.
$`_kf(k)={\displaystyle \frac{1}{2\pi }}{\displaystyle }{\displaystyle _\text{}^2}dxdyu_0(x,y)\mathrm{exp}(i(k+\overline{k})x(k^2\overline{k}^2)y)\times `$
$`\left[(ix2ky)\phi (x,y,k,0)+_k\phi (x,y,k,0)\right].`$
To evaluate the derivative $`_k\phi `$ we use the integral equation for the function $`\phi `$:
$`_k\phi ={\displaystyle \frac{1}{2i\pi }}V.P.{\displaystyle }{\displaystyle _{\text{}}}{\displaystyle \frac{dmd\overline{m}}{(km)^2}}F(\overline{m})\times `$
$`\mathrm{exp}(i(m+\overline{m})(m^2\overline{m}^2)y)\psi (x,y,m).`$
The functions $`F(\overline{m})`$ and $`\psi (x,y,m)`$ are smooth on the left and right-hand sides of complex plane of $`m`$ and one can show that the integral exists. Let us represent it into more convenient form. For this we rewrite the integral into sum of two integrals over left-hand side and right-hand side of complex plane. Let us integrate by parts these integrals. As a result we obtain
$`_k\phi ={\displaystyle \frac{1}{2i\pi }}{\displaystyle \underset{\pm }{}}{\displaystyle _{\mathrm{\Omega }^\pm }}{\displaystyle \frac{(\pm )d\overline{m}}{km}}[f(\overline{m})\times `$
$`\mathrm{exp}(i(m+\overline{m})x(m^2\overline{m}^2)y)\psi (x,y,m)]`$
$`{\displaystyle \frac{1}{2i\pi }}{\displaystyle \underset{\pm }{}}{\displaystyle }{\displaystyle _{\mathrm{\Omega }^\pm }}{\displaystyle \frac{dmd\overline{m}}{km}}[(ix2my)(\pm )f(\overline{m})\psi (x,y,m)+`$
$`(\pm )_mf(\overline{m})\psi (x,y,m)\phi (x,y,m)f(\overline{m})f(m)]\times `$
$`\times \mathrm{exp}(i(m+\overline{m})x(m^2\overline{m}^2)y).`$
Here $`\mathrm{\Omega }^\pm `$ is right-hand side (+) and left-hand side (-) of complex plane. We mean the integral over $`\mathrm{\Omega }^+`$ as sum of integral over right-hand side and left-hand side of the circle with center at origin of coordinates and with large radius (which tends into infinity) and of improper integral over imaginary line and of an integral over the circle with small radius $`\epsilon `$ with center at $`k=m`$. The direction of the integration over $`\mathrm{\Omega }^+`$ is determined by standard way.
Consider the sum of the integrals over $`\mathrm{\Omega }^\pm `$. Each of the integrals over large half-circles is equal to zero as $`R\mathrm{}`$ because the integrand decreases. The sum of the integrals over the imaginary axis gives doubled integral over the imaginary axis in positive direction. The integral over small circle equals zero as $`\epsilon 0`$.
Let us evaluate the derivative $`_mf(\overline{m})`$ of the integrand of the double integrals. For this we note that $`\phi (x,y,k,0)=\psi (x,y,\overline{k},0)`$, hence $`f(\overline{m})`$ has to be represented by the function $`\psi (x,y,m,0)`$ and then:
$$_mf(\overline{m})=\frac{1}{2\pi }_\text{}^2dxdyu_0(x,y)[f(m)\phi (x,y,m,0)+$$
$$+(ix2my)\psi (x,y,m,0)]\mathrm{exp}(i(m+\overline{m})x(m^2\overline{m}^2)y).$$
So we obtain a formula $`f_{10}(k^{})`$:
$$f_{10}(k^{})=f_{10}^{(1)}(k^{})+\text{sgn}[\text{Re}(k)]f_{10}^{(2)}(k^{}),$$
where
$`f_{10}^{(1)}(k^{})=f(\overline{k}^{})\psi (x,y,k^{},0){\displaystyle \frac{1}{2\pi i}}{\displaystyle }{\displaystyle _\text{}^2}dxdyu_0(x,y)\times `$
$`[(ix2ky)\varphi (x,y,k,0)+`$
$`({\displaystyle \frac{1}{2i\pi }}{\displaystyle \underset{\pm }{}}V.P.{\displaystyle _i\mathrm{}^i\mathrm{}}{\displaystyle \frac{d\overline{m}}{k^{}m}}(\pm )f(\overline{m})\psi (x,y,m,0)`$
$`{\displaystyle \frac{1}{2i\pi }}{\displaystyle \underset{\pm }{}}{\displaystyle }{\displaystyle _{\mathrm{\Omega }^\pm }}{\displaystyle \frac{dmd\overline{m}}{km}}[(ix2my)(\pm )f(\overline{m})\psi (x,y,m)+`$
$`(\pm )_mf(\overline{m})\psi (x,y,m)\phi (x,y,m)f(\overline{m})f(m)]\times `$
$`\mathrm{exp}(i(m+\overline{m})x(m^2\overline{m}^2)y))\times `$
$`\mathrm{exp}(i(k^{}+\overline{k}^{})x(k^2\overline{k}^2)y)];`$
$$f_{10}^{(2)}=\frac{1}{\pi }_\text{}^2𝑑x𝑑yu_0(x,y)f(\overline{k}^{})\psi (x,y,k^{},0).$$
An expression of $`f_{01}(k^{})=f_{01}^{(1)}(k^{})+\text{sgn}[\text{Re}(k)]f_{01}^{(2)}(k^{})`$ has the form:
$`f_{01}^{(1)}(k^{})`$ $`={\displaystyle \frac{1}{2i\pi }}{\displaystyle }{\displaystyle _\text{}^2}dxdyu_0(x,y)(ix+2\overline{k}^{}y)\phi (x,y,k^{},0)\times `$
$`\mathrm{exp}(i(k^{}+\overline{k}^{})x(k^2\overline{k}^2)y),`$
$`f_{01}^{(2)}(k^{})`$ $`={\displaystyle \frac{1}{2i\pi }}{\displaystyle }{\displaystyle _\text{}^2}dxdyu_0(x,y)\psi (x,y,k^{},0)\times `$
$`\times \mathrm{exp}(2i(k^{}+\overline{k}^{})x2(k^2\overline{k}^2)y).`$
The expressions of other coefficients of the expansion have to be evaluated by the same way. The lemma is proved.
## 4 Asymptotic solution of $`\overline{D}`$-problem
In this section we construct an asymptotic solution as $`t\mathrm{}`$ of the $`\overline{D}`$-problem:
$$\begin{array}{c}_{\overline{k}}\mu =\nu F(k)\mathrm{exp}(itS),\\ _k\nu =\mu F(\overline{k})\mathrm{exp}(itS);\end{array}\left(\begin{array}{c}\mu \\ \nu \end{array}\right)|_{|k|\mathrm{}}=\left(\begin{array}{c}1\\ 0\end{array}\right).$$
(9)
The solution of the problem (5) has to be obtained using the solution of (9) and the formula:
$$\left(\begin{array}{c}\phi \\ \psi \end{array}\right)=\left(\begin{array}{c}\mu (k,\xi ,\eta ,t)\\ \nu (k,\xi ,\eta ,t)\end{array}\right)+\left(\begin{array}{c}\nu (\overline{k},\xi ,\eta ,t)\\ \mu (\overline{k},\xi ,\eta ,t)\end{array}\right).$$
(10)
The asymptotic solution of the problem (9) is combined. Here we construct the asymptotics as $`t\mathrm{}`$ uniformly with respect to all parameters. The stationary points of the functions $`S(k,\overline{k},\xi ,\eta )`$ with respect to parameters $`k,\overline{k}`$ plays the important role in these constructions.
The asymptotic expansion of the solutions is constructed using the inverse powers of large parameter $`t`$ as a asymptotic sequence of asymptotic expansion outside of small neighborhoods of the stationary points of the function. Near the stationary points the asymptotic sequence has the form $`t^{n/2}`$, where $`n=0,1,2,\mathrm{}`$. Following the terminology of matching method , we call the asymptotics outside of small neighborhoods of stationary points by outer expansion and we call by interior expansion the asymptotics near the stationary points . The right-hand side in the system of equations (9) is discontinuous on the line $`\text{Re}(k)=0`$, therefore the asymptotic solution is differentiable out of the line $`\text{Re}(k)=0`$.
The domains of validity for interior and exterior asymptotic expansions are intersected. This fact is used by matching method in order to construct unique combined asymptotic expansions.
The phase function $`S`$ depends on two parameters $`(\xi ,\eta )`$ $`\text{}^2`$. On the curve $`12\xi +\eta ^2=0`$ the confluence of two stationary points of the function $`S`$ occurs. In this case we have one confluent stationary point. The structure of the asymptotic expansion of solution of (9) is changed here. As $`|12\xi +\eta ^2|1`$ the expansion is constructed on the powers of $`t^{n/3},`$ $`n=0,1,2,\mathrm{}`$ as a asymptotic sequence.
The formulas for the uniform asymptotic expansions of the solution of (9) are large and it seems convenient to formulate the results about general case in section 4.1 and about confluence case in section 4.2 for convenience .
### 4.1 Asymptotics in a general case
In this section the combined asymptotic solution of (9) is constructed when the phase function $`S`$ has nondegenerate stationary points $`k=k_1`$ and $`k=k_2`$. Here we suppose that $`k\text{},t^{1/3}\left|_k^2S|_{k=k_{1,2}}\right|1`$ as $`t\mathrm{}`$. This leads to restriction on values of the parameters $`\xi `$ and $`\eta `$, namely, $`t^{1/3}|12\xi +\eta ^2|1`$. The asymptotic expansion which is uniform with respect to $`k\text{}`$ is formulated in the end of this section.
To construct the combined asymptotic solution which is valuable as $`k\text{}`$ we obtain exterior and interior asymptotic expansions. These expansions are valid outside of small neighborhoods of $`k_j`$, $`j=1,2`$ and in the small neighborhoods of $`k_j`$ respectively. Denote: $`\theta =\sqrt{12\xi \eta ^2}`$. Then one can obtain an expressions for the stationary points of $`S`$ by using the parameter $`\theta `$: $`k_1=\frac{1}{12}(i\eta +\theta ),k_2=\frac{1}{12}(i\eta \theta )`$.
#### 4.1.1 The stationary points outside of the break line
Consider the case when $`\text{Re}(\theta )0`$, i.e. when the stationary points $`k_{1,2}`$ are outside of the discontinuity line $`\text{Re}(k)=0`$. Let us formulate a result of this section about the combined asymptotic solution:
###### Lemma 2
Let the system of the equations (9) have not the homogeneous solutions, $`F(k)C^2L_1`$ as $`\text{Re}(k)0`$ and the parameters $`\xi `$ and $`\eta `$ satisfy the inequality $`t^{1/3}|\theta |^21`$, then:
when $`\sqrt{t|\theta |}|kk_{1,2}|1`$ the formal asymptotic solution of the system (9) with respect to $`mod(O(t^2|_kS|^3))`$ has the form:
$$\stackrel{~}{\mu }=1+t^1\stackrel{1}{\mu }(k,\xi ,\eta ),$$
$$\stackrel{~}{\nu }=(t^1\underset{1}{\overset{1}{\nu }}(k,\xi ,\eta )+t^2\underset{1}{\overset{2}{\nu }}(k,\xi ,\eta ))\mathrm{exp}(itS)+t^1\underset{0}{\overset{1}{\nu }}(k,\xi ,\eta );$$
the functions $`\stackrel{1}{\mu },\underset{1}{\overset{1}{\nu }},\underset{1}{\overset{2}{\nu }},\underset{0}{\overset{1}{\nu }}`$ are defined by (19), (13), (17), (30);
when $`|\theta |^1|kk_j|1`$ the formal asymptotic solution of the system (9) with respect to $`mod(O(t^1|\theta |^1))`$ has the form:
$$\stackrel{~}{\mu }=1+t^1\stackrel{1}{M}(l_j,\xi ,\eta ),$$
$$\stackrel{~}{\nu }=\left(t^{1/2}\stackrel{1}{N}(l_j,\xi ,\eta )+t^1\stackrel{2}{N}(l_j,\xi ,\eta )\right)\mathrm{exp}(itS),$$
where $`l_j,j=1,2`$, are defined by formula:
$$l_j=\sqrt{t}\left(kk_j\right)\sqrt{\frac{_k^2S_j}{2}+\mathrm{\hspace{0.17em}4}(kk_j)},$$
the functions $`\stackrel{1}{M}(l_j,\xi ,\eta ),\stackrel{1}{N}(l_j,\xi ,\eta ),\stackrel{2}{N}(l_j,\xi ,\eta )`$ are defined by (32), (27), (31).
Proof. Let us construct the external asymptotic solution in the form:
$$\mu ^{ex}=1+t^1\stackrel{1}{\mu }(k,\xi ,\eta )+t^2\underset{1}{\overset{2}{\mu }}(k,\xi ,\eta )\mathrm{exp}(itS)+\mathrm{},$$
(11)
$`\nu ^{ex}=\left(t^1\underset{1}{\overset{1}{\nu }}(k,\xi ,\eta )+t^2\stackrel{2}{\nu }(k,\xi ,\eta )+\mathrm{}\right)\mathrm{exp}(itS)+`$
$`+t^1\underset{0}{\overset{1}{\nu }}(k,\xi ,\eta )+\mathrm{}.`$ (12)
Let us substitute the formulas (11) and (12) into (9), equate coefficients with equal power of $`t`$. As a result we obtain:
$`i_kS\underset{1}{\overset{1}{\nu }}(k,\xi ,\eta )=\text{sgn}(\text{Re}(k))f(\overline{k}),`$
$`_{\overline{k}}\stackrel{1}{\mu }(k,\xi ,\eta )+i_{\overline{k}}S\underset{1}{\overset{2}{\mu }}(k,\xi ,\eta )\mathrm{exp}(itS)=`$
$`\text{sgn}(\text{Re}(k))f(k)\underset{1}{\overset{1}{\nu }}(k,\xi ,\eta )+`$
$`\text{sgn}(\text{Re}(k))f(k)\underset{0}{\overset{1}{\nu }}(k,\xi ,\eta )\mathrm{exp}(itS),`$
$`i_kS\mathrm{exp}(itS)\underset{1}{\overset{2}{\nu }}(k,\xi ,\eta )+_k\underset{0}{\overset{1}{\nu }}(k,\xi ,\eta )=`$
$`\left(\text{sgn}(\text{Re}(k))f(\overline{k})\stackrel{1}{\mu }(k,\xi ,\eta )_k\underset{1}{\overset{1}{\nu }}(k,\xi ,\eta )\right)\mathrm{exp}(itS).`$
If we equate to zero the coefficients of oscillated terms and nonoscillated terms of the expansions (11) and (12) respectively , then we obtain formulas:
$$\underset{1}{\overset{1}{\nu }}=\frac{\text{sgn}(\text{Re}(k))f(\overline{k})}{i_kS};$$
(13)
$$_k\underset{0}{\overset{1}{\nu }}(k,\xi ,\eta )=0;$$
(14)
$$_{\overline{k}}\stackrel{1}{\mu }=\frac{f(\overline{k})f(k)}{i_kS};$$
(15)
$$\underset{1}{\overset{2}{\mu }}(k,\xi ,\eta )=\frac{\text{sgn}(\text{Re}(k))f(k)\underset{0}{\overset{1}{\nu }}(k,\xi ,\eta )}{i_{\overline{k}}S};$$
(16)
$$\stackrel{2}{\nu }=\frac{1}{i_kS}(\text{sgn}(\text{Re}(k))f(\overline{k})\stackrel{1}{\mu }_k\underset{1}{\overset{1}{\nu }}).$$
(17)
The function $`\underset{1}{\overset{1}{\nu }}`$ has a jump on the imaginary axis of $`k`$. This jump is $`\frac{2f(k)}{i_kS}`$. In order for the coefficients of the asymptotics of function $`\nu `$ as $`t^{}1`$ to be continuous, we add an analytic function of $`\overline{k}`$, which has the same jump on the imaginary axis in inverse direction, to $`\underset{1}{\overset{1}{\nu }}\mathrm{exp}(iSt)`$:
$$\underset{0}{\overset{1}{\nu ^{}}}(\overline{k})=\frac{1}{i\pi }_i\mathrm{}^i\mathrm{}\frac{dnf(n)}{(\overline{k}n)_nS}.$$
(18)
The function $`\underset{1}{\overset{1}{\nu ^{}}}`$ define an analytic function $`\underset{0}{\overset{1}{\nu }}`$ of $`\overline{k}`$ uncompletely. The rest terms will be defined when we will match the exterior and interior expansion.
A Cauchy-Green formula gives solution of (15) which is decreasing as $`|k|\mathrm{}`$ and bounded when $`k\text{}`$:
$$\stackrel{1}{\mu }=\frac{1}{2i\pi }_{\text{}}\frac{dpd\overline{p}}{kp}\frac{f(\overline{p})f(p)}{i_pS}.$$
(19)
We can obtain the domain of values of $`k`$, where the external expansion is valid, using the condition $`|t^1\underset{1}{\overset{1}{\nu }}/(t^2\stackrel{2}{\nu })|1`$. As a result of calculations we obtain:
$$\sqrt{t|\theta |}|kk_j|1.$$
Let us construct the interior asymptotic solution which is valid in the neighborhood of point $`k_j`$ as $`j=1,2`$. Denote new scaling variable by $`l_j`$:
$$l_j^2=t(kk_j)^2\frac{_k^2S_j}{2}+\mathrm{\hspace{0.17em}4}t(kk_j)^3.$$
(20)
When $`|l_j|`$ is not so large ($`t^{1/2}|_k^2S_j|^{3/2}|l_j|`$) an asymptotic formula is valid as $`t\mathrm{}`$
$$(kk_j)=\sqrt{\frac{2}{t_k^2S_j}}l_j\frac{8}{t(_k^2S_j)^2}l_j^2+\mathrm{}.$$
(21)
Rewrite the system (9) into terms of new variables $`l_j`$ and $`\overline{l}_j`$. Substitute the asymptotic expansion
$$\mu ^{in}=1+t^1\underset{j}{\overset{1}{M}}(l_j,\xi ,\eta )+\mathrm{},$$
(22)
$$\nu ^{in}=(t^{1/2}\underset{j}{\overset{1}{N}}(l_j,\xi ,\eta )+t^1\underset{j}{\overset{2}{N}}(l_j,\xi ,\eta )+\mathrm{})\mathrm{exp}(itS),$$
(23)
into system (9).
As a result one can obtain equations for the coefficients of expansions (22) and (23):
$$_{l_j}\underset{j}{\overset{1}{N}}2il_j\underset{j}{\overset{1}{N}}=\sqrt{\frac{2}{_k^2S_j}}\text{sgn}(\text{Re}(k_j))f(\overline{k}_j);$$
(24)
$$_{l_j}\underset{j}{\overset{2}{N}}2il_j\underset{j}{\overset{2}{N}}=\text{sgn}(\text{Re}(k_j))[(\frac{16}{(_k^2S_j)^2}f(\overline{k}_j)\frac{2}{_k^2S_j}f_{10}(\overline{k}_j))l_j$$
$$\frac{2\overline{l}_j}{|_k^2S_j|}f_{01}(\overline{k}_j)];$$
(25)
$$_{\overline{l}_j}\underset{j}{\overset{1}{M}}=\text{sgn}(\text{Re}(k_j))f(k_j)\underset{j}{\overset{1}{N}}\sqrt{\frac{2}{_{\overline{k}}^2S_j}}.$$
(26)
Construct the solutions of the equations (24)–(26). Since the external expansion has not the terms of order $`t^{1/2}`$, then the boundary condition for function $`\underset{1}{\overset{1}{N}}(l_j,\xi ,\eta )`$ has the form:
$$\underset{j}{\overset{1}{N}}(l_j,\xi ,\eta )|_{|l_j|\mathrm{}}=0.$$
The solution of the boundary condition for $`\underset{j}{\overset{1}{N}}(l_j,x,y)`$ is evaluated by formula:
$`\underset{j}{\overset{1}{N}}(l_j,\xi ,\eta )=\text{sgn}(\text{Re}(k_j)){\displaystyle \frac{\sqrt{2}f(\overline{k}_j)}{\sqrt{_k^2S_j}}}{\displaystyle \frac{\mathrm{exp}(i(l_j^2+\overline{l}_j^2))}{2i\pi }}\times `$
$`\times {\displaystyle }{\displaystyle _{\text{}}}{\displaystyle \frac{dnd\overline{n}}{\overline{l_jn}}}\mathrm{exp}(i(n^2+\overline{n}^2));`$ (27)
The solution of nonhomogeneous equation for the function $`\underset{j}{\overset{2}{N}}(l_j,\xi ,\eta )`$ has the form:
$`\underset{j}{\overset{s}{\stackrel{2}{N}}}(l_j,\xi ,\eta )=\text{sgn}(\text{Re}(k_j))[{\displaystyle \frac{8f(\overline{k}_j)}{(_k^2S_j)^2}}{\displaystyle \frac{f_{10}(\overline{k}_j)}{_k^2S_j}}`$
$`{\displaystyle \frac{2f_{01}(\overline{k}_j)}{|_k^2S_j|}}\mathrm{exp}(i(l_j^2+\overline{l}_j^2))`$
$`\overline{l}_j{\displaystyle \frac{2f_{01}(\overline{k}_j)}{|_k^2S_j|}}{\displaystyle \frac{\mathrm{exp}(i(l_j^2+\overline{l}_j^2))}{2i\pi }}{\displaystyle }{\displaystyle _{\text{}}}{\displaystyle \frac{dnd\overline{n}}{\overline{l_jn}}}\mathrm{exp}(i(n^2+\overline{n}^2))].`$ (28)
The partial solution of the equation for the function $`\underset{j}{\overset{1}{M}}(l_j,\xi ,\eta )`$ has to be written as four-multiple integral:
$$\underset{j}{\overset{s}{\stackrel{1}{M}}}=\frac{2f(k_j)f(\overline{k}_j)}{|_{\overline{k}}^2S_j|}J,$$
(29)
where
$`J={\displaystyle \frac{1}{2i\pi }}{\displaystyle }{\displaystyle _{\text{}}}{\displaystyle \frac{dnd\overline{n}}{l_jn}}{\displaystyle \frac{\mathrm{exp}(i(n^2+\overline{n}^2))}{2i\pi }}\times `$
$`{\displaystyle _{\text{}}\frac{dmd\overline{m}}{\overline{nm}}\mathrm{exp}(i(m^2+\overline{m}^2))}.`$
It is possible to reduce this four-multiple integral into two-multiple integral (see Appendix):
$`J=\overline{l}_j{\displaystyle \frac{\mathrm{exp}(i(l_j^2+\overline{l}_j^2))}{2i\pi }}{\displaystyle _{\text{}}\frac{dnd\overline{n}}{\overline{l_jn}}\mathrm{exp}(i(n^2+\overline{n}^2))}`$
$`\mathrm{exp}(i(l_j^2+\overline{l}_j^2)).`$
There exists an domain of values of complex parameter $`k`$, in which the external and internal asymptotic expansions of solution for the problem (9) are valid. In this domain these asymptotic expansions are equal up to the terms of order $`o(t^1)`$. In the domain, where $`t^{1/2}|kk_j|1`$ and $`|kk_j|t^{1/4}`$, the external and internal expansion are valid. We compute asymptotics of external expansion as $`kk_j`$ and internal expansion as $`|l_j|\mathrm{}`$.
Let us present the asymptotics of the functions $`\underset{j}{\overset{1}{N}}`$ and $`\underset{j}{\overset{s}{\stackrel{2}{N}}}`$ as $`|l_j|\mathrm{}`$:
$`\underset{j}{\overset{1}{N}}|_{|l_j|\mathrm{}}=\text{sgn}(\text{Re}(k_j))\sqrt{{\displaystyle \frac{2}{_k^2S_j}}}f(\overline{k}_j)\times `$
$`\left({\displaystyle \frac{1}{2il_j}}+{\displaystyle \frac{\mathrm{exp}(i(l_j^2+\overline{l}_j^2))}{\overline{l}_j}}+O(|l_j|^3)\right).`$
$`\underset{j}{\overset{2}{N}}(l_j,\xi ,\eta )|_{|l_j|\mathrm{}}=\text{sgn}(\text{Re}(k_j))[({\displaystyle \frac{8if(\overline{k}_j)}{(_k^2S_j)^2}}{\displaystyle \frac{if_{10}(\overline{k}_j)}{_k^2S_j}})`$
$`+{\displaystyle \frac{\overline{l}_j}{2il_j}}{\displaystyle \frac{2f_{01}(\overline{k}_j)}{|_{\overline{k}}^2S_j|}}+O(|l_j|^1)].`$
The matching condition for $`\stackrel{~}{\nu }`$ means that in the domain $`t^{1/3}|kk_j|1`$ and $`|kk_j|=o(1)`$:
$$\left(t^{1/2}\underset{j}{\overset{1}{N}}(l_j,\xi ,\eta )+t^1\underset{j}{\overset{2}{N}}(l_j,\xi ,\eta )\right)\mathrm{exp}(itS)$$
$$t^1(\underset{1}{\overset{1}{\nu }}(k,\xi ,\eta )\mathrm{exp}(itS)+\underset{0}{\overset{1}{\nu }}(k,\xi ,\eta ))=o(t^1).$$
Using the asymptotics of interior expansions of $`\underset{j}{\overset{2}{N}}`$, which are rewritten in the terms of external variable $`k`$ as $`kk_{1,2}`$, and external expansion of $`\underset{0}{\overset{1}{\nu }}`$, we can obtain:
$`\underset{0}{\overset{1}{\nu }}(k,x,y)`$ $`=\underset{0}{\overset{1}{\nu ^{}}}(\overline{k})+{\displaystyle \frac{2f(\overline{k}_1)\mathrm{exp}(itS_1)}{|_k^2S_1|\overline{(kk_1)}}}+`$ (30)
$`+{\displaystyle \frac{2f(\overline{k}_2)\mathrm{exp}(itS_2)}{|_k^2S_2|\overline{(kk_2)}}},`$
where the function $`\underset{0}{\overset{}{\stackrel{1}{\nu }}}(\overline{k})`$ is defined by (18);
$`\underset{j}{\overset{2}{N}}=\underset{j}{\overset{s}{\stackrel{2}{N}}}+\underset{0}{\overset{1}{\nu ^{}}}(k_j)\mathrm{exp}(i(l_j^2+\overline{l}_j^2))\mathrm{exp}(itS_j)+`$
$`+{\displaystyle \frac{2\text{sgn}(\text{Re}(k_m))f(\overline{k}_m)\mathrm{exp}(itS_m)}{|_k^2S_m|\overline{(k_jk_m)}}}\mathrm{exp}(i(l_j^2+\overline{l}_j^2)),`$ (31)
where $`mj`$, the function $`\underset{j}{\overset{s}{\stackrel{2}{N}}}`$ is defined by (28).
Let us construct $`\underset{j}{\overset{1}{M}}(l_j,\xi ,\eta )`$. The asymptotics of (29) as $`|l_j|\mathrm{}`$ has the form:
$$\underset{j}{\overset{s}{\stackrel{1}{M}}}|_{|l_j|\mathrm{}}=\frac{\overline{l}_j}{l_j}\sqrt{\frac{_k^2S_j}{_{\overline{k}}^2S_j}}\frac{i}{12\pi }\frac{f(\overline{k}_j)f(k_j)}{_k^2S}.$$
Evaluate the asymptotics of $`\underset{0}{\overset{1}{\mu }}`$ as $`kk_j`$:
$$\underset{0}{\overset{1}{\mu }}|_{kk_j}=\frac{\overline{(kk_j)}}{kk_j}\frac{f(\overline{k}_j)f(k_j)}{12i(k_jk_n)}\frac{f(k_j)f(\overline{k}_j)}{12i(k_jk_n)}+$$
$$+\frac{1}{2i\pi }_{\text{}}𝑑pd\overline{p}\frac{f(\overline{p})f(p)f(k_j)f(\overline{k}_j)}{12i(pk_j)^2(pk_n)},$$
here $`n1,2,nj`$.
The matching condition for $`\stackrel{~}{\mu }`$ when $`t^{1/3}|kk_j|1`$ and $`|kk_j|=o(t^{1/4})`$ has the form:
$$(1+t^1\underset{0}{\overset{1}{\mu }}(k,\xi ,\eta ))(1+t^1\underset{j}{\overset{1}{M}}(l_j,\xi ,\eta ))=o(t^1).$$
These condition allows to define $`\underset{j}{\overset{1}{M}}(l_j,\xi ,\eta )`$:
$$\underset{j}{\overset{1}{M}}(l_j,\xi ,\eta )=\underset{j}{\overset{s}{\stackrel{1}{M}}}(l_j,\xi ,\eta )+C_j(\xi ,\eta ).$$
(32)
Here the function $`\underset{j}{\overset{s}{\stackrel{1}{M}}}(l_j,\xi ,\eta )`$ is defined by (29), $`C_j(\xi ,\eta )`$ has the form:
$`C_j(\xi ,\eta )={\displaystyle \frac{f(k_j)f(\overline{k}_j)}{12i(k_jk_n)}}+`$
$`{\displaystyle \frac{1}{2i\pi }}{\displaystyle _{\text{}}𝑑p}d\overline{p}{\displaystyle \frac{f(\overline{p})f(p)f(k_j)f(\overline{k}_j)}{12i(pk_j)^2(pk_n)}},`$ (33)
where $`n1,2,nj`$.
Thus we have constructed the interior expansion in neighborhood of nondegenerate stationary point of $`S`$ as $`|12\xi +\eta ^2|t^{1/4}`$.
Lemma is proved.
Constructed asymptotic solutions are nonuniform with respect to $`k`$. But one can obtain an uniform asymptotic solution by using their combination. This uniform solution has the form (see, for example, ):
$`\left(\begin{array}{c}\widehat{\mu }\\ \widehat{\nu }\end{array}\right)=\left(\begin{array}{c}\stackrel{~}{\mu }\\ \stackrel{~}{\nu }\end{array}\right)+\left(\begin{array}{c}\stackrel{~}{M}_1\\ \stackrel{~}{N}_1\end{array}\right)+\left(\begin{array}{c}\stackrel{~}{M}_2\\ \stackrel{~}{N}_2\end{array}\right)`$ (42)
$`A_{1,k}\left(\begin{array}{c}\stackrel{~}{M}_1\\ \stackrel{~}{N}_1\end{array}\right)A_{1,k}\left(\begin{array}{c}\stackrel{~}{M}_2\\ \stackrel{~}{N}_2\end{array}\right).`$ (47)
The result of action of operator $`A_{n,k}`$ on the function $`\stackrel{~}{M}`$ is defined as follows (). Take the formula for $`\stackrel{~}{M}(l_j,\xi ,\eta ,t)`$, and change the dependence on variable $`l_j`$ into the dependence on variable $`k`$ using the formula (20). Rewrite the sum of all terms of asymptotic expansion up to $`t`$ with powers equal to $`m`$, where $`0mn`$. For example for the functions $`\stackrel{~}{M}(l_1,\xi ,\eta ,t)`$ and $`\stackrel{~}{N}(l_1,\xi ,\eta ,t)`$ this process leads to the formulas:
$$A_{1,k}(\stackrel{~}{M}_1)=1+t^1\frac{\overline{\sqrt{\theta (kk_1)^2+4(kk_1)^3}}}{\sqrt{\theta (kk_1)^2+4(kk_1)^3}}\frac{f(\overline{k}_1)f(k_1)}{2|_k^2S_1|};$$
$`A_{1,k}(\stackrel{~}{N}_1)=t^1[{\displaystyle \frac{f(\overline{k}_1)}{2i\theta \sqrt{\theta (kk_1)^2+4(kk_1)^3}}}`$
$`{\displaystyle \frac{f(\overline{k}_1)\mathrm{exp}(it(SS_1))}{|\theta |\overline{\sqrt{\theta (kk_1)^2+4(kk_1)^3}}}}+`$
$`+\underset{0}{\overset{1}{\nu ^{}}}{\displaystyle \frac{f_{01}(\overline{k}_1)}{|\theta |}}{\displaystyle \frac{\overline{\sqrt{\theta (kk_1)^2+4(kk_1)^3}}}{\sqrt{\theta (kk_1)^2+4(kk_1)^3}}}].`$
If we substitute (47) into (9) and evaluate a remainder, then we obtain:
###### Theorem 3
The formal asymptotic solution of the problem (9) with respect to $`mod\left(O\left((t|\theta |)^1\right)\right)`$, which is uniformly valuable when $`k\text{}`$ and $`\theta ^2t1`$, has the form (47).
#### 4.1.2 The stationary points on the imaginary axis
If $`\text{Re}(\theta )=0`$, then the stationary point of the phase function $`S`$ belongs to the line, where the coefficients of the equation (9) are discontinuous. In this case constructing of the formal asymptotic solution of problem (9) differs from the asymptotic solution which was constructed before. The main result about the combined asymptotic solution is:
###### Lemma 3
Let the system of equations (9) have no homogeneous solutions, $`F(k)C^2L_1`$ as $`\text{Re}(k)0`$; the parameters $`\xi `$ and $`\eta `$ are $`t^{2/3}(12\xi +\eta ^2)1`$, then:
when $`\sqrt{t|\theta |}|kk_{1,2}|1`$the formal asymptotic solution of (9) with respect to $`mod(O(t^2|_kS|^3))`$ has the form:
$$\stackrel{~}{\mu }=1+t^1\stackrel{1}{\alpha }(k,\xi ,\eta ),$$
$$\stackrel{~}{\nu }=(t^1\underset{1}{\overset{1}{\beta }}(k,\xi ,\eta )+t^2\underset{1}{\overset{2}{\beta }}(k,\xi ,\eta ))\mathrm{exp}(itS)+t^1\underset{0}{\overset{1}{\beta }}(k,\xi ,\eta );$$
the functions $`\stackrel{1}{\alpha },\underset{1}{\overset{1}{\beta }},\underset{1}{\overset{2}{\beta }},\underset{0}{\overset{1}{\beta }}`$ are defined by (50), (51), (54), (60);
when $`|\theta ||kk_j|1`$ the formal asymptotic solution of (9) with respect to $`mod(O(t^1|\theta |^1))`$ has the form:
$$\stackrel{~}{Y}=1+t^1\stackrel{1}{Y}(l_j,\xi ,\eta ),$$
$$\stackrel{~}{Z}=\left(t^{1/2}\stackrel{1}{Z}(l_j,\xi ,\eta )+t^1\stackrel{2}{Z}(l_j,\xi ,\eta )\right)\mathrm{exp}(itS),$$
where the functions $`\stackrel{1}{Y}(l_j,\xi ,\eta ),\stackrel{1}{Z}(l_j,\xi ,\eta ),\stackrel{2}{Z}(l_j,\xi ,\eta )`$ are defined by (62), (61).
The proof. Let us construct the asymptotics. The external expansion is constructed similarly as in 4.1.1. The main difference is that the stationary points of the function $`S`$ are on the line of discontinuity of the coefficients of system (9). It leads to sufficient modifications in the formulas for asymptotics. Let us find the external expansion in the form:
$$\mu ^{ex}=1+t^1\stackrel{1}{\alpha }(k,\xi ,\eta )+t^2\underset{1}{\overset{2}{\alpha }}(t,\xi ,\eta )\mathrm{exp}(itS)+\mathrm{},$$
(48)
$$\nu ^{ex}=t^1(\stackrel{1}{\beta }(k,\xi ,\eta )\mathrm{exp}(itS)+\underset{0}{\overset{1}{\beta }}(k,\xi ,\eta ))+t^2\stackrel{2}{\beta }\mathrm{exp}(itS)+\mathrm{}.$$
(49)
After substituting of (48) and (49) into the system (9) we equate the coefficients with identical powers of $`t`$ and with oscillated and nonoscillated terms correspondingly. As a result we obtain equations for the coefficients of asymptotic expansions (48) and (49). The obvious formulas are
$$\stackrel{1}{\alpha }=\frac{1}{2i\pi }_{\text{}}\frac{dnd\overline{n}}{kn}\frac{f(\overline{n})f(n)}{i_nS(n,\xi ,\eta )},$$
(50)
$$\underset{1}{\overset{1}{\beta }}=\frac{\text{sgn}(\text{Re}(k))f(\overline{k})}{i_kS}.$$
(51)
The function $`\underset{1}{\overset{1}{\beta }}`$ is discontinuous on the imaginary axis. We add an analytic function with the same jump on the imaginary axis in back direction then the coefficient of the expansion (49) at $`t^1`$ is continuous:
$$\underset{0}{\overset{1}{\beta ^{}}}=\frac{1}{\pi i_{\overline{k}}S}_i\mathrm{}^i\mathrm{}\frac{dnf(n)_{\overline{n}}S}{(\overline{k}n)_nS}=\frac{1}{\pi i_{\overline{k}}S}_i\mathrm{}^i\mathrm{}\frac{dnf(n)}{\overline{k}n}.$$
(52)
Unlike $`\underset{0}{\overset{1}{\nu ^{}}}`$, this function have singularities on the discontinuous line.
The formulas for the other coefficients of the expansions (48) and (49) have the forms:
$$\underset{1}{\overset{2}{\alpha }}=\frac{\text{sgn}(\text{Re}(k))f(k)\underset{0}{\overset{1}{\beta }}(k,\xi ,\eta )}{i_{\overline{k}}S},$$
(53)
$$\stackrel{2}{\beta }=\frac{1}{i_kS}(\text{sgn}(\text{Re}(k)f(\overline{k})\stackrel{1}{\alpha }_k\underset{1}{\overset{1}{\beta }}).$$
(54)
Here the analytic function $`\underset{0}{\overset{1}{\beta }}`$ is still undefined. We will finally define this function after matching of external and internal expansions.
The interior expansion in the neighborhood of the point $`k_j`$ depend on the scaling variable $`l_j`$. The asymptotics of the functions $`\mu `$ and $`\nu `$ has the same asymptotic sequence, but the equations for $`\underset{j}{\overset{1}{Z}}`$, $`\underset{j}{\overset{2}{Z}}`$ and $`\underset{j}{\overset{1}{Y}}`$ have discontinuous right-hand sides.
$`_{l_j}\underset{j}{\overset{1}{Z}}2il_j\underset{j}{\overset{1}{Z}}=\sqrt{{\displaystyle \frac{2}{_k^2S_j}}}f(\overline{k}_j)\text{sgn}(\text{Re}(\overline{l}\mathrm{exp}(i\pi /4)));`$ (55)
$`_{l_j}\underset{j}{\overset{2}{Z}}2il_j\underset{j}{\overset{2}{Z}}=\text{sgn}(\text{Re}[\overline{l}\mathrm{exp}(i\pi /4)])\times `$
$`\times \left({\displaystyle \frac{16}{(_k^2S_j)^2}}f_{00}(\overline{k}_j)+{\displaystyle \frac{2l_j}{_k^2S_j}}f_{01}^{(1)}(\overline{k}_j)+{\displaystyle \frac{2i\overline{l}_j}{|_k^2S_j|}}f_{10}^{(1)}(\overline{k}_j)\right)+`$
$`+{\displaystyle \frac{2l_j}{_k^2S_j}}f_{01}^{(2)}(\overline{k}_j)+{\displaystyle \frac{2i\overline{l}_j}{|_k^2S_j|}}f_{10}^{(2)}(\overline{k}_j).`$ (56)
$`_{\overline{l}_j}\underset{j}{\overset{1}{Y}}=\text{sgn}(\text{Re}(l\mathrm{exp}(i\pi /4)))f(k_j)\underset{j}{\overset{1}{Z}}\sqrt{{\displaystyle \frac{2}{_{\overline{k}}^2S_j}}}.`$ (57)
Continuous partial solutions of the equations (55), (56) and (57) are obtained by Cauchy-Green formula. The formulas for the partial solutions of the equations (55) and (56) are differing from the formulas obtained in section 4.1.1 by function $`\text{sgn}(\text{Re}(\overline{n}\mathrm{exp}(i\pi )))`$ under the sign of double integral:
$$\underset{j}{\overset{s}{\stackrel{1}{Z}}}=𝒥[\stackrel{1}{H_j}](l_j),\underset{j}{\overset{s}{\stackrel{2}{Z}}}=𝒥[\stackrel{2}{H_j}](l_j),$$
(58)
where $`\underset{j}{\overset{1}{H}}`$ and $`\underset{j}{\overset{2}{H}}`$ are the right-hand side of the equations (55) and (56). An operator $`𝒥`$ has the form:
$$𝒥[h](l_j)=\frac{\mathrm{exp}(i(l_j^2+\overline{l}_j^2))}{2i\pi }_{\text{}}\frac{dnd\overline{n}}{\overline{l_jn}}h(n)\mathrm{exp}(i(l_j^2+\overline{l}_j^2)).$$
The partial solution of (57) contains of the double integral the same as in (29) and the term, which is defined by integral over the discontinuous line. The function $`\underset{j}{\overset{s}{\stackrel{1}{Y}}}`$ has the form:
$$\underset{j}{\overset{s}{\stackrel{1}{Y}}}=\frac{2f(k_j)f(\overline{k}_j)}{|_{\overline{k}}^2S_j|}(J2J_+2J_+),$$
(59)
where:
$`J_\pm ={\displaystyle }{\displaystyle _\mathrm{\Omega }^{}}{\displaystyle \frac{dmd\overline{m}}{lm}}{\displaystyle \frac{\mathrm{exp}(i(m^2+\overline{m}^2))}{2i\pi }}\times `$
$`{\displaystyle _{\mathrm{\Omega }^\pm }\frac{dnd\overline{n}}{\overline{nm}}\mathrm{exp}(i(n^2+\overline{n}^2))}.`$
Here $`\mathrm{\Omega }^\pm =\{\pm \text{Re}[l\mathrm{exp}(i\pi /4)>0\}`$. Note that $`J_+(l)=J_+(l)`$. In the Appendix we show, that we can reduce the integral $`J_+`$ into double integral. As a result we obtain:
$`J_+={\displaystyle \frac{3}{4}}i\pi +`$
$`il{\displaystyle _{\mathrm{\Omega }^+}\frac{dmd\overline{m}}{il\overline{m}}\mathrm{exp}(i(m^2+\overline{m}^2))},\text{when}l\mathrm{\Omega }^+;`$
$`J_+=\overline{l}\mathrm{exp}(i(l^2+\overline{l}^2)){\displaystyle _{\mathrm{\Omega }^+}\frac{dmd\overline{m}}{il\overline{m}}\mathrm{exp}(i(m^2+\overline{m}^2))}`$
$`{\displaystyle \frac{1}{4}}i\pi {\displaystyle \frac{1}{2}}i\pi \mathrm{exp}(i(l^2+\overline{l}^2)),\text{when}l\mathrm{\Omega }^{};`$
Thus, the function $`\underset{j}{\overset{s}{\stackrel{1}{Y}}}`$ is represented by sum of double integrals.
Let us consider matching of the external and internal asymptotic expansions. For that we need the asymptotic behaviour of the partial solutions of (55) and (56) as $`|l|\mathrm{}`$. The asymptotic behaviour of the double integrals is evaluated in the Appendix. By using it we obtain:
$`\underset{j}{\overset{s}{\stackrel{1}{Z}}}|_{|l|\mathrm{}}=\text{sgn}\left[\text{Re}[\overline{l}_j\mathrm{exp}(i\pi /4)]\right]f(\overline{k}_j)\sqrt{{\displaystyle \frac{2}{_k^2S_j}}}({\displaystyle \frac{1}{2il_j}}+`$
$`+{\displaystyle \frac{\mathrm{exp}(i(l_j^2+\overline{l}_j^2))}{2\overline{l}_j}}+O(|l_j|^2));`$
$`\underset{j}{\overset{s}{\stackrel{2}{Z}}}|_{|l_j|\mathrm{}}=\text{sgn}\left[\text{Re}[\overline{l}_j\mathrm{exp}(i\pi /4)]\right]\times `$
$`[({\displaystyle \frac{\overline{l}_j}{2il_j}}{\displaystyle \frac{1}{2}}\mathrm{exp}(i(l_j^2+\overline{l}_j^2))){\displaystyle \frac{2if_{10}^{(1)}}{|_k^2S_j|^2}}+`$
$`+({\displaystyle \frac{1}{4}}{\displaystyle \frac{1}{4}}\mathrm{exp}(i(l_j^2+\overline{l}_j^2)))2f_{01}^{(1)}]+`$
$`{\displaystyle \frac{2if_{10}^{(2)}}{|_k^2S_j|}}\left({\displaystyle \frac{\overline{l}_j}{2il}}+{\displaystyle \frac{1}{2}}\mathrm{exp}(i(l_j^2+\overline{l}_j^2))\right)+if_{01}^{(2)}+O(|l_j|^1).`$
The asymptotic behaviour of the function $`\underset{0}{\overset{\prime \prime }{\stackrel{1}{\beta }}}`$ is:
$`\underset{0}{\overset{\prime \prime }{\stackrel{1}{\beta }}}|_{kk_j}={\displaystyle \frac{1}{12i\overline{(kk_j)(k_jk_m)}}}[f(k_j)\text{sgn}[\text{Re}(\overline{k})]+`$
$`+{\displaystyle \frac{1}{\pi i}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}dn{\displaystyle \frac{f(in)f(k_j)}{n\overline{ik_j}}}]+`$
$`{\displaystyle \frac{1}{12i(k_jk_m)^2}}[f(k_j)\text{sgn}[\text{Re}(\overline{k})]+`$
$`{\displaystyle \frac{1}{\pi i}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}dn{\displaystyle \frac{f(in)f(k_j)}{n\overline{ik_j}}}]+{\displaystyle \frac{1}{12i\overline{(k_jk_m)}}}[f_{01}^{(1)}(k_j)+`$
$`+f_{01}^{(2)}(k_j)\text{sgn}[\text{Im}(\overline{ik})]+f_{10}^{(1)}(k_j)+f_{10}^{(2)}(k_j)\text{sgn}[\text{Im}(\overline{ik})]+`$
$`+{\displaystyle \frac{1}{\pi i}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}dn{\displaystyle \frac{f(in)f(k_j)\overline{(inik_j)}f_{01}^{(2)}(k_j)(inik_j)f_{10}^{(2)}(k_j)}{(n\overline{ik_j})^2}}]`$
$`+o(1).`$
Using the matching conditions for the asymptotics of the external $`\nu ^{ex}`$ and internal $`\nu ^{in}`$ expansions we obtain:
$$\underset{0}{\overset{1}{\beta }}=\underset{0}{\overset{1}{\beta ^{}}}\frac{C_1^1}{12i\overline{(kk_1)(k_1k_2)}}\frac{C_2^1}{12i\overline{(kk_2)(k_2k_1)}},$$
(60)
where $`\underset{0}{\overset{}{\stackrel{1}{\beta }}}`$ is defined by formula (52),
$$C_j^1=\frac{1}{\pi i}_{\mathrm{}}^{\mathrm{}}𝑑n\frac{f(n)f(k_j)}{n\overline{ik_j}},j=1,2;$$
$$\underset{j}{\overset{1}{Z}}=\underset{j}{\overset{s}{\stackrel{1}{Z}}}\underset{j}{\overset{2}{Z}}=\underset{j}{\overset{s}{\stackrel{2}{Z}}}\mathrm{exp}(i(l_j^2+\overline{l}_j^2))C_j^2,$$
(61)
where functions $`\underset{j}{\overset{s}{\stackrel{1}{Z}}}`$ and $`\underset{j}{\overset{s}{\stackrel{2}{Z}}}`$ are defined by (58),
$`C_j^1={\displaystyle \frac{i}{_k^2S_j}}f_{10}^{(2)}(\overline{k}_j)`$
$`{\displaystyle \frac{1}{\pi i}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑n{\displaystyle \frac{f(n)f(\overline{k}_j)f_{01}^{(2)}(\overline{k}_j)\overline{(nik_j)}(nik_j)f_{10}^{(2)}(k_j)}{(n\overline{ik_j})^2}}.`$
The asymptotic behaviour of the function $`\underset{j}{\overset{s}{\stackrel{1}{Y}}}`$ is :
$$\underset{j}{\overset{s}{\stackrel{1}{Y}}}=(i\pi +\frac{i\overline{l}_j}{12\pi l_j})\frac{f(\overline{k}_j)f(k_j)}{|_k^2S_j|}.$$
The matching condition for the external expansion of $`\mu ^{ex}`$ and for the internal expansion of $`\mu ^{in}`$ gives:
$$\underset{j}{\overset{1}{Y}}=\underset{j}{\overset{s}{\stackrel{1}{Y}}}+C_j(\xi ,\eta )+i\pi \frac{f(\overline{k}_j)f(k_j)}{|_k^2S_j|},$$
(62)
where $`\underset{j}{\overset{s}{\stackrel{1}{Y}}}`$ is defined by (59), The function $`C_j(\xi ,\eta )`$ is defined by (33).
The lemma is proved.
Constructed external and internal expansions are irregular with respect to parameters $`k,\xi ,\eta `$. The uniform expansion with respect to $`k`$ when$`(12\xi +\eta ^2)t^{1/3}1`$ has to be constructed by the same way as in preceding section:
$$\left(\begin{array}{c}\widehat{\mu }\\ \widehat{\nu }\end{array}\right)=\left(\begin{array}{c}\stackrel{~}{\mu }\\ \stackrel{~}{\nu }\end{array}\right)+\left(\begin{array}{c}\stackrel{~}{Y}_1\\ \stackrel{~}{Z}_1\end{array}\right)+\left(\begin{array}{c}\stackrel{~}{Y}_2\\ \stackrel{~}{Z}_2\end{array}\right)A_{1,k}\left(\begin{array}{c}\stackrel{~}{Y}_1\\ \stackrel{~}{Z}_1\end{array}\right)A_{1,k}\left(\begin{array}{c}\stackrel{~}{Y}_2\\ \stackrel{~}{Z}_2\end{array}\right).$$
(63)
The operator $`A_{n,k}`$ was defined above. Let us substitute (63) into (9) and evaluate a remainder. As a result we obtain
###### Theorem 4
The formal asymptotic solution of the problem (9) with respect to $`mod\left(O\left((t|\theta |)^1\right)\right)`$, which is uniform when $`k\text{}`$ and $`\theta ^2t^{2/3}1`$, has the form (63).
### 4.2 Asymptotics in neighborhood of confluent stationary point
The system of equations in the problem (9) depends on two control parameters $`\xi `$ and $`\eta `$. On the parabola $`12\xi +\eta ^2=0`$ the degeneracy of stationary points occurs: $`k_1=k_2=k_0=i\eta `$. For this reason the asymptotics constructed in section 4.1 is invalid when parameter $`\theta =\sqrt{\eta ^2+12\xi }`$ is close to zero. For example, the asymptotics of $`\stackrel{1}{\mu }`$ as $`kk_0`$ and $`\theta 0`$ is discontinuous:
$`\left[\stackrel{1}{\mu }|_{kk_0}\right]|_{\theta 0}=({\displaystyle \frac{|kk_0|}{12i((kk_0)^2\theta ^2)}}{\displaystyle \frac{\overline{\theta }}{\theta }}{\displaystyle \frac{kk_0}{12i((kk_0)^2\theta ^2)}})\times `$
$`f(\overline{k}_0)f(k_0)+o(1).`$
This shows that we need new scaling of the parameter $`\theta `$. The scaled control parameter is:
$$v=t^{1/3}\frac{\theta }{\sqrt{12}}.$$
The external expansion constructed in sec. 4.1 becomes discontinuous at $`\theta =0`$ and the internal expansions constructed in sec. 4.1 become singular at the point $`\theta =0`$ and lose their asymptotic properties. Therefore here we must change the internal variable for the internal asymptotic expansion:
$$p=t^{1/3}(kk_0).$$
(64)
In this section we construct a formal asymptotic solution of the problem (9) with respect to $`mod(O(t^1))`$ when $`|\theta |1`$ uniform with respect to $`k\text{}`$. The result is formulated in the end of this section.
To construct the uniform asymptotic solution we need the external and internal asymptotics outside and inside a small neighborhood of the point $`k_0`$ respectively.
###### Lemma 4
Let the system of the equations (9)be have no the homogeneous nontrivial solutions, $`F(k)C^2L_1`$ and the parameters $`\xi `$ and $`\eta `$ satisfy the inequality $`|12\xi +\eta ^2|1`$, then:
when $`|kk_0|t^{1/3}1`$ the formal asymptotic solution of the system (9) with respect to $`mod(O(t^{2/3}/|kk_0|)+O(t^1))`$ has the form:
$$\stackrel{~}{\mu }=1+t^1\stackrel{1}{m}(k,\xi ,\eta ),$$
$$\stackrel{~}{\nu }=t^{2/3}\underset{0}{\overset{1}{n}}+t^1(\underset{1}{\overset{1}{n}}\mathrm{exp}(itS)+\underset{0}{\overset{2}{n}})$$
the function $`\stackrel{1}{m}`$ is defined by (90), the functions $`\underset{0}{\overset{1}{n}}`$ and $`\underset{2}{\overset{2}{n}}`$ are defined by (86) and (87) respectively, the function $`\underset{1}{\overset{1}{n}}`$ is defined by (77);
when $`|kk_0|1`$ the asymptotic solution of the system (9)with respect to $`mod(O(t^{2/3}|kk_0|+O(t^1))`$ has the form:
$$\stackrel{~}{\mu }=1+t^{2/3}\stackrel{1}{}+t^1\stackrel{2}{},$$
$$\stackrel{~}{\nu }=(t^{1/3}\stackrel{1}{𝒩}+t^{2/3}\stackrel{2}{𝒩}+t^1\stackrel{3}{𝒩})\mathrm{exp}(itS);$$
the functions $`\stackrel{j}{}`$, $`j=1,2`$ and $`\stackrel{j}{𝒩}`$, $`j=1,2,3`$ are defined by (89), and (88) respectively.
Proof. Let us find the internal formal asymptotic expansion for the solution of the system (9) in the form:
$$^{in}=1+t^{2/3}\stackrel{1}{}+t^1\stackrel{2}{}+\mathrm{},$$
(65)
$$𝒩^{in}=(t^{1/3}\stackrel{1}{𝒩}+t^{2/3}\stackrel{2}{𝒩}+t^1\stackrel{3}{𝒩}+\mathrm{})\mathrm{exp}(itS).$$
(66)
Change the variable $`k`$ into the variable $`p`$ in the system (9). The phase function of the exponent is depend on the variable $`p`$ as:
$$tS\omega (p)4(p^3+\overline{p}^3)v^2(p+\overline{p}),$$
where $`v=t^{1/3}\theta /\sqrt{12}`$. Substitute (65) and (66) into the system (9), equate the coefficients with equal power of $`t`$. As a result we obtain a sequence of equation for the coefficients of the expansions (65) and (66).
$`_p\stackrel{1}{𝒩}i(12p^2v^2)\stackrel{1}{𝒩}=\text{sgn}[\text{Re}(\overline{p})]f(\overline{k}_0),`$ (67)
$`_p\stackrel{2}{𝒩}i(12p^2v^2)\stackrel{2}{𝒩}=\text{sgn}[\text{Re}(\overline{p})](f_{01}^{(1)}(\overline{k}_0)\overline{p}+`$
$`+f_{10}^{(1)}(\overline{k})p)f^{(2)}_{01}(\overline{k}_0)\overline{p}f^{(2)}_{10}(\overline{k}_0)p.`$ (68)
$`_{\overline{p}}\stackrel{1}{}=\text{sgn}[\text{Re}(\overline{p})]f(k_0)\stackrel{1}{𝒩},`$ (69)
$`_p\stackrel{3}{𝒩}i(12p^2v^2)\stackrel{3}{𝒩}=`$
$`{\displaystyle \frac{1}{2}}\left[\text{sgn}\left[\text{Re}(\overline{p})\right]f_{20}^{(1)}(\overline{k}_0)+f_{20}^{(2)}(\overline{k}_0)\right]p^2`$
$`{\displaystyle \frac{1}{2}}\left[\text{sgn}\left[\text{Re}(\overline{p})\right]f_{02}^{(1)}(\overline{k}_0)+f_{02}^{(2)}(\overline{k}_0)\right]\overline{p}^2`$
$`\left[\text{sgn}\left[\text{Re}(\overline{p})\right]f_{11}^{(1)}(\overline{k}_0)+f_{11}^{(2)}(\overline{k}_0)\right]|p|^2`$
$`\text{sgn}\left[\text{Re}(\overline{p})\right]f(\overline{k}_0)\stackrel{1}{}.`$ (70)
$`_{\overline{p}}\stackrel{2}{}=\text{sgn}\left[\text{Re}(p)\right]f(k_0)\stackrel{2}{𝒩}+`$
$`\left[\text{sgn}\left[\text{Re}(\overline{p})\right]f_{10}^{(1)}(k_0)+f_{10}^{(2)}(k_0)\right]\stackrel{1}{𝒩}p+`$
$`+\left[\text{sgn}[\text{Re}(k)]f_{01}^{(1)}(k_0)+f_{01}^{(2)}(k_0)\right]\stackrel{1}{𝒩}\overline{p}.`$ (71)
These equations are obtained in a supposition that $`|p|t^{1/3}1`$. The uniform bounded partial solutions of the equations for the coefficients of the expansion (66) are obtained by using an integral operator:
$$𝒫[g]=\frac{\mathrm{exp}(i\omega (p))}{2i\pi }_{\text{}}\frac{drd\overline{r}}{\overline{pr}}\mathrm{exp}(i\omega (r))g(r),$$
(72)
where $`g(r)`$ is the right-hand side of corresponding equation.
Thus, the bounded partial solutions of the equations (67) and (68) (the functions $`\stackrel{s}{\stackrel{1}{𝒩}}`$ and $`\stackrel{s}{\stackrel{2}{𝒩}}`$) are represented by double integrals.
Bounded partial solutions of equations for the coefficients of the expansion (65) are constructed by using the Cauchy-Green formula.
The formula for the partial solution of the equation (69) has the form:
$$\stackrel{s}{\stackrel{1}{}}(p,\xi ,\eta )=f(k_0)f(\overline{k}_0)\left(J_1(p,v^2)2J_1^+2J_1^+\right).$$
(73)
Here
$$J_1=\frac{1}{2i\pi }_{\text{}}\frac{dnd\overline{n}}{pn}\frac{\mathrm{exp}(i\omega (n))}{2i\pi }_{\text{}}\frac{drd\overline{r}}{\overline{nr}}\mathrm{exp}(i\omega (r)).$$
$$J_1^\pm =\frac{1}{2i\pi }_\mathrm{\Omega }^{}\frac{dnd\overline{n}}{pn}\frac{\mathrm{exp}(i\omega (n))}{2i\pi }_{\mathrm{\Omega }^\pm }\frac{drd\overline{r}}{\overline{nr}}\mathrm{exp}(i\omega (r)).$$
In the Appendix the integrals $`J_1`$ and $`J_1^+`$ are reduced into double integrals. By the similar way we may represent the integral $`J_1^+`$ as the sum of double integrals. Therefore we may represent $`\stackrel{s}{\stackrel{1}{}}`$ as sum of double integrals also. Corresponded formula is very large and we don’t write it here. But we will use this formula to evaluate an asymptotic behaviour of $`\stackrel{s}{\stackrel{1}{}}`$ as $`|p|\mathrm{}`$.
The bounded partial solution of the equation (71) when $`p\text{}`$ has to be built similarly . We give following statement about $`\stackrel{2}{}`$.
###### Lemma 5
Continuous partial solution of the equation for $`\stackrel{2}{}`$ exists. This solution is uniformly bounded when $`p\text{}`$.
Sketch of proof. The partial solution of the equation for $`\stackrel{2}{}`$ has to be obtained as a result of using of the operator:
$$_{\overline{p}}^1\left[g\right]=\stackrel{s}{\stackrel{2}{}}\frac{1}{2i\pi }_{\text{}}\frac{dnd\overline{n}}{pn}g(n)$$
(74)
to the right-hand side of the equation. The integrals, which are obtained, are continuous by virtue of continuity of the integral operator with respect to the parameter $`p`$. The boundedness of these integrals with respect to $`p`$ when $`p\text{}`$ may be obtained by using the asymptotic behaviour of the right-hand side terms of the equation for $`\stackrel{2}{}`$ as $`|p|\mathrm{}`$.
Let us to construct the external expansion in the form:
$$m=1+t^1\underset{0}{\overset{1}{m}}(k,v)+t^{5/3}(\underset{1}{\overset{2}{m}}\mathrm{exp}(itS)+\underset{0}{\overset{2}{m}})+\mathrm{};$$
(75)
$`n`$ $`=(t^1\underset{1}{\overset{1}{n}}(k,v)+t^{5/3}\underset{1}{\overset{2}{n}}(k,v)+t^2\underset{1}{\overset{3}{n}}\mathrm{})\mathrm{exp}(itS)`$ (76)
$`+t^{2/3}\underset{0}{\overset{1}{n}}(k,v)+t^1\underset{0}{\overset{2}{n}}(k,v)+\mathrm{}.`$
Substitute (75) and (76) into (9). As a result we get the equations for coefficients of the expansions:
$$\underset{1}{\overset{1}{n}}=\text{sgn}\left[\text{Re}(\overline{k})\right]\frac{f(\overline{k})}{12i(kk_0)^2};\underset{1}{\overset{2}{n}}=\text{sgn}\left[\text{Re}(\overline{k})\right]\frac{v^2f(\overline{k})}{144i(kk_0)^4};$$
(77)
$$_{\overline{k}}\stackrel{1}{m}=\frac{f(k)f(\overline{k})}{12i(kk_0)^2};$$
(78)
$$_k\underset{0}{\overset{1}{n}}(k,\xi ,\eta )=0;_k\underset{0}{\overset{2}{n}}(k,\xi ,\eta )=0;$$
(79)
$`_{\overline{k}}\underset{0}{\overset{2}{m}}={\displaystyle \frac{f(k)f(\overline{k})v^2}{144i(kk_0)^4}};`$ (80)
$`\underset{1}{\overset{2}{m}}(k,\xi ,\eta )=\text{sgn}\left[\text{Re}(\overline{k})\right]f(k){\displaystyle \frac{\underset{0}{\overset{1}{n}}(k,\xi ,\eta )}{12i\overline{(}kk_))^2}};`$ (81)
$`\underset{1}{\overset{3}{n}}(k,\xi ,\eta )={\displaystyle \frac{1}{12i(kk_0)^2}}(\text{sgn}\left[\text{Re}(\overline{k})\right]f(\overline{k})\underset{0}{\overset{1}{m}}`$
$`+{\displaystyle \frac{2f(\overline{k})}{12i(kk_0)^3}}\text{sgn}[\text{Re}(\overline{k})]{\displaystyle \frac{f_{10}^{(1)}(\overline{k})\text{sgn}[\text{Re}(\overline{k})]+f_{10}^{(2)}(\overline{k})}{12i(kk_0)^2}}).`$ (82)
The formulas (77) and (82) define the functions $`\underset{1}{\overset{1}{n}}`$, $`\underset{1}{\overset{2}{n}}`$ and $`\underset{1}{\overset{3}{n}}`$. Using the formulas (79) we can see, that the functions $`\underset{0}{\overset{1}{n}}`$ $`\underset{0}{\overset{2}{n}}`$ are analytic of variable $`\overline{k}`$. The obvious form of this dependence is defined by two conditions. First one is the continuity of the asymptotics and second one is matching condition for the external and internal asymptotic expansions.
One can see, that the sufficient condition for the continuity of the coefficient of the asymptotics (76) as $`t^1`$ with respect to $`k`$ is an addition into $`\underset{0}{\overset{2}{n}}`$ of the term:
$$\underset{0}{\overset{2}{n^{}}}=\frac{1}{\pi i}\frac{1}{12i\overline{(kk_0)}^2}_i\mathrm{}^i\mathrm{}\frac{d\lambda }{\overline{k\lambda }}f(\overline{l}).$$
(83)
It follows from the formulas (81), that the coefficient of the asymptotics (75) as $`t^{5/3}`$ is defined after evaluating of $`\underset{0}{\overset{1}{n}}`$, i.e. after matching of the coefficients as $`t^{2/3}`$ of the external and internal expansions.
Consider the problem for defining of $`\stackrel{1}{m}`$ in detail. This function satisfies the boundary condition:
$$\underset{0}{\overset{1}{m}}|_{|k|\mathrm{}}=0.$$
The external solution is valid outside of small neighborhoods of the points $`k_j`$, $`j=0,1,2.`$ Therefore solutions of the equation (78), which are smooth, bounded and decreased as $`|k|\mathrm{}`$:
$`\underset{0}{\overset{s}{\stackrel{1}{m}}}(k,\theta )={\displaystyle \frac{1}{2i\pi }}{\displaystyle _{\text{}}\frac{drd\overline{r}}{kr}\frac{f(\overline{r})f(r)f(\overline{k}_0)f(k_0)}{12i(rk_0)^2}}+`$
$`+{\displaystyle \frac{\overline{(kk_0)}f(k_0)f(\overline{k}_0)}{12i(kk_0)^2}}`$ (84)
are defined to within an analytic function with respect to variable $`k`$, which has poles when $`t^{1/3}|kk_0|1`$. The full definition of $`\underset{0}{\overset{1}{m}}`$ will be done by matching of the external and internal expansions.
For matching process we need asymptotic behaviour of the partial solutions of equations (67)-(70) as $`|p|\mathrm{}`$. Constructing of these asymptotics is reduced to evaluating of an integrals with weak singularity in the integrand. Evaluating of these integrals is done in the Appendix. Here we write the asymptotic behaviour of partial solutions of equations (67)-(70) as $`|p|\mathrm{}`$.
$`\stackrel{s}{\stackrel{1}{𝒩}}(p,\xi ,\eta )={\displaystyle \frac{1}{\overline{p}}}\mathrm{exp}(i\omega (p))f(\overline{k}_0)\left[\varphi _{00}^+(v^2)\varphi _{00}^{}(v^2)\right]+`$
$`+{\displaystyle \frac{1}{\overline{p}^2}}\mathrm{exp}(i\omega (p))f(\overline{k})\left[\varphi _{01}^+(v^2)\varphi _{01}^{}(v^2)\right]+`$
$`+{\displaystyle \frac{1}{12i}}f(\overline{k}_0)\text{sgn}\left[\text{Re}(\overline{p})\right]\left({\displaystyle \frac{1}{p^2}}{\displaystyle \frac{\mathrm{exp}(i\omega (p))}{\overline{p}^2}}\right)+O(|v^2||p|^3).`$
$`\stackrel{s}{\stackrel{2}{𝒩}}(p,\xi ,\eta )=[{\displaystyle \frac{\mathrm{exp}(i\omega (p))}{\overline{p}}}(\varphi _{10}^+(v^2)\varphi _{10}^{}(v^2))+`$
$`({\displaystyle \frac{1}{p}}{\displaystyle \frac{\mathrm{exp}(i\omega (p))}{\overline{p}}}){\displaystyle \frac{\text{sgn}[\text{Re}(\overline{p})]}{12i}}]f_{10}^{(1)}(\overline{k}_0)+`$
$`\left[{\displaystyle \frac{1}{\overline{p}}}\mathrm{exp}(i\omega (p))\varphi _{10}(v^2)+{\displaystyle \frac{1}{12ip}}\right]f_{10}^{(2)}(\overline{k}_0)+`$
$`[{\displaystyle \frac{\mathrm{exp}(i\omega (p))}{\overline{p}}}(\varphi _{01}^+(v^2)\varphi _{01}^{}(v^2))+({\displaystyle \frac{\overline{p}}{p^2}}{\displaystyle \frac{\mathrm{exp}(\omega (p))}{\overline{p}}})\times `$
$`{\displaystyle \frac{\text{sgn}\left[\text{Re}(\overline{p})\right]}{12i}}]f_{01}^{(1)}(\overline{k}_0)+[{\displaystyle \frac{1}{\overline{p}}}\mathrm{exp}(i\omega (p))\varphi _{01}(v^2)+`$
$`{\displaystyle \frac{\overline{p}}{12ip^2}}]f^{(2)}_{01}(\overline{k}_0)+O(|p|^2)+O(|v|^2|p|^2),`$
$`\stackrel{s}{\stackrel{1}{}}(p,\xi ,\eta )=f(k_0)f(\overline{k}_0)[{\displaystyle \frac{\overline{p}}{12ip^2}}+`$
$`{\displaystyle \frac{1}{p}}\left(12i\varphi _{01}(v^2)\overline{\varphi _{10}(v^2)}+iv^2\varphi _{00}(v^2)\overline{\varphi _{00}(v^2)}\right)`$
$`{\displaystyle \frac{2}{p}}[iv^2(\varphi _{00}^+(v^2)\psi _{00}^{}(v^2)+\varphi _{00}^{}(v^2)\psi _{00}^+(v^2))+`$
$`+12i(\varphi _{01}^+(v^2)\psi _{01}^{}(v^2)+\varphi _{01}^{}(v^2)\psi _{01}^+(v^2)){\displaystyle \frac{1}{2}}\psi _{01}]]+`$
$`O(|v|^2|p|^2);`$
$`\stackrel{s}{\stackrel{2}{}}(p,\xi ,\eta )=f(k_0)f_{10}^{(1)}(\overline{k}_0){\displaystyle \frac{\overline{p}}{12ip}}+`$
$`f(k_0)f_{10}^{(2)}(\overline{k}_0)\left({\displaystyle \frac{\overline{p}}{12ip}}+{\displaystyle \frac{1}{12i}}\right)\text{sgn}[\text{Re}(p)]+`$
$`f(k_0)f_{01}^{(2)}(\overline{k}_0)\left({\displaystyle \frac{\overline{p}^2}{24ip^2}}{\displaystyle \frac{1}{24i}}\right)\text{sgn}[\text{Re}(p)]+`$
$`f(k_0)f_{01}^{(1)}(\overline{k}_0){\displaystyle \frac{\overline{p}^2}{24ip^2}}+\mathrm{\Phi }(v^2)+O(|v|^2|p|^1),`$
$`\stackrel{s}{\stackrel{3}{𝒩}}(p,\xi ,\eta )={\displaystyle \frac{1}{24i}}\left[\text{sgn}[\text{Re}(\overline{p})]\left(1\mathrm{exp}(i\omega (p))\right)f_{20}^{(1)}(\overline{k}_0)\right]+`$
$`\left[\text{sgn}[\text{Re}(\overline{p})]\left({\displaystyle \frac{\overline{p}^2}{12ip^2}}{\displaystyle \frac{\mathrm{exp}(i\omega (p))}{24i}}\right)f_{02}^{(1)}(\overline{k}_0)+{\displaystyle \frac{\overline{p}^2f_{02}^{(2)}(\overline{k}_0)}{24ip^2}}\right]+`$
$`\left[\text{sgn}[\text{Re}(\overline{p})]\left({\displaystyle \frac{\overline{p}}{12ip}}+{\displaystyle \frac{\mathrm{exp}(i\omega (p))}{12i}}\right)f_{11}^{(1)}(\overline{k}_0)+{\displaystyle \frac{\overline{p}}{12ip}}f_{11}^{(2)}(\overline{k}_0)\right]+`$
$`f_0^2(\overline{k}_0)f(k_0)\mathrm{exp}(i\omega (p))\mathrm{\Psi }(v^2)+O(|v||p|^1).`$
The functions $`\mathrm{\Phi }(v^2)`$ and $`\mathrm{\Psi }(v^2)`$ are smooth and uniformly bounded when $`v^2\text{}`$. Here we use notations:
$$\varphi _{mn}^\pm =_{\text{Re}(r)>0}𝑑rd\overline{r}r^m\overline{r}^n\mathrm{exp}(i\omega (r)),$$
$$\psi _{mn}^\pm =_{\text{Re}(r)>0}𝑑rd\overline{r}r^m\overline{r}^n\mathrm{exp}(i\omega (r)),$$
$$\varphi _{mn}=_{\text{}}𝑑rd\overline{r}r^m\overline{r}^n\mathrm{exp}(i\omega (r)),$$
$$\psi _{mn}=_{\text{}}𝑑rd\overline{r}r^m\overline{r}^n\mathrm{exp}(i\omega (r)).$$
(85)
Evaluate an asymptotic behaviours of the coefficients of the external expansion as $`kk_0`$ and $`_k^2S=o(1)`$.
$`\underset{1}{\overset{1}{n}}|_{kk_0}={\displaystyle \frac{f(\overline{k}_0)}{12i(kk_0)^2}}\text{sgn}[\text{Re}(\overline{k})]+`$
$`{\displaystyle \frac{1}{12i(kk_0)}}[f_{10}^{(1)}(\overline{k})+f_{01}^{(1)}(\overline{k}_0){\displaystyle \frac{\overline{kk_0}}{kk_0}}+`$
$`(f_{10}^{(2)}(\overline{k}_0)+f_{01}^{(2)}(\overline{k}_0){\displaystyle \frac{\overline{kk_0}}{kk_0}})\text{sgn}[\text{Re}(\overline{k})]]+`$
$`+{\displaystyle \frac{1}{2}}f_{20}^{(1)}(\overline{k})+f_{11}^{(1)}(\overline{k}_0){\displaystyle \frac{\overline{kk_0}}{kk_0}}+{\displaystyle \frac{1}{2}}f_{02}^{(1)}{\displaystyle \frac{\overline{(kk_0)^2}}{(kk_0)^2}}+`$
$`({\displaystyle \frac{1}{2}}f_{20}^{(2)}(\overline{k}_0)+f_{11}^{(2)}(\overline{k}_0){\displaystyle \frac{\overline{kk_0}}{kk_0}}+`$
$`{\displaystyle \frac{1}{2}}f_{02}^{(2)}(\overline{k}_0){\displaystyle \frac{\overline{(kk_0)^2}}{(kk_0)^2}})\text{sgn}[\text{Re}(\overline{k})]+o(1).`$
$`\underset{0}{\overset{}{\stackrel{2}{n}}}={\displaystyle \frac{1}{24\pi \overline{(kk_0)^2}}}\times `$
$`[\pi if(\overline{k}_0)\text{sgn}[\text{Re}(\overline{p})]+V.P.{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{d\lambda f(i\lambda )}{\lambda ik_0}}]+`$
$`{\displaystyle \frac{1}{12\pi \overline{(kk_0)}}}[\pi if_{01}^{(1)}(\overline{k}_0)+\pi if_{10}^{(1)}(\overline{k}_0)+(\pi if_{10}(2)(\overline{k}_0)+`$
$`+\pi if_{01}^{(2)}(\overline{k}_0))\text{sgn}[\text{Re}(\overline{k}_0)]+`$
$`+V.P.{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{d\lambda f(i\lambda )}{(\lambda ik_0)^2}}]+{\displaystyle \frac{1}{12i}}[\pi if_{20}^{(1)}(\overline{k}_0)+`$
$`+2\pi if_{11}^{(1)}(\overline{k}_0)+\pi if_{02}^{(1)}(\overline{k}_0)+(\pi if_{20}^{(2)}(\overline{k}_0)+2\pi if_{11}^{(2)}(\overline{k}_0)+`$
$`+\pi if_{02}^{(2)}(\overline{k}_0))\text{sgn}[\text{Re}(\overline{k})]+V.P.{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{d\lambda f(i\lambda )}{(\lambda ik_0)^3}}]+o(1).`$
$`\stackrel{s}{\stackrel{1}{m}}={\displaystyle \frac{\overline{(kk_0)}}{12i(kk_0)^2}}f(\overline{k}_0)f(k_0)`$
$`{\displaystyle \frac{\overline{kk_0}}{12i(kk_0)}}[(f_{10}^{(1)}(\overline{k}_0)\text{sgn}[\text{Re}(\overline{k})]+`$
$`+f_{10}^{(2)}(\overline{k}_0))f(k_0)+f(\overline{k}_0)(f_{10}^{(1)}(k_0)\text{sgn}[\text{Re}(k)]+f_{10}^{(2)}(k_0))]`$
$`{\displaystyle \frac{\overline{(kk_0)^2}}{12i(kk_0)^2}}[f(k_0)(f_{01}^{(1)}(\overline{k}_0)\text{sgn}[\text{Re}(\overline{k})]+f_{01}^{(2)}(\overline{k}_0))+`$
$`+f(\overline{k}_0)(f_{01}^{(1)}(k_0)\text{sgn}[\text{Re}(k)]+f_{01}^{(2)}(k_0))]+o(1).`$
Let us do matching of the external and internal asymptotic expansions of the function $`\stackrel{~}{\nu }`$. The matching condition for $`\stackrel{~}{\nu }`$ in domain $`t^{1/3}|kk_0|1`$ when $`|\theta |1`$ has the form:
$`\left(t^{2/3}\underset{0}{\overset{1}{n}}+t^1\underset{0}{\overset{2}{n}}+t^1\stackrel{1}{n_1}\mathrm{exp}(itS)\right)`$
$`(t^{1/3}\stackrel{1}{𝒩}+t^{2/3}\stackrel{2}{𝒩}+t^1\stackrel{3}{𝒩})=o(t^1).`$
Let us equate the coefficients with equal powers of the large parameter $`t`$. As a result we obtain:
$$\underset{0}{\overset{1}{n}}=\frac{1}{\overline{kk_0}}f(\overline{k}_0)[\varphi _{00}^+(v^2)\varphi _{00}^{}(v^2)],$$
(86)
$`\underset{0}{\overset{2}{n}}=`$ $`{\displaystyle \frac{1}{\overline{kk_0}^2}}f(\overline{k})[\varphi _{01}^+(v^2)\varphi _{01}^{}(v^2)]+`$ (87)
$`{\displaystyle \frac{1}{\overline{kk_0}}}f_{10}^{(1)}(\overline{k}_0)[\varphi _{10}^+(v^2)\varphi _{10}^{}(v^2)]+{\displaystyle \frac{1}{\overline{kk_0}}}f_{10}^{(2)}(\overline{k}_0)\varphi _{10}(v^2)+`$
$`{\displaystyle \frac{1}{\overline{kk_0}}}f_{01}^{(1)}(\overline{k}_0)[\varphi _{01}^+(v^2)\varphi _{01}^{}(v^2)]+`$
$`+{\displaystyle \frac{1}{\overline{kk_0}}}f_{01}^{(2)}(\overline{k}_0)\varphi _{01}(v^2),`$
where $`\varphi `$, $`\psi `$, $`\varphi ^\pm `$ $`\psi ^\pm `$ are defined by (85);
$`\stackrel{1}{𝒩}=𝒫[g_1],\stackrel{2}{𝒩}=𝒫[g_2],`$
$`\stackrel{3}{𝒩}=𝒫[g_3]f_0^2(\overline{k}_0)f(k_0)\mathrm{exp}(i\omega (p))\mathrm{\Psi }(v^2),`$ (88)
where $`g_j`$ are the right-hand sides of the equations for $`\stackrel{j}{𝒩}`$, the operator $`𝒫[g]`$ is defined by (72), the function $`\mathrm{\Psi }(v^2)`$ is smooth and uniformly bounded with respect to $`v^2`$ when $`v^2\text{}`$.
Matching condition for the function $`\stackrel{~}{\mu }`$ in the domain $`t^{1/3}|kk_0|1`$ when $`|\theta |1`$ has the form:
$$(1+t^1\stackrel{1}{m})(1+t^{2/3}\stackrel{1}{}+t^1\stackrel{2}{})=o(t^1).$$
As a result we obtain:
$$\stackrel{1}{}=\stackrel{s}{\stackrel{1}{}};\stackrel{2}{}=\stackrel{s}{\stackrel{2}{}}\mathrm{\Phi }(v^2),$$
(89)
where $`\stackrel{s}{\stackrel{1}{}}`$ is defined by (73), the function $`\stackrel{s}{\stackrel{2}{}}`$ is defined by (74) and the function $`\mathrm{\Phi }(v^2)`$ is smooth and uniformly bounded when $`v^2\text{}`$.
$`\stackrel{1}{m}=\stackrel{s}{\stackrel{1}{m}}{\displaystyle \frac{1}{kk_0}}(12i\varphi _{01}(v^2)\overline{\varphi _{10}(v^2)}+iv^2\varphi _{00}(v^2)\overline{\varphi _{00}(v^2)})`$
$`{\displaystyle \frac{2}{kk_0}}[iv^2(\varphi _{00}^+(v^2)\psi _{00}^{}(v^2)+\varphi _{00}^{}(v^2)\psi _{00}^+(v^2))+`$
$`+12i(\varphi _{01}^+(v^2)\psi _{01}^{}(v^2)+\varphi _{01}^{}(v^2)\psi _{01}^+(v^2)){\displaystyle \frac{1}{2}}\psi _{01}].`$ (90)
Here $`\stackrel{s}{\stackrel{1}{m}}`$ is defined by (84) and the functions $`\psi ^\pm `$, $`\varphi ^\pm `$, $`\varphi `$ and $`\psi `$ are defined by formulas (85).
Thus we have matched internal and external expansions of $`\stackrel{~}{\mu }`$ and $`\stackrel{~}{\nu }`$. The lemma is proved.
These expansions are ununiform with respect to $`k`$. Now we can construct uniform asymptotic expansion with respect to $`k\text{}`$. Following the matching method the uniform expansion is:
$$\left(\begin{array}{c}\widehat{\mu }\\ \widehat{\nu }\end{array}\right)=\left(\begin{array}{c}\stackrel{~}{m}\\ \stackrel{~}{n}\end{array}\right)+\left(\begin{array}{c}\stackrel{~}{}_1\\ \stackrel{~}{𝒩}_1\end{array}\right)A_{1,k}\left(\begin{array}{c}\stackrel{~}{}_1\\ \stackrel{~}{𝒩}_1\end{array}\right).$$
(91)
Here the operator $`A_{n,k}`$ processes on the function $`\stackrel{~}{}_1`$ in the formula for $`_1`$ by followed manner. One must change the variable $`p`$ into the variable $`k`$ using the formula (64) and write all terms of the asymptotic expansion with respect to $`t`$ with the powers are equal to $`m`$, where $`0mn`$. For example, one can obtain for the function $`\stackrel{~}{}(p,\xi ,\eta ,t)`$:
$`A_{1,k}[\stackrel{~}{}(p,\xi ,\eta ,t)]=1+t^1(f(k_0)f(\overline{k}_0)[{\displaystyle \frac{\overline{kk_0}}{12i(kk_0)^2}}+`$
$`+{\displaystyle \frac{1}{kk_0}}\left(12i\varphi _{01}(v^2)\overline{\varphi _{10}(v^2)}+iv^2\varphi _{00}(v^2)\overline{\varphi _{00}(v^2)}\right)`$
$`{\displaystyle \frac{2}{kk_0}}[iv^2(\varphi _{00}^+(v^2)\psi _{00}^{}(v^2)+\varphi _{00}^{}(v^2)\psi _{00}^+(v^2))+`$
$`+12i(\varphi _{01}^+(v^2)\psi _{01}^{}(v^2)+\varphi _{01}^{}(v^2)\psi _{01}^+(v^2)){\displaystyle \frac{1}{2}}\psi _{01}]]`$
$`f(k_0)f_{10}^{(1)}(\overline{k}_0){\displaystyle \frac{\overline{kk_0}}{12i(kk_0)}}+`$
$`f(k_0)f_{10}^{(2)}(\overline{k}_0)\left({\displaystyle \frac{\overline{kk_0}}{12i(kk_0)}}+{\displaystyle \frac{1}{12i}}\right)\text{sgn}[\text{Re}(k)]+`$
$`f(k_0)f_{01}^{(1)}(\overline{k}_0){\displaystyle \frac{\overline{kk_0}^2}{24i(kk_0)^2}}+`$
$`+f(k_0)f_{01}^{(2)}(\overline{k}_0)({\displaystyle \frac{\overline{kk_0}^2}{24i(kk_0)^2}}{\displaystyle \frac{1}{24i}})\text{sgn}[\text{Re}(k)]),`$
The obviously formula for $`A_{1,k}[\stackrel{~}{𝒩}]`$ is more large and doesn’t shown here.
###### Theorem 5
The formula (91) gives the asymptotics solution of the problem (9) with respect to $`mod(O(t^1))`$ as $`t\mathrm{}`$. This asymptotic solution is uniform with respect to $`k\text{}`$ and $`|\theta |1`$.
## 5 Justification of the asymptotics of the solution of $`\overline{D}`$-problem
In this section we prove that the remainder of the asymptotics has order by $`t^{4/3}`$ uniformly with respect to $`k\text{}`$ and this remainder has to be differentiable with respect to $`x`$. We call by the remainder of the asymptotics the difference between solution of the problem (9) and constructed asymptotic solutions (47) when $`\theta ^2t^{2/3}1`$, (63) when $`\theta ^2t^{2/3}1`$ and (91) when $`|\theta |1`$. The differentiability of the remainder will be important when we will construct an asymptotic behaviour of solution of the equation KP-2.
###### Theorem 6
Let $`^\alpha f(k,\overline{k})L_1C^1`$ when $`|\alpha |2`$ when $`k`$ is out of the imaginary axis and
$$\underset{z\text{}}{sup}\left|_{\text{}}\frac{dkd\overline{k}}{|kz|}|F(k)|\right|<\mathrm{\hspace{0.17em}2}\pi ,$$
then the solution of the problem (9) is:
$$\left(\begin{array}{c}\mu \\ \nu \end{array}\right)=\left(\begin{array}{c}\stackrel{~}{\mu }\\ \stackrel{~}{\nu }\end{array}\right)+O(t^{{}_{4}{}^{}/3}),$$
(92)
when $`k\text{},\xi ,\eta \text{}`$. The remainder of the asymptotics has to be differentiable with respect to $`x`$.
The proof. Let us write the system of differential equation for the remainder. Denote the remainder in (97) by $`V`$. Substitute (97) into (9). As a result we obtain:
$$\left(\begin{array}{cc}_{\overline{k}}& 0\\ 0& _k\end{array}\right)V=\left(\begin{array}{cc}0& F(\overline{k})\mathrm{exp}(itS)\\ F(k)\mathrm{exp}(itS)& 0\end{array}\right)V+f,$$
(93)
$$V|_{|k|\mathrm{}}=0.$$
(94)
We denote by the vector $`f`$ the residual which originates in (9) when we substitute the column into this equation $`(\stackrel{~}{\mu },\stackrel{~}{\nu })^T`$:
$$f_1=_{\overline{k}}\stackrel{~}{\mu }+F(\overline{k})\mathrm{exp}(itS)\stackrel{~}{\nu },$$
$$f_2=_k\stackrel{~}{\nu }+F(k)\mathrm{exp}(itS)\stackrel{~}{\mu }.$$
Let us denote by $`X`$ the space of bounded and continuous with respect to $`k`$ vector-functions with the norm:
$$W=\underset{k\text{},(\xi ,\eta )\text{}^2}{sup}|W_1|+\underset{k\text{},(\xi ,\eta )\text{}^2}{sup}|W_2|.$$
Consider a system of integral equations instead of the problem (93), (94):
$$V=G[F]V+H,$$
(95)
where $`G[F]`$ is the integral operator:
$`G[F]V={\displaystyle }{\displaystyle _m\text{}}dmd\overline{m}\times `$
$`\left(\begin{array}{cc}0& \frac{F(\overline{m})}{km}\mathrm{exp}(itS)\\ \frac{F(m)}{\overline{km}}\mathrm{exp}(itS)& 0\end{array}\right)V(m,\xi ,\eta ,t);`$
$$H=_m\text{}𝑑md\overline{m}\left(\begin{array}{c}\frac{f_1(m,\xi ,\eta ,t)}{km}\\ \frac{f_2(m,\xi ,\eta ,t)}{\overline{km}}\end{array}\right).$$
Using obvious forms of the functions $`f_{1,2}`$ one can proof, that
$$H=O(t^{4/3})$$
uniformly with respect to $`\xi ,\eta \text{}`$ when $`k\text{}`$.
The operator $`G[F]`$ is contracting in the space $`X`$, therefore the solution of the integral equation (95) exists in $`X`$ and is evaluated by $`O(t^{4/3})`$ uniformly with respect to $`\xi ,\eta \text{}`$.
Show that the remainder of the asymptotics is differentiable on $`x=t\xi `$. Differentiate with respect to $`x`$ the system of the equation for the remainder. Denote the derivative of the vector $`V`$ by $`\chi `$. Then we obtain:
$$\chi =G[F]\chi +_xG[F]V+_xH.$$
The terms $`_xG[F]V+_xH`$ may be evaluated by order $`O(t^{4/3})`$. The operator $`G[F]`$ is contracting, therefore one can obtain: $`\chi =O(t^{4/3})`$. The theorem is proved.
## 6 Solution of the equation KP-2
An asymptotics of the solution of the problem (5) as $`t\mathrm{}`$ may be written as:
$$\left(\begin{array}{c}\varphi \\ \psi \end{array}\right)=\left(\begin{array}{c}\stackrel{~}{\mu }\\ \stackrel{~}{\nu }\end{array}\right)(k,\overline{k},\xi ,\eta ,t)+\left(\begin{array}{c}\stackrel{~}{\nu }\\ \stackrel{~}{\mu }\end{array}\right)(\overline{k},k,\xi ,\eta ,t)+O(t^{4/3}).$$
(97)
The second term in this formula is the solution of the problem (5) with the boundary condition:
$$\left(\begin{array}{c}\stackrel{~}{\mu }\\ \stackrel{~}{\nu }\end{array}\right)|_{|k|\mathrm{}}=\left(\begin{array}{c}0\\ 1\end{array}\right).$$
The proof of the theorem 1. Let us substitute the function $`\psi (k,\overline{k},\xi ,\eta ,t)`$ into the formula for the solution of the equation KP-2(6). Differentiate the integrand with respect to $`x`$. The main terms of the integrand are the terms which appear after differentiating of the exponent. The derivatives with respect to $`x`$ of others factors of integrand are small because they depend on $`x`$ slowly. The integrals of such terms over the plane are evaluated by the order $`O(t^{4/3})`$ uniformly with respect to $`(x,y)\text{}`$. Rewrite the main term of the integral as the integral over real plane: $`(\kappa ,\lambda )\text{}^2,`$ where $`\kappa =\text{Re}(k),\lambda =\text{Im}(k)`$. As a result we obtain:
$`u(x,y,t)=4{\displaystyle }{\displaystyle _\text{}^2}d\kappa d\lambda |\kappa |f(\kappa +i\lambda )\times `$
$`\times \mathrm{exp}(it(8\kappa ^324\kappa \lambda ^2+2\kappa \xi +4\kappa \lambda \eta ))+O(t^{4/3}).`$ (98)
Thus the main term of the asymptotics of the solution of the Cauchy problem for the equation KP-2 is given by integral with fast oscillating exponent. Let us evaluate the asymptotic behaviour of this integral.
Let $`\eta ^2+12\xi >0`$ and $`t^{1/3}|\eta ^2+12\xi |1`$, then the stationary points of the exponent are: $`\lambda _{1,2}=\frac{\eta }{12}\pm \sqrt{\eta ^2+12\xi },\kappa _{1,2}=0`$. The stationary phase method (see for example ) gives:
$$u(x,y,t)=o(t^1).$$
Let $`\eta ^2+12\xi <0`$ and $`t^{1/3}|\eta ^2+12\xi |1`$, then $`\lambda _{1,2}=\frac{\eta }{12},\kappa _{1,2}=\pm \frac{1}{2}\sqrt{\eta ^212\xi }.`$ The asymptotic behaviour of the integral is:
$`u(x,y,t)=4t^1{\displaystyle \frac{\pi }{12i\sqrt{\eta ^212\xi }}}f_0({\displaystyle \frac{1}{2}}\sqrt{\eta ^212\xi }+{\displaystyle \frac{i\eta }{12}})\times `$
$`\times \mathrm{exp}(11it\sqrt{{\displaystyle \frac{y^2}{t^2}}12{\displaystyle \frac{x}{t}}})+c.c.+o(t^1).`$
To evaluate the main term of the asymptotics of the integral when $`|12\xi +\eta ^2|=o(1)`$, we substitute the scaled variables $`p_1=t^{1/3}\text{Re}(kk_0)`$, $`p_2=t^{1/3}\text{Im}(kk_0)`$ and the parameter $`v^2=t^{2/3}(\eta ^2+12\xi )/\sqrt{12}`$ into integral (98). As a result we obtain:
$`u(\xi ,\eta ,t)=4it^1f(k_0){\displaystyle }{\displaystyle _\text{}^2}dp_1dp_2\times `$
$`p_1\mathrm{exp}(i(8p_1^32v^2p_124p_1p_2^2))+o(t^1).`$
Let us integrate the internal integral with respect to the parameter $`p_2`$, use the even property of the integrand with respect to $`p_1`$. As a result we obtain:
$`u(x,y,t)=8it^1\sqrt{\pi }f(i\eta /12)({\displaystyle _0^{\mathrm{}}}dp_1\sqrt{p_1}\mathrm{cos}\left(8p_1(p_1^28v^2)\right)+`$
$`+{\displaystyle _0^{\mathrm{}}}dp_1\sqrt{p_1}\mathrm{sin}\left(8p_1(p_1^28v^2)\right))+o(t^1).`$
The theorem 1 is proved.
## Appendix A Asymptotic behaviour of double integral with weak singularity of the integrand
Here we obtain the asymptotic behaviour of integrals which are appeared when the asymptotic solution of (9) was studied.
Evaluating of the asymptotic behaviour of one-dimensional integrals with weak singular integrand and fast oscillated exponent was done in p.26 and p.332. The asymptotic behaviour of many-dimensional integrals with fast oscillating exponent was studied in , . The asymptotic behaviour of the Cauchy integrals with fast oscillating exponent in one-dimensional case was studied in and in many-dimensional case was studied in . The asymptotic behaviour of the two-dimensional for some integrals over all complex plane with weak singularity was studied in .
### A.1 Integrals over half-plane with general stationary point of fast oscillated exponent
Here we study an asymptotic behaviour of an integral:
$$I=_{\mathrm{\Omega }^+}\frac{dnd\overline{n}}{ln}\mathrm{exp}(i(n^2+\overline{n}^2)),$$
where $`|l|\mathrm{}`$ and the domain $`\mathrm{\Omega }^+=\{\text{Re}(l+\overline{l}i(l\overline{l}))>0\}`$.
###### Theorem 7
The asymptotic behaviour of the integral $`I`$ as $`|p|\mathrm{}`$ has the form:
$$I=2i\pi \frac{\mathrm{exp}(i(l^2+\overline{l}^2))}{2il}\frac{3i\pi }{2\overline{l}}+O(|l|^2)\text{ѳ }l\mathrm{\Omega }^+;$$
$$I=\frac{i\pi }{2\overline{l}}+O(|l|^2)\text{where}l\mathrm{\Omega }^+.$$
The proof. Let us suppose that $`l\mathrm{\Omega }^+`$. Divide the domain $`\mathrm{\Omega }^+`$ into three domains: first one is $`\mathrm{\Omega }_1^+=\{\mathrm{\Omega }^+\backslash \{[|n||l|/2][|nl|<\epsilon ]\}\}`$. This domain hasn’t stationary point of phase function of the exponent and the singularity of the integrand. Second one is $`\mathrm{\Omega }_2^+=\{\mathrm{\Omega }^+\backslash [|n||l|/2]\}`$. This domain contains the stationary point of the phase function. At last, third domain is $`\mathrm{\Omega }_3^+=\{|nl|<\epsilon \}`$. This domain contains the singularity of the integrand.
Let us integrate by parts over $`\mathrm{\Omega }_1^+`$. As a result we obtain:
$$_{\mathrm{\Omega }_1^+}\frac{dnd\overline{n}}{\overline{ln}}\mathrm{exp}(i(n^2+\overline{n}^2))=_{\mathrm{\Omega }_1^+}\frac{d\overline{n}}{\overline{ln}}\frac{\mathrm{exp}(i(n^2+\overline{n}^2))}{2in}$$
$$_{\mathrm{\Omega }_1^+}\frac{dnd\overline{n}}{\overline{ln}}\frac{\mathrm{exp}(i(n^2+\overline{n}^2))}{2in^2}.$$
(99)
The boundary of the domain $`\mathrm{\Omega }_1^+`$ includes a large half-circle of radius $`R`$ as $`R\mathrm{}`$, a circle of radius $`\epsilon `$ over the point $`l=n`$, a half-circle of radius $`|l|/2`$ and two segments: $`[R\mathrm{exp}(3i\pi /4),|l|\mathrm{exp}(3i\pi )/2]`$ and $`[|l|\mathrm{exp}(i\pi /4)/2,R\mathrm{exp}(3i\pi )]`$.
Consider the integrals over boundary of the domain $`\mathrm{\Omega }_1^+`$. The integral over large half-circle is equals to zero as $`R\mathrm{}`$. The integral over half-circle of radius $`|l|/2`$ has order $`|l|^3`$ (because of oscillations and small value of the integrand). The integral over the circle at $`n=l`$ equals to a residue of the integrand multiplied by $`2i\pi `$ as $`\epsilon 0`$. The integral over the segments $`[R\mathrm{exp}(3i\pi /4),|l|\mathrm{exp}(3i\pi )/2]`$ and $`[|l|\mathrm{exp}(i\pi /4)/2,R\mathrm{exp}(3i\pi )]`$ (let us denote their union by $`L`$) is reduced to the form:
$$_{nL}\frac{d\overline{n}}{2in\overline{(ln)}}=\frac{\mathrm{exp}(3i\pi /4)}{2i\overline{l}}_{|\lambda ||l|/2}\frac{d\lambda }{\overline{l}\lambda \mathrm{exp}(i\pi /4)}.$$
The second term of (99) has to be evaluated as by $`O(|l|^3)`$.
Let us consider the integral over $`\mathrm{\Omega }_2^+`$. Represent:
$$\frac{1}{\overline{ln}}=\frac{1}{\overline{l}}\left(1+\frac{\overline{n}}{\overline{ln}}\right).$$
Then the integral over $`\mathrm{\Omega }_2^+`$ has to be written as:
$`I_2={\displaystyle _{\mathrm{\Omega }_2^+}\frac{dnd\overline{n}}{\overline{ln}}\mathrm{exp}(i(n^2+\overline{n}^2))}=`$
$`{\displaystyle \frac{1}{\overline{l}}}{\displaystyle _{\mathrm{\Omega }_2^+}\frac{dnd\overline{n}}{\overline{ln}}\mathrm{exp}(i(n^2+\overline{n}^2))}+`$
$`{\displaystyle \frac{1}{\overline{l}}}{\displaystyle _{\mathrm{\Omega }_2^+}\frac{dnd\overline{n}}{\overline{ln}}\overline{n}\mathrm{exp}(i(n^2+\overline{n}^2))}.`$
Here we integrate by parts the second term. After evaluations we obtain:
$`I_2={\displaystyle \frac{1}{\overline{l}}}{\displaystyle _{\mathrm{\Omega }^+}𝑑n}d\overline{n}\mathrm{exp}(i(n^2+\overline{n}^2))+`$
$`{\displaystyle \frac{1}{2i\overline{l}}}{\displaystyle _{|l|\mathrm{exp}(3i\pi /4)/2}^{|l|\mathrm{exp}(i\pi /4)/2}}{\displaystyle \frac{dn}{\overline{ln}}}+O(|l|^2).`$
The integral over $`\mathrm{\Omega }_3^+`$ as $`\epsilon 0`$ equals to zero.
Let us sum the obtained asymptotics:
$$I=2i\pi \frac{\mathrm{exp}(i(l^2+\overline{l}^2))}{2il}+\frac{\mathrm{exp}(i\pi /2)}{2i\overline{l}}_{\mathrm{}}^{\mathrm{}}\frac{d\lambda }{\overline{l}\mathrm{exp}(i\pi /4)\lambda }+$$
$$\frac{1}{2i\overline{l}}_{\mathrm{\Omega }^+}𝑑nd\overline{n}\mathrm{exp}(i(n^2+\overline{n}^2))+O(|l|^2).$$
Thus the first statement of the theorem is proved. The second statement has to be proved by the same way.
The theorem is proved.
### A.2 Asymptotic behaviour of the integral <br>with confluent phase function
In this section we obtain an asymptotic behaviour of an integral as $`|p|\mathrm{}`$
$$W^+=_{\mathrm{\Omega }^+}\frac{drd\overline{r}}{\overline{pr}}\mathrm{exp}(i\omega (r)).$$
(100)
###### Theorem 8
The asymptotic behaviour of the integral (100) where $`\omega (p)`$ $`=4(p^3+\overline{p}^3)v^2(p+\overline{p})`$ as $`|p|\mathrm{}`$ and $`|p\pm \frac{v}{\sqrt{12}}|0`$ has the form:
$$W^+=[\begin{array}{c}\frac{1}{\overline{p}}\varphi _{00}^+(v^2)+\frac{1}{\overline{p}^2}\varphi _{01}^+(v^2)+\frac{\pi i}{12\overline{p}^2}+2\pi i\frac{\mathrm{exp}(i\omega (p))}{12p^2}+\\ +O(|p|^3+|v|^2|p|^2),\text{when}p\mathrm{\Omega }^+;\\ \\ \frac{1}{\overline{p}}\varphi _{00}^+(v^2)+\frac{1}{\overline{p}^2}\varphi _{01}^+(v^2)\frac{\pi i}{12\overline{p}^2}+\\ +O(|p|^3+|v|^2|p|^2),\text{when}p\mathrm{\Omega }^+.\end{array}$$
(101)
Let us prove the theorem. Represent the integral in the form:
$`W^+={\displaystyle \frac{1}{\overline{p}}}{\displaystyle _{\mathrm{\Omega }^+}𝑑n}d\overline{n}\mathrm{exp}(i\omega (n))+`$
$`+{\displaystyle \frac{1}{\overline{p}^2}}{\displaystyle _{\mathrm{\Omega }^+}𝑑n}d\overline{n}\mathrm{exp}(i\omega (n))\overline{n}+`$
$`+{\displaystyle \frac{v^2}{12\overline{p}^2}}{\displaystyle _{\mathrm{\Omega }^+}\frac{dnd\overline{n}}{\overline{pn}}\overline{n}^2\mathrm{exp}(i\omega (n))}+`$
$`+{\displaystyle \frac{1}{12\overline{p}^2}}{\displaystyle _{\mathrm{\Omega }^+}\frac{dnd\overline{n}}{\overline{pn}}(12\overline{n}^2v^2)\mathrm{exp}(i\omega (n))}.`$
Integrate by parts the last term. As a result we obtain:
$`W^+={\displaystyle \frac{1}{\overline{p}}}\varphi _{00}^+(v^2)+{\displaystyle \frac{1}{\overline{p}^2}}\varphi _{01}^+(v^2)+`$
$`+{\displaystyle \frac{1}{12\overline{p}^2}}{\displaystyle _{\mathrm{\Omega }^+}}{\displaystyle \frac{d\overline{n}}{\overline{pn}}}{\displaystyle \frac{(12\overline{n}^2v^2)\mathrm{exp}(i\omega (n)))}{(i)(12n^2v^2)}}+`$
$`O(|p|^3+|v|^2|p|^2).`$
Consider the integral over the boundary of the domain $`\mathrm{\Omega }^+`$. This integral may be considered as a sum of integrals over imaginary axis $`\text{Re}(r)=0`$, over half-circle of radius $`R`$ as $`R\mathrm{},\text{Re}(r)>0`$ and over the circle $`|pr|=\epsilon `$ as $`\epsilon 0`$ (if $`p\mathrm{\Omega }^+`$). We can see that the integral over the large half-circle tends to zero. As a result we obtain the statement of the theorem.
One more double integral which we need is:
$$U^+=_{\mathrm{\Omega }^+}\frac{drd\overline{r}}{pr}\overline{r}\mathrm{exp}(i\omega (r)).$$
Its asymptotics can be evaluated by the same way as the asymptotics of the integral (100) as $`|p|\mathrm{}`$ $`|p\pm v|0`$:
$`U=[\begin{array}{c}\frac{1}{\overline{p}}\varphi _{10}^+(v^2)+2i\pi \frac{1}{12ip}\mathrm{exp}(i\omega (p))+\frac{\pi i}{12\overline{p}}+2\pi i\frac{\mathrm{exp}(i\omega (p))}{12p^2}+\\ +O(|p|^2+|v|^2|p|^2),\text{when}p\mathrm{\Omega }^+;\\ \\ \frac{1}{\overline{p}}\varphi _{10}^+(v^2)\frac{\pi i}{12\overline{p}}+O(|p|^2+|v|^2|p|^2),\text{when}p\mathrm{\Omega }^+.\end{array}`$
## Appendix B Reducing of four-multiply integral into double integral
This section is pure technical. Here we show as the four-multiply integrals have to be reduced into double integrals over half-plane and all complex plane.
### B.1 Four-multiply integral with nondegenerate phase of the exponent.
Let us show that the four-multiply integral may be written as a sum of double integrals. Change the variable: $`nm=r`$, then the integral $`J`$ has the form:
$`J={\displaystyle \frac{1}{2i\pi }}{\displaystyle }{\displaystyle _{\text{}}}{\displaystyle \frac{dnd\overline{n}}{l_jn}}\times `$
$`{\displaystyle \frac{1}{2i\pi }}{\displaystyle _{\text{}}\frac{drd\overline{r}}{\overline{r}}\mathrm{exp}(i(r^2+\overline{r}^2))\mathrm{exp}(2i(rn+\overline{rn}))}.`$
Integrate by part over $`\overline{n}`$. As a result we obtain:
$`J={\displaystyle \frac{1}{4\pi ^2}}\underset{R\mathrm{}}{lim}{\displaystyle _{|n|=R}}{\displaystyle \frac{dn\overline{n}}{l_jn}}\times `$
$`{\displaystyle _{\text{}}\frac{drd\overline{r}}{\overline{r}}\mathrm{exp}(i(r^2+\overline{r}^2))\mathrm{exp}(2i(rn+\overline{rn}))}+`$
$`{\displaystyle \frac{1}{2i\pi }}\underset{\epsilon 0}{lim}{\displaystyle _{|ln|=\epsilon }}{\displaystyle \frac{dn\overline{n}}{l_jn}}\times `$
$`{\displaystyle \frac{1}{2i\pi }}{\displaystyle _{\text{}}\frac{drd\overline{r}}{\overline{r}}\mathrm{exp}(i(r^2+\overline{r}^2))\mathrm{exp}(2i(rn+\overline{rn}))}`$
$`{\displaystyle \frac{1}{2i\pi }}{\displaystyle }{\displaystyle _{\text{}}}{\displaystyle \frac{dnd\overline{n}}{l_jn}}\overline{n}\times `$
$`{\displaystyle \frac{1}{2i\pi }}{\displaystyle _{\text{}}𝑑r}d\overline{r}(2i)\mathrm{exp}(i(r^2+\overline{r}^2))\mathrm{exp}(2i(rn+\overline{rn})).`$
First term is equal to zero because of fast oscillating of the integrand, second term is equal to residue with sign minus of the integrand at $`n=l`$, the internal integral in the third term has to be evaluated. Finally we obtain:
$$J=\overline{l}_j\frac{\mathrm{exp}(i(l_j^2+\overline{l}_j^2))}{2i\pi }_{\text{}}\frac{dnd\overline{n}}{\overline{l_jm}}\mathrm{exp}(i(n^2+\overline{n}^2))\mathrm{exp}(i(l_j^2+\overline{l}_j^2)).$$
### B.2 Four-multiply integral over different half-planes with nondegenerate phase function
Let us consider an integral:
$`J_+={\displaystyle }{\displaystyle _\mathrm{\Omega }^{}}{\displaystyle \frac{dnd\overline{n}}{ln}}{\displaystyle \frac{\mathrm{exp}(i(n^2+\overline{n}^2))}{2i\pi }}\times `$
$`{\displaystyle _{\mathrm{\Omega }^+}\frac{dmd\overline{m}}{\overline{nm}}\mathrm{exp}(i(m^2+\overline{m}^2))}.`$
A partially integrating over $`\overline{n}`$ gives:
$`J_+={\displaystyle _\mathrm{\Omega }_{}}{\displaystyle \frac{dn}{ln}}\overline{n}{\displaystyle \frac{\mathrm{exp}(i(n^2+\overline{n}^2))}{2i\pi }}\times `$
$`{\displaystyle _{\mathrm{\Omega }^+}\frac{dmd\overline{m}}{\overline{nm}}\mathrm{exp}(i(m^2+\overline{m}^2))}`$
$`{\displaystyle }{\displaystyle _\mathrm{\Omega }^{}}{\displaystyle \frac{dnd\overline{n}}{ln}}\overline{n}{\displaystyle \frac{}{\overline{n}}}{\displaystyle \frac{\mathrm{exp}(i(n^2+\overline{n}^2))}{2i\pi }}\times `$
$`{\displaystyle _{\mathrm{\Omega }^+}\frac{dmd\overline{m}}{\overline{nm}}\mathrm{exp}(i(m^2+\overline{m}^2))}.`$
The derivative is:
$`J_+={\displaystyle _\mathrm{\Omega }^{}}{\displaystyle \frac{dn}{ln}}{\displaystyle \frac{\overline{n}\mathrm{exp}(i(n^2+\overline{n}^2))}{2i\pi }}\times `$
$`{\displaystyle _{\mathrm{\Omega }^+}\frac{dmd\overline{m}}{\overline{nm}}\mathrm{exp}(i(m^2+\overline{m}^2))}`$
$`{\displaystyle \frac{3}{4}}{\displaystyle _\mathrm{\Omega }^{}}{\displaystyle \frac{dn}{ln}}\mathrm{exp}(i(n^2+\overline{n}^2)).`$
Let us consider integral over part of the boundary in first term :
$`=\underset{R\mathrm{}}{lim}{\displaystyle _{R\mathrm{exp}(i\pi /4)}^{R\mathrm{exp}(3i\pi /4)}}{\displaystyle \frac{dn}{ln}}{\displaystyle \frac{\overline{n}\mathrm{exp}(i(n^2+\overline{n}^2))}{2i\pi }}\times `$
$`{\displaystyle _{\mathrm{\Omega }^+}\frac{dmd\overline{m}}{\overline{nm}}\mathrm{exp}(i(m^2+\overline{m}^2))}.`$
Change the order of integration over $`m,\overline{m}`$ and $`n`$. The internal integral over $`n`$ has to be evaluate:
$$\underset{R\mathrm{}}{lim}_{R\mathrm{exp}(i\pi /4)}^{R\mathrm{exp}(3i\pi /4)}dn\frac{\overline{n}}{(ln)\overline{(nm)}}=[\begin{array}{c}\frac{1}{2}+\frac{\overline{m}}{il\overline{m}},\text{when}l\mathrm{\Omega }^+;\\ 1/2\text{when}l\mathrm{\Omega }^+.\end{array}$$
In last expression of $`J_+`$ the integrals over the large half of circle $`|n|=R`$ as $`R\mathrm{}`$ are equal to zero.
Final formulas for $`J_+`$ have to be written as:
$$J_+=[\begin{array}{c}\frac{5}{4}i\pi +il_{\mathrm{\Omega }^+}𝑑md\overline{m}\frac{\mathrm{exp}(i(m^2+\overline{m}^2))}{il\overline{m}},\text{when}l\mathrm{\Omega }^+;\\ \\ \frac{5}{4}i\pi +\frac{3}{2}i\pi \mathrm{exp}(i)l^2+\overline{l}^2))\overline{l}\mathrm{exp}(i(l^2+\overline{l}^2))\times \\ _{\mathrm{\Omega }^+}𝑑md\overline{m}\frac{\mathrm{exp}(i(m^2+\overline{m}^2))}{\overline{lm}},\text{when}l\mathrm{\Omega }^{}.\end{array}$$
### B.3 Four-multiply integral with confluent <br>phase function
Let us reduce an four-multiply integral
$$J_1=_{\text{}}\frac{dnd\overline{n}}{pn}\mathrm{exp}(i\omega (n))_{\text{}}\frac{dmd\overline{m}}{\overline{nm}}\mathrm{exp}(i\omega (m))$$
into sum of double integrals. Denote:
$$V(n,v^2)=\frac{\mathrm{exp}(i\omega (n))}{2i\pi }_{\text{}}\frac{drd\overline{r}}{\overline{nr}}\mathrm{exp}(i\omega (r)).$$
The function $`V(p,v^2)`$ has to be written as:
$$\frac{V}{\overline{n}}=12i\overline{n}\mathrm{exp}(i\omega (n))\varphi _{00}(v^2)+\mathrm{\hspace{0.17em}12}i\mathrm{exp}(i\omega (n))\varphi _{01}(v^2).$$
We can obtain this formula by changing variable in the integrand $`\rho =nr`$ and differentiating the obtained expression of $`V(n,v^2)`$ with respect to $`\overline{n}`$ and, finally changing the variable back: $`r=n\rho `$.
Partial differentiating of the double integral over $`\overline{n}`$ and $`n`$ gives:
$$_{\text{}}\frac{dnd\overline{n}}{pn}V(n,v^2)=\underset{R\mathrm{}}{lim}_{|n|=R}\frac{dn\overline{n}V(n,v^2)}{pn}+$$
$$+\underset{\epsilon 0}{lim}_{|pn|=\epsilon }\frac{dn\overline{n}}{pn}V(n,v^2)_{\text{}}\frac{dnd\overline{n}}{pn}\overline{n}\frac{V(n,v^2)}{\overline{n}}.$$
(103)
Asymptotic behaviour of $`V(n,v^2)`$ as $`|n|\mathrm{}`$ was obtained above. Using this asymptotics for the first part of the formula (103) gives zero. Second term of the right-hand side of (103) gives the residue of the integrand at $`p`$ multiplied by $`2i\pi `$. Then we obtain:
$$J_1(p,v^2)=\overline{p}V(p,v^2)\varphi _{00}(v^2)\mathrm{exp}(i\omega (p))$$
$$\mathrm{\hspace{0.17em}12}i\varphi _{01}(v^2)\frac{1}{2i\pi }_{\text{}}\frac{dnd\overline{n}}{pr}\overline{n}\mathrm{exp}(i\omega (n))+$$
$$+iv^2\varphi _{00}(v^2)\frac{1}{2i\pi }_{\text{}}\frac{dnd\overline{n}}{pr}\mathrm{exp}(i\omega (n)).$$
On the same way one can evaluate an integral
$$J_1^+=_\mathrm{\Omega }^{}\frac{dnd\overline{n}}{pn}\mathrm{exp}(i\omega (n))_{\mathrm{\Omega }^+}\frac{dmd\overline{m}}{\overline{nm}}\mathrm{exp}(i\omega (m)).$$
The difference between the integral $`J_1`$ and the integral $`J_1^+`$ consists in the addition terms over boundaries of the domains $`\mathrm{\Omega }^+`$ and $`\mathrm{\Omega }^{}`$ in the result. The finally forms are:
when $`p\mathrm{\Omega }^+`$:
$`J_1^+=\varphi _{00}^+(v^2)\mathrm{exp}(i\omega (p))+`$
$`+\overline{p}\mathrm{exp}(i\omega (p)){\displaystyle _{\mathrm{\Omega }^+}\frac{dmd\overline{m}}{\overline{pm}}\mathrm{exp}(i\omega (m))}`$
$`{\displaystyle \frac{1}{2}}{\displaystyle _\mathrm{\Omega }^{}\frac{dnd\overline{n}}{pn}\overline{n}\mathrm{exp}(i\omega (n))}+`$
$`+iv^2\varphi _{00}^+(v^2){\displaystyle \frac{1}{2i\pi }}{\displaystyle _\mathrm{\Omega }^{}\frac{dnd\overline{n}}{pn}\mathrm{exp}(i\omega (n))}`$
$`12i\varphi _{01}^+(v^2){\displaystyle \frac{1}{2i\pi }}{\displaystyle _\mathrm{\Omega }^{}\frac{dnd\overline{n}}{pn}\overline{n}\mathrm{exp}(i\omega (n))};`$
when $`p\mathrm{\Omega }^+`$:
$$J_1^+=\varphi _{00}^+(v^2)+p_{\mathrm{\Omega }^+}\frac{dmd\overline{m}}{p+\overline{m}}\mathrm{exp}(i\omega (m))$$
$`{\displaystyle \frac{1}{2}}{\displaystyle _\mathrm{\Omega }^{}\frac{dnd\overline{n}}{pn}\overline{n}\mathrm{exp}(i\omega (n))}`$
$`iv^2\varphi _{00}^+(v^2){\displaystyle \frac{1}{2i\pi }}{\displaystyle _\mathrm{\Omega }^{}\frac{dnd\overline{n}}{pn}\mathrm{exp}(i\omega (n))}`$
$`12i\varphi _{01}^+(v^2){\displaystyle \frac{1}{2i\pi }}{\displaystyle _\mathrm{\Omega }^{}\frac{dnd\overline{n}}{pn}\overline{n}\mathrm{exp}(i\omega (n))}.`$
Acknowledgements.
I am grateful to S.G. Glebov, A.M. Il’in, L.A. Kalyakin and M.M. Shakir’yanov for stimulated discussions.
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# Untitled Document
THE TWO-EXPONENTIAL LIOUVILLE THEORY AND
THE UNIQUENESS OF THE THREE-POINT FUNCTION
L. O’Raifeartaigh<sup>1</sup> e-mail: lor@stp.dias.ie, J. M. Pawlowski<sup>2</sup> e-mail: jmp@stp.dias.ie, and V. V. Sreedhar<sup>3</sup> e-mail: sreedhar@stp.dias.ie
School of Theoretical Physics
Dublin Institute for Advanced Studies
10 Burlington Road, Dublin 4, Ireland
Abstract
It is shown that in the two-exponential version of Liouville theory the coefficients of the three-point functions of vertex operators can be determined uniquely using the translational invariance of the path integral measure and the self-consistency of the two-point functions. The result agrees with that obtained using conformal bootstrap methods. Reflection symmetry and a previously conjectured relationship between the dimensional parameters of the theory and the overall scale are derived.
Introduction: Although the quantisation of two-dimensional Liouville theory with a potential of the form $`V_b=\mu e^{2b\stackrel{~}{\varphi }(x)}`$, where $`\stackrel{~}{\varphi }(x)`$ is a scalar field and $`\mu `$ and $`b`$ are constants, has been widely investigated , it still presents some problems. What is perhaps the most disturbing problem is the following: The three point functions of vertex operators $`\mathrm{exp}[2\alpha _I\stackrel{~}{\varphi }(x_I)]`$ play a central role in the theory, in the sense that all $`N`$-point functions can be obtained from these by integration, and have the following form
$$𝒢_3=𝐂_3\underset{I=1}{\overset{3}{}}|\frac{x_Ix_J}{L}|^{\mathrm{\Delta }_{(IJ)}(\alpha )}\text{where}𝐂_3=\left(g_b(\xi )Z_0^b(\xi )\right)\left(\frac{K(\xi ,\alpha _I)}{k^{}(\xi )}\right)$$
$`(1)`$
Here $`L`$ is a constant scale, the $`\mathrm{\Delta }_{(IJ)}`$ are known combinations of constant conformal weights, $`\xi =q_{I=1}^3\alpha _I`$, and $`q=b+1/b`$, while $`Z_0^b(\xi )`$ is related to the zero-mode integration and can be computed explicitly. $`K`$ and $`k`$ are the functions defined in (26). The problem is that $`𝐂_3`$ can actually be computed only at the points $`\xi =mb`$ for $`mZ_+`$ and the extrapolation from these points to general values of $`\xi `$ leaves a factor $`g_b(\xi )`$ in (1) undetermined. So far it has not been possible to determine this factor from first principles. Instead, what has been done is to make the so-called DOZZ Ansatz that
$$𝐂_3=\frac{K(\xi ,\alpha _I)}{k(\xi )}h_b(\xi )g_bZ_0^b\frac{k}{k^{}}=1$$
$`(2)`$
and check that the resulting three-point function satisfies some reasonable, but extra, conditions such as reflection symmetry and crossing symmetry. The latter check is the decisive one, but uses some special four-point functions .
In a previous paper it was suggested that the gap in the direct derivation of $`h(\xi )`$ could be closed by using a potential of the form $`V(\varphi )=\mu _be^{2b\stackrel{~}{\varphi }}+\mu _ce^{2c\stackrel{~}{\varphi }}`$, where $`bc=1`$, – this being the most general potential one can use in the path integral context for a conformal field theory of a single real scalar field without derivative interactions – rather than the standard form $`V_b`$. As in the one-exponential theory, the $`𝐂_3`$ in the two-exponential theory can be calculated only at a discrete, but much larger, set of points, namely $`\xi _{mn}=mb+nc`$ and the equation corresponding to (1) takes the form
$$𝐂_3=\left(g(\xi )Z_0(\xi )\right)\left(\frac{K(\xi ,\alpha _I)}{k^{}(\xi )}\right)$$
$`(3)`$
where $`g(\xi )`$ is a priori undetermined, the functions $`K`$ and $`k`$ are exactly the same as in (1) but $`Z_0`$ is not. It was shown in that, subject to two conditions, the factor $`g(\xi )`$ could be fixed and thus the DOZZ Ansatz (2) could be derived. The conditions were (a) that the dimensional parameters $`\mu _b`$ and $`\mu _c`$ were related to the overall scale $`L`$ of the system by an equation of the form $`\mathrm{\Omega }(\mu _b,\mu _c,L)=1`$, where $`\mathrm{\Omega }`$ is the function defined in (25), and (b) that $`h(\xi )`$, defined as $`h(\xi )=g(\xi )Z_0(\xi )k(\xi )/k^{}(\xi )`$, had no singular points.
In the present paper we refine these results considerably and extend the analysis of the two-exponential theory. In particular we show that, subject to the mild technical condition given in (37), the DOZZ Ansatz follows from the translational invariance of the path integral measure and the self-consistency of the two-point functions. We also present an extrapolation of the fluctuating part of the path integral from the lattice points $`\xi _{mn}`$, for which the $`Z_0`$ part of the path integral is automatically
$$Z_0(\xi )=\frac{k^{}(\xi )}{k(\xi )}$$
$`(4)`$
This means that the zero mode integral for the two-exponential theory produces exactly the factor $`k^{}/k`$ that was postulated in the DOZZ Ansatz.
As in , we use the path-integral formalism; and as the symmetries of the path-integral and the associated sum rules are of interest in their own right, a secondary purpose of the paper is to present these in a systematic way. The most important symmetries are those connected with the translational invariance of the measure and conformal (Weyl) invariance.
A by-product is an analysis of the two-point function $`𝒢_2`$. Although $`𝒢_2`$ cannot be defined directly because of conformal invariance, it can be defined both as a limit of $`𝒢_3`$ when one of the $`\alpha `$’s becomes zero, and as a volume integral of $`𝒢_3`$ when the corresponding $`\alpha `$ is $`b`$ or $`c`$. The compatibility of the two definitions and the sum rule mentioned above lead to a simple linear homogeneous sum rule for the quantity $`h(\xi )`$. Together with the boundary conditions $`h(\xi _{mn})=\mathrm{\Omega }^{m+n}`$, obtained by direct computation, this sum rule fixes $`h(\xi )=1`$ uniquely. It also fixes $`\mathrm{\Omega }=1`$ which is the relation between the dimensional constants $`\mu _b`$ and $`\mu _c`$ and the overall scale $`L`$ that was conjectured in an earlier paper. The corresponding sum rule and boundary conditions in the one-exponential theory would imply only that $`h(\xi )`$ is periodic. If $`𝒢_2`$ is interpreted as an inner-product of states, then $`h(\xi )=1`$ implies that for each conformal weight a zero-norm state decouples, so that there is only one physical state. The decoupling is equivalent to reflection invariance, which therefore emerges as an output.
The generating functional: The generating functional of the two-exponential Liouville theory is defined as
$$Z[J]=[d\stackrel{~}{\varphi }]e^{{\scriptscriptstyle d^2x\sqrt{g(x)}\left[{\scriptscriptstyle \frac{1}{4\pi }}\stackrel{~}{\varphi }\mathrm{\Delta }\stackrel{~}{\varphi }+{\scriptscriptstyle \frac{q}{4\pi }}\stackrel{~}{\varphi }+\mu _be^{2b\stackrel{~}{\varphi }}+\mu _ce^{2c\stackrel{~}{\varphi }}J(x)\stackrel{~}{\varphi }(x)\right]}}$$
$`(5)`$
where $`\mathrm{\Delta }`$ is the Laplace-Beltrami operator, $``$ is the Ricci scalar, the coupling constants $`\mu _b,\mu _c`$ have dimensions of mass squared, and $`b`$ and $`c`$ are dimensionless constants. For the $`N`$-point functions of vertex operators we have $`e^{{\scriptscriptstyle d^2x\sqrt{g}J\stackrel{~}{\varphi }}}=_{I=1}^N\nu _I\sqrt{g}e^{2\alpha _I\stackrel{~}{\varphi }}`$, but we need not yet specialise to this case. The $`\nu _I`$ are constants which, like the $`\mu _b`$ and $`\mu _c`$, have to be renormalised because of the short distance singularities of the Green’s function defined by the equation $`\mathrm{\Delta }G(x,y)=\frac{\pi }{\sqrt{g}}\delta ^2(xy)`$. As explained in detail in , the Green’s function for the non-coincident and coincident arguments is defined as:
$$G(x,y)=\frac{1}{2}\text{ln}\frac{xy}{L}+𝒪(V^1)\text{and}G(x,x+dx)\frac{1}{2}\mathrm{ln}[\frac{ds}{L}]+\frac{1}{4}\mathrm{ln}\sqrt{g(x)}$$
$`(6)`$
respectively, where $`ds`$ is the infinitesimal geodesic separation, $`L`$ sets the overall scale, and $`V`$ is the volume of the space.
As the Green’s function $`G(x,x)`$ occurs only in the form $`\alpha ^2G(x,x)`$ at each vertex, its infinite part may be absorbed by renormalising $`\mu _b`$, $`\mu _c`$ and $`\nu _I`$. Let us define the renormalisations of $`\mu _b`$, $`\mu _c`$ and $`\nu _I`$ to be
$$\mu _b\mu _b(\mathrm{\Lambda }\frac{ds}{L})^{2(b^2qb)},\mu _c\mu _c(\mathrm{\Lambda }\frac{ds}{L})^{2(c^2qc)},\nu _I\nu _I(\mathrm{\Lambda }\frac{ds}{L})^{2(\alpha _I^2q\alpha _I)}$$
$`(7)`$
where $`\mathrm{\Lambda }`$ is the dimensionless ultra-violet renormalisation scale. Anticipating the result that $`b+c=q,bc=1`$ from Weyl invariance, it is easy to see that $`\mu _b`$ and $`\mu _c`$ have the naive scaling dimensions with respect to $`\mathrm{\Lambda }`$.$`^1`$ There is arbitrariness in the definition of renormalisation associated with the translational invariance of the path integral measure. However, this arbitrariness does not affect equation (8). Clearly the $`b^2`$, $`c^2`$, $`\alpha ^2`$ powers of $`ds`$ in the above equation cancel the infinities coming from the Green’s functions at coincident points. Counting the powers of the remaining factors of $`ds`$ and $`\mathrm{\Lambda }`$, we find, after a little algebra, that the renormalised functional integral is proportional to
$$[\frac{ds}{L}]^{2q^2}\underset{I=1}{\overset{N}{}}\mathrm{\Lambda }^{2(\alpha _I^2q\alpha _I)}$$
$`(8)`$
Since the factor $`[ds/L]^{q^2}`$ may be absorbed in the Polyakov conformal anomaly term which is proportional to $`q^2`$, we see that the $`\mathrm{\Lambda }`$ dependence of the renormalised functional integral cancels except for the contribution coming from the external sources.
The conformal weights $`\mathrm{\Delta }_\alpha `$ of the fields exp \[$`2\alpha _I\stackrel{~}{\varphi }(x_I)`$\] may be read off directly by varying the path integral with respect $`\sqrt{g(x_I)}`$. Taking the $`\sqrt{g}`$ dependence of $`\alpha _I^2G(x_I,x_I)`$ and the cross-term $`q\alpha _Id^2xd^2y\sqrt{g(x)}(x)G(x,y)\delta (yx_I)`$ in the Gaussian integration into account, this yields $`\mathrm{\Delta }_\alpha =\alpha (q\alpha )`$. From (7) it is clear that the scaling dimensions of $`\nu _I`$ are minus the conformal weight of corresponding field. Thus fields of the same conformal weight are renormalised in the same way.
A sum rule for the generating functional: If we now specialise to the case where the external current is of the form $`_{I=1}^N\nu _Ie^{{\scriptscriptstyle d^2x\sqrt{g}J\stackrel{~}{\varphi }}}=_{I=1}^N\nu _I\sqrt{g}e^{2\alpha _I\stackrel{~}{\varphi }}`$, the renormalised functional integral is defined by
$$Z[J]=[d\stackrel{~}{\varphi }]e^{{\scriptscriptstyle d^2x\sqrt{g(x)}\left[{\scriptscriptstyle \frac{1}{4\pi }}\stackrel{~}{\varphi }\mathrm{\Delta }\stackrel{~}{\varphi }+{\scriptscriptstyle \frac{q}{4\pi }}\stackrel{~}{\varphi }+\sqrt{g}^{b^2}\mu _be^{2b\stackrel{~}{\varphi }}+\sqrt{g}^{c^2}\mu _ce^{2c\stackrel{~}{\varphi }}\right]}}\underset{I=1}{\overset{N}{}}\sqrt{g}^{\alpha _I^2}\nu _Ie^{2\alpha _I\stackrel{~}{\varphi }(x_I)}$$
$`(9)`$
where we are using the renormalisation prescription in (7) and we have absorbed, for convenience, the ln$`\sqrt{g}`$ part of $`G(x,x)`$ by letting $`e^{2\alpha \varphi }\sqrt{g}^{\alpha ^2}e^{2\alpha \varphi }`$. It is then understood that $`G_R(x,x)=0`$.
The translation invariance of the path integral measure can now be formulated as follows:
$$[d\stackrel{~}{\varphi }]\frac{\delta }{\delta \stackrel{~}{\varphi }(x)}\left[e^{S[\stackrel{~}{\varphi }]+{\scriptscriptstyle d^2x\sqrt{g}J(x)\stackrel{~}{\varphi }(x)}}\right]=0$$
$`(10)`$
where $`\frac{\delta }{\delta \stackrel{~}{\varphi }(x)}`$ is the generator of translations on the space of fields. From (10) we derive the quantum equations of motion for $`Z[J]`$:
$$\frac{1}{4\pi }\mathrm{\Delta }\frac{\delta Z[J]}{\delta J(x)}=\left(b\mu _b\sqrt{g}^{b^2}Z[J_{b,x}]+c\mu _c\sqrt{g}^{c^2}Z[J_{c,x}]\right)+\frac{1}{2}\left(J(x)\frac{q}{4\pi }(x)\right)Z[J]$$
$`(11)`$
where
$$J_{b,x}(y)=J(y)+2b\frac{1}{\sqrt{g(y)}}\delta ^2(yx),J_{c,x}(y)=J(y)+2c\frac{1}{\sqrt{g(y)}}\delta ^2(yx)$$
$`(12)`$
Carrying out the same procedure for the constant, zero-mode, part $`\varphi _0`$ of $`\stackrel{~}{\varphi }`$, we obtain the integrated form of (11), namely
$$d^2x\sqrt{g}\left(b\mu _b\sqrt{g}^{b^2}Z[J_{b,x}]+c\mu _c\sqrt{g}^{c^2}Z[J_{c,x}]\right)=\frac{1}{2}\left(J_0q\chi \right)Z[J]$$
$`(13)`$
where
$$J_0=d^2x\sqrt{g}J(x)\text{and}\chi =\frac{1}{4\pi }d^2x\sqrt{g}(x)$$
$`(14)`$
$`\chi `$ being the Euler number of the underlying manifold. Eq. (13) embodies the sum rule for the generating functional.
Weyl transformations: A local Weyl transformation can be performed by varying the generating functional with respect to $`\sqrt{g(x)}`$ where $`xx_I`$, the external points. For $`Z[J]`$ in (9), we find
$$\frac{\delta Z}{\delta \sqrt{g}}=(1+b^2)\sqrt{g}^{b^2}\mu _bZ[J_{b,x}]+(1+c^2)\sqrt{g}^{c^2}\mu _cZ[J_{c,x}]+\frac{q}{4\pi }\mathrm{\Delta }\frac{\delta Z[J]}{\delta J(x)}$$
$`(15)`$
The third term may be eliminated using (11) to get
$$\left(1+b^2qb\right)\sqrt{g}^{b^2}\mu _bZ[J_{b,x}]+\left(1+c^2qc\right)\sqrt{g}^{c^2}\mu _cZ[J_{c,x}]\left(\frac{q}{4\pi }(x)J(x)\right)\frac{q}{2}Z(J)$$
$`(16)`$
The Weyl condition is that the variation $`\frac{\delta Z}{\delta \sqrt{g}}`$ should be proportional to the external current namely $`\frac{q}{4\pi }J`$. Since this has to be valid for all currents $`J`$, the condition for Weyl invariance is that the first two terms must vanish and we have
$$q=(b+c)\text{and}bc=1,$$
$`(17)`$
The above approach may be contrasted with the one in where (17) was derived only after the path integral was evaluated.
The $`N`$-point functions: For the computation of general $`N`$-point functions of vertex operators, we let the underlying manifold be a two dimensional sphere. In that case, $`\chi =2`$ and there is only one zero-mode for $`\stackrel{~}{\varphi }`$, namely the constant $`\varphi _0`$. As explained in detail in , the expression for the $`N`$-point function may be simplified by using a Sommerfeld-Watson transform for the exponential of an integrated vertex operator. With $`\stackrel{~}{\varphi }=\varphi _0+\varphi `$ the resulting expression for the $`N`$-point function takes the form
$$𝒢_N=𝑑\varphi _0\frac{dudv}{\mathrm{\Gamma }(1+iu)\mathrm{\Gamma }(1+iv)}\frac{e^{2(i(bu+cv)\xi _N)\varphi _0}}{\text{sinh}\pi u\text{sinh}\pi v}\times 𝐂_N(iu,iv)$$
$`(18)`$
where
$$𝐂_N(iu,iv)=𝑑\varphi U_b^{iu}U_c^{iv}e^{{\scriptscriptstyle d^2x\sqrt{g}\left[{\scriptscriptstyle \frac{1}{4\pi }}\varphi \mathrm{\Delta }\varphi +{\scriptscriptstyle \frac{q}{4\pi }}R\varphi \right]}}𝚷_N\text{with}\xi _N=q\underset{I=1}{\overset{N}{}}\alpha _I$$
$`(19)`$
and
$$𝚷_N=\underset{I=1}{\overset{N}{}}\sqrt{g}^{\alpha _I^2}\psi _Ie^{2\alpha _I\stackrel{~}{\varphi }(x_I)},U_b(\varphi )=\mu _bd^2x(\sqrt{g})^{qb}e^{2b\varphi }$$
$`(20)`$
and similarly for $`bc`$. The integral $`𝐂_N(iu,iv)`$ is a Gaussian path integral for the fluctuations $`\varphi `$ which, for $`iu=m`$ and $`iv=n`$, $`m`$ and $`n`$ being positive integers, can be done in a straightforward manner and produces an ordinary multiple integral.
In terms of the $`N`$-point functions $`𝒢_N(x_I,\alpha _I)`$, the sum rule (13) takes the form
$$d^2x\sqrt{g}\left[b\mu _b𝒢_{N+1}(x_I,x,\alpha _I,b)+c\mu _c𝒢_{N+1}(x_I,x,\alpha _I,c)\right]=\xi _N𝒢_N(x_I,\alpha _I)$$
$`(21)`$
and thus relates the (integrated) $`N+1`$-point function to the $`N`$-point function. Note that the above equation requires that $`\xi _N0`$ because $`𝒢_N\mathrm{}\text{as}\xi _N0`$, making the right hand side indefinite.
The Three-point function: As is well-known, the three-point function is the lowest $`N`$-point function for which conformal invariance does not require the extraction of an infinite group volume factor. If we choose $`\xi `$ to be pure imaginary and integrate over the zero-mode $`\varphi _0`$ in (18) we obtain a delta function $`\delta (\xi bucv)`$, in which case the coefficient $`𝐂_3`$ may be written as
$$𝐂_3(\xi ,i(u+v))=𝑑\varphi e^{\frac{1}{4\pi }{\scriptscriptstyle d^2x\sqrt{g}[\varphi \mathrm{\Delta }\varphi +qR\varphi ]}}U_b^{iu}U_c^{iv}𝚷_\mathrm{𝟑}$$
$`(22)`$
Apart from the spectator variables $`\alpha _I\alpha _J`$, we see that, due to the delta-function, $`𝐂_3`$ is a function of only two variables, chosen as $`\xi `$ and $`u+v`$ for convenience. This is the great advantage of using the Sommerfeld-Watson transform. We then have in the infinite volume limit,
$$𝒢_3=\frac{dudv}{\mathrm{\Gamma }(1+iu)\mathrm{\Gamma }(1+iv)}\frac{\delta (\xi ibuicv)}{\text{sinh}\pi u\text{sinh}\pi v}𝐂_3(\xi ,i(u+v))$$
$`(23)`$
The problem is that $`𝐂_3`$ can only be computed at the points $`u,v=im,in`$ for $`m,nZ_+`$, where, as shown in , it is given by, for $`\xi _{mn}bm+cn`$,
$$𝐂_3(\xi _{mn},m+n)=(1)^{m+n}m!n!\lambda ^{\xi _{mn}}\mathrm{\Omega }^{m+n}\left(\frac{K(\xi _{mn},\alpha _I)}{k^{}(\xi _{mn})}\right)\underset{I=1}{\overset{3}{}}\frac{x_{IJ}}{L}^{\mathrm{\Delta }_{IJ}(\alpha )}$$
$`(24)`$
where
$$\lambda =\left(\frac{\mu _b\mathrm{\Phi }_b}{\mu _c\mathrm{\Phi }_c}\right)^{\frac{1}{bc}},\mathrm{\Omega }^{cb}=\frac{(\mu _bL^2\mathrm{\Phi }_b)^c}{(\mu _cL^2\mathrm{\Phi }_c)^b},\mathrm{\Phi }_b=\pi \gamma (b^2)(b^2)^{2qb},\gamma (x)=\frac{\mathrm{\Gamma }(x)}{\mathrm{\Gamma }(1x)}$$
$`(25)`$
The functions $`k(\xi )`$ and $`K(\xi ,\alpha _I)`$ are defined by the following equations:
$$\mathrm{ln}k(\xi )=_0^{\mathrm{}}\frac{dt}{t}\left((\frac{q}{2}x)^2e^{2t}\frac{\text{sinh}^2(\frac{q}{2}x)t}{\text{sinh}bt\text{sinh}ct}\right)\text{and}K(\xi ,\alpha _I)=k^{}(0)\underset{I=1}{\overset{3}{}}\frac{k(2\alpha _I)}{k(\xi +2\alpha _I)}$$
$`(26)`$
To extrapolate $`𝐂_3`$ to other values of $`\xi `$ we note that although the arguments of $`k`$’s are constrained by the relation $`_I\alpha _I=q\xi _{mn}`$, they range over the whole real axis for fixed $`\xi _{mn}`$. Hence the only reasonable extrapolation is the obvious one, $`k(\xi _{mn},\alpha _I)k(\xi ,\alpha _I)`$, in which case $`K(\xi _{mn},\alpha _I)K(\xi ,\alpha _I)`$. Unfortunately, this argument is not valid for the rest of $`𝐂_3(\xi _{mn},m+n)`$, which depends only on $`\xi _{mn}`$. Thus the most general extrapolation of (24) is
$$𝐂_3(iu,iv,\xi )=\text{cosh}\pi u\text{cosh}\pi v\mathrm{\Gamma }(1+iu)\mathrm{\Gamma }(1+iv)\lambda ^\xi \mathrm{\Omega }^{i(u+v)}\frac{K(\xi ,\alpha _I)}{k^{}(\xi )}f(\xi ,u+v)\frac{x_{IJ}}{L}^{\mathrm{\Delta }_{(IJ)}}$$
$`(27)`$
where the cosh terms take care of the $`(1)^{m+n}`$ terms and $`f(\xi ,u+v)`$ is an arbitary function with $`f(\xi =mb+nc,m+n)=1`$. Then $`𝒢_3`$ becomes
$$𝒢_3=\lambda ^\xi h(\xi )\frac{K(\xi ,\alpha _I)}{k(\xi )}\underset{I=1}{\overset{3}{}}\frac{x_{IJ}}{L}^{\mathrm{\Delta }_{(IJ)}(\alpha _I)}$$
$`(28)`$
where
$$h(\xi )=\frac{k(\xi )}{k^{}(\xi )}Z_0(\xi )\text{and}Z_0(\xi )=\frac{dudv\delta (\xi ibuicv)}{\text{tanh}\pi u\text{tanh}\pi v}\mathrm{\Omega }^{i(u+v)}f[\xi ,i(u+v)]$$
$`(29)`$
Since the function $`f`$ is unknown, we cannot proceed from (29). Hence we take an alternative route using the two-point function.
Uniqueness (An Application of the Sum Rule): For two and three point functions the sum rule (21) in the infinite volume limit is
$$b\mu _bd^2x𝒢_3(x_I,x,\alpha _I,b)+c\mu _cd^2x𝒢_3(x_I,x,\alpha _I,c)=\xi _2𝒢_2(x_I,\alpha _I),\xi _20$$
$`(30)`$
This is not useful unless we have an alternative definition for the two-point function. Such a definition may be obtained by regarding it as twice$`^2`$ Actually any constant independent of $`\xi `$ is allowed a priori, but the requirement that the sum rule be satisfied at the points $`\xi =mb+nc`$ fixes the constant to be two. the limit of the three-point function as $`\alpha _30`$ and $`x_3\mathrm{}`$. It is easy to see that the limit is non-zero only if $`\mathrm{\Delta }_1=\mathrm{\Delta }_2`$ i.e if $`\alpha _1=\alpha _2`$ or $`\alpha _1=q\alpha _2`$. Since for scattering states $`\alpha _I=\frac{q}{2}+i\beta _I`$, the quantities $`\alpha _1\alpha _2`$ and $`\xi _2=q\alpha _1\alpha _2`$ are pure imaginary, and we obtain
$$𝒢_2(\alpha _1,\alpha _2;x_{12})=4\pi \lambda ^{\xi _2}\frac{x_{12}}{L}^{(\mathrm{\Delta }_1+\mathrm{\Delta }_2)}h(\xi _2)R(\xi _2)[\delta (\beta _1\beta _2)+\delta (\beta _1+\beta _2)]$$
$`(31)`$
where $`R(\xi )=k(\xi )/k(\xi )`$.
To compute the left hand side of (30) we note that the first integral is
$$\lambda ^{\xi _2b}h(\xi _2b)\frac{K(\alpha _1,\alpha _2,b)}{k(\xi _2+b)}|\frac{x_{12}}{L}|^{2(\mathrm{\Delta }_b\mathrm{\Delta }_1\mathrm{\Delta }_2)}\times \zeta $$
$`(32)`$
where
$$\zeta =d^2x_3|\frac{x_{31}}{L}|^{2(\mathrm{\Delta }_b+D)}|\frac{x_{23}}{L}|^{2(\mathrm{\Delta }_bD)}=2\pi ^2\frac{x_{12}}{L}^{2\mathrm{\Delta }_b}\delta (D)$$
$`(33)`$
and $`\mathrm{\Delta }_b=1`$ and $`D=\mathrm{\Delta }_1\mathrm{\Delta }_2=\beta _1^2\beta _2^2`$. The delta functions in these equations and the condition $`\xi _20`$ in (30) mean that we only need the coefficient for $`\alpha _1=\alpha _2\alpha `$, which is easily computed to be
$$\frac{K(\alpha ,\alpha ,b)}{k(\xi _2+b)}=\xi _2^2k^{}(0)R(\xi _2)\frac{k(2b)}{k^2(b)}=\xi _2^2R(\xi _2)\frac{\mathrm{\Phi }_b}{\pi b},\xi _2=q2\alpha $$
$`(34)`$
Inserting these formulae into the sum rule (30) and using the identity $`\lambda ^b=\mu _b\mathrm{\Phi }_b\mathrm{\Omega }`$, and similarly for $`c`$, we get the sum rule
$$h(\xi _2+b)+h(\xi _2+c)=2\mathrm{\Omega }h(\xi _2)\text{with}h(mb+nc)=\mathrm{\Omega }^{m+n}$$
$`(35)`$
for the unknown function $`h(\xi _2)`$. The second equation is obtained by explicit computation from (24). The question is whether the equations in (35) determine $`h(\xi _2)`$ uniquely. We show that they do subject to a mild technical assumption to be introduced presently. First let us restrict ourselves to the case when $`b`$ is fractional i.e. $`b=r/s`$ where $`r>s`$ are positive integers with no common factor and lowest common multiple $`rs`$. Next we rewrite the sum rule (35) in terms of new variables defined as follows:
$$ye^\xi h(e^{\frac{r}{s}}y)+h(e^{\frac{s}{r}}y)=2\mathrm{\Omega }h(y)$$
$`(36)`$
We now make the assumption<sup>3</sup> It is probable that this could be actually proved directly from the sum rule (35) which rules out a whole class of functions.
$$\underset{y0}{lim}\frac{h\left(e^{(\frac{r}{s})^{\pm 1}}y\right)}{h(y)}C_\pm <\mathrm{}$$
$`(37)`$
Considering the restriction
$$y_t=e^{\frac{t}{rs}}h(e^{\frac{tr^2}{rs}})+h(e^{\frac{ts^2}{rs}})=2\mathrm{\Omega }h(e^{\frac{t}{rs}}),\text{where}tZ_+$$
$`(38)`$
we see that the symmetry group we need to implement is the dilatation group. Thus we may expand a solution of (38) as follows:
$$h(y_t)=\underset{p=0}{\overset{N}{}}C_py_t^{\sigma _p}\text{where}e^{r^2\omega \sigma _p}+e^{s^2\omega \sigma _p}2\mathrm{\Omega }=0\text{and}\omega =\frac{1}{rs}$$
$`(39)`$
$`N(0Nr^2)`$ being a finite number. This follows because the sum rule restricts the number of $`C_p`$s to be the dimension of the solution space of the polynomial equation in (39). This is of course true only on the first sheet of the covering of the $`y`$ variable. In the general case we would also have a sum over the infinite number of coverings.
$$h(y)=\underset{p=0}{\overset{N}{}}\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}C_{p_n}y^{\sigma _{p_n}}\text{where}\sigma _{p_n}=\sigma +2\pi inrs$$
$`(40)`$
We now use the second equation in (35) to write, for the special values $`\xi =mb+nct=mr^2+ns^2`$,
$$\underset{p=0}{\overset{N}{}}C_py_t^{\sigma _p}=\mathrm{\Omega }^{m+n}C_N=1,\sigma _N=0,C_{pN}=0\mathrm{\Omega }=1$$
$`(41)`$
The latter follows as the only consistent solution since the finite sum on the left hand side cannot in general match the right hand side which can be made arbitrarily large by letting $`m`$ and $`n`$ tend to infinity. Note that the covering does not play a role in the discussion of (41). It then follows that
$$h(\frac{t}{rs})=1tZ_+$$
$`(42)`$
Thus for $`b=r/s`$ the function $`h(\xi )`$ is unity when restricted to the subset of points $`\xi =t/rs`$.
This is as far as we can go for a fixed $`b=r/s`$. However we now invoke the fact that (29) implies that $`h(\xi )`$ is continuous in $`b`$ and $`\xi `$. We then define $`\stackrel{~}{b}=\stackrel{~}{r}/\stackrel{~}{s}`$ where $`\stackrel{~}{r}=lr+1`$ and $`\stackrel{~}{s}=ls`$. It follows that $`\stackrel{~}{r}`$ and $`\stackrel{~}{s}`$ have no common factor, $`\stackrel{~}{b}b=1/ls`$ and, by applying the above result to the tilde-variables, $`h(t/l^2sr(1+1/lr)=1`$. But this means that in any $`1/l`$ neighbourhood of $`b`$ there is a $`\stackrel{~}{b}`$ for which $`h(\xi )=1`$ at points which are separated by distances of order $`1/l^2`$. As these distances tend to zero as $`l`$ tends to infinity, we see that this is compatible with the continuity of $`h(\xi )`$ in $`b`$ and $`\xi `$ only if $`h(\xi )=1`$ for all $`\xi `$.
Reflection symmetry: Once $`h(\xi )=1`$ it follows that the denominator in the three-point function is invariant under the reflection $`\alpha _Iq\alpha _I`$ for each $`I`$ and thus the three-point function is covariant with respect to reflection symmetry in the sense that
$$𝒢_3(q\alpha _1,\alpha _2,\alpha _3)=R(q2\alpha _1)𝒢_3(\alpha _1,\alpha _2,\alpha _3)$$
$`(43)`$
where the prefactor depends only on the reflected parameter $`\alpha _1`$. Thus in the two-exponential theory, reflection covariance is an output rather than an input.
It is interesting to note how this reflection covariance expresses itself in terms of the two-point function defined. If we interpret the two-point function normalised by the volume factor $`4\pi `$, as the inner product of primary states
$$<\alpha _1,\alpha _2>\underset{x0}{lim}\frac{1}{4\pi }x^{2\mathrm{\Delta }_\alpha }𝒢_2(q\alpha _1,\alpha _2;x)$$
$`(44)`$
we have
$$\left(\begin{array}{cc}<\alpha ,\alpha >& <\alpha ,q\alpha >\\ <q\alpha ,\alpha >& <q\alpha ,q\alpha >\end{array}\right)=\left(\begin{array}{cc}1& R^1(q2\alpha )\\ R(q2\alpha )& 1\end{array}\right)\delta (0)$$
$`(45)`$
It is clear that the matrix in (45) is hermitian and has zero determinant. Hence one linear combination of the states, namely $`|\alpha >R(q2\alpha )|q\alpha >`$, has zero norm and decouples. Thus effectively,
$$|\alpha >=R(q2\alpha )|q\alpha >$$
$`(46)`$
which means that there is actually only one physical state for each conformal weight $`\mathrm{\Delta }_\alpha `$. This may seem surprising but from (43) it is seen to be a manifestation of the reflection covariance.
Comparison of the one and two-exponential path integrals: In order to compare the one and two exponential theories, we begin by recalling that $`𝐂_3`$ in both the theories is defined in terms of the correlation functions of vertex operators in a free field theory. In the two-exponential theory, the relevant integral is given by
$$𝑑\varphi e^{S(\varphi )}U_b^{iu}U_c^{iv}𝚷_3=\mathrm{\Gamma }(1+iu)\mathrm{\Gamma }(1+iv)\frac{K(\xi ,\alpha _I)}{k^{}(\xi )}f[\xi ,i(u+v)]$$
$`(47)`$
where $`S(\varphi )`$ is the Action for the free theory. The corresponding equation, in the one-exponential theory, is obtained by letting $`v0`$ and takes the form
$$𝑑\varphi e^{S(\varphi )}U_b^{iu}𝚷_3=\mathrm{\Gamma }(1+iu)\frac{K(\xi ,\alpha _I)}{k^{}(\xi )}\times f[\xi ,iu]$$
$`(48)`$
The integral corresponding to the zero-mode integral in (29) in the one-exponential theory for an a priori arbitrary $`f[\xi ,iu]`$ can be performed to yield $`f[\xi ,c\xi ]/\text{tanh}\pi c\xi `$. If we assume that the one-exponential theory leads to the DOZZ Ansatz, then the function $`f[\xi ,iu]`$ is determined uniquely to be $`b\text{tan}(\pi u)k^{}(ibu)/k(ibu)`$ for $`\xi =ibu`$. We may now ask, what choice, if any, for the function $`f[\xi ,i(u+v)]`$ in the two-exponential theory will produce the DOZZ Ansatz. It is easy to see that if we choose
$$f[\xi ,i(u+v)]=\frac{k^{}(\xi )}{k(\xi )}[b\text{tanh}\pi u+c\text{tanh}\pi v]$$
$`(49)`$
the $`Z_0`$ integral in (29) produces $`k^{}/k`$ in accordance with the DOZZ Ansatz. This choice has the virtue that it reduces to the one-exponential result in the limit $`v0`$. However, one can easily convince oneself that there exist other choices for the function $`f[\xi ,i(u+v)]`$. If we choose, for example, $`f[\xi ,i(u+v)]=\text{sgn}[i(u+v)]`$ (the sgn-factor is necessary to preserve the symmetry of the path integral under a change of sign of $`b,c`$ and the $`\alpha `$’s) and integrate over $`u`$ and $`v`$ first, we obtain,
$$Z_0(\xi )=_0^{\mathrm{}}𝑑t\frac{\text{sinh}((q2\xi )t)}{\text{sinh}(bt)\text{sinh}(ct)},t|\varphi _0|$$
$`(50)`$
Up to a regulating term,$`^4`$ The regulating term can be obtained naturally by modifying the path-integral measure $`d\varphi _0`$ to $`lim_{ϵ0}d\varphi _0\left[1\text{exp}(|\varphi _0|/ϵ)\right]`$. this will be recognized, from the definition of $`k(\xi )`$ in (26), as $`k^{}(\xi )/k(\xi )`$. Since the choice for the extrapolation function is not unique we conclude that any result that is based on an extrapolation is not conclusive. It is therefore desirable to have a direct proof of the DOZZ result for the three-point function without any reference to the extrapolation. This is exactly what was achieved in this paper by showing that the sum rule in the two-exponential theory has a unique (constant) solution. In contrast, in the one-exponential theory, if the specific extrapolation leading to the DOZZ Ansatz is not made, the final result can only be obtained up to a periodic function of $`\xi `$.
Acknowledgements: We thank I. Sachs, J. Teschner and P. Watts for useful discussions.
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# Deriving the Nonlinear Cosmological Power Spectrum and Bispectrum from Analytic Dark Matter Halo Profiles and Mass Functions
## 1 Introduction
Two conceptual pictures of galaxy clustering have been examined in the literature, the continuous hierarchical clustering model and the power-law cluster model (Peebles 1980, §61). In the hierarchical clustering model, which has emerged as the accepted model over the past two decades, galaxy clustering is characterized by power-law correlation functions: the $`N`$-point correlation function $`\xi _N`$ scales with configuration size as $`\xi _Nr^{\gamma _N}\xi _2^{(N1)}`$, where $`\gamma _N=(N1)\gamma `$ and the two-point correlation function goes as $`\xi _2=\xi r^\gamma `$. The hierarchical model is motivated by the observed power-law behavior $`\gamma 1.8`$ of galaxy correlations (Groth & Peebles 1977; Fry & Peebles 1978), with a theoretical basis in a self-similar, scale-invariant solution to the equations of motion (Davis & Peebles 1977).
The alternative power-law cluster model has an even longer history (Neyman & Scott 1952; Peebles 1974, 1980; McClelland & Silk 1977; Scherrer & Bertschinger 1991; Sheth & Jain 1997; Valageas 1998; Yano & Gouda 1999). In this model, galaxies are placed in spherical clumps that are assumed to follow a power-law density profile $`\rho (r)r^ϵ`$, with the centers of the clumps distributed randomly. The resulting two-point correlation function is also a power law with a logarithmic slope $`\gamma =2ϵ3`$. While it is possible to reproduce the observed two-point function by an appropriate choice of the power index $`ϵ=(3+\gamma )/22.4`$, Peebles and Groth (1975) pointed out that this model produces a three-point function that is too steep to be consistent with observations in the Zwicky and Lick catalogs.
In an earlier paper (Ma & Fry 2000a), we have shown that in the nonlinear regime, the three-point correlation function $`\zeta =\xi _3`$ of the cosmological mass density field does not exactly follow the prediction $`\zeta \xi ^2`$ of the hierarchical clustering model. These conclusions are drawn from study of high resolution numerical simulations of a cold dark matter (CDM) model with cosmological constant and of a model with scale-free initial conditions $`P(k)k^n`$ with $`n=2`$. In experiments replacing simulation dark matter halos with power-law density profiles, $`\rho (r)r^ϵ`$, we have demonstrated that the behavior of the correlation functions in the nonlinear regime are determined by the halo profiles, but that it is not possible to match both the two- and three-point correlations with a single slope $`ϵ`$. These results differ from the predictions of both of these two conceptual models.
In this paper, we expand our previous study of the nonlinear two- and three-point correlation functions by investigating a new prescription that takes into account the non-power-law profiles of halos, the distribution of halo masses, and the spatial correlations of halo centers. Each of these ingredients has been well studied in the literature. We find that this halo model provides a good description of the two- and three-point correlation functions in both the $`n=2`$ and CDM simulations over the entire range of scales from the weak clustering, perturbative regime on large length scales, to the strongly nonlinear regime on small length scales. Our result is approximately hierarchical over an intermediate range of scales, thus uniting the two pictures. An independent recent study by Seljak (2000), which appeared during completion of this work, has also examined the two-point power spectrum in a similar construction and has found that this type of approach can reproduce the power spectrum in the CDM model. The analytic model proposed here can be used to compute the two- and three-point correlation functions and their Fourier transforms, the power spectrum and bispectrum, over any range of scale where the input halo properties are valid.
In a subsequent paper (Ma & Fry 2000c), we study the predictions of this analytic halo model for the asymptotic nonlinear behavior of the $`N`$-point correlation functions and the pairwise velocities and examine the conditions required for stable clustering.
The outline of this paper is as follows. In §2 we describe the three input ingredients of the model: halo density profiles, halo mass functions, and halo-halo correlations. In §3 we assemble these ingredients and construct analytic expressions for the two-point correlation function $`\xi (r)`$ and the power spectrum $`P(k)`$. In §4 we do the same for the three-point correlation function $`\zeta (r_1,r_2,r_3)`$ and its Fourier transform, the bispectrum $`B(k_1,k_2,k_3)`$. In §5 we test the validity of this new model by comparing its predictions with results from numerical simulations of an $`n=2`$ scale free model and a low-density CDM model with a cosmological constant ($`\mathrm{\Lambda }`$CDM). We also present results of the synthetic halo replacement technique used to enhance the numerical resolution. In §6 we discuss further the physical meanings and implications of the model. In particular, we elaborate on two important implications of this model: deviations from the common assumptions of stable clustering and hierarchical clustering. Section 7 is a summary.
## 2 Model Ingredients
### 2.1 Halo Mass Density Profile
It has been suggested recently that the mass density profiles of cold dark matter halos have a roughly universal shape, generally independent of cosmological parameters (Navarro, Frenk, & White 1996, 1997)
$$\frac{\rho (r)}{\overline{\rho }}=\overline{\delta }u(r/R_s),$$
(1)
where $`\overline{\delta }`$ is a dimensionless density amplitude, $`R_s`$ is a characteristic radius, and $`\overline{\rho }`$ is the mean background density. We consider two functional forms for the density profiles
$`u_I(x)`$ $`=`$ $`{\displaystyle \frac{1}{x^p(1+x)^{3p}}},`$
$`u_{II}(x)`$ $`=`$ $`{\displaystyle \frac{1}{x^p(1+x^{3p})}}.`$ (2)
Both forms have asymptotic behaviors $`x^p`$ at small $`x`$ and $`x^3`$ at large $`x`$, but they differ in the transition region. The first form $`u_I(x)`$ with $`p=1`$ is found to provide a good fit to simulation halos by Navarro et al. (1996, 1997), whereas the second form $`u_{II}(x)`$ with a steeper inner slope $`p=3/2`$ is favored by Moore et al. (1998, 1999). Some independent simulations have produced halos that are well fit by the shallower $`p=1`$ inner slope (e.g., Hernquist 1990; Dubinski & Carlberg 1991; Huss, Jain, & Steinmetz 1999), and others the steeper $`p>1`$ slope (e.g., Fukushige and Makino 1997). Jing & Suto (2000) have recently reported a mass-dependent inner slope, with $`p1.5`$ for galactic-mass halos and $`p1`$ for cluster-mass halos. Many of these authors find that the outer profile scales as $`r^3`$, but steeper outer profiles have also been suggested (Hernquist 1990; Dubinski & Carlberg 1991). Given these uncertainties, we will consider in this paper both types of profiles in equation (2).
The parameters $`R_s`$ and $`\overline{\delta }`$ in equation (1) are generally functions of the halo mass $`M`$. A concentration parameter,
$$c=\frac{R_{200}}{R_s},$$
(3)
can be used to quantify the central density of a halo (Navarro et al. 1997), where $`R_{200}`$ is the radius within which the average density is 200 times the mean density of the universe. Using $`M=800\pi \overline{\rho }R_{200}^3/3`$, we can relate $`R_s`$ and $`\overline{\delta }`$ to $`M`$ and $`c`$, where the scale radius $`R_s`$ is
$$R_s=\frac{1}{c}\left(\frac{3M}{800\pi \overline{\rho }}\right)^{1/3}=\frac{1.63\times 10^5}{\mathrm{\Omega }_m^{1/3}c}\left(\frac{M}{h^1M_{}}\right)^{1/3}h^1\mathrm{Mpc},$$
(4)
and the density amplitude $`\overline{\delta }`$ is
$`\overline{\delta }_I`$ $`=`$ $`{\displaystyle \frac{200c^3}{3[\mathrm{ln}(1+c)c/(1+c)]}},p=1,`$
$`\overline{\delta }_{II}`$ $`=`$ $`{\displaystyle \frac{100c^3}{\mathrm{ln}(1+c^{3/2})}},p={\displaystyle \frac{3}{2}}.`$ (5)
Typical values of $`c`$ are in the range of a few to ten for type I and perhaps a factor of three smaller for type II. There is a weak dependence on mass, such that less massive halos have a larger central density (e.g., Cole & Lacey 1996; Tormen, Bouchet, & White 1997; Navarro et al. 1996, 1997; Jing & Suto 2000). This is understood in general terms as reflecting the mean density at the redshift $`z_f`$ when the halo initially collapsed, $`\overline{\delta }(1+z_f)^3`$. For $`\mathrm{\Omega }=1`$ this is $`c\sigma (M)`$, or $`cM^{(3+n)/6}`$ in a scale-free model.
### 2.2 Halo Mass Function
The number density of halos with mass $`M`$ within a logarithmic interval is often approximated by the prescription of Press & Schechter (1974),
$$\frac{dn}{d\mathrm{ln}M}=\sqrt{\frac{2}{\pi }}\frac{d\mathrm{ln}\sigma ^1}{d\mathrm{ln}M}\frac{\overline{\rho }}{M}\nu e^{\nu ^2/2},\nu =\frac{\delta _c}{\sigma (M)},$$
(6)
where $`\delta _c`$ is a parameter characterizing the linear overdensity at the onset of gravitational collapse, and $`\sigma `$ is the linear rms mass fluctuations in spheres of radius $`R`$
$$\sigma ^2(M)=_0^{\mathrm{}}\frac{4\pi k^2dk}{(2\pi )^3}P(k)W^2(kR),$$
(7)
where $`W(x)=3(\mathrm{sin}xx\mathrm{cos}x)/x^3`$ is the Fourier transform of a real-space tophat window function. The mass $`M`$ is related to $`R`$ by $`M=4\pi \overline{\rho }R^3/3`$. For scale free models with a power law initial power spectrum $`Pk^n`$, this is $`\sigma =(M/M_{})^{(3+n)/6}`$. The parameter $`M_{}`$ characterizes the mass scale at the onset of nonlinearity, $`\sigma (M_{})=1`$, and is related to the nonlinear wavenumber $`k_{\mathrm{nl}}`$ (defined as \]$`_0^{k_{\mathrm{nl}}}4\pi k^2𝑑kP(k)/(2\pi )^3=1`$) by
$$M_{}=\frac{4\pi \overline{\rho }}{3}\frac{B(n)}{k_{\mathrm{nl}}^3}=\frac{4\pi \overline{\rho }}{3}R_{}^3,$$
(8)
where
$`B(n)`$ $`=`$ $`(k_{\mathrm{nl}}R_{})^3=\left[(n+3){\displaystyle _0^{\mathrm{}}}𝑑xx^{n+2}W^2(x)\right]^{3/(n+3)},`$
$`B^{(3+n)/3}`$ $`=`$ $`\mathrm{sin}\left[(n+2){\displaystyle \frac{\pi }{2}}\right]\mathrm{\Gamma }(n+2){\displaystyle \frac{9(2^n)(3+n)}{(n)(1n)(3n)}}`$ (9)
(defined for $`3n<1`$). Various modifications to the Press-Schechter mass function have been suggested (e.g., Sheth & Tormen 1999; Lee & Shandarin 1999; Jenkins et al. 2000) to improve the accuracy of the original formula.
### 2.3 Halo-Halo Correlations
Dark matter halos do not cluster in the same way as the mass density field. On large scales, a bias parameter $`b`$ is typically used to quantify this difference. Let $`\xi _{\mathrm{halo}}(r;M,M^{})`$ be the two-point correlation function of halos with masses $`M`$ and $`M^{}`$, $`\xi _{\mathrm{lin}}(r)`$ be the linear correlation function for the mass density field, and $`P_{\mathrm{halo}}`$ and $`P_{\mathrm{lin}}`$ be the corresponding power spectra. On large length scales, we assume a linear bias and write
$`\xi _{\mathrm{halo}}(r;M,M^{})`$ $`=`$ $`b(M)b(M^{})\xi _{\mathrm{lin}}(r),`$
$`P_{\mathrm{halo}}(k;M,M^{})`$ $`=`$ $`b(M)b(M^{})P_{\mathrm{lin}}(k).`$ (10)
Based on the peak and the Press-Schechter formalism, Mo & White (1996) developed a model for the linear bias $`b(M)`$, which is later modified by Jing (1998) to be
$$b(M)=\left(1+\frac{\nu ^21}{\delta _c}\right)\left(\frac{1}{2\nu ^4}+1\right)^{0.060.02n},\nu =\frac{\delta _c}{\sigma (M)}.$$
(11)
The original formula for $`b(M)`$ by Mo & White includes only the first factor above; the second factor, dependent on the primordial spectral index $`n`$, is obtained empirically for an improved fit to simulation results at the lower mass end (Jing 1998). In this bias model, $`b(M)`$ is below unity for $`MM_{}`$ (where $`\sigma (M_{})=1`$) and reaches $`0.5`$ for $`M0.01M_{}`$. Small dark matter halos are therefore anti-biased relative to the mass density. For $`MM_{}`$, $`b(M)`$ increases monotonically with the halo mass and reaches $`b10`$ at $`M100M_{}`$. Nonlinear effects on the bias have been studied (Kravtsov & Klypin 1999 and references therein), but they are unimportant in our model because the halo-halo correlation terms contribute significantly only on large length scales in the linear regime (see §3 and 4).
Similarly, we use higher order bias parameters to relate the higher-order correlation functions for halos and mass density. In this paper we examine the three-point correlation function $`\zeta (r_1,r_2,r_3)`$ and its Fourier transform, the bispectrum $`B(k_1,k_2,k_3)`$ (see §4 for a more detailed discussion). On large length scales where the amplitude of $`\delta `$ is small, perturbation theory can be used to relate the lowest order contribution to the bispectrum of the mass density to the linear power spectrum $`P_{\mathrm{lin}}`$ (Fry 1984):
$`B^{(0)}(k_1,k_2,k_3)`$ $`=`$ $`F_{12}P_{\mathrm{lin}}(k_1)P_{\mathrm{lin}}(k_2)+F_{23}P_{\mathrm{lin}}(k_2)P_{\mathrm{lin}}(k_3)+F_{31}P_{\mathrm{lin}}(k_3)P_{\mathrm{lin}}(k_1),`$
$`F_{ij}`$ $`=`$ $`{\displaystyle \frac{10}{7}}+(k_i/k_j+k_j/k_i)(\widehat{\text{k}}_i\widehat{\text{k}}_j)+{\displaystyle \frac{4}{7}}(\widehat{\text{k}}_i\widehat{\text{k}}_j)^2.`$ (12)
Using this perturbative result and the results of Mo, Jing, & White (1997), we can write the halo bispectrum as
$`B_{\mathrm{halo}}(k_1,k_2,k_3;M,M^{},M^{\prime \prime })`$ $`=`$ $`\left[b(M)b(M^{})b(M^{\prime \prime })F_{12}+b(M)b(M^{})b_2(M^{\prime \prime })\right]P_{\mathrm{lin}}(k_1)P_{\mathrm{lin}}(k_2)`$
$`+`$ $`\left[b(M)b(M^{})b(M^{\prime \prime })F_{23}+b(M)b_2(M^{})b(M^{\prime \prime })\right]P_{\mathrm{lin}}(k_2)P_{\mathrm{lin}}(k_3)`$
$`+`$ $`\left[b(M)b(M^{})b(M^{\prime \prime })F_{31}+b_2(M)b(M^{})b(M^{\prime \prime })\right]P_{\mathrm{lin}}(k_3)P_{\mathrm{lin}}(k_1),`$
where $`b(M)`$ is given by equation (11), and the quadratic bias parameter $`b_2(M)`$ is
$$b_2(M)=\frac{8}{21}\frac{(\nu ^21)}{\delta _c}+\left(\frac{\nu }{\delta _c}\right)^2(\nu ^23).$$
(14)
For the special equilateral case of $`k_1=k_2=k_3=k`$, equation (LABEL:Bbias) simplifies to
$`B_{\mathrm{halo}}^{\mathrm{eq}}(k;M,M^{},M^{\prime \prime })`$ $`=`$ $`[{\displaystyle \frac{12}{7}}b(M)b(M^{})b(M^{\prime \prime })+b(M)b(M^{})b_2(M^{\prime \prime })`$ (15)
$`+`$ $`b(M)b_2(M^{})b(M^{\prime \prime })+b_2(M)b(M^{})b(M^{\prime \prime })]P_{\mathrm{lin}}^2(k).`$
In practice, the terms involving $`b_2(M)`$ in equations (LABEL:Bbias) and (15) make only a small net contribution. For simplicity, we will therefore not include this term in the subsequent derivations and calculations.
## 3 Two-Point Statistics: $`\xi (r)`$ and $`P(k)`$
We now construct our analytic halo model for the two-point correlation function $`\xi (r)`$ and the power spectrum $`P(k)`$. The two-point correlation function of the cosmological mass density field $`\delta =\delta \rho /\overline{\rho }`$ is
$$\xi (\text{r})=\delta (\text{x})\delta (\text{x}+\text{r}).$$
(16)
The Fourier transform of $`\xi (r)`$ is the mass power spectrum $`P(k)=d^3re^{i\text{k}\text{r}}\xi (r)`$, which is related to the density field in $`k`$-space by $`\delta (\text{k}_1)\delta (\text{k}_2)=P(k_1)(2\pi )^3\delta _D(\text{k}_1+\text{k}_2)`$, where $`\delta _D`$ is the Dirac delta-function.
The two-point correlation function measures the excess probability above the Poisson distribution of finding a pair of objects with separation $`r`$ (Peebles 1980). The objects can be taken to be dark matter particles, most of which cluster gravitationally in the form of dark matter halos. One should therefore be able to express $`\xi `$ for the density field in terms of properties of dark matter halos. In this picture, we can write the contributions to $`\xi `$ as two separate terms, one from particle pairs in the same halo, and the other from pairs that reside in two different halos. In realistic cosmological models, dark matter halos exhibit a spectrum of masses that can be characterized by a distribution function $`dn/dM`$, and the halo centers are spatially correlated. Taking these factors into consideration, we can write the two-point correlation function for $`\delta `$ in terms of the halo density profile $`u(x)`$, halo mass function $`dn/dM`$, and halo-halo correlation function $`\xi _{\mathrm{halo}}`$ discussed in §2. We write
$$\xi (r)=\xi _{1h}(r)+\xi _{2h}(r),$$
(17)
where the subscripts “$`1h`$” and “$`2h`$” denote contributions from particle pairs in “1-halo” and “2-halos”, respectively, and
$`\xi _{1h}(r)`$ $`=`$ $`{\displaystyle d^3r^{}𝑑M\frac{dn}{dM}\overline{\delta }^2u(r^{}/R_s)u(|\text{r}^{}+\text{r}|/R_s)}`$
$`\xi _{2h}(r)`$ $`=`$ $`{\displaystyle d^3r^{}d^3r^{\prime \prime }𝑑M^{}\frac{dn}{dM^{}}\overline{\delta }^{}u(r^{}/R_s^{})𝑑M^{\prime \prime }\frac{dn}{dM^{\prime \prime }}\overline{\delta }^{\prime \prime }u(r^{\prime \prime }/R_s^{\prime \prime })\xi _{\mathrm{halo}}(|\text{r}^{}\text{r}^{\prime \prime }+\text{r}|)}`$ (18)
$`=`$ $`{\displaystyle d^3r^{}d^3r^{\prime \prime }\left[𝑑M\frac{dn}{dM}\overline{\delta }u(r^{}/R_s)b(M)\right]\left[𝑑M\frac{dn}{dM}\overline{\delta }u(r^{\prime \prime }/R_s)b(M)\right]}`$
$`\times \xi _{\mathrm{lin}}(|\text{r}^{}\text{r}^{\prime \prime }+\text{r}|).`$
These expressions arise from averaging over displacements $`\text{r}^{}`$, $`\text{r}^{\prime \prime }`$ of halo centers from the particle positions $`\text{r}_1`$, $`\text{r}_2`$, where $`r=|\text{r}_1\text{r}_2|`$. In the last expression above, we have used the bias model of equation (10) to relate the halo-halo correlation function $`\xi _{\mathrm{halo}}`$ to the linear correlation function $`\xi _{\mathrm{lin}}`$ of the mass density field.
As we will show in §5, the dominant contribution to the two-point correlation function in the nonlinear regime on small length scales is from the first, 1-halo term $`\xi _{1h}`$ for particle pairs that reside in the same halos. This makes intuitive sense, because closely spaced particle pairs are most likely to be found in the same halo. This term is determined by the convolution of the dimensionless density profile with itself,
$$\lambda (x)=d^3yu(y)u(|\text{x}+\text{y}|).$$
(19)
For many forms of $`u(x)`$, the angular integration in this equation is analytic, and $`\lambda `$ can be reduced to a simple one-dimensional integral over $`y`$. For some special cases, $`\lambda `$ can even be reduced to an analytic expression. We leave the detailed results for $`\lambda `$ to the Appendix.
In $`k`$-space, the convolutions in equation (18) for $`\xi (r)`$ become simple products. Using $`\stackrel{~}{u}(q)`$ to denote the Fourier transform of $`u(x)`$, where $`\stackrel{~}{u}(q)=d^3xu(x)e^{i\text{q}\text{x}}`$, we can readily transform equation (18) into expressions for the mass power spectrum:
$$P(k)=P_{1h}(k)+P_{2h}(k),$$
(20)
where the 1-halo and 2-halo terms are
$`P_{1h}(k)`$ $`=`$ $`{\displaystyle 𝑑M\frac{dn}{dM}[R_s^3\overline{\delta }\stackrel{~}{u}(kR_s)]^2}`$
$`P_{2h}(k)`$ $`=`$ $`{\displaystyle 𝑑M\frac{dn}{dM}R_s^3\overline{\delta }\stackrel{~}{u}(kR_s)𝑑M^{}\frac{dn}{dM^{}}R_s^3\overline{\delta }^{}\stackrel{~}{u}(kR_s^{})P_{\mathrm{halo}}(k)}`$
$`=`$ $`\left[{\displaystyle 𝑑M\frac{dn}{dM}R_s^3\overline{\delta }\stackrel{~}{u}(kR_s)b(M)}\right]^2P_{\mathrm{lin}}(k).`$
To arrive at the last expression above, we have again used the bias model of equation (10). For computational efficiency, we find that the algebraic expressions
$`\stackrel{~}{u}_I(q)`$ $`=`$ $`{\displaystyle \frac{4\pi \{\mathrm{ln}(e+1/q)\mathrm{ln}[\mathrm{ln}(e+1/q)]/3\}}{(1+q^{1.1})^{(2/1.1)}}},p=1`$
$`\stackrel{~}{u}_{II}(q)`$ $`=`$ $`{\displaystyle \frac{4\pi \{\mathrm{ln}(e+1/q)+0.25\mathrm{ln}[\mathrm{ln}(e+1/q)]\}}{1+0.8q^{1.5}}},p={\displaystyle \frac{3}{2}}`$ (22)
provide excellent fits for the profiles of Navarro et al. (1997) and Moore et al. (1999), with less than 4% rms error for form I and less than 1% rms error for form II. The functional form is chosen to reproduce the asymptotic behaviors: $`\stackrel{~}{u}4\pi \mathrm{ln}q`$ at small $`q`$ (with no radial cutoff), and $`\stackrel{~}{u}q^2`$ (type I) and $`\stackrel{~}{u}q^{3/2}`$ (type II) at large $`q`$.
The two-point $`\xi (r)`$ and $`P(k)`$ can now be computed analytically from equations (18) and (3). The inputs are equation (2) or (22) for the halo density profile $`u(x)`$ or $`\stackrel{~}{u}(q)`$, equations (4) and (5) for $`R_s`$ and $`\overline{\delta }`$, equation (6) for the halo mass function $`dn/dM`$, and equation (10) for the halo-halo correlation function. Since the halo density profile appears to have a nearly universal form regardless of background cosmology, $`\xi (r)`$ and $`P(k)`$ depend on cosmological parameters mainly through $`\sigma (M)`$ of equation (7) and the halo concentration $`c(M)`$ or central density $`\overline{\delta }(M)`$. (See Ma & Fry 2000c for a more detailed discussion of $`c(M)`$.)
## 4 Three-Point Statistics: $`\zeta `$ and $`B`$
Here we construct our analytic halo model for the three-point correlation function $`\zeta `$ and the bispectrum $`B`$. The joint probability of finding three objects in volume elements $`dV_1,dV_2`$, and $`dV_3`$ is given by
$$dP=[1+\xi (r_1)+\xi (r_2)+\xi (r_3)+\zeta (r_1,r_2,r_3)]\overline{n}^3dV_1dV_2dV_3,$$
(23)
where $`\xi (r)`$ and $`\zeta (r_1,r_2,r_3)`$ are the two- and three-point correlation functions, respectively, $`\overline{n}`$ is the mean number density of objects, and $`r_1,r_2`$ and $`r_3`$ are the lengths of the sides of the triangle defined by the three objects (Peebles 1980). The Fourier transform of the three-point correlation function $`\zeta (r_1,r_2,r_3)`$ is the bispectrum $`B(k_1,k_2,k_3)`$, which is related to the density field in $`k`$-space by $`\delta (\text{k}_1)\delta (\text{k}_2)\delta (\text{k}_3)=B(k_1,k_2,k_3)(2\pi )^3\delta _D(\text{k}_1+\text{k}_2+\text{k}_3)`$. The bispectrum depends on any three parameters that define a triangle in $`k`$-space. A particular simple configuration to study is the equilateral triangle ($`k_1=k_2=k_3=k`$), and in this case the bispectrum $`B^{\mathrm{eq}}`$ depends only on a single wavenumber.
Similar to the two-point halo model of §3, we can write the contributions to the three-point correlation function $`\zeta `$ of the mass density as three separate terms, each term representing particle triplets that reside in a single halo, two distinct halos, or three distinct halos. Taking into account the halo mass distribution and halo-halo correlations discussed in §2, we obtain
$$\zeta (r_1,r_2,r_3)=\zeta _{1h}(r_1,r_2,r_3)+\zeta _{2h}(r_1,r_2,r_3)+\zeta _{3h}(r_1,r_2,r_3),$$
(24)
where the separate 1-halo, 2-halo, and 3-halo terms are
$`\zeta _{1h}(r_1,r_2,r_3)`$ $`=`$ $`{\displaystyle d^3r𝑑M\frac{dn}{dM}\overline{\delta }^3u(r/R_s)u(|\text{r}+\text{r}_1\text{r}_2|/R_s)u(|\text{r}+\text{r}_1\text{r}_3|/R_s)}`$
$`\zeta _{2h}(r_1,r_2,r_3)`$ $`=`$ $`{\displaystyle d^3rd^3r^{}𝑑M\frac{dn}{dM}\overline{\delta }^2u(r/R_s)u(|\text{r}+\text{r}_1\text{r}_2|/R_s)𝑑M^{}\frac{dn}{dM^{}}\overline{\delta }^{}u(r^{}/R_s^{})}`$ (25)
$`\times \xi _{\mathrm{halo}}(|\text{r}\text{r}^{}+\text{r}_1\text{r}_3|)+\text{sym.(1,2,3)}`$
$`\zeta _{3h}(r_1,r_2,r_3)`$ $`=`$ $`{\displaystyle d^3rd^3r^{}d^3r^{\prime \prime }𝑑M\frac{dn}{dM}\overline{\delta }u(r/R_s)𝑑M^{}\frac{dn}{dM^{}}\overline{\delta }^{}u(r^{}/R_s^{})}`$
$`\times {\displaystyle }dM^{\prime \prime }{\displaystyle \frac{dn}{dM^{\prime \prime }}}\overline{\delta }^{\prime \prime }u(r^{\prime \prime }/R_s^{\prime \prime })\zeta _{\mathrm{halo}}(\text{r}+\text{r}_1,\text{r}^{}+\text{r}_2,\text{r}^{\prime \prime }+\text{r}_3).`$
The dominant contribution to the three-point correlation function in the nonlinear regime is from the first term $`\zeta _{1h}`$, which comes from particle triplets that reside in the same halo. This term is determined by the convolution $`\gamma (\text{x}_1,\text{x}_2)=d^3yu(y)u(|\text{y}+\text{x}_1|)u(|\text{y}+\text{x}_2|)`$ of three factors of the density profile $`u(x)`$, and is analogous to the convolution $`\lambda `$ in equation (19) for the one-halo term $`\xi _{1h}`$ in the two-point correlation function.
The bispectrum of the mass density field $`\delta `$ in $`k`$-space can be obtained by Fourier transforming the equations above. We find
$$B(k_1,k_2,k_3)=B_{1h}(k_1,k_2,k_3)+B_{2h}(k_1,k_2,k_3)+B_{3h}(k_1,k_2,k_3),$$
(26)
where
$`B_{1h}(k_1,k_2,k_3)`$ $`=`$ $`{\displaystyle 𝑑M\frac{dn}{dM}[R_s^3\overline{\delta }\stackrel{~}{u}(k_1R_s)][R_s^3\overline{\delta }\stackrel{~}{u}(k_2R_s)][R_s^3\overline{\delta }\stackrel{~}{u}(k_3R_s)]}`$
$`B_{2h}(k_1,k_2,k_3)`$ $`=`$ $`{\displaystyle 𝑑M\frac{dn}{dM}[R_s^3\overline{\delta }\stackrel{~}{u}(k_1R_s)][R_s^3\overline{\delta }\stackrel{~}{u}(k_2R_s)]}`$ (27)
$`\times {\displaystyle }dM^{}{\displaystyle \frac{dn}{dM^{}}}R_s^3\overline{\delta }^{}\stackrel{~}{u}(k_3R_s^{})P_{\mathrm{halo}}(k_3;M,M^{})+\text{sym.(1,2,3)}`$
$`B_{3h}(k_1,k_2,k_3)`$ $`=`$ $`{\displaystyle 𝑑M\frac{dn}{dM}R_s^3\overline{\delta }\stackrel{~}{u}(k_1R_s)𝑑M^{}\frac{dn}{dM^{}}R_s^3\overline{\delta }^{}\stackrel{~}{u}(k_2R_s^{})}`$
$`\times {\displaystyle }dM^{\prime \prime }{\displaystyle \frac{dn}{dM^{\prime \prime }}}R_s^{\prime \prime 3}\overline{\delta }^{\prime \prime }\stackrel{~}{u}(k_3R_s^{\prime \prime })B_{\mathrm{halo}}(k_1,k_2,k_3;M,M^{},M^{\prime \prime }).`$
The halo-halo power spectrum $`P_{\mathrm{halo}}(k)`$ and bispectrum $`B_{\mathrm{halo}}(k_1,k_2,k_3)`$ are related to the linear mass power spectrum $`P_{\mathrm{lin}}(k)`$ by equations (10) and (LABEL:Bbias).
The expressions for the mass bispectrum above simplify considerably for the equilateral triangle configuration, and
$$B^{\mathrm{eq}}(k)=B_{1h}^{\mathrm{eq}}(k)+B_{2h}^{\mathrm{eq}}(k)+B_{3h}^{\mathrm{eq}}(k),$$
(28)
where
$`B_{1h}^{\mathrm{eq}}(k)`$ $`=`$ $`{\displaystyle 𝑑M\frac{dn}{dM}[R_s^3\overline{\delta }\stackrel{~}{u}(kR_s)]^3}`$
$`B_{2h}^{\mathrm{eq}}(k)`$ $`=`$ $`3\left[{\displaystyle 𝑑M\frac{dn}{dM}[R_s^3\overline{\delta }\stackrel{~}{u}(kR_s)]^2b(M)}\right]\left[{\displaystyle 𝑑M\frac{dn}{dM}R_s^3\overline{\delta }\stackrel{~}{u}(kR_s)b(M)}\right]P_{\mathrm{lin}}(k)`$ (29)
$`B_{3h}^{\mathrm{eq}}(k)`$ $`=`$ $`\left[{\displaystyle 𝑑M\frac{dn}{dM}R_s^3\overline{\delta }\stackrel{~}{u}(kR_s)b(M)}\right]^3{\displaystyle \frac{12}{7}}P_{\mathrm{lin}}^2(k).`$
Here we have written out explicitly the bias factors $`b(M)`$ using equations (10) and (15), and we have neglected terms with $`b_2(M)`$ as discussed in §2.3.
## 5 $`N`$-body Experiments and Numerical Results
In this section we compare the predictions of our analytical model described in §2, 3, and 4 with results from cosmological $`N`$-body simulations. We examine two cosmological models: an $`n=2`$ scale-free model and a low-density $`\mathrm{\Lambda }`$CDM model. These are the same simulations studied in Ma & Fry (2000a). The $`n=2`$ simulation has $`256^3`$ particles and a Plummer force softening length of $`L/5120`$, where $`L`$ is the box length. The $`\mathrm{\Lambda }`$CDM model is spatially flat with matter density $`\mathrm{\Omega }_m=0.3`$ and cosmological constant $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$. This run has $`128^3`$ particles and is performed in a $`(100\mathrm{Mpc})^3`$ comoving box with a comoving force softening length of $`50\mathrm{kpc}`$ for Hubble parameter $`h=0.75`$. The baryon fraction is set to zero for simplicity. The primordial power spectrum has a spectral index of $`n=1`$, and the density fluctuations are drawn from a random Gaussian distribution. The gravitational forces are computed with a particle-particle particle-mesh (P<sup>3</sup>M) code (Ferrell & Bertschinger 1994). We compute the density field $`\delta `$ on a grid from particle positions using the second-order triangular-shaped cloud (TSC) interpolation scheme. A fast Fourier transform is then used to obtain $`\delta `$ in $`k`$-space. The $`k`$-space TSC window function is deconvolved to correct for smearing in real space due to the interpolation, and shot noise terms are subtracted to correct for discreteness effects. We then compute the second and third moments of the density amplitudes in Fourier space.
We show results for the power spectrum as the dimensionless variance $`\mathrm{\Delta }(k)4\pi k^3P(k)/(2\pi )^3`$. A useful dimensionless three-point statistic is the hierarchical three-point amplitude
$$Q(k_1,k_2,k_3)\frac{B(k_1,k_2,k_3)}{P(k_1)P(k_2)+P(k_2)P(k_3)+P(k_3)P(k_1)}.$$
(30)
The three-point amplitude $`Q`$ has the convenient feature that for the lowest nonvanishing result in perturbation theory, $`Q`$ is independent of time and the overall amplitude of $`P`$; for scale-free models with a power-law $`P`$, $`Q`$ is independent of overall scale as well. To lowest order, it follows from equation (12) that the equilateral bispectrum has a particularly simple form, $`B^{(0)}(k)=\frac{12}{7}P_{\mathrm{lin}}^2(k)`$, and we have $`Q^{(0)}(k)=\frac{4}{7}`$, independent of the power spectrum.
### 5.1 Synthetic Halo Replacement
To investigate the numerical effects of limited resolution in the simulations, we have experimented with the distribution of matter in halos identified in the simulations. In these experiments, we keep the locations and masses of the halos unchanged but redistribute the subset of particles which lies within the virial radius $`R_{200}`$ (the radius within which the mean overdensity is 200) of each halo according to a prescribed density profile. We then recompute the two- and three-point statistics $`\mathrm{\Delta }`$ and $`Q`$ from the redistributed particle positions as well as the original non-halo particles, which remain at their original positions. By using density profiles obtained empirically from higher-resolution simulations of individual halos, this recipe allows us to model accurately the inner regions of the halos on scales below the numerical softening length scale while at the same time preserving all the large-scale information available in the large parent simulation. This technique should also be useful for other studies that are sensitive to the inner halo density profiles, for example the ray-tracing method in gravitational lensing.
Ma & Fry (2000a) have used this replacement technique to experiment with synthetic halos that follow a pure power-law profile $`ur^ϵ`$. It is found that $`\mathrm{\Delta }(k)`$ and $`Q(k)`$ at high-$`k`$ indeed obey $`\mathrm{\Delta }(k)k^{2ϵ3}`$ and $`Q(k)k^{3ϵ}`$ as predicted by the simple power-law model of Peebles (1974). The scaling works even in the presence of the full distribution of matter outside the halo cores. Here we extend this replacement technique to more realistic halo profiles of equation (2). Figures 1 and 2 illustrate the effects on the matter power spectrum and bispectrum when the original halos in large cosmological simulations are replaced by synthetic halos with the density profile $`u_{II}=1/(x^{3/2}+x^3)`$ of equation (2). For the $`n=2`$ scale-free model, the concentration parameter is taken to be $`c(M)=3(M_{}/M)^{1/6}`$, which is consistent with Navarro et al. (1997) and has the expected scaling with mass, $`cM^{(3+n)/6}`$, in a scale-free model. For $`\mathrm{\Lambda }`$CDM models, we use $`c(M)=5(M_{}/M)^{1/6}`$ as suggested by Figure 3 of Moore et al. (1999). We note, however, that $`c(M)`$ from various recent simulations has shown a large scatter, and its functional form depends on the exact form of the density profile used. For the $`\mathrm{\Lambda }`$CDM model and form $`u_{II}=1/(x^{3/2}+x^3)`$, for example, a flatter and smaller $`c(M)=3(M_{}/M)^{0.084}`$ appears to be preferred by Jing & Suto (2000) and Navarro et al. (1997). The results of Tormen et al. (1997) and Cole & Lacey (1996) are also only marginally consistent with each other. A more detailed investigation of the different forms of $`c(M)`$ can be found in Ma & Fry (2000c).
In Figures 1 and 2, the agreement at low values of $`k`$ between the original and synthetic halos is excellent, confirming that the correlation functions on larger length scales are insensitive to the spatial distribution of particles in the halo cores. The only significant difference between the simulation and synthetic halos appears at small length scales, where the coarser resolution of the simulation blurs out the structure of the inner halo and results in an inner profile flatter than in equation (2). This effect is manifested in the bending over of the dashed curves for $`P(k)`$ in Figures 1 and 2 at high $`k`$, and is corrected for when the synthetic halos are used.
### 5.2 $`N`$-body Results vs. Analytic Halo Model
We now proceed to compare the predictions of the analytic model of §2 – §4 with the numerical results from cosmological simulations. Figures 3 and 4 show the $`k`$-space density variance $`\mathrm{\Delta }(k)`$ (upper panel) and the three-point amplitude $`Q_{\mathrm{eq}}(k)`$ for equilateral triangles for the $`n=2`$ scale-free model and the $`\mathrm{\Lambda }`$CDM model. The solid black curves are the model predictions computed from equations (3) and (29). The contribution from the single-halo and multiple-halo terms are shown separately as dashed curves. For the density profile, we use the same $`u_{II}=1/(x^{3/2}+x^3)`$ and concentration parameters as in Figures 1 and 2. For the mass function, we use the Press-Schechter formula but reduce its overall amplitude by 25%, which we find necessary in order to match the halo mass functions for our numerical simulations. This overestimation of halo numbers with $`MM_{}`$ by Press-Schechter is a well known result reported in many other studies (see Jenkins et al. 2000 and references therein). The mass limits for the integrals in equations (3) and (29) do not significantly affect the model predictions for the total $`\mathrm{\Delta }`$ or $`Q`$. Raising the lower mass limit does reduce the contribution from lower mass halos and hence lower the high-$`k`$ amplitudes of the multiple halo terms $`\mathrm{\Delta }_{2h}`$, $`Q_{2h}`$, and $`Q_{3h}`$, but these terms make negligible contributions to the total $`\mathrm{\Delta }`$ and $`Q`$.
As discussed in §3 and 4, the nonlinear parts of both the two- and three-point statistics are determined by the dominant 1-halo term because the closely spaced particle pairs and triplets mostly reside in the same halos. The multiple-halo terms are therefore significant only on larger length scales comparable to the separation between halos. Their inclusion, however, is necessary for the transition into the linear regime.
For the $`n=2`$ model in Figure 3, we plot the results against the scaled $`k/k_{\mathrm{nl}}`$, where $`k_{\mathrm{nl}}`$ characterizes the length scale that is becoming nonlinear and is defined by $`_0^{k_{\mathrm{nl}}}d^3kP_{\mathrm{lin}}(a,k)/(2\pi )^3=1`$. Three time outputs are shown, where the expansion factor (1 initially) and $`k_{\mathrm{nl}}`$ (in units of $`2\pi /L`$) are: $`(a,k_{\mathrm{nl}})=(13.45,29),(19.03,14.5)`$, and $`(26.91,7.25)`$ (from left to right). For the two-point $`\mathrm{\Delta }(k)`$, the agreement between the model prediction and the simulations is excellent. The three simulation outputs also overlap well, indicating that self-similarity is obeyed, as reported in Jain & Bertschinger (1998). For the three-point $`Q_{\mathrm{eq}}`$, however, self-similar scaling does not hold as rigorously (Ma & Fry 2000a). It is interesting to note that the analytic prediction agrees most closely with the earliest output $`(a,k_{\mathrm{nl}})=(13.45,29)`$ (green curve). This provides further evidence to the suggestion of Ma & Fry (2000a) that the later outputs of the $`n=2`$ simulation may be affected by the finite volume of the simulation box. For the $`\mathrm{\Lambda }`$CDM model in Figure 4, the analytic model again provides a good match to the $`N`$-body results within the fluctuations among the simulations. We illustrate the numerical effects due to box sizes by showing results from two runs with volume (100 Mpc)<sup>3</sup> and (640 Mpc)<sup>3</sup>. The model predictions extend well beyond the resolution of the simulations.
The real-space two-point correlation function for the $`n=2`$ and $`\mathrm{\Lambda }`$CDM models is shown in Figures 5 and 6. For the halo model predictions, we have chosen to show only the results for the 1-halo term $`\xi _{1h}`$ because this term dominates the interesting nonlinear portion of $`\xi `$. The agreement between the halo model (dashed curves) and the simulations (symbols) is again excellent. For the 2-halo terms $`\xi _{2h}`$, the computation can be done more easily in $`k`$-space as shown in Figures 3 and 4, so we do not include them here.
For comparison, we plot in Figures 3–6 the results from the commonly used fitting formulas for the nonlinear power spectrum (Hamilton et al. 1991; Jain et al. 1995; Peacock & Dodds 1996; Ma 1998; Ma et al. 1999). While the formulas provide a good approximation to $`\mathrm{\Delta }(k)`$ up to $`k/k_{\mathrm{nl}}50`$ for the $`n=2`$ model and $`k20h`$ Mpc<sup>-1</sup> for the $`\mathrm{\Lambda }`$CDM model, the figures show that significant deviations occur at higher $`k`$, and the fitting formula and our current model predict different high-$`k`$ slopes for $`\mathrm{\Delta }(k)`$. Since the high-$`k`$ behavior of the fitting formulas has been constructed to obey the stable clustering prediction, this discrepancy has an important implication for the validity of stable clustering, which we discuss briefly in the next section and at length in Ma & Fry (2000c).
## 6 Discussion
We have constructed a physical model for the correlation functions of the mass density field in which the correlations are derived from properties of dark matter halos. We have described in detail the input, construction, and results of this model in §2 – §5. We now examine more closely its physical meanings and implications in three separate regimes.
On scales larger than the size of the largest halo, the contributions from separate halos dominate, and (by design) the model reproduces the results of perturbation theory. On intermediate scales, $`1/R_{}k1/R_s(M_{})`$, because of the exponential cutoff in the mass function $`dn/dM`$ at the high mass end, the contribution to the volume integrals in equation (18) is dominated by the large-$`r`$ regime where the halo profiles are roughly $`r^3`$. The correlation functions therefore behave approximately as predicted by the power-law model with $`ϵ=3`$, i.e., $`\mathrm{\Delta }k^{2ϵ3}k^3`$ and $`Qk^{3ϵ}`$ constant. This is why $`Q`$ exhibits an approximately flat plateau at intermediate $`k`$ in the bottom panels of Figures 3 and 4.
On the smallest and most nonlinear scales, the correlation functions probe the innermost regions of the halos. Intriguingly, the halo model predicts on these scales a behavior that is different from either the frequently-assumed stable clustering result of $`\mathrm{\Delta }(k)k^\gamma `$ with $`\gamma =(9+3n)/(5+n)`$ (Davis & Peebles 1977), or the power-law profile result of $`\gamma =2ϵ3`$. The implication of departure from stable clustering is significant because all the fitting formulas for the nonlinear $`P(k)`$ in the literature (see §5.2) have been constructed to approach the stable clustering limit at high $`k`$. A more detailed study on the criteria for stable clustering in this model is given in a separate paper (Ma & Fry 2000c).
The origin of the deviation from stable clustering in the model at high-$`k`$ can be understood as follows. For the two-point function, as $`k`$ becomes large, the one-halo integral $`P_{1h}(k)`$ in equation (3) converges before the exponential cutoff, and is dominated by contributions near the mass scale for which $`kR_s=1`$. The behavior now depends on the mass distribution function. The various mass functions discussed in §2.2 have the same general behavior of $`dn/dMM^2\nu ^\alpha e^{\nu ^2/2}`$, where $`\nu =\delta _c/\sigma `$. The Press-Schechter form assumes $`\alpha =1`$ (see eq. ), while others (e.g., Sheth & Tormen 1999; Jenkins et al. 2000) suggest a flatter slope of $`\alpha 0.4`$ for the lower mass halos. Since the scale radius $`R_s`$ depends on mass as $`R_s=R_{200}/cM^{1/3}/M^{(3+n)/6}M^{(5+n)/6}`$, and $`R_s^3\overline{\delta }M`$ (up to logarithmic factors), we find from equation (3) that the power spectrum at high $`k`$ goes as
$$\mathrm{\Delta }(k)\mathrm{\Delta }_{1h}(k)k^3𝑑M\nu ^\alpha \stackrel{~}{u}^2(kR_s).$$
(31)
Changing variables to $`y=kR_sk(M/M_{})^{(5+n)/6}`$, we see that
$$\mathrm{\Delta }(k)k^\gamma ,\gamma =\left(\frac{9+3n}{5+n}\right)\alpha \left(\frac{3+n}{5+n}\right),$$
(32)
where the first term in $`\gamma `$ is the prediction of stable clustering. The departure arises from the factor $`\nu ^\alpha `$ in the mass function, and would vanish only if $`\alpha =0`$ or $`n=3`$. This is the origin of the difference in $`\mathrm{\Delta }(k)`$ at high $`k`$ between the model prediction (solid curves) and the fitting formula (dotted curves) shown in Figures 3 and 4.
For the three-point function, the one-halo integral $`B_{1h}^{\mathrm{eq}}`$ in equation (29) converges (barely, for $`p=\frac{3}{2}`$ and $`n=2`$), giving
$`B^{\mathrm{eq}}(k)`$ $``$ $`k^{\gamma _36},\gamma _3=2\left({\displaystyle \frac{9+3n}{5+n}}\right)\alpha \left({\displaystyle \frac{3+n}{5+n}}\right)`$
$`Q^{\mathrm{eq}}(k)`$ $``$ $`k^{\alpha (3+n)/(5+n)}`$ (33)
This again disagrees with the prediction of stable clustering that $`Q`$ is constant, but it appears to be consistent with numerical simulations as shown in Figures 3 and 4.
For yet higher order correlations, details of the halo profile begin to matter. For $`p=1`$, the pattern of equations (32) and (33) persists to all orders, but for $`p=\frac{3}{2}`$ they apply only for the two- and three-point functions; for four-point and higher functions the nonlinear scale $`M_{}`$ and $`\gamma _n=np3`$ for $`n4`$. Thus there seems to be some potentially interesting behavior that is tested only in the four-point function and higher.
## 7 Summary
We have presented an analytic model for the two- and three-point correlation functions $`\xi (r)`$ and $`\zeta (r_1,r_2,r_3)`$ of the cosmological mass density field and their Fourier transforms, the mass power spectrum $`P(k)`$ and the bispectrum $`B(k_1,k_2,k_3)`$. In this model, the clustering statistics of the density field are derived from a superposition of dark matter halos with a given set of input halo properties. These input ingredients include realistic halo density profiles of equation (2), halo mass distribution of equation (6), and halo-halo spatial correlations of equations (10) and (15). The main results of the model are given by equations (18) and (3) for the two-point statistics $`\xi `$ and $`P`$, and by equations (25) and (29) for the three-point statistics $`\zeta `$ and $`B`$. This model provides a rapid way to compute the correlation functions over all length scales where the model inputs are valid; it also gives a physical interpretation of the clustering process of matter in the universe.
We have tested the validity of this model by comparing its predictions with results from cosmological simulations of an $`n=2`$ scale-free model and a $`\mathrm{\Lambda }`$CDM model. As Figures 3 – 6 illustrate, the model describes well the simulation results spanning the entire range of behavior from the perturbative regime on large scales to the strongly nonlinear regime on small scales. To probe the critical high-$`k`$ range in the deeply nonlinear regime, we have used a halo replacement technique to increase the resolution of the large parent simulations. As Figures 1 and 2 illustrate, this method of replacing the original halos that suffer from numerically softened cores with synthetic halos of analytic profiles is a reasonable way to improve the resolution of numerical simulations. By using density profiles obtained empirically from higher-resolution simulations of individual halos, this recipe allows us to model accurately the inner regions of the halos on scales below the numerical softening length scale, while at the same time preserving all the large-scale information available in the large parent simulation. This technique should also be useful for other studies that depend on the inner halo density profiles, for example, the ray-tracing method in gravitational lensing.
Given that dark matter halos in simulations (and presumably in nature) are not perfectly spherical, cleanly delineated objects, it is intriguing that the model constructed in this paper works as well as it does at matching the simulation results. Nevertheless, this analytic model provides a good qualitative and quantitative description over the entire range of scales covered by the simulation, and it can be used to make predictions beyond these scales. This is the first model prescription that successfully reproduces both two- and three-point mass correlations. We believe that it will prove to be a generally useful framework.
We have enjoyed stimulating discussions with John Peacock and David Weinberg. We thank Edmund Bertschinger for valuable comments and for providing the $`n=2`$ scale-free simulation. Computing time for this work is provided by the National Scalable Cluster Project and the Intel Eniac2000 Project at the University of Pennsylvania. C.-P. M. acknowledges support of an Alfred P. Sloan Foundation Fellowship, a Cottrell Scholars Award from the Research Corporation, a Penn Research Foundation Award, and NSF grant AST 9973461.
## Appendix A Appendix
In this Appendix we display analytic forms for the convolution of the dimensionless profile shape
$$\lambda (x)=d^3yu(y)u(|\text{x}+\text{y}|)$$
(A1)
discussed in §3. These analytic expressions are useful for computing the nonlinear two-point correlation function $`\xi `$ of the mass density field, which is dominated by the 1-halo term $`\xi _{1h}`$ in equation (18) and is related to $`\lambda `$ by
$$\xi (r)\xi _{1h}(r)=𝑑M\frac{dn}{dM}\overline{\delta }^2R_s^3\lambda (r/R_s),\mathrm{for}\xi 1.$$
(A2)
For the type-I profile $`u_I`$ of equation (2), the angular integration in equation (A1) is analytic, and $`\lambda `$ is reduced to a simple integral
$$\lambda _I(x)=\frac{2\pi }{(2p)x}_0^{\mathrm{}}\frac{ydy}{y^p(1+y)^{3p}}\left[\frac{(x+y)^{2p}}{(1+x+y)^{2p}}\frac{|xy|^{2p}}{(1+|xy|)^{2p}}\right].$$
(A3)
For the special case $`p=1`$, this integral can be further reduced to the analytical form
$$\lambda _I(x)=\frac{8\pi }{x^2(x+2)}\left[\frac{(x^2+2x+2)\mathrm{ln}(1+x)}{x(x+2)}1\right],p=1.$$
(A4)
For $`u_{II}`$ of equation (2), we are able to simplify $`\lambda `$ to
$$\lambda _{II}(x)=\frac{2\pi }{x}_0^{\mathrm{}}\frac{ydy}{y^p(1+y^{3p})}F_p(x,y),$$
(A5)
where the function $`F_p(x,y)`$ represents the angular part of the integration in equation (A1) and
$$F_p(x,y)=_{|xy|}^{x+y}\frac{zdz}{z^p(1+z^{3p})}.$$
(A6)
The integral in $`F_p`$ can be reduced to analytic forms for special values of $`p`$. Here we display the six cases $`p=0`$, $`1/2`$, 1, $`3/2`$, 2, and $`5/2`$:
$`F_0`$ $`=`$ $`{\displaystyle \frac{1}{6}}\left\{2\sqrt{3}\mathrm{tan}^1\left[{\displaystyle \frac{1+2(x+y)}{\sqrt{3}}}\right]+\mathrm{ln}\left[{\displaystyle \frac{1(x+y)+(x+y)^2}{1+2(x+y)+(x+y)^2}}\right]\right\}`$ (A7)
$`{\displaystyle \frac{1}{6}}\left\{\text{replace }(x+y)\text{ above with }|xy|\right\}`$
$`F_{1/2}`$ $`=`$ $`{\displaystyle \frac{1}{10}}\{2\sqrt{10+2\sqrt{5}}\mathrm{tan}^1\left({\displaystyle \frac{1+\sqrt{5}4\sqrt{x+y}}{\sqrt{102\sqrt{5})}}}\right)`$ (A8)
$`2\sqrt{102\sqrt{5}}\mathrm{tan}^1\left({\displaystyle \frac{1+\sqrt{5}+4\sqrt{x+y}}{\sqrt{10+2\sqrt{5})}}}\right)`$
$`+4\mathrm{ln}\left(1+\sqrt{x+y}\right)(1+\sqrt{5})\mathrm{ln}\left[1+{\displaystyle \frac{1}{2}}(1+\sqrt{5})\sqrt{x+y}+x+y\right]`$
$`(1\sqrt{5})\mathrm{ln}[1{\displaystyle \frac{1}{2}}(1+\sqrt{5})\sqrt{x+y}+x+y]\}`$
$`{\displaystyle \frac{1}{10}}\left\{\text{replace }(x+y)\text{ above with }|xy|\right\}`$
$`F_1`$ $`=`$ $`\mathrm{tan}^1(x+y)\mathrm{tan}^1(|xy|)`$ (A9)
$`F_{3/2}`$ $`=`$ $`{\displaystyle \frac{1}{3}}\left\{2\sqrt{3}\mathrm{tan}^1\left[{\displaystyle \frac{1+2\sqrt{x+y}}{\sqrt{3}}}\right]+\mathrm{ln}\left[{\displaystyle \frac{1+2\sqrt{x+y}+x+y}{1\sqrt{x+y}+x+y}}\right]\right\}`$ (A10)
$`{\displaystyle \frac{1}{3}}\left\{\text{replace }(x+y)\text{ above with }|xy|\right\}`$
$`F_2`$ $`=`$ $`\mathrm{ln}\left[{\displaystyle \frac{x+y}{1+x+y}}\right]\mathrm{ln}\left[{\displaystyle \frac{|xy|}{1+|xy|}}\right]`$ (A11)
$`F_{5/2}`$ $`=`$ $`{\displaystyle \frac{2}{\sqrt{|xy|}}}{\displaystyle \frac{2}{\sqrt{x+y}}}+\mathrm{ln}\left[{\displaystyle \frac{(1+2\sqrt{x+y}+x+y)|xy|}{(1+2\sqrt{|xy|}+|xy|)(x+y)}}\right]`$ (A12)
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# On the theory of system administration
## 1 Introduction
System administration is the realm of computer science which deals with the planning, configuration and maintenance of computer systems. It is presently a discipline founded mainly on the anecdotal experiences of system managers. To date, no formal (mathematical) analyses of system administration have been undertaken, with the aim of making more scientific studies. This makes it difficult to express objective truths about the field, avoiding marketing assertions and the vested interests of companies and individuals, which are common in the commercial sector.
The aim of the present work is to establish a formal basis for the field, a way of formulating a framework for objective discussions about computer management. It will hopefully serve as a bridge between mathematical disciplines and system administrators. In this respect, the paper may be viewed mainly as a commentary, laying some foundations for future work, rather than providing immediate solutions.
In previous work, it has been shown how the average behaviour of systems of computers and users can be approximated by a blend of statistical models and thermodynamical ideas. That work allows us to form a mathematical model of computer systems which can be used as a basis for modelling system administration. The study of computer behaviour has much in common with the physics of thermodynamics. From a coarse mathematical viewpoint, system administration can be viewed in much the same way as thermodynamical experiments with pistons and engines, i.e. moving information and resources around in such a way as to change the state of the system. However this viewpoint is mainly useful in a calculational setting. System administration also has much in common with medicine. In many ways, system administration is medical science for computers: a somewhat simpler problem than that of human physiology, but nonetheless involving many of the same themes: nutrition, regulation, immunity and repair.
What then should a theory of system administration be about? The task of elucidating this sounds straightforward, but it is a slippery business. System administration, in reality, is based on mainly qualitative, high level concepts, which mix technical and sociological issues at many levels. Although it is clear to system administrators that there is a body of technical principles involved in the discipline, it remains somewhat intangible from the viewpoint of a scientist. It is hard to find anything of general, reproducible value on which to base a more quantitative theory.
One of the obstacles to formulating such a theory is the complexity of interaction between humans and computers. There are many variables in a computer system, which are controlled at distributed locations. Computer systems are complex in the sense of having many embedded causal relationships and controlling parameters. Computer behaviour is strongly affected by human social behaviour, which is often unpredictable. The task of identifying and completely specifying the ideal state is therefore a non-trivial one. It is nonetheless this task which this paper attempts to address. Can one formulate a quantitative theory of system administration, which is general enough to be widely applicable, but which is specific enough to admit analysis?
If this, already significant problem can be addressed in sufficient terms, one might then aim to look further towards general regulatory systems and approach more ambitious questions. It is not difficult to see many analogous questions in other areas of science, which could be applicable to system administration. For instance: what is the effectiveness of generalized immunity and repair systems (automatic repair and regulation)? Is there an optimal strategy for error detection and correction? Is a system administrator’s human mind (playing the role of doctor/surgeon) better or worse than a mechanistic response or immune system? This last point is often a bone of contention in the system administration community. Should tasks be automated? Or should a human lawgiver always remain in manual control? What is more efficient? Biological systems point to the need for both types of management: at any given moment, a doctor’s intelligence and superior human cognition can compensate for a lack of adaptation in our programmed immune responses, but the automatic immune response is both faster and more capable than a doctor when its program is sufficient. Certainly the empirical evidence in biological information systems is compelling: after billions of years of evolution, nature has established immune systems in all vertebrates larger than a tadpole. Of course, this is no indication that the solution is optimal. No acceptable analysis has been used to demonstrate this yet. It could be that vertebrate evolution is merely poised on some plateau between minima of much deeper importance.
The aim of this paper then is to elevate system administration from an expression of subjective opinion to a more objective, scientific level, hopefully without inflating it meaninglessly into pseudo-science or philosophy. In order to limit the length of this paper, solutions of the models and constraints will be kept to a minimum here. However, it will be possible to draw a few general conclusions, even without reference to specific models.
The outline of this paper is as follows. To begin the discussion it is necessary to establish some basic axioms. It is important to restrict the scope of what a theory of system administration encapsulates; without such a restriction, one ends up with either many disjointed pieces or only vague hand-waving notions. Having determined the ground rules, it is then appropriate to identify the basic operations which can be carried out within that scope. This identification is required in order to formulate a discussion of strategies for system management. Once this level of formality has been attained, strategies can be formulated, based on types of action and timing and the task of administrating a computer system can be described in precise game theoretical terms. This is the primary goal of this work.
## 2 The scope of system administration
One of the first obstacles in discussing the theory of system administration is defining its scope. System administrators are called upon to perform all manner of tasks as part of their duties. This battery of skills has no particular cohesion or structure to it, so it resists formalization. We must improve on this situation if we are to make progress in forming a theory of system administration. In particular, we must restrict its scope to encompass only core activities. These core activities will include insuring availability, efficiency, and security for all users, and finally fault diagnosis of the system. This includes issues such as software installation and upgrades, which can be classified under availability and efficiency. It also includes user management to a certain extent, though it will not be useful to address the issue of creation of user accounts in this context.
## 3 On scales
A well known feature of descriptions of complex systems is that a complete understanding is best organized as a unification of the partial understanding of the system at several different levels or scales. Complex systems are often so disparate at different scales that quite different descriptions are required to capture the full essence. A theory of system behaviour at, say the microscopic level of system calls, need not resemble a theory for the behaviour at a macroscopic scale of larger entities, such as patterns of user behaviour. Both are needed in order to understand the whole hierarchy of things going on.
If one is only interested in high level phenomena, then the details of low level phenomena are seldom directly relevant, to a good approximation. This is the principle of separation of scales. The principle states that, as one moves from microscopic to macroscopic scales, new behaviour can emerge as collective phenomena, which often depends only weakly on the microscopic details of the levels below. This is a simple idea, which is quite intuitive, but which has far reaching consequences. It can easily be appreciated with the help of a couple of examples.
A bridge, for instance, has the property of spanning a distance and carrying weight, regardless of whether it is made from steel or copper or wood. The choice of material and the microscopic arrangement of atoms in the metal or wood, of course, tells us something about the strength of the bridge, but perhaps not as much as the structure of the bridge at the scale of the whole thing. In other words, the construction of the bridge at the scale of the users of the bridge is far more important to its function than the microscopic construction of its pieces under a microscope.
Similarly, to an acceptable approximation, the behaviour and operation of a sales database, at the level of information transactions (their order and type), is more important to an information retrieval system than how those transactions are implemented through system calls (e.g. whether the system runs on Unix or on NT). The ability to retrieve information does not depend on whether the storage medium is, an IDE or a SCSI disk. The same job will be done regardless.
To summarize, a description of system behaviour at a high level is, for many purposes, independent of specific details of the lower levels. Computer systems can be modelled by generic computer systems with certain high level characteristics; similarly users can be modelled as idealized users, also with common characteristics. A theory of system administration will be most successful if it appeals to such generalities, rather than delving into unnecessary specifics.
## 4 Axioms of system administration
To begin a formal discussion, we need to establish a frame of reference, i.e. the ground rules for the discussion. In this section, a basic fundament is proposed with the aim of striking a balance between reality and suitability for analysis. It is also necessary to partially limit the scope of the discussion to avoid unnecessary complication. Although the aim of this presentation is not precise mathematical rigor, it is the aim to indicate that such a rigor is possible and to indicate how it can be provided. A secondary aim is to communicate the key elements of the discussion to a more theoretical audience; for these reasons, the language adopted is one which is meant to build bridges between system administration and more mathematical disciplines. Readers are asked to keep an open mind with regard to use of terms however, since technical disciplines often use words in meanings which are specific to those disciplines, and this could lead to confusion.
A computer system is analogous to a community composed of many interacting and competing players: i.e. users and administrators. It can only properly be discussed in terms of the aims and activities of this collective and of individual members of that community. Not all the members of a community share the same objectives, as a general rule. Traugott and Huddleston have pointed out that it is often pertinent to view a local computer community as a single virtual machine, rather than as a conglomeration of individual hosts. In this paper, the term computer system will be used to refer to the collective hosts of a local domain, or some appropriate logical unit of networked computers. It is taken for granted that there may be internal competition for resources and even conflict between competing parties.
In order to formulate a theory of system administration we must establish a set of possible goals, procedures and obstructions and state them in formal terms. The aims and intentions of each computer system are different; usually they are prescribed by a system policy, i.e. a formal statement of intent and allowed practice. The aim is then to postulate or derive strategies which best achieve those goals, given the essential constraints. From this viewpoint, one expects the language of constrained competition to play a role in a theory of system administration. Even if one could frame such a theory in formal terms, what would be the purpose of such an exercise? The principal benefit of such an attempt is to create a rigid protocol for discussing system administration, which is general enough to cover most of the actual problems and possibilities, but which is stringent enough to prevent its perversion by parties with vested interests in proving a certain point of view.
There is a number of stages in this programme. To begin with, one needs some basic axioms which all parties agree on, propositions which define the aims of system administration. Next one needs to abstract a model of a computer system which is sufficient to capture the dynamical interaction between all of the players, but which is sufficiently simple to be surrendered for analysis. Here we shall suppose that a computer’s resources (memory, CPU, disk etc) are divided into two parts,
$`R=R_cR_m,`$ (1)
i.e. a part which determines the behavioural configuration parameters $`C`$ of the system working resources, and a remainder part (the working resources themselves) which users of the system can change as a normal part of their interaction with the system. This remainder part can be observed over an appropriate time scale, giving a set of measurements $`\overline{M}`$ which indicate how the system is being used.
The different possible configurations of the system resources $`\{C\}`$ are made up from the independent operation types $`\{T\}`$ which lead to these configurations. From this definition one needs to be sure that a unique description is possible, i.e. eliminate points of contention about the description itself. Finally, one must be sure that the description is sufficiently complete, i.e. that there exists a mapping between policy and system configuration which is as complete as the problem itself. The purpose of this section is to introduce the key players in this description, in advance of a fuller description in the coming sections.
In order to state the purpose of system administration, we may take the basic tenet or principle to be the following:
###### Basic assumption 1
The requirements and constraints of any computer system are defined at any time $`t`$ by an implementable system policy $`P(t)`$. This policy determines the actions or rules of play for a system administrator, but not necessarily the actions of users. It includes a specification of which and how many users are allowed to access the system.
The policy $`P(t)`$ is not usually a continuous function of time, but may change catastrophically (in the mathematical sense) over a time scale which is much longer than the time scale over which users act and make changes to the system. The nature of this policy is not yet determined.
In order to make a policy implementable, it must be possible to relate it to a complete configuration instruction for the system $`C(t)`$, through rules and constraints. These rules and constraints could be issued verbally to users, or could be programmed into configuration files of software components which form the system. A single complete configuration instruction for the system can be thought of as being a sum of two parts:
$`C=C_rC_u,`$ (2)
a specification of resource configurations $`C_r`$ which describe how the software and hardware landscape is configured, and a specification of user configurations $`C_u`$, which describes who is allowed to do what with the resources (this includes remote, network users who access services through local agents). A specification of user configurations $`C_u`$ (numbers of users and their rights to resources) could easily be separated from system policy conceptually, but it is convenient to view the policy as a complete specification of the system plus its intended and actual usage. The meaning of the symbol $``$ is that of a heuristic union: configuration specifications take many forms (are objects of many types). They are most easily thought of as sets of more primitive objects, in which case the addition of sets implies their strict union.
A complete configuration instruction can be thought of geometrically as a point in a vector space, which is found by adding together instructions of linearly independent (orthogonal) types. One does this by introducing a set of primitive configuration instruction types $`\{T^i\}`$, and writing the complete configuration as a linear combination of these:
$`C={\displaystyle \underset{i}{}}c_iT^i,`$ (3)
with set-valued coefficients $`c_i`$. The basis of primitive configuration operations will be described later. A complete configuration usually contains instructions for the operators of the system also. The system administrator can also be viewed as part of this system for the sake of abstraction.
The set of all configurations $`\{C\}`$ contains much redundancy. Let us imagine that the mapping of complete configuration instructions to the final state of the system is many to one, and that this multiplicity can be represented by a group of permutations and transformations $`𝒢`$<sup>1</sup><sup>1</sup>1The nature of this group could be fairly complicated and is not particularly important to the discussion. The fact that the redundancy, in principle, may be represented by a mathematical group is an idealization which is attainable in in theory. It is not an expression of the current state of affairs in the world of computers.. Thus equivalent configurations could be formed by permuting configuration instructions, if ordering is unimportant, or by exchanging (transforming) one component of the configuration for another. An example of this is the following: the configuration of a World Wide Web service might be possible with several equivalent software systems, each with equivalent configuration files: in this case, these would form an equivalence class. Conversely, if even a minor detail distinguishes them, then they are inequivalent configurations.
Let us define an implementable policy $`P(t)`$ as being any representative member of the set of equivalent configurations $`\{C(t)\}`$
$`P(t)\{C(t)\}/𝒢.`$ (4)
A policy could naturally mean more than a configuration (or computer and its operators), but as long as other aspects of the policy cannot be implemented by either machine or human, they are irrelevant to the system. Having related policy to configuration instructions, the path is clear to define the state $`S(t)`$ of the system.
Let a state of the system describe a single configuration of users $`C_u`$, of system resources $`C_r`$ and a set of measured average metrics $`\overline{M}`$ which summarizes the average usage of the system in relation to the users. The metrics $`\overline{M}`$ represent a first order response (feedback interaction) between users and resources. The state is written, again, as a direct sum
$`\overline{S}_p\left({\displaystyle \frac{C_rC_u}{𝒢}}\right)\overline{M}(C_r,C_u)`$ (5)
where the behavioural metrics $`\overline{M}(C_r,C_u)`$ are functions of resources and user activity. One may now state the following provable hypothesis.
###### Theorem 1
Any sufficiently complete system policy $`P(t)`$ specifies, by implication, a representative average ideal state $`\overline{S}_p(t)`$, from an equivalence class of ideal states $`\{\overline{S}\}`$ under $`𝒢`$, for the computer system concerned, over user time-scales $`T_u`$, provided that the rate of change of policy $`dP/dt`$ is much smaller than changes in user behaviour $`d\overline{M}/dt`$, i.e. the policy changes on the order of weeks or months rather than hours or days.
Corollary: The ideal state can only be identified on average, since interactions with unpredictable user activity are constantly causing fluctuations $`\delta S(t)`$ in the state of the system. These fluctuations also occur at a rate $`d\overline{M}/dtdP/dt`$.
The existence of an ideal state has already been used in designing the author’s site Configuration Engine (cfengine), but it has not previously been explained at length. The proof of this theorem is straightforward, from the definitions. Every computer system has a finite set of resources and configuration objects which is completely prescribed by a total configuration $`C_p`$. Each resource object may be in a state described by a finite length bit string, describing a distinct configuration $`s_iC_r`$. There is therefore a mapping from the configuration $`C_p`$ to the actual state
$`C_p\overline{S}_p+\delta S.`$ (6)
This mapping is one to many, since $`\delta S`$ is a stochastic variable. The averaging operation eliminates the non-uniqueness by extinguishing $`\delta S`$, provided the averaging process is defined, i.e., $`d\delta S/dtdP/dt`$. Thus any complete set of average measurements contributes to the average state in a well-defined manner.
The meaning of ‘sufficiently complete system policy’ is now clear. The covering of the policy domain must be as large as the domain of state one wishes to cover, since it follows from the above definitions that the association of policies to states is now one to one, after one factors out the equivalences $`𝒢`$. The uniqueness is secured by making the configuration instruction itself a part of the state. Without this, there would still be ambiguity, since there is no guarantee that a measurement $`\overline{M}(C_r,C_u)`$ is a unique function of its arguments. This completes the proof.
This sufficiency referred to above has the corollary that an incomplete system policy $`P_1`$ cannot determine a unique state for the whole system, only a part of it. An incomplete policy divides the system into two or more parts, since the total policy is still in one to one correspondence with the states.
$`P_1+P_2+\mathrm{}S_{p1}+S_{p2}+\mathrm{}`$ (7)
By the virtue of the fundamental theorem, we have the important conclusion that the necessary and sufficient condition for implementation a policy $`P(t)`$ (i.e. the ability to map it onto a system configuration over a period of time) is that the total average state $`\overline{S}(t)=\overline{S}_p(t)`$.
Let us take a moment to understand the structure and meaning of average ideal state. It is tempting to think of the system as being in an ideal state at some time $`t_0`$ and then deviating from it at later times. The precise state of the system at some reference time might seem to characterize an ideal to our subjective judgement, but the ideal state of configuration must change with time, since the computer system is, by nature, influenced by users whose activities are not completely secured by a policy. To freeze one’s view of the ideal in time, is to place unreasonable restrictions on the use of the system (we shall see this later in examples connected to the use of fixed disk quotas). A specification of resource and user boundary conditions is not the same as a specification of the ideal dynamical behaviour of the system, if users are allowed to act on the resources.
Given that the policy and user configurations are stable over the prescribed time-scales, one may take the average value (or distribution of values) for each metric which characterizes the response of system over shorter time-scales $`\overline{M}(C_r,C_u)`$ as being representative of the state at time $`t`$. This summarizes the effect of feedback of users on resources, or the statistical interaction between the users and the system. Since we have prescribed every bit string affecting the dynamics of the system at the outset as policy, and we have measured the average result of those prescriptions at time $`t`$, we have a complete description of the system in terms of an implementable policy.
$`\overline{S}_p=P\overline{M}(P).`$ (8)
Not surprisingly, this expression is directly analogous to linear response theory in the physics of time-varying systems. The policy plays the role of a constraint of the motion, while the statistical metric $`\overline{M}`$, has the role of the integrated response of the evolved state at time $`t`$. The ‘equations of motion’ which lead to the evolution of a system also have an analogue here: they are the operations carried out by the system software on the resources.
The permutation or invariance group $`𝒢`$ is of no concern to this paper except as a matter of principle for the most pedantic. It is a heuristic representation of all of the involved details which are irrelevant to a theoretical formulation, but which occur in practice<sup>2</sup><sup>2</sup>2The fact that such details do indeed represent a set, indeed a group of permutations and transformations is clear from the empirical facts. The factoring of redundancy means picking only one representative member of each configuration which gives the same results, in the same manner that factor groups are formed in group theory. Indeed, it is a trivial technical point that the sets and equivalence classes in a computer system may all be represented by operations on a single binary string (a computing machine does precisely this on a finite, possibly disjointed binary string). The existence of transformations with closure is assured by extending a binary string to encompass all possibiluties; the existence of an inverse is trivial for permutations, as is the existence of a null operation and associativity. That the formal factor group exists follows from the existence of heuristic equivalences with respect to system function.. Nevertheless, it has a theoretical implication: the ideal state $`\overline{S}`$ has a number of equivalent representations, e.g. those formed by permuting or swapping configuration details whose ordering or equivalence is unimportant. This multiplicity could be a benefit or a hazard to the task of implementation of policy. This remains to be determined.
Having established the existence of an ideal state, the second basic assumption is:
###### Basic assumption 2
The long term aim of system administration is to optimize the policy $`P(t)`$ for maximum productivity, insofar as this is allowed by local constraints. The short term aim is to keep the system as close to the resulting ideal state $`\overline{S}_p(t)`$ as possible, i.e. to minimize fluctuations $`\delta S=\overline{S}S`$.
One expects the average state $`\overline{S}_p`$ to exhibit persistent behaviour however, i.e. be invariant for periods approaching the duration of the system policy. Note that, what we are essentially doing, by making the assertion of an ideal state, is to separate slowly varying changes from quickly varying changes.
$$S(t)=\overline{S}_p(t)+\delta S(t).$$
(9)
Errors and misconfigurations (fluctuations $`\delta S`$) can accumulate over short periods of time, shorter than the time scale over which the average or ideal state changes. In terms of the relative rates of change:
$$\mathrm{max}\left|\frac{1}{\delta S}\frac{d}{dt}\delta S\right|\mathrm{max}\left|\frac{1}{\overline{S}_p}\frac{d}{dt}\overline{S}_p\right|.$$
(10)
The business of system administration is therefore a problem in regulation, or in minimizing the effect of $`\delta S`$.
We arrive at the following: a theory of system administration would attempt to answer the questions:
> Is there an optimal strategy for keeping the system as close as possible to its ideal state, and maximize its productivity?
To answer these questions, we must understand more deeply the meaning of the abstract formulation above. To begin, we backtrack and re-examine the underpinning concepts.
## 5 A generic computer system
In order to elucidate the goals of computer configuration and maintenance, it will be necessary to identify the main characteristics of computer systems at a suitable level of abstraction. These include finding:
* Relevant variables,
* Invariances,
* Persistent structures,
* Sources of information loss or entropy,
which affect the principal goals. Several studies of computer systems have attempted to identify such qualities and it is hypothesized that a suitably abstracted description can be built on the few simple principles identified by these authors.
The basic model of a computer is that of a dynamical community of processes and resources, coupled to an external environment (an external source or force). The source includes the stochastic influences of all of the users of the system, and any other computer systems which communicate with hosts within the perimeter of our own system. As pointed out in ref. , the issue of networking does not increase the complexity of the administration problem, only its localization and perhaps its magnitude. A set of networked hosts, sending external messages, is simply a single virtual host with internal inter-process communication.
The variables, important in characterizing the usage of a computer system, are measures of average behaviour, such as rate of work, numbers of processes, network connections and so forth. Other measures, such as average service latencies, affect the system only at the level of the network. Latencies are very complex phenomena and are unlikely to be predictable by any simple model.
Invariances refer to the independence of qualities and values to changes. In the long run, there are no features of a computer system which are fully constant, but for long periods of time, certain things can be considered invariant. For instance, the software tools one uses to edit a file usually make no difference to the outcome, thus the outcome of an editing operation may be considered invariant with respect to differences in software used; the CPU efficiency of the software used makes no difference to the result in most cases. Invariance could also mean that a particular piece of software never changes (is never upgraded), or that the content of a configuration file is fixed with respect to other changes. In the space of changes, such invariances may be considered to be ignorable coordinates.
Persistent structures are, like invariances, values or qualities which do not change over appreciable periods of time. This includes checksums of important software, kernel profiles of software; it might also include numbers of user accounts. Persistent structures are not expected to change. Changes in these structures might be considered anomalous behaviour.
An important characteristic of computer systems is that they are strongly coupled to human users’ behaviour patterns. The majority of human users follow strict daily and weekly work patterns and this is reflected in many measurements of system resource behaviour. A consequence of this is that measurements which are periodically constrained are distributed according to a Planck spectrum. The Planck spectrum can therefore be considered a general characteristic of computer statistics in many cases.
## 6 The scope of a theory of system administration
Even a limited theory of system administration should cover some key aspects of the problem:
* Policy determination,
* Strategic decisions about resource usage,
* Productivity considerations (the economics of the system),
* Empirical verification of strategies and policies,
* Efficiency of policy and of policy implementation,
* Efficiency of the system in doing its job.
More pragmatic details such as the need for software installation and upgrade have to be tackled at an abstract level, in terms of productivity, probability of failure, resource usage and so on. Software bugs can be addressed in terms of productivity or security. Security, in turn can be viewed as a contest for resources at the level of the system.
The benefits of automation versus human incursion are often discussed in system administration, sometimes as a bone of contention. This is one area that a theory of system administration can address objectively and have a real prospect of answering once and for all. An aspect of this will be discussed later as an example.
### 6.1 Measures and characters
As an empirical science, system administration suffers from many shortcomings. It has all of the problems associated with the social sciences: statistical measures are seldom forthcoming, experimental repeatability is a luxury, and sufficient repetition to obtain statistically meaningful samples is a near impossibility. The conditions under which measurements are made are constantly changing. The situation is somewhat analogous to that of non-equilibrium statistical mechanics in physics, but markedly less controlled.
The characteristics which are of interest to us refer to the actions and results which inter-weave in the dynamical behaviour of the system. These include the quality of actions of the system administrator, in relation to the prescribed policy, a typical characterization of the environment which affects the system. The measurements which are most useful are those based on persistent variables, since these have a stable value. Other fluctuating values can be treated stochastically or averaged out into persistent values.
The following measures will be useful in formulating ‘pay-off’ matrices for administration models, as in the example to follow below. The accuracy with which a policy is implemented by an agent of system management (human or automatic system) can be gauged with the following ratio:
$`\mathrm{Accuracy}={\displaystyle \frac{\mathrm{Number}\mathrm{of}\mathrm{policy}\mathrm{actions}}{\mathrm{All}\mathrm{actionsperformed}}}`$ (11)
i.e. the fraction of work which is within prescribed guidelines. In algebraic terms:
$`\alpha ={\displaystyle \frac{N_p(t)}{N(t)}}={\displaystyle \frac{{\displaystyle \underset{(aP)}{}}N_a}{{\displaystyle \underset{(a)}{}}N_a}}`$ (12)
For humans $`\alpha 1`$. For any bug-free automatic system, $`\alpha =1`$. Similarly, one may define the efficiency of a system by its use of resources (memory and CPU share):
$`\mathrm{Efficiency}=\mathrm{Accuracy}\times \left(1{\displaystyle \frac{\mathrm{Resources}\mathrm{used}}{\mathrm{Resources}\mathrm{available}}}\right)`$ (13)
In algebraic terms:
$`\epsilon =\alpha \left(1{\displaystyle \frac{{\displaystyle \underset{(aP)}{}}r_a}{{\displaystyle \underset{(a)}{}}r_a}}\right)`$ (14)
i.e. the more resources which are consumed in implementing a policy, the less efficient it can be considered to be.
Other measures are more useful for describing the relationship of a computer system to its environment, or the influential forces which steer its dynamical evolution. The response of a computer system to its users is characterized by averages which fluctuate in time. Human society’s diurnal work pattern imposes a twenty four hour periodic character on these measurements and a also a weekly work pattern, which is dominant during weekdays and slight at weekends (at least in the Western world). The periodic topology implies that the distribution of resource usage takes on the special form of a Planck distribution with a Gaussian component, by analogy with statistical physics at temperature $`T`$:
$`D(\lambda )=Ae^{\left(\frac{(\lambda \overline{\lambda })^2}{2\sigma ^2}\right)}+{\displaystyle \frac{B}{(\lambda \lambda _0)^3(e^{1/(\lambda \lambda _0)T}1)}}.`$ (15)
$`\lambda `$ is the deviation of a measurement from its average value over a period. The values of the constants $`A`$, $`B`$, $`\lambda _0`$ and $`T`$ may be chosen to fit the behaviour of any variable which is strongly coupled to periodic usage. Their absolute values have no significance, since there is no ‘standard candle’ computer system to compare to, but changes relative to the local norm could be interpreted as anomalies. Non-zero $`A`$ allows for the presence of additional Gaussian noise in some measurements.
### 6.2 Interactions of time scales
The identification of suitable time-scales is of crucial importance to any dynamical problem. Time-scales control rates of competition which lead to balance, and also rates of change.
It is easy to show that human administrators only compete with automatic systems in speed and efficiency at times of the day when they have nothing pressing to do. Indeed, it is always possible to arrange for an automatic system to beat a human, provided it can run in overlapping instantiations. A straightforward comparison of the time-scales involved in automated maintenance, to those of manual human maintenance can be made for any operation which is programmable in an automatic system with available technology.
Alarm systems which merely notify humans of errors and then rely on a human response are intrinsically slower than automatic systems which repair errors, provided the alarms represent errors which can be corrected with current automation.
The response time $`t_{\mathrm{auto}}`$ of a automatic machine system $`M`$, falls between two bounds (see figure 1)
$`nT_p+T_e(A)t_{\mathrm{auto}}T_e(A)`$ (16)
where $`T_p`$ is the scheduling period for regular execution of the system (e.g. the cron interval, typically half-hour to an hour), $`T_e(A)`$ is the execution time of the automatic system (typically seconds). The integer $`n0`$ since the number of iterations of the automatic system required to fix a problem might be greater than one. The time required to make a decision about the correct course of action $`T_d(A)`$ is negligible for the automatic system.
For a human being, making a decision based on a predecided policy, the response time $`t_{\mathrm{human}}`$ falls between the limits:
$`\mathrm{}t_{\mathrm{human}}T_w(H)+T_d(H)+T_e(H).`$ (17)
$`T_d(H)`$ is again the decision time, or time required to determine the correct policy response (typically seconds to minutes). $`T_e(H)`$ is the time required for a human to execute the required remedy (typically seconds to minutes). $`T_w(H)`$ is the time for which the human is busy or unavailable to respond to the request, i.e. the wait-time. The availability of human beings is limited by social and physiological considerations. In a simple way, one can expect this to follow a pattern in which the response time is greatest during the night; simplistically, if one assumes that humans sleep 8 hours,
$$T_w(H)>4(1+\mathrm{sin}(t/24)),$$
(18)
where time is measured in hours, whereas
$`T_w(A)0.`$ (19)
We can note that human response times are usually much longer than the corresponding machine response times,
$`T_d(A)T_d(H)`$
$`T_e(A)T_e(H)`$ (20)
and that the periodic interval of execution of the automatic system is generally taken to be greater than the execution time of the automatic system
$`T_pT_e,`$ (21)
thus avoiding overlapping executions (though this is not necessarily a problem, see the discussion of adaptive locks) It is always possible to choose the scheduling interval to be arbitrarily close to $`T_e(A)`$ (i.e. as short as one likes). Then provided,
$`T_w(H)>T_e(A)`$ (22)
the automatic system can always win over a human. This last inequality requires qualification however, since very long jobs (such as backups or file tree parses) increase exponentially in time with the size of the file tree concerned. This makes a prediction: it tells us that one should always arrange to allow such long jobs to be run last in a sequence of maintenance tasks, and also in overlappable threads. This means that long jobs will not hinder the rapid execution of a maintenance program.
Cfengine allows overlapping runs using its scheme of adaptive locks. Thus, by scheduling long jobs last in a cfengine program, it is virtually always possible for cfengine to beat a human, unless it is prevented from running, or the human is given the chance to respond with a head-start; this seldom happens by chance.
## 7 Primitive moves
Having identified the principle aims and methods of system administration, one is free to represent a model for these in any convenient calculational scheme. Almost immediately, one is confounded by the multiplicity, or non-uniqueness of the mapping between problem and solution: It is common-lore amongst system administrators, and it is to be expected logically in any causal web, that
* One problem can have several solutions.
* Several problems can be solved with a single solution.
How should one classify such mappings? By coarse-graining? Some degree of coarse classification is inevitable to make the analysis tractable, but it needs to be performed in a well-defined way. To some extent, we have already dealt with this problem in the factoring out of redundant expression in section 2. However the same problem returns in specifying the actions required to maintain the ideal state.
In order to unravel this situation as far as possible, it is reasonable to try to express problems and solutions in terms of linear combinations of primitive actions. The analysis of primitive operations has already been considered by Burgess in ref. in developing automated approach to system administration. There is little to add here, except to say that it is required that every implementable policy be decomposable as a combination of these primitives.
The available channels for action, i.e. the possible moves which a ‘player of the computer system game’ (user or administrator) can choose from, form a huge set if one views them at the level of the user. Formulating generic activity would be an intractable problem if one chose to consider every nuance of the system, viewed from a user perspective. Fortunately it is possible to break down the variety of activities available to users into a number of primitive actions. Any task can be considered as some linear combination of these few basic actions. The actions are:
| Primitive type $`T^i`$ | Comments/Examples |
| --- | --- |
| Create file | |
| Delete file | Tidy garbage |
| Rename file | Disable |
| Edit file | Used in configuration |
| Access control | Permissions |
| Request resource | Read/Mount |
| Copy file | Read/write |
| Process control | Start/stop |
| Process priority | Nice |
| Configure device | |
We should be careful to distinguish between how functions are implemented and how they can be decomposed. The method of implementation is not necessarily relevant to the analysis. What is important is that there exists a finite number of primitive actions which can be used to express all others in combinations.
Are these primitives sufficient in themselves? Could we implement the following policy, for instance: downloading of pornographic material between the hours of 9:00 and 17:00 is forbidden? If such a policy is implementable by an automatic system, it must be possible to filter content-specific data. Such a filter would need a configuration file which would need to be edited. The time-dependent behaviour could be handled by a scheduler, also configured by a text file. These configuration details are all implementable with file editing and process control. The ability of software to perform the task has to be assumed. This has nothing to do with management of the system. If the same job is to be carried out by a human, then the model of the management system must be extended to include humans, in which case job control and job definition require the analogous concepts to file editing and process control, for human brains. In other words, when humans are involved in a theory of manual work, they must be considered a part of the computer system.
## 8 The ideal average state
In order to have a chance of repairing damage, or maintaining a detailed balance of resources, we need to be able to trace the development or history of the system, from an ideal average state at an initial time, to a less than ideal state at a later time. In accordance with the axioms lain out at the beginning of the paper, it is assumed that the ideal state is determined as a matter of policy, by local considerations.
Many minor changes take place all the time in a computer system; these are healthy. Programs are started and stopped, files are created and destroyed: this is part of the work done by the system. However, certain features of the system should not change greatly (they should be persistent, at least on average). For instance, resources like disks and network services should be available to users at all times. If a crucial service falls out, then it affects other changes in the system.
Some changes are important to operation of the system, others are unimportant. For instance, it would be unimportant if one swapped the process ID’s of two programs. The process ID is just a label which has no bearing on the performance of the system or the productivity of users. However, if one process stopped running prematurely, this would be a change of state.
If we know what changes have taken place to move the system away from the ideal state, it should be possible to undo them, provided these do not involve the destruction of useful work. To accomplish this tracking of changes, in formal terms, we need to quantify the state of the system with respect to specific changes<sup>3</sup><sup>3</sup>3Note that, while one is interested in tracking changes in principle, in order to formulate the theory of system changes, this does not imply memorizing changes in a system is a desirable thing to do. Some system administration tools attempt to do this, often unsuccessfully, but as a counter-example one has cfengine which simply acts as a generic counter-force, pushing the system towards the ideal state, regardless of what specific historical chain the system follow.. Suppose one considers system administration as a game, framed on a lattice of $`n`$-dimensions, and suppose that the system has an ideal state located at the origin of this lattice, based on a policy and described in terms of primitive system variables. Each node of the lattice is a new state of the system. Let us suppose that the aim of the game is to remain as close as possible to the ideal state, i.e. the origin of this discrete space. How can one formulate such a game? How many dimensions does the lattice extend into, and what do they represent? These questions are central to formulating an analysis.
In a general sense, a computer system is a dynamical system like any other, and it must follow the same basic principles as any set of variables which changes in time. Let $`\varphi _i(t)`$, where $`i=1,2,\mathrm{}N`$ be the set of measurable variables which can be associated with a computer system. A canonically complete dynamical system can be associated with the set of phase-space variables,
$$q_i(t),\dot{q}_i(t),$$
i.e. the variables and their time derivatives. Not all variables can be considered differentiable functions of time, but it will be possible to give the derivative a meaning even for discrete variables, so this may be regarded symbolically for the present. Given that the values of these variables can change statistically with time (the nature of this variation will be qualified later), at any time $`t`$, we can decompose the value of $`q(t)`$ into a local average and a fluctuating piece.
$`q(t)=\overline{q}(t)+\delta q(t).`$ (23)
This means essentially decomposing $`q(t)`$ into fast and slowly changing variables. The average value $`\overline{q}(t)`$ varies only slowly with time, but many rapid changes $`\delta q(t)`$ fluctuate about the average value. The average may be defined by
$`\overline{q}_i(t)={\displaystyle \frac{1}{tt_i}}{\displaystyle _{t_i}^t}q_i(t^{})𝑑t^{},`$ (24)
where $`tt_i`$ is the interval over which the average is taken, and it is assumed that
$`tt_i<tt_0,`$ (25)
where $`t_0`$ is the ‘zeroth’ time at which the system was in the ideal state. The rate at which variables are changing $`\dot{q}(t)`$ can also be measured. A similar procedure can be implemented for the $`N`$ derivatives and their local average values.
For the sake of characterizing the state of the system, one is interested in change in the average values since some ideal zeroth time $`t_0`$:
$`d_i\{\overline{q}_i(t)\overline{q}_i(t_0),\overline{\dot{q_i}}(t)\overline{\dot{q_i}}(t_0)\}.`$ (26)
In terms of the deviations $`d_i`$ in key system variables, one may postulate a $`2N`$-dimensional lattice whose independent, orthogonal axes are the $`n=2N`$ variables of the phase space $`d_i`$, for $`i=1\mathrm{}n`$. Positions on this lattice are denoted by the vector of these component deviations. It is collectively denoted $`\stackrel{}{d}`$.
Suppose now that the system has deviated from the ideal state at $`\stackrel{}{0}`$ and has reached a point $`\stackrel{}{d}`$ on the lattice (see figure 2). The number of equivalent paths $`H(\stackrel{}{d})`$ back to the ideal state, is
$`H(\stackrel{}{d})={\displaystyle \frac{({\displaystyle \underset{j=1}{\overset{n}{}}}d_j)!}{{\displaystyle \underset{k=1}{\overset{n}{}}}(d_k!)}}`$ (27)
This grows rapidly with the Euclidean distance $`|\stackrel{}{d}|`$
$`|\stackrel{}{d}|d=\sqrt{{\displaystyle \underset{i=1}{\overset{n}{}}}(d_i)^2}.`$ (28)
$`H(\stackrel{}{d})`$ may be considered as a measure of the entropy, or disorder in the system. The entropy may be thought of as measuring the ‘hopelessness’ of finding the original route which led to the deviation. If all the paths are equivalent, i.e. the particular route by which the current state was achieved was not important, then it measures the number of equivalent ways in which the deviation can be fixed.
If the path is important then a different interpretation is more appropriate. In common with its analogue from physics, $`H`$ may be thought of as a measure of the amount of potential work has been lost to the system as a result of its deviation from the ideal state. Or conversely, here it may be considered a measure of the amount of work which would have to be expended in order to return the system to its ideal state. To gauge how quickly this grows with distance, one may compute the rate of increase in numbers of paths as $`\stackrel{}{d}`$ increases. Define
$`\stackrel{d_i}{}H(\stackrel{}{d})`$ $`=`$ $`{\displaystyle \frac{H(\stackrel{}{d}+\stackrel{}{\mathrm{\Delta }d})H(\stackrel{}{d})}{|\stackrel{}{\mathrm{\Delta }d}|}}`$ (29)
$`=`$ $`H(d_1,..,d_i+1,..d_n)H(d_1,..,d_i,..d_n)`$
Thus we define the rate of increase on the discrete lattice by,
$`{\displaystyle \frac{\stackrel{d_i}{}H(\stackrel{}{d})}{H(\stackrel{}{d})}}={\displaystyle \frac{1}{(d_i+1)}}{\displaystyle \underset{ji}{\overset{n}{}}}d_j1.`$ (30)
This shows that the increase is in fact approximately proportional to the distance. In other words, the rate of increase is approximately exponential. Clearly, this simple quantification of cumulative system error indicates that deviations from an ideal state should be dealt with as quickly as possible, since it becomes increasingly difficult to make corrections as the errors are compounded.
The ideal state itself needs to be characterized in terms of reasonable tolerances in system variables. The important variables include the availability of resources (ability to create new files and processes) as well as the level of activity. In dynamical terms one considers a set of variables and their rates of change:
$`q(t),{\displaystyle \frac{dq(t)}{dt}}.`$ (31)
If the system is a complete characterization of every possible influence and change in the system, then these form a simplectic algebra and the behaviour of the system is, at least in principle, completely deterministic, if not exactly predictable. In most cases there are influences which are not completely known, or may be regarded as random. In that case, one moves from simple mechanical systems into to realm of statistical mechanics and non-equilibrium studies.
These underlying variables are only indirectly linked to the ideal state, through averaging.
The above view is quite simplistic. In reality there might not be only one ideal state, but a set of equivalent ideal states. These can all be formulated as direct sums or quotients of a simply-connected state space however, so these need not be of concern to the principle of the argument. Having identified an ideal state as a point in a vector space, or lattice, one is now free to discuss how changes in the forces or influences on the system lead to movements through the lattice.
## 9 Game theory and the contest for the ideal state
There are two separable issues in the ideal-state view of system administration. The distinction concerns the perceived intelligence behind the changes which lead to a degradation of the ideal state. We may classify changes as either random (stochastic) or as intentional (strategic) depending on the nature of the adversary.
This distinction is partly artificial: all changes can be traced back to the actions of humans at some level, but it is not always pertinent to do so. Not all users act in response to a specific provocation, or with a specific aim in mind. It just happens that their actions lead to a general degradation of the ideal state, no malice intended. This strikes back to the fundamental principle of detail, namely that high level effects wash out the specifics of low-level origins. Thus there is a part of the spectrum of changes which averages out to a kind of faceless background noise. The details of who did what are of no concern. Random influences have been analyzed in ref. and are found to follow a number of well-known statistical distributions. Their study is part of the problem to be solved, but not all of it.
The other part of the problem is the case of actions which may be regarded as being more carefully calculated, or following a systematic behavioural pattern. These are caused by conflicts of interest between system policy and user wishes. A suitable framework for analyzing conflicts of interest, in a closed system, is the theory of games. Game theory is about introducing players, with goals and aims, into a scheme of rules and then analyzing how much a player can win, according to those restrictions. Each move in a game affords the player a characteristic value, often referred to as the ‘payoff’. Game theory has been applied to warfare, to economics (commercial warfare) and many other situations. In this case, the game takes place on the $`n`$-dimensional board, spanned by the $`\stackrel{}{d}`$ vectors.
There are many types or classifications of game. Some games are trivial: one-person games of chance, for example, are not analyzable in terms of strategies, since the actions of the player are irrelevant to the outcome. In a sense, these are related to the first kind of deviation referred to above. Some situations in system administration fit this scenario. More interesting, is the case in which the outcome of the game can be determined by a specific choice of strategy on the part of the players. The most basic model for such a game is that of a two-person zero-sum game, or a game in which there are two players, and where the losses of one player are the gains of the other. ‘Zero sum’ is the law of conservation of currency (current).
Many games can be stated in terms of this basic model, although this is often a simplification of reality. Games in complex systems are rarely true zero-sum games: energy leaks out, money gets burned or printed and thus there is no exact zero-sum conservation.
### 9.1 Models
The basic valuables of system administration are the system resources: file space, CPU share, memory share and network share. The theory of system administration can be viewed as a competition for these resources and for user privileges. The central obstacle in formulating a scenario in terms of game theory is the classification of strategies and their evaluation in terms of a characteristic (payoff) matrix.
* As a zero sum, two person game system administration is a game between the collective users and the system administrator. The aim of the users is to consume all of the system resources, while the aim of the administrator is to keep the system as close as possible to its ideal state. Ideally, the system administrators strategies should always bring the system closer to the ideal state. This is the property of convergence referred to in ref. . The ideal outcome of this game is a stalemate, or equilibrium somewhere close to the ideal state.
This game is often one with perfect information since all the important moves are visible to both players, however both sides can engage in bluffing. Clearly the administrator can win, either by limiting or reducing the consumption of resources and by extending the resources of the system. A user can ‘win’ in a certain pessimistic sense by moving the ideal state so far from the ideal that the system crashes and thus the game ends.
* A more optimistic variant of the above, is to view the aim of users as being to produce as much useful work as possible. This is a more complicated aim, since users can now impede their own progress by consuming too many resources, thus impairing the system as a whole and preventing themselves from being able to work (users need to be environmentally friendly). Experience from reality shows that most users do not concern themselves with this aspect however; they see it as the system administrators job to deal with such problems when they arise.
* As a zero sum, $`N`$-person game one can make a more detailed model, in which users compete against one another in addition to the system administrator. The system administrator’s task then becomes to act as a kind of Robin Hood character, preventing any one user’s consumption of all resources, trying to distribute resources fairly. Again, the aim of the administrator is to maximize the duration of the game by keeping the system as close to the ideal state as possible.
### 9.2 Payoffs and work
The next obstacle concerns the level at which we decide to address the behaviour of the system. Appropriate measures can be defined at various levels.
In order to formulate the characteristic matrix (often called the pay-off matrix) we must identify the book-keeping parameters and aims by which one hopes to win the game. What is the currency of this system? In social systems one has money as the book-keeping parameter for transactions. In physical systems, one has energy as the book-keeping parameter. These quantities count resources, in some well-defined sense. An analogous quantity is needed in system administration.
* The aim of the system administrator is to keep the system alive and running so that users can perform useful work.
* The aim of benign users is to produce useful work using the system. The aim of malicious users is often to maximize their control over system resources.
In a community, games are not necessarily cut and dried zero-sum engagements. We are faced with a Nash problem, or prisoner’s dilemma, which often ends in a Nash equilibrium.
> A user of the system who pursues solely private interests, does not necessarily promote the best interest of the community as a whole.
In other words, users can shoot themselves in the proverbial foot by using up all the available resources on a finite system. This affects them as much as anyone else. The empirical evidence suggests that, on average, users consume resources at a rate which is periodic and polynomial in time.
$`W(t)\mathrm{sin}(\mathrm{\Omega }t){\displaystyle \underset{n}{}}c_nt^n.`$ (32)
A definition of work is required in order to quantify the production of useful work in a non-prejudicial manner. Clearly the term ‘useful work’ spans a wide variety of activities. Clearly work can increase and decrease (work can be lost through accidents), but this is not really germane to the problem at hand. The work generated by a user (physical and mental work and then computationally assisted results) is a function of the information input into the system by the user. Since the amount of computation resulting from a single input might be infinite, in practice, the function is an unknown.
In general, the pay-off in not just a scalar value, but a vector. This indicates that a game might actually be decomposable into a number of parallel but interacting games.
What is the value of a game? How much can a user or an attacker hope to win? The system administrator, or embodiment of system policy, is not interested in winning the game, but rather in confounding the game for users who gain too much control. The system administrator plays a similar role to that of a police force. In some vague sense, the administrator’s jobs is to make sure that resources are distributed fairly, according to the policies laid down for the computer society.
### 9.3 Strategy expression
In a realistic situation one expects both parties in the two-person game to use mixed strategies. The formulation of the game theoretical pay-off matrix requires one to consider the strategies which the players can adopt. Again, the number of possible strategies is huge and the scope for strategic contrivance is almost infinite. In order to limit the formulation of the problem, it is necessary to break down strategies into linear combinations of primitives again. What is a strategy?
* A set of operations
* A schedule of operations
* Rules for counter-moves
In addition to simple strategies, there can be meta-strategies, or long-term goals. For instance, a nominal community strategy might be to:
* Maximize productivity or generation of work.
* Gain the largest feasible share of resources.
An attack strategy might be to
* Consume as many resources as possible.
* Destroy key resources.
Other strategies for attaining intermediate goals might include covert strategies such as bluffing (falsely naming files). Defensive strategies might involve taking out an attacker, counter attacking, or evasion (concealment), exploitation, trickery, antagonization, incessant complaint (spam), revenge etc. Security and privilege, levels of access, integrity and trust must be woven into algebraic measures for the pay-off. A means of expressing these devices must be formulated within a language which can be understood by system administrators, but which is primitive enough to enable the problem to be analyzed in an unambiguous fashion.
### 9.4 Stable and dominant strategies
It has been argued here, and in earlier papers, that computer systems can be viewed as fluctuating around statistically stable configurations, for the most part. This assumes that both users and system administration mechanisms are in approximate balance. Game theory is suited to finding equilibria, or stable superiorities in a set of strategies. Let us consider how game theory can be used to frame system behaviour as a contest for control of the system’s resources.
The simplest case of a two-person, zero-sum game is chosen.
We are interested in determining whether any optimal strategies can be adopted by the system (and its administrator) in order to maintain control of the system, i.e. in order to prevent users from winning control of the system. This situation is analogous to the analysis of dominant evolutionary strategies, considered by Hamilton and Maynard-Smith. These so-called Evolutionary Stable Strategies are the winning strategies favoured by natural selection mechanisms in the animal or plant kingdom. In our case, we are simply interested in strategies which are clear winners over all other strategies. If we consider the characteristic matrix, or pay-off matrix, as a function of strategies for attack and defense $`\pi (\sigma _a,\sigma _d)`$, then one may characterize a dominant attack-strategy $`\sigma _a^{}`$ by the criterion:
$`\pi (\sigma _a^{},\sigma _d)`$ $`>`$ $`\pi (\sigma _a,\sigma _d)`$ (33)
i.e. $`\sigma _a^{}`$ must be a better move than any other strategy against and arbitrary counter-move $`\sigma _d`$. If this is the case, then there exists at least one pure strategy which is optimal for the attacker. Similarly, an optimal defensive strategy $`\sigma _d^{}`$ is characterized by:
$`\pi (\sigma _a,\sigma _d^{})`$ $`>`$ $`\pi (\sigma _a,\sigma _d)`$ (34)
A more general situation is that one can find a winning mixture of strategies $`\mathrm{\Sigma }`$ (a linear combination of pure strategies)
$`\mathrm{\Sigma }={\displaystyle \frac{1}{N}}{\displaystyle \underset{i}{\overset{N}{}}}c_i\sigma _i.`$ (35)
Then if the dominant mixture of strategies $`\mathrm{\Sigma }_a^{}`$ satisfies,
$`\pi (\mathrm{\Sigma }_a^{},\mathrm{\Sigma }_d)`$ $`>`$ $`\pi (\mathrm{\Sigma }_a,\mathrm{\Sigma }_d)`$ (36)
then the attacker must win, but if some optimal mixture of strategies $`\mathrm{\Sigma }_d^{}`$ satisfies,
$`\pi (\mathrm{\Sigma }_a,\mathrm{\Sigma }_d^{})`$ $`>`$ $`\pi (\mathrm{\Sigma }_a,\mathrm{\Sigma }_d),`$ (37)
then the defender must prevail. It is this final solution which one hopes to find in order to secure a stable computer environment.
To illustrate this idea, consider an example of some importance, namely the issue of garbage collection. The need for forced garbage collection has been argued on several occasions, but the value of this strategy to system rule has not been analyzed previously.
The first issue is to determine the currency of this game. What payment will be transferred from one player to the other in play? Here, there are three relevant measurements to take into account: (i) the amount of resources consumed by the attacker (or freed by the defender), and sociological rewards: (ii) ‘goodwill’ or (iii) ‘privilege’ which are conferred as a result of sticking to the policy rules. These latter rewards can most easily be combined into an effective variable ‘satisfaction’. Then the player who can’t get no satisfaction is the poorer one. A satisfaction measure is needed in order to balance the situation in which the system administrator prevents users from using any resources at all. This is clearly not a defensible use of the system, thus the system defenses should be penalized for restricting users too much. The characteristic matrix now has two contributions,
$`\pi =\pi _r(\mathrm{resources})+\pi _s(\mathrm{satisfaction}).`$ (38)
It is convenient to define
$`\pi _r\pi (\mathrm{resources})={\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\mathrm{Resources}\mathrm{won}}{\mathrm{Total}\mathrm{resources}}}\right).`$ (39)
Satisfaction $`\pi _s`$ is assigned arbitrarily from values from plus to minus one half, such that,
$`{\displaystyle \frac{1}{2}}\pi _r+{\displaystyle \frac{1}{2}}`$
$`{\displaystyle \frac{1}{2}}\pi _s+{\displaystyle \frac{1}{2}}`$
$`1\pi +1.`$ (40)
The pay-off is related to the movements made through the lattice $`\stackrel{}{d}`$. The different strategies can now be regarded as duels, or games of timing.
| Users/System | Ask to tidy | Tidy by date | Tidy above | Quotas |
| --- | --- | --- | --- | --- |
| | | | Threshold | |
| Tidy when asked | $`\pi (1,1)`$ | $`\pi (1,2)`$ | $`\pi (1,3)`$ | $`\pi (1,4)`$ |
| Never tidy | $`\pi (2,1)`$ | $`\pi (2,2)`$ | $`\pi (2,3)`$ | $`\pi (2,4)`$ |
| Conceal files | $`\pi (3,1)`$ | $`\pi (3,2)`$ | $`\pi (3,3)`$ | $`\pi (3,4)`$ |
| Change timestamps | $`\pi (4,1)`$ | $`\pi (4,2)`$ | $`\pi (4,3)`$ | $`\pi (4,4)`$ |
The elements of the characteristic matrix must now be modelled by suitable algebraic or constant terms. The rate at which users produce files may be written
$`r_u={\displaystyle \frac{n_br_b+n_gr_g}{n_b+n_g}},`$ (41)
where $`r_b`$ is the rate for bad users and $`r_g`$ is the rate for good users. The total number of users $`n_u=n_b+n_g`$. From the authors experience, the ratio $`n_b/n_g`$ is about one percent. The rate can be expressed as a scaled number between zero and one, for convenience, so that $`r_b=1r_g`$.
The payoff in terms of the consumption of resources by users, to the users themselves, is then
$`\pi _u={\displaystyle \frac{1}{2}}{\displaystyle _0^T}𝑑t{\displaystyle \frac{r_u(\mathrm{sin}(2\pi t/24)+1)}{R_{\mathrm{tot}}}},`$ (42)
where the factor of 24 is the human daily rhythm, measured in hours, and $`R_{\mathrm{tot}}`$ is the total amount of resources to be consumed. Note that, by considering only good user or bad users, one has a corresponding expression for $`\pi _g`$ and $`\pi _b`$, with $`r_u`$ replaced by $`r_g`$ or $`r_b`$ respectively. An automatic garbage collection system results in a negative pay-off to users, i.e. a pay-off to the system administrator. This may be written
$`\pi _a={\displaystyle \frac{1}{2}}{\displaystyle _0^T}𝑑t{\displaystyle \frac{r_a(\mathrm{sin}(2\pi t/T_p)+1)}{R_{\mathrm{tot}}}},`$ (43)
where $`T_p`$ is the period of execution for the automatic system, considered earlier. This is typically hourly or more often, so the frequency of the automatic cycle is some twenty times greater than that of the human cycle. The rate of resource-freeing $`r_a`$ is also greater than $`r_u`$, since file deletion takes little time compared to file creation, and also an automated system will be faster than a human. The quota payoff yields a fixed allocation of resources, which are assumed to be distributed equally amongst users and thus each quota slice assumed to be unavailable to other users. The users are nonchalant, so $`\pi _s=0`$ here, but the quota yields
$`\pi _q=+{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{1}{n_b+n_g}}\right).`$ (44)
The matrix elements are expressed in terms of these.
* Here $`\pi _s=\frac{1}{2}`$ since the system administrator is maximally satisfied by the users’ behaviour. $`\pi _r`$ is the rate of file creation by good users $`\pi _g`$, i.e. only legal files are produced. Comparing the strategies, it is clear that $`\pi (1,1)=\pi (1,2)=\pi (1,3)`$.
* Here $`\pi _s=0`$ since the users are dissatisfied by the quotas, but the system administrator must be penalized for restricting the functionality of the system. With fixed quotas, users cannot generate large temporary files. $`\pi _q`$ is the fixed quota payoff, a fair slice of the resources. Clearly $`\pi (4,1)=\pi (4,2)=\pi (4,3)=\pi (4,4)`$. This tells us that quotas put a straight-jacket on the system. The game has a fixed value if this strategy is adopted by system administrators. However, it does not mean that this is the best strategy, according to the rules of the game, since the system administrator loses points for restrictive practices. This is yet to be determined.
* Here $`\pi _s=\frac{1}{2}`$ since the system administrator is maximally dissatisfied with users’ refusal to tidy their files. The pay-off for users is also maximal in taking control of resources, since the system administrator does nothing to prevent this, thus $`\pi _r=\pi _u`$. Examining the strategies, one find that $`\pi (2,1)=\pi (3,1)=\pi (3,2)=\pi (3,3)=\pi (4,1)=\pi (4,2)`$.
* Here $`\pi _s=\frac{1}{2}`$ since the system administrator is maximally dissatisfied with users’ refusal to tidy their files. The pay-off for users is now mitigated by the action of the automatic system which works in competition, thus $`\pi _r=\pi _u\pi _a`$. The automatic system is invalidated by user bluffing (file concealment).
* Here $`\pi _s=\frac{1}{2}`$ since the system administrator is maximally dissatisfied with users’ refusal to tidy their files. The pay-off for users is mitigated by the automatic system, but this does not activate until some threshold time is reached, i.e. until $`t>t_0`$. Since changing the date cannot conceal files from the automatic system, when they are tidied above threshold, we have $`\pi (2,3)=\pi (4,3)`$.
Thus, in summary, the characteristic matrix is given by:
$`\pi (u,s)=\left(\begin{array}{cccc}\frac{1}{2}+\pi _g(t)& \frac{1}{2}+\pi _g(t)& \frac{1}{2}+\pi _g(t)& \pi _q\\ \frac{1}{2}+\pi _u(t)& \frac{1}{2}+\pi _u(t)+\pi _a(t)& \frac{1}{2}+\pi _u(t)+\pi _a(t)\theta (t_0t)& \pi _q\\ \frac{1}{2}+\pi _u(t)& \frac{1}{2}+\pi _u(t)& \frac{1}{2}+\pi _u(t)& \pi _q\\ \frac{1}{2}+\pi _u(t)& \frac{1}{2}+\pi _u(t)& \frac{1}{2}+\pi _u(t)+\pi _a(t)\theta (t_0t)& \pi _q\end{array}\right),`$ (49)
where the step function is defined by,
$`\theta (t_0t)=\{\begin{array}{c}1(tt_0)\\ 0(t<t_0)\end{array},`$ (52)
and represents the time-delay in starting the automatic tidying system in the case of tidy-above-threshold.
It is possible to make several remarks about the relative sizes of these contributions. The automatic system works at least as fast as any human so, by design, in this simple model we have
$`{\displaystyle \frac{1}{2}}|\pi _a||\pi _u||\pi _g|0,`$ (53)
for all times. In addition , for short times $`\pi _q>\pi _u`$, but users can quickly fill their quota and overtake this. In a zero-sum game, the automatic system can never tidy garbage faster than users can create it, so the first inequality is always saturated. From the nature of the cumulative pay-offs, we can also say that
$`({\displaystyle \frac{1}{2}}+\pi _u)({\displaystyle \frac{1}{2}}+\pi _u+\pi _a\theta (t_0t))({\displaystyle \frac{1}{2}}+\pi _u+\pi _a),`$ (54)
and
$`|{\displaystyle \frac{1}{2}}+\pi _u||\pi _g{\displaystyle \frac{1}{2}}|.`$ (55)
Let us now apply these results to a modest strategy of automatic tidying, of garbage, once per day, in order to illustrate the utility of the game formulation. The first step is to compute the pay-off rate contributions. Referring to figure 4, one sees that the automatic system can always match users’ moves. As drawn, the daily ripples of the automatic system are in phase with the users’ activity. This is not realistic, since tidying would normally be done at night when user activity is low, however such details need not concern us in this illustrative example.
The policy we have created in setting up the rules of play for the game, penalizes the system administrator for employing strict quota shares. Even so, users do not gain much from this, because quotas are constant for all time. A quota is a severe handicap to users in the game, except for very short times before users reach their quota limits. Quotas could be considered cheating in such a game, since they determine the outcome even before play commences. There is no longer a contest. Moreover, comparing the values in the figure, it is possible to see how resource inefficient quotas are. Users cannot create temporary files which exceed these hard and fast quotas. An immunity type model which allows fluctuations is a considerably more resource efficient strategy, since it allows users to span all the available resources for short periods of time, without consuming them for ever.
Any two-person zero-sum game has a solution, either in terms of a pair of optimal pure strategies or as a pair of optimal mixed strategies. This result is known as the minimax theorem and was proved by Von Neumann. The solution is found as the balance between one player’s attempt to maximize his pay-off and the other player’s attempting to minimize the opponent’s result.
In general one can say of the pay-off matrix that
$`\underset{}{\mathrm{max}}\underset{}{\mathrm{min}}\pi _{rc}\underset{}{\mathrm{min}}\underset{}{\mathrm{max}}\pi _{rc},`$ (56)
where the arrows refer to the directions of increasing rows ($``$) and columns ($``$). The left hand side is the least users can hope to win (or conversely the most that the system administrator can hope to keep) and the right is the most users can hope to win (or conversely the least the system admin can hope to keep). If we have
$`\underset{}{\mathrm{max}}\underset{}{\mathrm{min}}\pi _{rc}=\underset{}{\mathrm{min}}\underset{}{\mathrm{max}}\pi _{rc},`$ (57)
it implies the existence of a pair of single, pure strategies $`(r^{},c^{})`$ which are optimal for both players, regardless of what the other does. If the equality is not satisfied, then the minimax theorem tells us that there exist optimal mixtures of strategies, where each player selects at random from a number of pure strategies with a certain probability weight.
The situation for our time-dependent example matrix is different for small $`t`$ and for large $`t`$. The distinction depends on whether users have had time to exceed fixed quotas or not; thus ‘small $`t`$’ refers to times when users are not impeded by the imposition of quotas.
For small $`t`$, we have:
$`\underset{}{\mathrm{max}}\underset{}{\mathrm{min}}\pi _{rc}`$ $`=`$ $`\underset{}{\mathrm{max}}\left(\begin{array}{c}\pi _g\frac{1}{2}\\ \frac{1}{2}+\pi _u+\pi _a\\ \frac{1}{2}+\pi _u\\ \frac{1}{2}+\pi _u+\pi _a\theta (t_0t)\end{array}\right)`$ (62)
$`=`$ $`{\displaystyle \frac{1}{2}}+\pi _u.`$ (63)
The ordering of sizes in the above minimum vector is:
$`{\displaystyle \frac{1}{2}}+\pi _u{\displaystyle \frac{1}{2}}+\pi _u+\pi _a\theta (t_0t)\pi _u+\pi _a\theta (t_0t)\pi _g{\displaystyle \frac{1}{2}}.`$ (64)
This is useful to know, if we should examine what happens when certain strategies are eliminated. For the opponent’s endeavours we have
$`\underset{}{\mathrm{min}}\underset{}{\mathrm{max}}\pi _{rc}`$ $`=`$ $`\underset{}{\mathrm{min}}({\displaystyle \frac{1}{2}}+\pi _u,{\displaystyle \frac{1}{2}}+\pi _u,{\displaystyle \frac{1}{2}}+\pi _u,\pi _q)`$ (65)
$`=`$ $`{\displaystyle \frac{1}{2}}+\pi _u.`$
This indicates that the equality in eqn. (57) is satisfied and there exists at least one pair of pure strategies which is optimal for both players. In this case, the pair is for users to conceal files, and for the system administrator to tidy by any means (these all contribute the same weight in eqn (65). Thus for small times, the users are always winning the game if we assume that they are allowed to bluff by concealment. If the possibility of concealment or bluffing is removed (perhaps through an improved technology used by the administrator), then the next best strategy is for users to bluff by changing the date. In that case, the best system administrator strategy is to tidy at threshold.
These results make qualitative sense and tally well with the author’s experience. The result also makes a prediction for system administration tools like cfengine. Cfengine must be able to see through attempts at bluffing if it is to be an effective opponent against the worst users.
For large times (when system resources are becoming or have become scarce), then the situation looks different. In this case one finds that
$`\underset{}{\mathrm{max}}\underset{}{\mathrm{min}}\pi _{rc}=\underset{}{\mathrm{min}}\underset{}{\mathrm{max}}\pi _{rc}=\pi _q.`$ (66)
In other words, the quota solution determines the outcome of the game for any user strategy. As already commented, this might be considered cheating or poor use of resources, at the very least. If one eliminates quotas from the game, then the results for small times hold also at large times.
This simple example of system administration as a strategic game between users and administrators was not intended to be as realistic as possible, rather it was intended as an illustration of the principles involved. Nevertheless, it is already clear that user bluffing and system quotas are strategies which are to be avoided in an efficient system. By following this basic plan, it should be possible to analyze more complex situations in future work.
### 9.5 The policy $`P(t)`$ and the pay-off matrix $`\pi (t)`$
At the beginning of this paper, we referred to a central axiom which involved the changing system policy $`P(t)`$. The characteristic (pay-off) matrix $`\pi _{rc}(t)`$ must clearly be related to this policy.
Let us suppose that the pay-off matrix is a $`u\times s`$ matrix, with $`u`$ user strategies and $`s`$ system strategies. The administrators strategies are limited by the policy, and the rewards are also limited, so both the dimension $`s`$ and form of the pay-off matrix are functions of the policy. The user’s strategies cannot be assumed to be limited by policy however, since ‘criminal’ users will ignore policy for personal gain. Although one may think of the dimension $`s[P(t)]`$ as being a functional of the policy, it would not be correct to think of $`u`$ as being a functional of the policy, since there can be no restriction on what users will try, simply as a result of law-giving. User’s actions can only be restricted by applying counter-measures within the
$`\pi =\pi _{rs[P(t)]}(P(t)).`$ (67)
It should be noted, however, that there can be no unique mapping between policy and pay-off matrix.
### 9.6 Change and future models
Expressing deterministic changes in generic computer systems would be a huge undertaking unless one restricted ambitions to general features and trends. Dynamical systems are difficult to trace, even in the simplest of cases, so one cannot expect to get very far without making significant simplifications. The aim of considering a dynamical theory is thus to characterize the significant trends of change which might occur, owing to idealized influences. A full discussion of this topic is beyond the scope of the present paper, however based on the axioms and deliberations presented here, it is possible to outline the way forward in studying them.
The expression of strategies in the previous section is too general to be useful for a fully general, dynamical theory. Taking account of every strategic detail would be a vast undertaking. Instead, one can analyze the development at the level of a generic computer system undergoing generic changes as a matter of principle. The purpose of such a vague preliminary investigation is to elucidate the relationship between the system administration game and the lattice description of the ideal state, presented in section 8.
Once a strategy mixture has been decided, one must address the fact that, in real-world games, the speed of information is finite. It will take a finite amount of time for a response to develop after a strategy is implemented. Moves and counter-moves do not follow a rigid time-plan as in games like chess. This kind of delay leads to races and duels for superiority between competing players. Delay is the province of linear response theory.
The aim, then, is to express the causal structure of system development in the foregoing mathematical language. In order to reduce the dynamical game to algebra we must express each of these in terms of basic primitives. Causality is about relating actions to outcomes, or changes of state $`\delta S`$. A general action $`A(t)`$ is built up from a number of primitive action-types $`T_a`$ (called the generators for the action transformation) in a linear combination
$`A(t)={\displaystyle \underset{i}{}}a_i(t)T_i`$ (68)
where $`a_i(t)`$ are functions of time (not necessarily differentiable, often step-like) and $`i`$ takes values which number the full spectrum of primitive actions. The $`T^i`$ are orthogonal vectors or matrices (indices suppressed), one for each primitive action type, which span an abstract vector space. This vector space is the chequerboard on which the game takes place.
Each complete action $`A(t)`$, results in a change in the state of the system, which may be denoted $`\delta S`$. An action can also be a causal chain of sub-actions, characterizing a sequence of changes in the state. This type of causal relationship is summarized by a Green function, propagator, or response-function formulation:
$`\delta S(t)={\displaystyle 𝑑t^{}G(t,t^{})A(t^{})},`$ (69)
where $`G(t,t^{})`$ is the two-point response function, as yet unspecified. If the rules of the game are independent of time, then $`G(t,t^{})=G(tt^{})`$; if the rules change over time, then $`G(t,t^{})=G(tt^{},t+t^{})`$. In this language of dynamical systems, an action plays the role of a source or driving force for the system. The equation above may be inverted to provide an inhomogeneous differential equation for the changing state of the system. If one formally introduces a differential operator $`D_t`$ which is the inverse of the response function:
$`{\displaystyle 𝑑t^{}D_tG(t,t^{})}=1,`$ (70)
then the differential equation may be written, schematically:
$`D_tS(t)=A(t),`$ (71)
where, as ad hoc an example one might have,
$`D_t{\displaystyle \frac{d^2}{dt^2}}+i\gamma {\displaystyle \frac{d}{dt}}+\omega _0^2,`$ (72)
for an approximately periodic system which degrades over time, like a damped harmonic oscillator. Each action $`A(t)`$ thus leads to a response or change of state; this in turn implies that the state of the system must be a linear combination of the same action types:
$`S(t)={\displaystyle \underset{i}{}}s_i(t)T_i.`$ (73)
The state is thus defined on the same lattice, or chequerboard as the actions themselves. Differential (difference) characterizations of state have been studied in ref. ; this type of description is interesting, since it leads often to rich dynamics. Alternating periods of change and stability (riffles and pools in the flow of the system) might be best described by a difference representation.
Returning to the idea of the contest as a game, one writes a strategy as a statistical mixture of actions (i.e. moves in the game) $`A(t)`$, applied over an interval of time. This stochastic mixture specifies the boundary conditions under which the actions are applied. It may be formed as a linear combination of basic actions $`A_n`$, with probability weights $`w_n`$:
$`J(t)={\displaystyle \underset{n}{}}w_nA_n(t)={\displaystyle \underset{i}{}}p_iT_i,`$ (74)
The strategy vector $`J_i`$ is the vector of probabilities for each primitive action, given the chosen mixture of full actions $`A_n`$ for $`J`$. In other words, $`J_i`$ are the components of the decomposition of the strategy $`J(t)`$ on the space of primitive actions.
$`p_i={\displaystyle \underset{n}{}}w_n.`$ (75)
It is easy to normalize these so as to be actual probabilities which sum to unity
$`{\displaystyle \underset{i}{}}p_i1.`$ (76)
Notice that the specific representation of basis generators $`T_i`$ does not affect the strategy vector, since it only serves to label the lattice-work of independent actions. There is no unique labelling. The components with respect to the basis must be related by a response function $`\mathrm{\Pi }_{ij}(t,t^{})`$
$`\delta S_i={\displaystyle 𝑑t^{}\mathrm{\Pi }_{ij}(t,t^{})J_j(t^{})}.`$ (77)
The matrix value distribution is related to the pay-off matrix, and a basis of so-called ladder operators, also called creation and annihilation operator. The represent do-action/undo-action operations of the users and system administrator.
$`\mathrm{\Pi }(t,t^{})\pi _{ij}\stackrel{}{d}|\widehat{S}_+(t)\widehat{S}_{}(t^{})|\stackrel{}{d}`$ (78)
where $`\widehat{S}_\pm `$ are operators which annihilate a configuration state at $`t^{}`$ and create a new configuration at time $`t`$. This is the generic mechanism by which the system develops. This form of description might seem unnecessarily formal, but it is actually highly useful, since the continuous generalization of this kind of dynamical system has been widely studied in statistical field theory. By picking out universal features of statistical models and restricting the scope of the computer system, there is a real chance of being able to build toy models which have qualitative, predictive power. However, this is no trivial undertaking and will be considered in a later paper.
Since the actions which configure a computer system form a lattice, and these primitive action types do not necessarily commute with one another, one concludes that a suitable idealization of the system administration’s stochastic dynamics is found in non-Abelian, statistical field theories. This line of study would be suitable for modelling resource availabilities for large numbers of users, in which all users behave approximately equally on average (like an ideal gas). This approach promises therefore to be relevant to the problem of anomaly detection and will be returned to in later work.
## 10 Summary
The aim of this paper has been to formulate a trustworthy framework for analyzing models of system administration. There is good cause to view computers as dynamical systems, approximated by mechanistic rules developing in time, with idealized properties which can be summarized by a finite state lattice. The theory of games has been employed in order to select between alternative strategies in a contest for machine resources, moving the state of the system through the lattice, as if on a chequerboard. It has been shown that it is possible to see system administration as the effort to keep the system close to an ideal state, by introducing countermeasures in the face of competitive resource consumption. This is the formal basis which opens the way for objective analyses in the field.
It is important to understand that, even an answer obtained with the assistance of a mathematical formalism is not necessarily the last word on the subject. Mathematics is only a tools for relating assumptions to conclusions, in an impartial way. With a mathematical approach, it becomes easier to see through personal opinions and vested interests when assumptions and methods are clearly and rigorously appraised. However, one can only distinguish between those possibilities which are taken into account. That means that every relevant strategy, or alternative, has to be considered, or else one could miss the crucial combination which wins the game. This is the limitation of game theory. It is not generally possible to determine strategies without creative input; this means that human intelligence will be required for the foreseeable future. There can be no zero-maintenance computer system. With this caution, how can one know that the ideal state of a system can be reached? How can one know that the system will not run away in an unstable spiral to catastrophe?
Two things are clear from the limited analysis here. The first is that purely dumb automatic systems are inadequate to perform every task in system administration today. Intelligent incursions are required to solve complex problems, to extend or adjust the strategies of the automatic system. Interestingly, this is the approach by which evolution has solved the immunity problem: the automatic responses of lymphocytes only go so far; the emergence of intelligence in humans has enabled us to develop medical research and develop drugs and other treatments against damage and disease. It seems naive to believe that any simple mechanistic system would be able to do any better than this; we can expect to require the assistance of humans at least until alternative machine intelligences have been developed.
The second point is that the use of quotas is a highly inefficient way of counteracting the effects of selfish users. A quota strategy can never approach the same level of productivity as one which is based on competitive counterforce. The optimal strategies for garbage collection are rather found to lie in the realm of the immunity model. However, it is a sobering thought that a persistent user, who is able to bluff the immune system into disregarding it, (like a cancer) will always win against the resource battle. The need for new technologies which can see through bluffs will be an ever present reality in the future. With the ability of encryption and compression systems to obscure file contents, this is a contest which will not be easily won by system administrators.
There is plenty of work to be done on the theory of system administration. This paper is merely a small push in the direction of progress.
I am grateful to Trond Reitan for a useful discussion about evolutionary stable strategies and to Hårek Haugerud, Lars Kristiansen and particularly Sigmund Straumsnes for their critical readings of the manuscript.
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# Probing flavor changing interactions in hadron collisions 11footnote 1This work was supported by National Natural Science Foundation of China.
Abstract
The subprocess $`ggt\overline{c}+\overline{t}c`$ in the two-Higgs-doublet model with flavor-changing scalar couplings is examined at the one loop level. With perturbative QCD factorization theorem, the corresponding cross sections for hadron-hadron collisions are computed numerically. The results are applicable to the whole mass range of the weakly coupled Higgs bosons. In case we could efficiently exclude the severe backgrounds of the $`t\overline{c}\left(\overline{t}c\right)`$ production signal, probing the flavor-changing top-charm-scalar vertex at hadron colliders would be very promising and accessible experimentally.
PACS number(s):13.85.Ni, 11.30.Hv, 12.60.Fr, 14.65.Ha
I. Introduction
There are stringent experimental constraints against the existence of tree level neutral flavor changing interactions, especially for light quarks. This leads to the naturalness to suppress the flavor changing neutral current interactions (FCNCs) for all kinds of model building studies, that is realized in terms of the Glashow-Iliopoulos-Maiani (GIM) mechanism in the standard model(SM).
The two-Higgs-doublet models(THDMs) are extensions of the SM in which one more scalar doublet is added. In order to forbid possible tree level FCNCs to appear in the models, Glashow and Weinberg proposed a neutral flavor conservation (NFC) condition by imposing discrete symmetries on the models. The common THDMs with the NFC condition can be divided into two categories, i.e., Model I and Model II. In Model I, both the up- and down-type quarks couple to the same one of the Higgs doublets, but in the Model II the up- and down-type quarks couple to the two Higgs doublets respectively. In fact, the suppression facts of the FCNC processes in the cases involving the light down-type quark sector have been observed experimentally, whereas those relevant to the up-type quark sector have not been established so well. Therefore the so-called Model-III of the THDM is proposed and it allows the FCNC Yukawa couplings which has the character that the couplings are related to the masses of the coupled flavors even at tree level.
As pointed out by Cheng, Sher and other authors , the Yukawa couplings are typically proportional to the masses of the coupled fermions at the vertices, it is rather natural to imagine having such Yukawa couplings for the FCNC interactions instead of placing the constraints on the theory. In this case, low energy limits on the FCNCs may be evaded, because the flavor changing couplings to the light quarks are small, and the suppression of the FCNC processes involving light quarks can be automatically satisfied. Then the imposition of discrete symmetries to obtain the NFC condition, therefore, is unnecessary, which is normally invoked in common THDMs to prevent the FCNC interactions at tree level. In the literature such the THDM without the NFC condition is called as the ‘third’ THDM (i.e., THDM III) , where the up-type and down-type quarks are allowed simultaneously to couple to more than one scalar doublet. In the framework of the model the effects of the FCNC interactions involving the heavy quark will be enhanced.
Now the top quark has been observed and a unexpected very large mass $`m_t=173.8\pm 5.2GeV`$ (world average value) has been recognized. This ‘extraordinary’ mass scale of the top quark may have many important implications pertaining to many outstanding issues. One of them is to test the FCNC processes requested in the Model-III. Namely, it is a ‘good place’ to observe the facts of the Model-III for the FCNC interactions. The measurement of FCNC processes involving top quark would provide an important test for the discrimination of various models. In the THDM III one would expect that large effects of the FCNC interactions could manifest themselves in the cases involving the massive top quark. Therefore testing the existence of the flavor changing scalar interactions involving top quarks are clearly important. D. Atwood et al. presented results of a calculation for the process $`e^+e^{}t\overline{c}`$(or $`\overline{t}c`$)in the THDM III, and they obtained $`R^{tc}/\lambda ^4`$ to be in the order of $`10^5`$ with proper parameters. Recently, Jiang et al.studied the production rates of $`e^+e^{}\gamma \gamma t\overline{c}+\overline{t}c`$ for the NLC and found that this process is more promising than the straight production via $`e^+e^{}`$ collisions for probing the FCNC interactions. Abraham et al. also investigated the anomalous $`\overline{t}q\gamma `$ couplings via single top quark production processes by considering effective Lagrangian of the lowest dimension with $`\gamma \gamma `$ collisions, and found the processes would be observable with suitable strength of anomalous coupling as long as b-tagging and suitable kinematic cuts are taken properly.
After the termination of the running of the LEP2, the hadron colliders Tevatron and LHC will be the only machines in searching for the FCNC processes. It is believed that more experimental events involving top quark will be collected in these hadronic machines. It will give a good chance to study the physics relevant to the FCNC processes of the Model-III. In this paper we are to study the problem and present complete one-loop calculation for the subprocess $`ggt\overline{c}`$(or $`\overline{t}c`$) to the order $`O(m_tm_c/m_W^2)`$ for the THDM Model-III. In fact, the obtained results in the paper are applicable to the whole mass range for weakly coupled Higgs bosons. The production cross sections of $`pp(\overline{p})ggt\overline{c}+\overline{t}c+X`$ are also given for the Tevatron and LHC. The paper is organized as follows: The details of the calculation are given in Sec. II. In Sec. III numerical results, discussions and a short summary are presented. Finally, the explicit expressions used in the paper are collected in appendix.
II. Calculation
In the third type of the two-Higgs-doublet model, the up-type and down-type quarks are allowed simultaneously to couple to more than one scalar doublet. Since there is no global symmetry that distinguishes the two doublets in the model, we will assume that only one of them $`(\varphi _1)`$ develops a vacuum expectation value and the second one $`(\varphi _2)`$ remains unbroken, i.e.
$$<\varphi _1>=\left(\begin{array}{c}0\hfill \\ v/\sqrt{2}\hfill \end{array}\right),<\varphi _2>=0$$
$`\left(1\right)`$
where $`v246GeV`$. The physical spectrum of Higgs bosons consists of two scalar neutral bosons $`h^0`$ and $`H^0`$, one pseudoscalar neutral boson $`A^0`$ and two charged Higgs $`H^\pm `$,
$$\begin{array}{c}H^0=\sqrt{2}\left[\left(Re\varphi _1^0v\right)\mathrm{cos}\alpha +Re\varphi _2^0\mathrm{sin}\alpha \right],\hfill \\ h^0=\sqrt{2}\left[\left(Re\varphi _1^0v\right)\mathrm{sin}\alpha +Re\varphi _2^0\mathrm{cos}\alpha \right],\hfill \\ A^0=\sqrt{2}\left(Im\varphi _2^0\right).\hfill \end{array}$$
$`\left(2\right)`$
The masses of the five neutral and charged Higgs bosons and the mixing angle $`\alpha `$ are free parameters of the model. The Yukawa couplings to quarks are,
$$_Y^Q=\lambda _{ij}^U\overline{Q_i}\stackrel{~}{\varphi _1}U_j+\lambda _{ij}^D\overline{Q_i}\varphi _1D_j+\xi _{ij}^U\overline{Q_i}\stackrel{~}{\varphi _2}U_j+\xi _{ij}^D\overline{Q_i}\varphi _2D_j,$$
$`\left(3\right)`$
where the first two terms give masses of the quark mass eigenstates, and $`\xi _{ij}^U`$ and $`\xi _{ij}^D`$ are the $`3\times 3`$ matrices which give the strength of the flavor changing neutral scalar vertices. The $`\xi `$s are all free parameters and can be constrained by the experimental data. If we neglect CP violation, the $`\xi `$s are all real. We will use the Cheng-Sher Ansatz(CSA) and let
$$\xi _{ij}\frac{\sqrt{m_im_j}}{v}.$$
And we can parametrize the Yukawa couplings as
$$\xi _{ij}=g\frac{\sqrt{m_im_j}}{m_W}\lambda .$$
$`\left(4\right)`$
Comparing it with the usual gauge couplings of $`SU(2)\times U(1)`$, one has $`\lambda =\frac{1}{\sqrt{2}}`$. In our calculation we use $`\lambda =\frac{1}{\sqrt{2}}`$ and note that there is no theoretical bound on the coupling factor $`\lambda `$.
The subprocess producing $`t\overline{c}(\overline{t}c)`$ via gluon-gluon collisions,
$$ggt\overline{c}\left(\overline{t}c\right)$$
can be induced through one-loop diagrams at the lowest order, and the Feynman diagrams are drawn in figure 1(a) and figure 1(b), where the contributions from the one-loop diagrams involving neutral and charged Higgs bosons are presented, respectively. The contributions from the figures of Fig.1(b) with the charged Higgs boson in loops being replaced by W-boson, is much smaller due to the GIM suppression and Yukawa coupling. We can omitted this part in our calculation. The diagrams exchanging the two external gluon-gluon lines are not shown, but are numbered in Fig.1(a) and Fig.1(b). Fig.1(a)(1 $``$ 12) and Fig. 1(b)(1 $``$ 6) are the self-energy diagrams, Fig. 1(a)(13 $``$ 20) and Fig. 1(b)(7 $``$ 10) are the vertex correction diagrams, Fig. 1(a)(25 $``$ 34) and Fig. 1(b) (13 $``$ 15) are the s-channel diagrams, Fig. 1(a)(21 $``$ 24) and Fig. 1(b)(11 $``$ 12) are the box diagrams. Note that in the present case at one-loop level the ultraviolet divergence would be canceled automatically, if all the one-loop diagrams at the $`O(m_tm_c/m_W^2)`$ order in the THDM III are included. In this work, we perform the calculation in the t’Hooft-Feynman gauge.
To simplify the calculation we set $`\alpha =0`$ and let all scalar bosons be degenerate, i.e., $`m_{h^0}=m_{A^0}=m_{H^\pm }=M_s`$ where $`M_s`$ is the common scalar mass. The contribution from the coupling involving $`H^0`$ is suppressed due to $`\alpha =0`$.
In the calculation for the s-channel diagrams(Fig.1.(a)(25 $``$ 28)), we take into account the width effects of the $`h^0`$ and $`A^0`$ propagators. The decays of $`h^0`$ to WW and ZZ are suppressed, because of the factor $`\mathrm{sin}\alpha `$ in the $`h^0WW`$ and $`h^0ZZ`$ couplings, and $`h^0`$ decay to $`A^0A^0`$ is also forbidden due to the case of the degenerate masses of $`h^0`$ and $`A^0`$. Note that the pseudoscalar $`A^0`$ does not couple with gauge boson pair. Therefore only the decays of $`h^0`$ and $`A^0`$ to final states $`q_i\overline{q_j}`$ need to be considered, where $`q_i`$ and $`q_j`$ represent quarks of flavor i and j respectively. The decay width for the scalar $`h^0`$ can be written as
$$\mathrm{\Gamma }\left(h^0q\overline{q}\right)=\frac{3g^2m_{h^0}}{32\pi M_W^2}m_q\left(14m_q^2/m_{h^0}^2\right)^{3/2}$$
and
$$\mathrm{\Gamma }\left(h^0t\overline{c}+\overline{t}c\right)=\frac{3g^2m_{h^0}}{32\pi M_W^2}2m_tm_c\left[1\left(m_t+m_c\right)^2/m_{h^0}^2\right]^{3/2}\times \left[1\left(m_tm_c\right)^2/m_{h^0}^2\right]^{1/2}.$$
$`\left(5\right)`$
The decay width for the pseudoscalar $`A^0`$ boson can be represented by exchanging exponents $`3/21/2`$ and $`m_{h^0}m_{A^0}`$ in Eq.(5). When $`m_t+m_c<M_s<2m_t`$, the dominant decay modes of $`h^0`$ and $`A^0`$ are $`h^0,A^0c\overline{c},b\overline{b},t\overline{c}+\overline{t}c`$, whereas when $`M_s>2m_t`$, the final state $`t\overline{t}`$ decay channel is open, and their decay widths are rather large due to the large masses of $`M_s`$ and $`m_t`$.
We denote $`\theta `$ as the scattering angle between one of the gluons and the final top quark. Then we express all the four-momenta of the initial and final particles in the center-of-mass(CMS) by means of the total energy $`\sqrt{\widehat{s}}`$ and the scattering angle $`\theta `$. The four-momenta of top quark and charm quark are $`p_1`$ and $`p_2`$ respectively and are read
$$\begin{array}{c}p_1=(E_t,\sqrt{E_t^2m_t^2}sin\theta ,0,\sqrt{E_t^2m_t^2}cos\theta ),\hfill \\ p_2=(E_c,\sqrt{E_c^2m_c^2}sin\theta ,0,\sqrt{E_c^2m_c^2}cos\theta ),\hfill \end{array}$$
$`\left(6\right)`$
where
$$E_t=\frac{1}{2}\left(\sqrt{\widehat{s}}+\left(m_t^2m_c^2\right)/\sqrt{\widehat{s}}\right),E_c=\frac{1}{2}\left(\sqrt{\widehat{s}}\left(m_t^2m_c^2\right)/\sqrt{\widehat{s}}\right).$$
$`\left(7\right)`$
$`p_3`$ and $`p_4`$ are the four-momenta of the initial gluons and are expressed as
$$p_3=(\frac{1}{2}\sqrt{\widehat{s}},0,0,\frac{1}{2}\sqrt{\widehat{s}}),p_4=(\frac{1}{2}\sqrt{\widehat{s}},0,0,\frac{1}{2}\sqrt{\widehat{s}}).$$
$`\left(8\right)`$
The corresponding matrix element for all the diagrams in figure 1(a) and figure 1(b) is written as
$$M=Tr\left(T^aT^b\right)\delta _{jl}M^{\widehat{s}_1}+\left(f_{abc}T_{jl}^c\right)M^{\widehat{s}_2}+\left(T_{jm}^aT_{ml}^b\right)M^{\widehat{t}}+\left(T_{jm}^bT_{ml}^a\right)M^{\widehat{u}}$$
$`\left(9\right)`$
The upper indexes $`\widehat{s}_1`$, $`\widehat{s}_2`$, $`\widehat{t}`$ and $`\widehat{u}`$ represent the amplitudes corresponding to the s-channel diagrams Fig. 1(a) (25$``$ 28), s-channel diagrams Fig. 1(a) (29 $``$ 34) and Fig. 1(b) (13 $``$ 15), t-channel and u-channel diagrams in figure 1(a) and figure 1(b) respectively. The $`T^a(a=18)`$ are the $`SU(3)_c`$ generators introduced by Gell-Mann and $`f_{abc}`$ are the antisymmetric $`SU(3)_c`$ structure constants. The subscripts $`j,l(j,l=1,2,3)`$ of $`T^a`$ represent the color of final state top quark and charm quark respectively. The variables $`\widehat{s}`$, $`\widehat{t}`$ and $`\widehat{u}`$ are usual Mandelstam variables in the center of mass system of gluon-gluon. Their definitions are:
$$\begin{array}{c}\widehat{s}=\left(p_1+p_2\right)^2=\left(p_3+p_4\right)^2,\widehat{t}=\left(p_1p_3\right)^2=\left(p_2p_4\right)^2,\hfill \\ \widehat{u}=\left(p_1p_4\right)^2=\left(p_2p_3\right)^2.\hfill \end{array}$$
$`\left(10\right)`$
We collect all the explicit expressions of the amplitudes appearing in equation (9) in the appendix. The total cross section for $`ggt\overline{c}+\overline{t}c`$ can be written in the form
$$\widehat{\sigma }\left(\widehat{s}\right)=\frac{2}{16\pi \widehat{s}^2}_{\widehat{t^{}}}^{\widehat{t^+}}𝑑\widehat{t}\left|\overline{M}\right|^2$$
$`\left(11\right)`$
where $`|\overline{M}|^2`$ is the initial spin-averaged matrix element squared and $`\widehat{t^\pm }=1/2(m_t^2+m_c^2\widehat{s})\pm \sqrt{E_t^2m_t^2}\sqrt{\widehat{s}}`$. The cross section for $`ppggt\overline{c}+\overline{t}c+X`$ is conveniently written in terms of the rapidities $`y_1`$ and $`y_2`$ of the two jets (finial states) and their common transverse momentum $`p_T`$. Here we neglect the intrinsic transverse momentum carried by partons. It is
$$\frac{d\sigma }{dy_1dy_2dp_T}=\frac{\pi \tau p_T}{\widehat{s}}f_g(x_1,Q^2)f_g(x_2,Q^2)\widehat{\sigma }\left(ggt\overline{c}+\overline{t}c\mathrm{at}\widehat{s}=\tau s\right),$$
$`\left(12\right)`$
where $`\sqrt{s}`$ and $`\sqrt{\widehat{s}}`$ denote the proton-proton and gluon-gluon c.m. energies respectively and $`\widehat{s}=s\tau `$. $`f_g(x_i,Q^2)`$ is the distribution function of gluon in proton.
Defining
$$y^{}=\frac{1}{2}\left(y_1y_2\right)$$
$`\left(13\right)`$
and
$$y_{boost}=\frac{1}{2}\left(y_1+y_2\right).$$
$`\left(14\right)`$
We may write
$$\tau =\frac{4p_T^2}{s}\mathrm{cosh}^2y^{}$$
$`\left(15\right)`$
and
$$x_1=\sqrt{\tau }e^{y_{boost}},x_2=\sqrt{\tau }e^{y_{boost}}.$$
$`\left(16\right)`$
In our numerical calculation we adopt the MRS set G parton distribution function $`f_g(x_i,Q^2)`$ and let the factorization scale $`Q^2=\widehat{s}`$. The numerical calculation is carried out around the Tevatron and LHC energy ranges.
III. Numerical Results and Discussions
In the numerical calculation we take the input parameters as $`m_b=4.5GeV`$, $`m_c=1.35GeV`$, $`m_t=175GeV`$, $`M_W=80.2226GeV`$, $`G_F=1.166392\times 10^5(GeV)^2`$ and $`\alpha =1/128`$. We adopt a simple one-loop formula for the running strong coupling constant $`\alpha _s`$ as
$$\alpha _s\left(\mu \right)=\frac{\alpha _s\left(m_Z\right)}{1+\frac{332n_f}{6\pi }\alpha _s\left(m_Z\right)\mathrm{ln}\left(\frac{\mu }{m_Z}\right)}.$$
$`\left(14\right)`$
where $`\alpha _s(m_Z)=0.117`$ and $`n_f`$ is the number of active flavors at energy scale $`\mu `$.
Figure 2 shows the cross sections for $`ggt\overline{c}+\overline{t}c`$ as a function of the masses of the Higgs bosons $`M_s`$. The cross sections are displayed with the three values of the gluon-gluon CMS energy 200 GeV, 400 GeV and 500 GeV respectively. Because there is no stringent bound on the Higgs bosons masses, we choose $`M_s`$ in the range from 50 GeV to 800 GeV. The peak of each curve comes from s-channel resonance effects, where $`M_s=m_{h^0}=m_{A^0}\sqrt{\widehat{s}}`$. From these curves we find that the cross section can be obviously enhanced when $`M_s`$ gets close to $`\sqrt{\widehat{s}}`$.
Figure 3 shows the cross sections of $`ggt\overline{c}+\overline{t}c`$ as a function of $`\sqrt{\widehat{s}}`$, and the three curves correspond to the $`M_s`$ values 100 GeV, 250 GeV and 500 GeV, respectively. For $`M_s=100GeV`$, the effects of the widths of the Higgs bosons are not obvious, and s-channel resonance effects are suppressed, since $`\sqrt{\widehat{s}}`$ is far beyond the Higgs boson masses $`M_s`$. Therefore the curve of its cross section is relative flat with the increasing of $`\sqrt{\widehat{s}}`$. When $`\sqrt{\widehat{s}}`$ approaches the value of $`M_s`$, such as $`M_s=250GeV`$ and $`M_s=500GeV`$, the cross sections will be enhanced by the s-channel resonance effects, and the width effects become stronger, since the $`h^0,A^0t\overline{c}+\overline{t}c`$ channels are opened. In Fig.3, we can see that the curve for $`M_s=500GeV`$ shows a sharp peak around the position at $`\sqrt{\widehat{s}}500GeV`$ due to the s-channel resonance effects and large width effects of $`h^0`$ and $`A^0`$.
In figure 4 and figure 5 we show the transverse momentum spectrum of the top quark at the Tevatron and LHC energies respectively, where we assume the off-line analysis will require at least the event selection criterion of involving one isolated high $`p_T`$ track with the cut of pseudorapidity $`|\eta |<2`$. Again, the peaks on the curves of $`M_s=250`$ and $`M_s=500`$ show the s-channel resonance effects, where $`M_s=m_{h^0}=m_{A^0}\sqrt{\widehat{s}}`$.
In figure 6 and 7 we show the cross section of $`ppggt\overline{c}+\overline{t}c+X`$ as a function of center-of-mass energy of electron-positron system $`\sqrt{s}`$. The cross section may reach 0.83 femtobarn when $`M_s=100GeV`$ and $`\sqrt{s}=2TeV`$ at the Tevatron and 131 femtobarn when $`M_s=100GeV`$ and $`\sqrt{s}=14TeV`$ at the LHC. For the Tevatron at 2 TeV we can expect about 4 raw events when $`M_s=100GeV`$ if we assume $`5fb^1`$ integrated luminosity, and for the LHC at 14 TeV we can expect about $`1.3\times 10^4`$ raw events if we assume $`100fb^1`$ integrated luminosity. Since the cross section of this process is roughly scaled by $`\lambda ^4`$, if we let $`\lambda 1`$, the cross section will be 4 times larger. There are several potentially severe backgrounds from the SM to the signal. A top quark with a mass about $`174GeV`$, decays dominantly to $`tW^+b`$. In the $`t\overline{c}W^+b\overline{c}l^+\nu b\overline{c}`$ detection mode, the backgrounds are mainly from $`t\overline{t}W+jets`$ and $`t\overline{t}WWb\overline{b}l^+\nu q\overline{q}^{^{}}b\overline{b}`$ processes. The calculation shows that the cross section of top pair production will reach about $`5pb`$ at the Tevatron for $`\sqrt{s}=1.8TeV`$ and about $`10^2pb`$ at the LHC for $`\sqrt{s}=14TeV`$. If we use the $`t\overline{c}+\overline{t}clepton+jets`$ detection mode, the production rate of $`t\overline{c}W^+b\overline{c}l^+\nu b\overline{c}`$ can reach about $`0.1fb`$ at the Tevatron and $`14fb`$ at the LHC, while the possible background from $`t\overline{t}WWb\overline{b}l^+\nu q\overline{q}^{^{}}b\overline{b}`$ would be about $`0.8pb`$ at the Tevatron and some dozens picobarn at the LHC. The reduction of the these backgrounds is possible through various kinematics cuts on the transverse energy, on the rapidity of jets and leptons, or involving b-tagging. Due to the very small production rate for the signal of $`ppt\overline{c}+\overline{t}c`$, it is not so easy to suppress its backgrounds. Therefore the further precise analyses are necessary to exclude these backgrounds.
In summary, from our calculation, we can conclude that if we could efficiently exclude the severe backgrounds of the $`t\overline{c}(\overline{t}c)`$ production signal, it would be possible at the Tevatron and the LHC that the process $`ppggt\overline{c}+\overline{t}c`$ could be used to probe the flavor changing interactions in the context of the THDM III.
Acknowledgement:
This work was supported in part by the National Natural Science Foundation of China(project numbers: 19675033, 19875049), the Youth Science Foundation of the University of Science and Technology of China(USTC) and a grant from the Research Fund for the Doctoral Program of Higher Education(RFDP) of China.
Appendix
We adopt the same definitions of one-loop A, B, C and D integral functions as in Ref. and the references therein. The dimension $`D=4ϵ`$. The integral functions are defined as
$$A_0\left(m\right)=\frac{\left(2\pi \mu \right)^{4D}}{i\pi ^2}d^Dq\frac{1}{\left[q^2m^2\right]},$$
$$\{B_1;B_\mu ;B_{\mu \nu }\}(p,m_1,m_2)=\frac{\left(2\pi \mu \right)^{4D}}{i\pi ^2}d^Dq\frac{\{1;q_\mu ;q_{\mu \nu }\}}{\left[q^2m_1^2\right]\left[\left(q+p\right)^2m_2^2\right]},$$
$$\{C_0;C_\mu ;C_{\mu \nu };C_{\mu \nu \rho }\}(p_1,p_2,m_1,m_2,m_3)=\frac{\left(2\pi \mu \right)^{4D}}{i\pi ^2}$$
$$\times d^Dq\frac{\{1;q_\mu ;q_{\mu \nu };q_{\mu \nu \rho }\}}{\left[q^2m_1^2\right]\left[\left(q+p_1\right)^2m_2^2\right]\left[\left(q+p_1+p_2\right)^2m_3^2\right]},$$
$$\{D_0;D_\mu ;D_{\mu \nu };D_{\mu \nu \rho };D_{\mu \nu \rho \alpha }\}(p_1,p_2,p_3,m_1,m_2,m_3,m_4)=\frac{\left(2\pi \mu \right)^{4D}}{i\pi ^2}$$
$$\times d^Dq\{1;q_\mu ;q_{\mu \nu };q_{\mu \nu \rho };q_{\mu \nu \rho \alpha }\}$$
$$\times \left\{\left[q^2m_1^2\right]\left[\left(q+p_1\right)^2m_2^2\right]\left[\left(q+p_1+p_2\right)^2m_3^2\right]\left[\left(q+p_1+p_2+p_3\right)^2m_4^2\right]\right\}^1.$$
In our calculation we take the strange quark mass $`m_s=0`$. The $`M^{\widehat{s}_1}`$ and $`M^{\widehat{s}_2}`$ in equation (9) can be written as
$$\begin{array}{ccc}M^{\widehat{s}_1}\hfill & =\hfill & \frac{i\alpha _s^2g^2}{16\pi ^2M_W^2}m_t\sqrt{m_tm_c}ϵ_\mu \left(p_3\right)ϵ_\nu \left(p_4\right)\overline{u}\left(p_1\right)\hfill \\ & & \{2a_{h^0}m_t^2(C_0+4C_{22}4C_{23})[p_3,p_1p_2,m_t,m_t,m_t](p_1^\mu p_1^\nu +p_1^\mu p_2^\nu +p_1^\nu p_2^\mu +p_2^\mu p_2^\nu )\hfill \\ & & +2a_{h^0}m_t^2g^{\mu \nu }(B_0[p_1p_2,m_t,m_t]((p_1+p_2)p_3C_0+4C_{24})[p_3,p_1p_2,m_t,m_t,m_t]\hfill \\ & & +2ia_{A^0}m_t^2C_0[p_3,p_1p_2,m_t,m_t,m_t]ϵ^{\mu \nu \alpha \beta }\gamma _5(p_1^\alpha p_3^\beta +p_2^\alpha p_3^\beta )\left\}v\right(p_2),\hfill \end{array}$$
$`\left(A.1\right)`$
where
$$a_{h^0}=\frac{1}{\widehat{s}m_{h^0}^2+i\mathrm{\Gamma }_{h^0}m_{h^0}},a_{A^0}=\frac{1}{\widehat{s}m_{A^0}^2+i\mathrm{\Gamma }_{A^0}m_{A^0}}.$$
$$\begin{array}{ccc}M^{\widehat{s}_2}\hfill & =\hfill & \frac{i\alpha _s^2g^2}{128\pi ^2M_W^2\widehat{s}}m_t\sqrt{m_tm_c}ϵ_\mu \left(p_3\right)ϵ_\nu \left(p_4\right)\overline{u}\left(p_1\right)\hfill \\ & & \{16m_t(C_{11}C_{12}+C_{21}C_{23})[p_1,p_1+p_2,M_s,m_t,m_t](p_1^\mu p_2^\nu p_1^\nu p_2^\mu )\hfill \\ & & +(8(2C_{24}+m_t^2(C_0C_{11}+C_{12}C_{21}C_{22}+2C_{23})2(p_1p_2)(C_{22}+C_{23}))\hfill \\ & & [p_1,p_1+p_2,M_s,m_t,m_t]+4m_t^2B_1[p_1,m_t,M_s]\hfill \\ & & 4m_tm_cB_1[p_2,m_t,M_s]+4m_t^2B_1[p_1,m_b,M_s]+\hfill \\ & & 4m_tm_cB_1[p_2,m_b,M_s]+m_t^2B_1[p_1,m_t,M_s]\hfill \\ & & m_tm_cB_1[p_2,m_t,M_s]m_t^2B_1[p_1,m_b,M_s]\hfill \\ & & +m_cm_tB_1[p_2,m_b,M_s]\left)\right(\gamma ^\nu p_1^\mu \gamma ^\mu p_1^\nu +\gamma ^\nu p_2^\mu \gamma ^\mu p_2^\nu +g^{\mu \nu }\text{/}p_3)\hfill \\ & & +(4(2m_tC_{24}2m_t^3C_0+m_t(m_t^2+2(p_1p_2)4(p_1p_3))(C_{11}C_{12}+C_{21}+C_{22}\hfill \\ & & 2C_{23})+4(p_2p_3)(m_cC_{12}m_tC_{22}+m_tC_{23})\left)\right[p_1,p_1+p_2,M_s,m_t,m_t]\hfill \\ & & (m_cm_t^2C_0+2m_tC_{24}+m_t^3(C_{11}C_{12}+C_{21}+C_{22}2C_{23})+2m_t((p_1p_2)\hfill \\ & & 2(p_1p_2)\left)\right(C_{11}C_{12}+C_{21}+C_{22}2C_{23})+4(p_2p_3\left)\right(m_cC_{12}\hfill \\ & & +m_t(C_{23}C_{22})\left)\right[p_1,p_1+p_2,M_s,m_b,m_b]m_t^3B_1[p_1,m_t,M_s]+\hfill \\ & & m_t^2m_cB_1[p_2,m_t,M_s]m_t^3B_1[p_1,m_b,M_s]+(m_cm_t^2+\hfill \\ & & 2m_cp_1p_24m_cp_1p_3)B_1[p_2,m_b,M_s])g^{\mu \nu }4m_t(C_{11}\hfill \\ & & C_{12}+C_{21}C_{23}\left)\right[p_1,p_1+p_2,M_s,m_b,m_b\left]\right(p_1^\mu p_2^\nu p_1^\nu p_2^\mu +\gamma _5p_1^\nu p_2^\mu \gamma _5p_1^\mu p_2^\nu )\hfill \\ & & +(2(2C_{24}+m_cm_tC_0m_t^2(C_{11}C_{12}+C_{21}+C_{22}2C_{23})2(p_1p_2)(C_{22}C_{23}))\hfill \\ & & [p_1,p_1+p_2,M_s,m_b,m_b]+4m_t^2B_1[p_1,m_b,M_s]\hfill \\ & & 4m_cm_tB_1[p_2,m_b,M_s]\left)\right(\gamma ^\nu p_1^\mu \gamma ^\mu p_1^\nu +\gamma ^\nu p_2^\mu \gamma ^\mu p_2^\nu +g^{\mu \nu }\text{/}p_3)(m_cm_t^2C_02m_tC_{24}\hfill \\ & & m_t^3(C_{11}C_{12}+C_{21}+C_{22}2C_{23})2m_t((p_1p_2)2(p_1p_2))(C_{11}C_{12}+C_{21}\hfill \\ & & +C_{22}2C_{23})+4(p_2p_3\left)\right(m_cC_{12}m_t(C_{23}C_{22})\left)\right[p_1,p_1+p_2,M_s,m_b,m_b]\gamma _5g^{\mu \nu }\hfill \\ & & +((4C_{24}2m_cm_tC_02m_t^2(C_{11}C_{12}+C_{21}+C_{22}2C_{23})+4(p_1p_2)(C_{23}\hfill \\ & & C_{22}))[p_1,p_1+p_2,M_s,m_b,m_b])+2m_t^2B_1[p_1,m_b,M_s]\hfill \\ & & 2m_cm_tB_1[p_2,m_b,M_s])\gamma _5(\gamma ^\nu p_1^\mu \gamma ^\mu p_1^\nu +\gamma ^\nu p_2^\mu \gamma ^\mu p_2^\nu g^{\mu \nu }\text{/}p_3)\left\}v\right(p_2),\hfill \end{array}$$
$`\left(A.2\right)`$
The amplitude of $`M^{\widehat{t}}`$ can be written as
$$\begin{array}{ccc}M^{\widehat{t}}\hfill & =\hfill & \frac{i\alpha _s^2g^2}{64\pi ^2M_W^2}m_t\sqrt{m_tm_c}ϵ_\mu \left(p_3\right)ϵ_\nu \left(p_4\right)\overline{u}\left(p_1\right)(f_1p_1^\mu p_1^\nu +f_2p_1^\mu p_2^\nu +f_3p_1^\nu p_2^\mu +f_4p_2^\mu p_2^\nu \hfill \\ & & +f_5\gamma ^\nu p_1^\mu +f_6\gamma ^\mu p_1^\nu +f_7\gamma ^\nu p_2^\mu +f_8\gamma ^\mu p_2^\nu +f_9\gamma ^\mu \gamma ^\nu +f_{10}\gamma ^\nu \gamma ^\mu +f_{11}\text{/}p_3p_1^\mu p_1^\nu \hfill \\ & & +f_{12}\text{/}p_3p_1^\mu p_2^\nu +f_{13}\text{/}p_3p_1^\nu p_2^\mu +f_{14}\text{/}p_3p_2^\mu p_2^\nu +f_{15}\text{/}p_3\gamma ^\nu p_1^\mu +f_{16}\text{/}p_3\gamma ^\mu p_1^\nu +f_{17}\text{/}p_3\gamma ^\nu p_2^\mu \hfill \\ & & +f_{18}\text{/}p_3\gamma ^\mu p_2^\nu +f_{19}\text{/}p_3\gamma ^\mu \gamma ^\nu +f_{20}\text{/}p_3\gamma ^\nu \gamma ^\mu +f_1^{}\gamma _5p_1^\mu p_1^\nu +f_2^{}\gamma _5p_1^\mu p_2^\nu +f_3^{}\gamma _5p_1^\nu p_2^\mu \hfill \\ & & +f_4^{}\gamma _5p_2^\mu p_2^\nu +f_5^{}\gamma _5\gamma ^\nu p_1^\mu +f_6^{}\gamma _5\gamma ^\mu p_1^\nu +f_7^{}\gamma _5\gamma ^\nu p_2^\mu +f_8^{}\gamma _5\gamma ^\mu p_2^\nu +f_9^{}\gamma _5\gamma ^\mu \gamma ^\nu \hfill \\ & & +f_{10}^{}\gamma _5\gamma ^\nu \gamma ^\mu +f_{11}^{}\gamma _5\text{/}p_3p_1^\mu p_1^\nu +f_{12}^{}\gamma _5\text{/}p_3p_1^\mu p_2^\nu +f_{13}^{}\gamma _5\text{/}p_3p_1^\nu p_2^\mu +f_{14}^{}\gamma _5\text{/}p_3p_2^\mu p_2^\nu \hfill \\ & & +f_{15}^{}\gamma _5\text{/}p_3\gamma ^\nu p_1^\mu +f_{16}^{}\gamma _5\text{/}p_3\gamma ^\mu p_1^\nu +f_{17}^{}\gamma _5\text{/}p_3\gamma ^\nu p_2^\mu +f_{18}^{}\gamma _5\text{/}p_3\gamma ^\mu p_2^\nu \hfill \\ & & +f_{19}^{}\gamma _5\text{/}p_3\gamma ^\mu \gamma ^\nu +f_{20}^{}\gamma _5\text{/}p_3\gamma ^\nu \gamma ^\mu \left)v\right(p_2),\hfill \end{array}$$
$`\left(A.3\right)`$
where the $`f_is`$ and $`f_i^{}s`$ are expressed explicitly as,
$$\begin{array}{ccc}f_1^{}\hfill & =\hfill & 4\left(m_c\left(D_{38}D_{310}\right)+m_t\left(D_{32}D_{36}\right)\right)[p_3,p_1,p_2,m_b,m_b,M_s,m_b]\hfill \end{array}$$
$`\left(A.4\right)`$
$$\begin{array}{ccc}f_1\hfill & =\hfill & f_1^{}\left(m_tm_t\right)+8m_t\left(D_{11}D_{12}+D_{21}D_{24}D_{25}+D_{26}\right)[p_1,p_3,p_4,M_s,m_t,m_t,m_t]\hfill \end{array}$$
$`\left(A.5\right)`$
$$\begin{array}{ccc}f_2^{}\hfill & =\hfill & 4a_1\left(m_cC_{22}m_tC_{12}m_tC_{23}\right)[p_4,p_2,m_b,m_b,M_s]\hfill \\ & & +4\left(m_c\left(D_{39}D_{310}\right)m_t\left(D_{22}D_{26}+D_{36}D_{38}\right)\right)[p_3,p_1,p_2,m_b,m_b,M_s,m_b]\hfill \end{array}$$
$`\left(A.6\right)`$
$$\begin{array}{ccc}f_2\hfill & =\hfill & f_2^{}\left(m_tm_t\right)+8a_1m_tC_{11}[p_2,p_4,M_s,m_t,m_t]\hfill \\ & & 8m_t\left(D_{12}+D_{24}D_{26}\right)[p_1,p_3,p_4,M_s,m_t,m_t,m_t]\hfill \end{array}$$
$`\left(A.7\right)`$
$$\begin{array}{ccc}f_3^{}\hfill & =\hfill & 4\left(m_c\left(D_{39}D_{37}\right)m_t\left(D_{310}D_{38}\right)\right)[p_3,p_1,p_2,m_b,m_b,M_s,m_b]\hfill \end{array}$$
$`\left(A.8\right)`$
$$f_3=f_3^{}\left(m_tm_t\right)8m_t\left(D_{25}D_{26}\right)[p_1,p_3,p_4,M_s,m_t,m_t,m_t]$$
$`\left(A.9\right)`$
$$f_4^{}=4\left(m_c\left(D_{33}D_{37}\right)+m_t\left(D_{23}D_{26}D_{310}+D_{39}\right)\right)[p_3,p_1,p_2,m_b,m_b,M_s,m_b]$$
$`\left(A.10\right)`$
$$f_4=f_4^{}\left(m_tm_t\right)+8m_tD_{26}[p_1,p_3,p_4,M_s,m_t,m_t,m_t]$$
$`\left(A.11\right)`$
$$\begin{array}{ccc}f_5^{}\hfill & =\hfill & 2a_2a_3m_t^2B_1[p_1,m_b,M_s]2a_1a_2\left(m_t^22\left(p_1p_3\right)\right)B_1[p_1+p_3,m_b,M_s]\hfill \\ & & 2a_1a_3m_cm_tB_1[p_2,m_b,M_s]+2a_2\left(m_cm_t\left(C_0+2C_{11}+C_{21}\right)2C_{24}\right)[p_1,p_3,M_s,m_t,m_t]\hfill \\ & & +2a_1(2C_{24}m_b^2C_0+m_c^2(C_{22}2C_{23})m_cm_tC_{12}+m_t^2(C_{11}+C_{21})\hfill \\ & & +2\left(p_1p_2p_2p_3\right)\left(C_{11}C_{12}+C_{21}C_{23}\right)\hfill \\ & & 2(p_1p_3)(C_{11}+C_{21})\left)\right[p_4,p_2,m_b,m_b,M_s]\hfill \\ & & +2(2D_{27}+2D_{312}+m_b^2D_0m_c^2D_{23}m_cm_t(D_{22}D_{26})\hfill \\ & & +2(p_2p_3)(D_{25}D_{26})\left)\right[p_3,p_1,p_2,m_b,m_b,M_s,m_b]\hfill \end{array}$$
$`\left(A.12\right)`$
$$\begin{array}{ccc}f_5\hfill & =\hfill & f_5^{}\left(m_tm_t\right)+4a_2a_3m_t^2B_0[p_1,m_t,M_s]\hfill \\ & & +4a_1a_2m_t^2B_0[p_1+p_3,m_t,M_s]4a_1a_3m_t^2B_0[p_2,m_t,M_s]\hfill \\ & & 4a_2m_t^2\left(C_0+C_{11}\right)[p_1,p_3,M_s,m_t,m_t]4a_1m_t^2[p_2,p_4,M_s,m_t,m_t]\hfill \\ & & +4m_t^2\left(D_0+D_{11}D_{13}\right)[p_1,p_3,p_4,M_s,m_t,m_t,m_t]\hfill \end{array}$$
$`\left(A.13\right)`$
$$\begin{array}{ccc}f_6^{}\hfill & =\hfill & 2(4(D_{311}D_{312})+m_b^2(D_{11}D_{12})m_c^2(D_{37}D_{39})\hfill \\ & & +m_t^2\left(D_{32}D_{36}\right)+2\left(p_1p_2\right)\left(D_{38}D_{310}\right)2\left(p_1p_3\right)\left(D_{22}D_{24}D_{34}+D_{36}\right)\hfill \\ & & +2(p_1p_2)(D_{35}D_{310})\left)\right[p_3,p_1,p_2,m_b,m_b,M_s,m_b]\hfill \end{array}$$
$`\left(A.14\right)`$
$$f_6=f_6^{}$$
$`\left(A.15\right)`$
$$\begin{array}{ccc}f_7^{}\hfill & =\hfill & 2(4D_{313}m_b^2D_{13}+m_c^2D_{33}+m_cm_tD_{23}+m_t^2(D_{26}+D_{38})+2(p_1p_2)(D_{23}+D_{39})\hfill \\ & & 2(p_1p_3)(D_{25}+D_{310}+D_{23}+D_{37})\left)\right[p_3,p_1,p_2,m_b,m_b,M_s,m_b]\hfill \end{array}$$
$`\left(A.16\right)`$
$$f_7=f_7^{}\left(m_tm_t\right)4m_t^2D_{13}[p_1,p_3,p_4,M_s,m_t,m_t,m_t]$$
$`\left(A.17\right)`$
$$\begin{array}{ccc}f_8^{}\hfill & =\hfill & 4a_1\left(p_1p_3\right)\left(C_{12}+C_{23}\right)[p_4,p_2,m_b,m_b,M_s]\hfill \\ & & +2(2(D_{27}+2D_{311}+D_{313})+m_b^2(D_0+D_{11})m_c^2(D_{23}+D_{37})\hfill \\ & & m_t^2\left(D_{22}+D_{36}\right)[p_3,p_1,p_2,m_b,m_b,M_s,m_b]\hfill \\ & & \left(p_1p_2\right)\left(D_{26}+D_{310}\right)+2\left(p_1p_3\right)\left(D_{12}D_{13}+4D_{24}D_{26}+D_{34}\right)\hfill \\ & & +2(p_2p_3)(D_{25}+D_{35})\left)\right[p_3,p_1,p_2,m_b,m_b,M_s,m_b]\hfill \end{array}$$
$`\left(A.18\right)`$
$$f_8=f_8^{}$$
$`\left(A.19\right)`$
$$\begin{array}{ccc}f_9^{}\hfill & =\hfill & a_2a_3m_t^3B_1[p_1,m_b,M_s]\hfill \\ & & +2a_1a_2m_c(p_1p_3)B_1[p_1+p_3,m_b,M_s]a_2m_t(2C_{24}\hfill \\ & & m_t^2(C_0+2C_{11}+C_{21})+2(p_1p_3)(C_0+C_{11}+C_{12}+C_{23})\left)\right[p_1,p_3,M_s,m_t,m_t]\hfill \\ & & 2a_1m_c(p_1p_3)C_{12}[p_4,p_2,m_b,m_b,M_s]+(4m_c(D_{313}\hfill \\ & & +m_b^2D_{13}m_c^2D_{33})+m_t(2D_{27}+4D_{312}+m_b^2(D_0+D_{12})m_c^2(D_{23}+D_{39})\hfill \\ & & m_t(m_cD_{38}m_tD_{22}m_tD_{32}))2(p_1p_2\left)\right(m_cD_{39}\hfill \\ & & +m_tD_{26}+m_tD_{38})+2(p_1p_3\left)\right(m_cD_{12}m_cD_{13}+m_cD_{310}+m_tD_{24}+m_tD_{36})\hfill \\ & & +2(p_2p_3)(m_cD_{37}+m_tD_{25}+m_tD_{310})\left)\right[p_3,p_1,p_2,m_b,m_b,M_s,m_b]\hfill \end{array}$$
$`\left(A.20\right)`$
$$\begin{array}{ccc}f_9\hfill & =\hfill & f_9^{}\left(m_tm_t\right)2a_2a_3m_t^3B_0[p_1,m_t,M_s]\hfill \\ & & 4a_1a_2m_t\left(p_1p_3\right)B_0[p_1+p_3,m_t,M_s]\hfill \\ & & +4a_2m_t\left(p_1p_3\right)C_0[p_1,p_3,M_s,m_t,m_t]\hfill \\ & & +4a_1m_t\left(p_1p_3\right)C_0[p_2,p_4,M_s,m_t,m_t]\hfill \\ & & +4m_t\left(D_{27}+D_{311}D_{313}\right)[p_1,p_3,p_4,M_s,m_t,m_t,m_t]\hfill \end{array}$$
$`\left(A.21\right)`$
$$f_{10}^{}=2\left(m_cD_{313}+m_tD_{312}\right)[p_3,p_1,p_2,m_b,m_b,M_s,m_b]$$
$`\left(A.22\right)`$
$$f_{10}=f_{10}^{}\left(m_tm_t\right)4m_tD_{27}[p_1,p_3,p_4,M_s,m_t,m_t,m_t]$$
$`\left(A.23\right)`$
$$f_{11}^{}=4\left(D_{22}D_{24}D_{34}+D_{36}\right)[p_3,p_1,p_2,m_b,m_b,M_s,m_b]$$
$`\left(A.24\right)`$
$$f_{11}=f_{11}^{}$$
$`\left(A.25\right)`$
$$f_{12}^{}=4\left(D_{13}D_{12}2D_{24}+2D_{26}+D_{310}D_{34}\right)[p_3,p_1,p_2,m_b,m_b,M_s,m_b]$$
$`\left(A.26\right)`$
$$f_{12}=f_{12}^{}$$
$`\left(A.27\right)`$
$$f_{13}^{}$$
$`\left(A.28\right)`$
$$f_{13}=f_{13}^{}$$
$`\left(A.29\right)`$
$$f_{14}^{}=4\left(D_{23}D_{25}D_{35}+D_{37}\right)[p_3,p_1,p_2,m_b,m_b,M_s,m_b]$$
$`\left(A.30\right)`$
$$f_{14}=f_{14}^{}$$
$`\left(A.31\right)`$
$$\begin{array}{ccc}f_{15}^{}\hfill & =\hfill & 2a_2m_t\left(C_{11}C_{12}+C_{21}C_{23}\right)[p_1,p_3,M_s,m_t,m_t]\hfill \\ & & +2\left(m_c\left(D_{13}D_{12}+D_{26}\right)+m_t\left(D_{22}D_{24}\right)\right)[p_3,p_1,p_2,m_b,m_b,M_s,m_b]\hfill \end{array}$$
$`\left(A.32\right)`$
$$\begin{array}{ccc}f_{15}\hfill & =\hfill & f_{15}^{}\left(m_tm_t\right)+4a_2m_tC_{11}[p_1,p_3,M_s,m_t,m_t]\hfill \\ & & 4m_t\left(D_{11}D_{13}\right)[p_1,p_3,p_4,M_s,m_t,m_t,m_t]\hfill \end{array}$$
$`\left(A.33\right)`$
$$f_{16}^{}=2\left(m_c\left(D_{25}D_{26}\right)m_t\left(D_{22}D_{24}\right)\right)[p_3,p_1,p_2,m_b,m_b,M_s,m_b]$$
$`\left(A.34\right)`$
$$f_{16}=f_{16}^{}\left(m_tm_t\right)+4m_t\left(D_{11}D_{12}\right)[p_1,p_3,p_4,M_s,m_t,m_t,m_t]$$
$`\left(A.35\right)`$
$$f_{17}^{}=2\left(m_cD_{23}m_t\left(D_{25}D_{26}\right)\right)[p_3,p_1,p_2,m_b,m_b,M_s,m_b]$$
$`\left(A.36\right)`$
$$f_{17}=f_{17}^{}\left(m_tm_t\right)+4m_tD_{13}[p_1,p_3,p_4,M_s,m_t,m_t,m_t]$$
$`\left(A.37\right)`$
$$\begin{array}{ccc}f_{18}^{}\hfill & =\hfill & 2a_1\left(m_cC_{22}m_tC_{12}m_tC_{23}\right)[p_4,p_2,m_b,m_b,M_s]\hfill \\ & & +2\left(m_c\left(D_{25}D_{23}\right)+m_t\left(D_{12}D_{13}+D_{24}D_{26}\right)\right)[p_3,p_1,p_2,m_b,m_b,M_s,m_b]\hfill \end{array}$$
$`\left(A.38\right)`$
$$\begin{array}{ccc}f_{18}\hfill & =\hfill & f_{18}^{}\left(m_tm_t\right)+4a_1m_tC_{11}[p_2,p_4,M_s,m_t,m_t]\hfill \\ & & 4m_tD_{12}[p_1,p_3,p_4,M_s,m_t,m_t,m_t]\hfill \end{array}$$
$`\left(A.39\right)`$
$$\begin{array}{ccc}f_{19}^{}\hfill & =\hfill & a_2a_3m_t^2B_1[p_1,m_b,M_s]\hfill \\ & & a_1a_2\left(m_t^22\left(p_1p_3\right)\right)B_1[p_1+p_3,m_b,M_s]\hfill \\ & & a_1a_3m_cm_tB_1[p_2,m_b,M_s]a_2(2C_{24}m_t(m_cC_0+m_tC_{11}\hfill \\ & & +m_tC_{21})+2(C_{12}+C_{23})(p_1p_3)\left)\right[p_1,p_3,M_s,m_t,m_t]\hfill \\ & & +a_1(2C_{24}m_b^2C_0+m_c^2(C_{22}2C_{23})m_cm_tC_{12}\hfill \\ & & +m_t^2\left(C_{11}+C_{21}\right)+2\left(\left(p_1p_2\right)\left(p_2p_3\right)\right)\left(C_{11}C_{12}+C_{21}C_{23}\right)\hfill \\ & & 2(p_1p_3)(C_{11}+C_{21})\left)\right[p_4,p_2,m_b,m_b,M_s]\hfill \\ & & +(4(D_{27}+D_{311})+2m_b^2(D_0+D_{11})2m_c^2(D_{23}+D_{37})D_{13}m_cm_t\hfill \\ & & m_t^2(D_{12}+D_{22}+D_{36})2(p_1p_2)(D_{12}+2D_{26}+D_{310})+2(p_1p_3)(D_{12}\hfill \\ & & +2D_{24}+D_{34})+2(p_2p_3)(D_{13}+2D_{25}+D_{35})\left)\right[p_3,p_1,p_2,m_b,m_b,M_s,m_b]\hfill \end{array}$$
$`\left(A.40\right)`$
$$\begin{array}{ccc}f_{19}\hfill & =\hfill & f_{19}^{}\left(m_tm_t\right)+2a_2a_3m_t^2B_0[p_1,m_t,M_s]\hfill \\ & & +2a_1a_2m_t^2B_0[p_1+p_3,m_t,M_s]2a_1a_2m_t^2B_0[p_2,m_t,M_s]\hfill \\ & & 2a_2m_t^2C_0[p_1,p_3,M_s,m_t,m_t]2a_1m_t^2C_0[p_2,p_4,M_s,m_t,m_t]\hfill \\ & & +2m_t^2D_0[p_1,p_3,p_4,M_s,m_t,m_t,m_t]\hfill \end{array}$$
$`\left(A.41\right)`$
$$f_{20}^{}=2D_{311}[p_3,p_1,p_2,m_b,m_b,M_s,m_b]$$
$`\left(A.42\right)`$
$$f_{20}=f_{20}^{}$$
$`\left(A.43\right)`$
where
$$a_1=\frac{1}{\widehat{t}m_t^2},a_2=\frac{1}{\widehat{t}m_c^2}anda_3=\frac{1}{m_t^2m_c^2}.$$
$$M^{\widehat{u}}=M^{\widehat{t}}(p_3p_4,\mu \nu ,\widehat{t}\widehat{u})\left(A.44\right)$$
Figure Captions
Fig.1(a)(b) The Feynman diagrams of the subprocess $`ggt\overline{c}`$.
Fig.2 Total cross sections of the subprocess $`ggt\overline{c}+\overline{t}c`$ as function of $`M_s`$. The solid curve is for $`\sqrt{\widehat{s}}=200GeV`$, the dashed curve is for $`\sqrt{\widehat{s}}=400GeV`$ and the dotted curve is for $`\sqrt{\widehat{s}}=500GeV`$.
Fig.3 Total cross sections of the subprocess $`ggt\overline{c}+\overline{t}c`$ as function of $`\sqrt{\widehat{s}}`$. The solid curve is for $`M_s=100GeV`$, the dashed curve is for $`M_s=250GeV`$ and the dotted curve is for $`M_s=500GeV`$.
Fig.4 The transverse momentum spectrum of top quark for the Tevatron at 2 TeV. The solid curve is for $`M_s=100GeV`$, the dashed curve is for $`M_s=250GeV`$ and the dotted curve is for $`M_s=500GeV`$.
Fig.5 The transverse momentum spectrum of top quark for the LHC at 14 TeV. The solid curve is for $`M_s=100GeV`$, the dashed curve is for $`M_s=250GeV`$ and the dotted curve is for $`M_s=500GeV`$.
Fig.6 Total cross sections of the process $`ppggt\overline{c}+\overline{t}c`$ as function of $`\sqrt{s}`$ at the Tevatron. The solid curve is for $`M_s=100GeV`$, the dashed curve is for $`M_s=250GeV`$ and the dotted curve is for $`M_s=500GeV`$.
Fig.7 Total cross sections of the process $`ppggt\overline{c}+\overline{t}c`$ as function of $`\sqrt{s}`$ at the LHC. The solid curve is for $`M_s=100GeV`$, the dashed curve is for $`M_s=250GeV`$ and the dotted curve is for $`M_s=500GeV`$.
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# Some remarks on the free fields realization of the bosonic string on 𝐴𝑑𝑆₃.
## 1 Introduction.
The last year has seen an increasing and renewed interest in string theory propagating on $`AdS`$ backgrounds. A particular attention has been dedicated to the $`AdS_3`$ background with NSNS flux, this because it is the only non trivial model which has been possible to treat exactly at the quantum level. Nevertheless and despite its apparent simplicity our understanding is far from being complete. In particular the question of which is its spectrum is not satisfactorily settled. It was already noted in the old days that with a naive quantization the Virasoro constraints were not sufficient to eliminate the ghost from the spectrum (). In order to remedy to this unfortunate circumstance two different proposals have been put forward<sup>1</sup><sup>1</sup>1 Some recent papers have done some steps forwards in increasing our understanding of the quantum theory (,,). Especially () has given a nice interpretation of the old results by () on the construction of a modular invariant partition function.
1. To truncate the spectrum to the unitary part: this was the approach first advocated by Petropolous () and independently by Mohammedi (), developed by Hwang () and collaborators. Last year Evans, Gaberdiel and Perry () showed that the free spectrum is actually ghost free;
2. To introduce some new hidden dof in form of zero momenta of the fields used to bosonize the KM currents. This way was first undertaken by Bars () and developed by Satoh () along slightly different lines.
We think that both proposals have some weak sides but the main criticism is due to their philosophical attitude to the problem (see also () for a recent review of the open problems): string theory is so consistent that there should not be necessary to introduce new elements in the game since the string itsself should give the answer.
Our approach is different from both the previous ones even if it is closer to the second one: we find the classical canonical transformation from the string fields to the Wakimoto fields and then we discuss how and in what extent the Wakimoto fields are free.
## 2 The Classical Bosonic String Theory on $`AdS_3`$.
We first exam the equations of motion and boundary conditions for a classical bosonic string with a Kalb-Ramond background thought of as a WZW theory plus Virasoro constraints not relying on the KM symmetry structure present in the theory and we derive the general solution. We set up the canonical formalism which we use to discuss the classical ’canonical’ transformation to the classical Wakimoto fields. finally we discuss the boundary conditions and their consequences. It turns out that the Wakimoto fields while satisfying free fields equations of motion are generically not free because of the boundary conditions; we discuss how to circumvent the problem.
### 2.1 The action and the constraints of the bosonic string on $`AdS_3`$.
The string action is given by<sup>2</sup><sup>2</sup>2 Spacetime coordinates: $`x^\pm =\frac{x^1\pm x^0}{\sqrt{2}}`$ WS metric $`\eta _{00}=\eta _{11}=1`$; $`\eta _+=\frac{1}{2}`$ $`\eta ^+=2`$ WS coordinates: $`\xi ^\pm =\tau \pm \sigma `$; $`\sigma [0,2\pi ]`$; $`d^2\xi =2d\tau d\sigma `$; $`_\pm =\frac{1}{2}\left(_\tau \pm _\sigma \right)`$ Sigma matrices: $`\sigma _3=\left(\begin{array}{cc}1& \\ & 1\end{array}\right)`$$`\sigma _+=\left(\begin{array}{cc}& 1\\ 0& \end{array}\right)`$$`\sigma _{}=\left(\begin{array}{cc}& 0\\ 1& \end{array}\right)`$ $`sl(2,R)`$ algebra: $`\sigma _a\sigma _b=ϵ_{abc}\sigma ^c+\eta _{ab}1`$ with $`ϵ_{+3}=\frac{1}{2}`$ and $`\eta _+=\frac{1}{2}`$ $`\eta _{33}=1`$
$`S_{wzw}`$ $`=`$ $`{\displaystyle \frac{|k|}{4\pi }}{\displaystyle d^2\xi tr(\omega _+\omega _{})}`$ (2.1)
$`+{\displaystyle \frac{k}{12\pi }}{\displaystyle tr(\omega ^3)}`$
where $`\omega =g^1dg=\omega _+d\xi ^++\omega _{}d\xi ^{}`$ is the pullback on the string worldsheet of the left invariant one form on the group $`SL(2,R)`$. When we use the explicit expression for $`g`$ in the Gauss form<sup>3</sup><sup>3</sup>3 This expression only covers half $`AdS_3`$ , even if we perform an analytic continuation letting $`e^\rho R`$ they do not cover the whole manifold since points like $`\left(\begin{array}{cc}\alpha & \beta \\ \frac{1}{\beta }& 0\end{array}\right)`$ are left out.
$$g=\left(\begin{array}{cc}1& \\ x^{}& 1\end{array}\right)\left(\begin{array}{cc}e^\rho & \\ & e^\rho \end{array}\right)\left(\begin{array}{cc}1& x^+\\ & 1\end{array}\right)=\left(\begin{array}{cc}e^\rho & e^\rho x^+\\ e^\rho x^{}& e^\rho +e^\rho x^+x^{}\end{array}\right)$$
(2.2)
and take $`k=|k|`$ the previous expression for the action becomes
$$S_{wzw}=\frac{|k|}{2\pi }d^2\xi _+\rho _{}\rho +e^{2\rho }_{}x^+_+x^{}$$
(2.3)
We can interpret this as a string action in the conformal gauge if we add the Virasoro constraints
$`T_{++}`$ $`=`$ $`|k|\left[\left(_+\rho \right)^2+e^{2\rho }_+x^+_+x^{}\right]=0`$
$`T_{}`$ $`=`$ $`|k|\left[\left(_{}\rho \right)^2+e^{2\rho }_{}x^+_{}x^{}\right]=0`$
We want now to derive the the equations of motion and the allowed boundary conditions associated with the action (2.3) in the same way of we proceed with the usual string action in Minkowski space. The equations of motion read
$`_{}\left(e^{2\rho }_+x^{}\right)=_+\left(e^{2\rho }_{}x^+\right)`$ $`=`$ $`0`$ (2.4)
$`_{}_+\rho +e^{2\rho }_{}x^+_+x^{}`$ $`=`$ $`0`$ (2.5)
while from the request of the cancellation of the surface terms obtained from the fields variation we get the boundary conditions
$`\delta \rho |_{\sigma =0}=\delta \rho |_{\sigma =2\pi }`$ $``$
$`\rho ^{}|_{\sigma =0}`$ $`=`$ $`\rho ^{}|_{\sigma =2\pi }`$
$`\delta x^{}|_{\sigma =0}=\delta x^{}|_{\sigma =2\pi }`$ $``$
$`e^{2\rho }_{}x^+|_{\sigma =0}`$ $`=`$ $`e^{2\rho }_{}x^+|_{\sigma =2\pi }`$ (2.6)
$`\delta x^+|_{\sigma =0}=\delta x^+|_{\sigma =2\pi }`$ $``$
$`e^{2\rho }_+x^{}|_{\sigma =0}`$ $`=`$ $`e^{2\rho }_+x^{}|_{\sigma =2\pi }`$ (2.7)
Anticipating the expressions for the KM currents (2.25,2.29) we can rewrite the last two conditions (2.6,2.7) as
$`J^{}|_{\sigma =0}=J^{}|_{\sigma =2\pi }`$ $`\overline{J}^+|_{\sigma =0}=\overline{J}^+|_{\sigma =2\pi }`$
### 2.2 The general solution of the equations of motions.
We are now ready to discuss the general solution of eq.s (2.4-2.5) . We can write the general solution as
$`x^+`$ $`=`$ $`a(\xi ^+)+e^{2c(\xi ^+)}{\displaystyle \frac{\overline{b}(\xi ^{})}{1+\overline{b}(\xi ^{})b(\xi ^+)}}`$ (2.8)
$`x^{}`$ $`=`$ $`\overline{a}(\xi ^{})+e^{2\overline{c}(\xi ^{})}{\displaystyle \frac{b(\xi ^+)}{1+\overline{b}(\xi ^{})b(\xi ^+)}}`$ (2.9)
$`\rho `$ $`=`$ $`\mathrm{lg}\left(1+\overline{b}(\xi ^{})b(\xi ^+)\right)`$ (2.10)
$`+c(\xi ^+)+\overline{c}(\xi ^{}).`$
from the knowledge of the general solution of the equations of motion associated with a WZW action, i.e<sup>4</sup><sup>4</sup>4 The reason why we choose such a parametrization is because we want the canonical momenta associated with $`x^\pm `$ be symmetric in the exchange of the barred and unbarred quantities. .
$`g(\xi ^+,\xi ^{})`$ $`=`$ $`g_R^T(\xi ^{})g_L(\xi ^+)`$ (2.11)
$`g_R^T(\xi ^{})`$ $`=`$ $`\left(\begin{array}{cc}e^{\overline{c}}& e^{\overline{c}}\overline{b}\\ e^{\overline{c}}\overline{a}& e^{\overline{c}}\overline{a}\overline{b}+e^{\overline{c}}\end{array}\right)`$ (2.14)
$`g_L(\xi ^+)`$ $`=`$ $`\left(\begin{array}{cc}e^c& e^ca\\ e^cb& e^cab+e^c\end{array}\right)`$ (2.17)
As it is well known the solution (2.11) does not fix completely $`g_R,g_L`$ (2.17) which are determined up to a redefinition
$$g_Lg_0g_Lg_Rg_0^Tg_R$$
(2.18)
we can (partially) fix this invariance by choosing a canonical form for the monodromies. This invariance is also connected to the possibility of using different charts: our parametrization is not global and therefore the group has to be covered with charts where one patch is parametrized as in (2.17) and the others can be chosen to be
$`g_{(1)L}(\xi ^+)`$ $`=`$ $`\left(\begin{array}{cc}e^{c_{(1)}}b_{(1)}& e^{c_{(1)}}a_{(1)}b_{(1)}+e^{c_{(1)}}\\ e^{c_{(1)}}& e^{c_{(1)}}a_{(1)}\end{array}\right)`$
$`g_{(3)L}(\xi ^+)`$ $`=`$ $`\left(\begin{array}{cc}e^{c_{(3)}}b_{(3)}& e^{c_{(3)}}a_{(3)}b_{(3)}+e^{c_{(3)}}\\ e^{c_{(3)}}& e^{c_{(3)}}a_{(3)}\end{array}\right)`$
$`g_{(2)L}(\xi ^+)`$ $`=`$ $`\left(\begin{array}{cc}e^{c_{(2)}}& e^{c_{(2)}}a_{(2)}\\ e^{c_{(2)}}b_{(2)}& e^{c_{(2)}}a_{(2)}b_{(2)}+e^{c_{(2)}}\end{array}\right)`$
with transition function given by $`\mathrm{\Omega }=\left(\begin{array}{cc}& 1\\ 1& \end{array}\right)`$ (i.e. $`g_{(i+1)L}=\mathrm{\Omega }g_{(i)L}`$ with $`imod\mathrm{\hspace{0.25em}4}`$) for $`g_L`$ and similarly for $`g_R`$. If we do not want to use charts we have to use singular functions as it happens with the Dirac monopole.
### 2.3 Canonical formalism.
Since we want to to discuss canonical transformations from interacting fields to Wakimoto ones we need to set up the canonical formalism. This is easily done and we find the momenta
$`𝒫`$ $`=`$ $`{\displaystyle \frac{|k|}{2\pi }}\dot{\rho }𝒫_+={\displaystyle \frac{|k|}{2\pi }}e^{2\rho }_+x^{}𝒫_{}={\displaystyle \frac{|k|}{2\pi }}e^{2\rho }_{}x^+`$
along with the classical hamiltonian
$$=\frac{\pi }{|k|}𝒫^2+\frac{|k|}{4\pi }\rho ^2+\frac{4\pi }{|k|}e^{2\rho }𝒫_+𝒫_{}𝒫_{}x^{}+𝒫_+x^+$$
Moreover we can write the non vanishing Poisson brackets as
$`\{x^+(\sigma ),𝒫_+(\sigma ^{})\}`$ $`=`$ $`\{x^{}(\sigma ),𝒫_{}(\sigma ^{})\}=\delta (\sigma \sigma ^{})`$
$`\{\rho (\sigma ),𝒫(\sigma ^{})\}`$ $`=`$ $`\delta (\sigma \sigma ^{})`$ (2.22)
Obviously this expressions are not very useful because we cannot use them to deduce the commutation relations between the “oscillators” $`a,b,c`$ and $`\overline{a},\overline{b},\overline{c}`$ due to the highly non linear way they enter the expressions for $`x^\pm ,`$ $`\rho `$, explicitly
$`𝒫_+`$ $`=`$ $`{\displaystyle \frac{|k|}{2\pi }}e^{2c}_+b𝒫_{}={\displaystyle \frac{|k|}{2\pi }}e^{2\overline{c}}_{}\overline{b}`$ (2.23)
### 2.4 The KM algebra and the energy-momentum tensor.
From the standard classical expression for the left/right KM currents $`J=|k|g^1_+g`$ ($`\overline{J}=|k|_{}gg^1`$) we can compute the classical KM currents which read
$`J^{}`$ $`=`$ $`|k|e^{2\rho }_+x^{}`$ (2.24)
$`J^3`$ $`=`$ $`|k|\left(_+\rho x^+e^{2\rho }_+x^{}\right)`$ (2.25)
$`J^+`$ $`=`$ $`|k|\left(_+x^++2x^+_+\rho x^{+2}e^{2\rho }_+x^{}\right)`$
and
$`\overline{J}^{}`$ $`=`$ $`|k|\left(_{}x^{}+2x^{}_{}\rho x^2e^{2\rho }_{}x^+\right)`$
$`\overline{J}^3`$ $`=`$ $`|k|\left(_{}\rho x^{}e^{2\rho }_{}x^+\right)`$ (2.28)
$`\overline{J}^+`$ $`=`$ $`|k|e^{2\rho }_{}x^+`$ (2.29)
The previous currents can be rewritten in the canonical formalism in the following way
$`J^{}`$ $`=`$ $`2\pi 𝒫_+`$
$`J^3`$ $`=`$ $`{\displaystyle \frac{|k|}{2}}\left({\displaystyle \frac{2\pi }{|k|}}𝒫+\rho ^{}\right)2\pi x^+𝒫_+`$ (2.30)
$`J^+`$ $`=`$ $`|k|x^++|k|x^+\left({\displaystyle \frac{2\pi }{|k|}}𝒫+\rho ^{}\right)`$
$`2\pi x^{+2}𝒫_++2\pi e^{2\rho }𝒫_{}`$
$`\overline{J}^{}`$ $`=`$ $`|k|x^{}+|k|x^{}\left({\displaystyle \frac{2\pi }{|k|}}𝒫\rho ^{}\right)`$ (2.32)
$`2\pi x^2𝒫_{}+2\pi e^{2\rho }𝒫_+`$
$`\overline{J}^3`$ $`=`$ $`{\displaystyle \frac{|k|}{2}}\left({\displaystyle \frac{2\pi }{|k|}}𝒫\rho ^{}\right)2\pi x^{}𝒫_{}`$
$`\overline{J}^+`$ $`=`$ $`2\pi 𝒫_{}`$ (2.33)
while the momentum-energy tensor reads
$`T_{++}`$ $`=`$ $`{\displaystyle \frac{|k|}{4}}\left({\displaystyle \frac{2\pi }{|k|}}𝒫+\rho ^{}\right)^2+{\displaystyle \frac{4\pi ^2}{|k|}}e^{2\rho }𝒫_+𝒫_{}`$
$`+2\pi 𝒫_+x^+`$
$`T_{}`$ $`=`$ $`{\displaystyle \frac{|k|}{4}}\left({\displaystyle \frac{2\pi }{|k|}}𝒫\rho ^{}\right)^2+{\displaystyle \frac{4\pi ^2}{|k|}}e^{2\rho }𝒫_+𝒫_{}`$
$`2\pi 𝒫_{}x^{}`$
It is then an easy matter to verify that they satisfy the following classical Virasoro (with vanishing central charge)
$`\{T_{++}(\sigma ),T_{++}(\sigma ^{})\}`$ $`=`$
$`2\pi \left[T_{++}(\sigma )+T_{++}(\sigma ^{})\right]_\sigma \delta (\sigma \sigma ^{})`$
$`2\pi {\displaystyle \frac{c}{12}}_\sigma ^3\delta (\sigma \sigma ^{})`$
$`\{T_{}(\sigma ),T_{}(\sigma ^{})\}`$ $`=`$
$`2\pi \left[T_{}(\sigma )+T_{}(\sigma ^{})\right]_\sigma \delta (\sigma \sigma ^{})`$
$`+2\pi {\displaystyle \frac{c}{12}}_\sigma ^3\delta (\sigma \sigma ^{})`$
with $`c=0`$ and Kac-Moody algebra (of level $`|k|`$ )
$`\{J^a(\sigma ),J^b(\sigma ^{})\}`$ $`=`$ $`2\pi ϵ_{..c}^{ab}J^c\delta (\sigma \sigma ^{})`$ (2.35)
$`+\pi |k|\eta ^{ab}\delta ^{}(\sigma \sigma ^{})`$
$`\{\overline{J}^a(\sigma ),\overline{J}^b(\sigma ^{})\}`$ $`=`$ $`2\pi ϵ_{..c}^{ab}\overline{J}^c\delta (\sigma \sigma ^{})`$ (2.36)
$`\pi |k|\eta ^{ab}\delta ^{}(\sigma \sigma ^{})`$
It is not difficult to check the classical Sugawara construction, i.e. $`T=\frac{1}{|k|}\eta _{ab}J^aJ^b`$.
### 2.5 The classical canonical transformation to the Wakimoto fields.
We are now ready to discuss the classical Wakimoto canonical fields. In order to get a hint on how they are related to our starting canonical variables we evaluate $`T`$ on the general solution of the equations of motion and we get
$$T_{++}=|k|(_+c)^2+|k|e^{2c}_+b_+a$$
Remembering the value of $`𝒫_+`$ given in eq. (2.23) the previous expression suggests that $`c`$ has something to do with a “free” field while $`a`$ could be proportional to the field canonically conjugate to $`𝒫_+`$. Starting from this observation it is not difficult to show that
$`F=\sqrt{2|k|}(c+\overline{c})`$ (2.37)
$`\beta =2\pi i𝒫_+,\gamma =a`$ (2.38)
$`\overline{\beta }=2\pi i𝒫_{},\overline{\gamma }=\overline{a}`$ (2.39)
which satisfy the following canonical Poisson brackets
$`\{F(\sigma ),\dot{F}(\sigma ^{})\}`$ $`=`$ $`4\pi \delta (\sigma \sigma ^{})`$ (2.40)
$`\{\beta (\sigma ),\gamma (\sigma ^{})\}`$ $`=`$ $`2\pi i\delta (\sigma \sigma ^{})`$ (2.41)
$`\{\overline{\beta }(\sigma ),\overline{\gamma }(\sigma ^{})\}`$ $`=`$ $`2\pi i\delta (\sigma \sigma ^{})`$ (2.42)
reproduce eq.s (2.3).
We can now use these new canonical variables to express the classical energy momentum tensor
$$T_{++}=\frac{1}{2}\left(_+F\right)^2i\beta _+\gamma $$
and the classical left $`sl(2,R)`$ KM generators
$`J^{}`$ $`=`$ $`i\beta J^3=\sqrt{{\displaystyle \frac{|k|}{2}}}_+F+i\beta \gamma `$ (2.43)
$`J^+`$ $`=`$ $`|k|_+\gamma +\sqrt{2|k|}\gamma _+F+i\beta \gamma ^2`$ (2.44)
and the classical right $`sl(2,R)`$ ones
$`\overline{J}^+`$ $`=`$ $`i\overline{\beta }\overline{J}^3=\sqrt{{\displaystyle \frac{|k|}{2}}}_{}F+i\overline{\beta }\overline{\gamma }`$ (2.45)
$`\overline{J}^{}`$ $`=`$ $`|k|_{}\overline{\gamma }\sqrt{2|k|}\overline{\gamma }_{}F+i\overline{\beta }\overline{\gamma }^2`$ (2.46)
## 3 Monodromy matrices and boundary conditions.
Which are the boundary conditions to be imposed, here we follow the approach of ()? Usually and naively we would take $`g_R,g_LSL(2,R)`$ but in this case proceeding in this way we would get some very ugly results or miss two different physical sectors as we are going to explain, we take therefore $`g_R,g_LSL(2,C)`$ with the further condition dictated by reality of $`g`$
$$g_L^{}=g_{CC}g_Lg_R^{}=g_Rg_{CC}^1g_{CC}SL(2,C).$$
(3.1)
We can now ask which is the most general boundary conditions we can impose on $`g_{L,R}`$ compatible with eq.s (2.1-2.7). The simplest answer is the periodicity in $`g`$ which can be achieved by imposing the following boundary conditions
$`g_L(\xi ^++2\pi )`$ $`=`$ $`g_Pg_L(\xi ^+)`$
$`g_R(\xi ^{}2\pi )`$ $`=`$ $`g_R(\xi ^{})g_P^1`$ (3.2)
A less obvious answer which is nevertheless allowed by the boundary conditions is
$`g_L(\xi ^++2\pi )`$ $`=`$ $`g_Pg_L(\xi ^+)g_{L0}`$
$`g_R(\xi ^{}2\pi )`$ $`=`$ $`g_{R0}g_R(\xi ^{})g_P^T`$ (3.3)
where $`g_{L0}=\left(\begin{array}{cc}1& 2\pi w\\ 0& 1\end{array}\right)`$ and $`g_{R0}=\left(\begin{array}{cc}1& 2\pi \overline{w}\\ 0& 1\end{array}\right)`$ as it can be checked from (2.1-2.7). From the explicit form of the Wakimoto fields it turns out that both $`w`$ and $`\overline{w}`$ have vanishing Poisson brackets with all the fields: this is not strange as it can appear because the same happens in the flat limit.
Obviously eq. (3.1) has to be compatible with eq. (3.3), i.e. $`g_{CC}`$ and $`g_P`$ have to satisfy
$$g_{CC}g_P=g_P^{}g_{CC}$$
(3.4)
As far as the periodicity is concerned there are three different equivalence classes: $`g_P`$ can be either hyperbolic, parabolic or elliptic. Such classes correspond to tachionic, massless and massive string excitations in the flat limit $`k\mathrm{}`$ .
#### 3.0.1 Hyperbolic sector.
Let us start with the hyperbolic class is the most natural with the coordinates associated with the Gauss decomposition (2.2). This case can be described, up to conjugation by a constant element, by taking
$$g_P=\left(\begin{array}{cc}e^{2\pi p}& \\ & e^{2\pi p}\end{array}\right)p>0$$
where the constraint $`p>0`$ is due to the symmetry $`g_L\mathrm{\Omega }g_L`$ with $`\mathrm{\Omega }=\left(\begin{array}{cc}& 1\\ 1& \end{array}\right)`$. This form of the periodicity matrix does not fix completely the invariance (2.18) which can, for example, be generically fixed by the further constraint $`\overline{b}(0)=b(0)`$; analogous considerations apply for the other sectors. This $`g_P`$ is compatible with $`g_{CC}=1`$.
We can write the explicit expansions of the functions entering the general solution with the given boundary conditions as
$`a=w\xi ^++a_{periodic},b=e^{2p\xi ^+}b_{periodic},c=p\xi ^++c_{periodic}`$
$`F=p\tau +F_{periodic},\beta =\beta _{periodic},\gamma =w\xi ^++\gamma _{periodic}`$
similarly for the barred quantities with $`\overline{p}=p`$ but with $`\overline{w}`$ independent of $`w`$. The presence of these two constants $`w`$ and $`\overline{w}`$ is allowed by the equality of the variations of $`x^\pm `$ at $`\sigma =0,2\pi `$ in particular setting $`w=\overline{w}=Rn`$ ($`nZ`$) is equivalent to compactify the $`x^1`$ with radius $`R`$, i.e. to choose an extremal BH as background (). It is important to notice that both $`w`$ and $`\overline{w}`$ have vanishing Poisson brackets with everything as it can be verified from the free field representation. Analogous considerations apply to the other sectors.
It is interesting to notice that this is the only sector considered in (,) as it can be seen from the $`F`$ expansion.
#### 3.0.2 Elliptic sector.
Let us now consider the elliptic case which describes massive excitations in the flat limit where
$$g_P=\left(\begin{array}{cc}\mathrm{cos}2\pi p& \mathrm{sin}2\pi p\\ \mathrm{sin}2\pi p& \mathrm{cos}2\pi p\end{array}\right)\left(\begin{array}{cc}e^{i2\pi p}& \\ & e^{i2\pi p}\end{array}\right)\mathrm{\hspace{0.25em}\hspace{0.25em}\hspace{0.25em}\hspace{0.25em}\hspace{0.25em}0}<p<\frac{1}{2}$$
where the first expression is the natural one when restricting the attention to real $`g_{L,R}`$ , i.e. when $`g_{CC}=1`$ while the second one is the most natural when considering complex $`g_{L,R}`$ with $`g_{CC}=\pm \left(\begin{array}{cc}& i\\ i& \end{array}\right)`$.
If we insist to use the real fields we get strange and ugly boundary conditions such as
$$c(\xi ^++2\pi )=c(\xi ^+)+\mathrm{log}\left(\mathrm{cos}2\pi p+b(\xi ^+)\mathrm{sin}2\pi p\right)$$
which propagate to strange, non free boundary conditions for the Wakimoto fields while using complex fields we get nice and free boundary conditions since the complex fields have analogous expansion as (3.0.1) with the substitution $`pi(p+k)`$ (where $`kZ`$) but it obliges us to impose the constraints
$$a^{}=a+\frac{1}{be^{2c}},b^{}=\frac{1}{b},c^{}=c+\mathrm{log}\left(\pm ib\right)$$
They can be imposed in a better way by requiring the reality of the KM currents.
Notice that if we restrict the momentum $`p`$ to the first Block wave the constraint $`0<p<\frac{1}{2}`$ implies $`0<j=p_F<\frac{k}{2}`$ which is equivalent to the unitary spectrum truncation at the quantum level.
#### 3.0.3 Parabolic sector.
Let us now consider the parabolic case where
$$g_P=\left(\begin{array}{cc}1& 0\\ 2\pi p& 1\end{array}\right)pR$$
and $`g_{CC}=1`$, then the fields in the 0th patch can be expanded as
$`a=w\xi ^++a_{periodic},b=p\xi ^++b_{periodic},`$
$`c=c_{periodic}F=F_{periodic},`$
$`\beta =\beta _{periodic},\gamma =\gamma _{periodic}`$
However life is not so easy since the fields in the 1st patch satisfy
$`a_{(1)}(\xi ^++2\pi )`$ $`=`$ $`a_{(1)}(\xi ^+){\displaystyle \frac{e^{2c_{(1)}}}{12\pi pb_{(1)}(\xi ^+)}}`$
$`{\displaystyle \frac{1}{b_{(1)}(\xi ^++2\pi )}}`$ $`=`$ $`{\displaystyle \frac{1}{b_{(1)}(\xi ^+)}}2\pi p`$
$`c_{(1)}(\xi ^++2\pi )`$ $`=`$ $`c_{(1)}(\xi ^+)+\mathrm{log}\left(12\pi pb_{(1)}(\xi ^+)\right)`$
There is apparently not an easy way to impose these constraints on the Wakimoto fields however it is necessary to use patches in this sector. A comment is now necessary in order to explain why it is necessary to work with two patches in order the canonical transformation work fine, exactly as it happens for the Liouville case. . If $`p_F0`$ we can express $`b`$ using the new canonical variables $`F`$ and $`\beta `$ but if we try to solve for $`b`$ when $`p_F=0`$ then we cannot recover $`b_0`$ (the constant mode of $`b_{periodic}`$ ) from the expression for $`\beta `$. This can clearly be avoided if we use two patches.
## 4 Conclusions.
We have shown that a free fields approach to string propagating on $`AdS_3`$ requires a lot of attention and that we must work with different charts, as done in () for the Liouville theory, if we want to treat the parabolic sector correctly. As byproducts of this analysis we have shown that the spectrum proposed by (,) is unnaturally truncated to the hyperbolic sector and that it is possible to describe a string propagating on an extremal BH background without much effort.
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# KUNS-1654hep-th/0003231 Noncommutative Monopole from Nonlinear Monopole
## 1 Introduction
Recently non-commutative gauge theory has received much attention for its origin in string theory. The effective action of D-brane in string theory with a constant NS 2-form $`B_{ij}`$ is the non-commutative Born-Infeld theory when the point splitting regularization is adopted . On the other hand, if we adopt the Pauli-Villars regularization we obtain the ordinary Born-Infeld theory. Since the method of regularization should not change the physical S-matrices, the two descriptions are argued to be related by a field redefinition (called the Seiberg-Witten map).
This relation has been investigated intensively from various aspects, and among other things, from various BPS solutions . Since the constant NS 2-form serves as a uniform magnetic field, if we view the monopole solution as a D-string ending on a D3-brane, we expect the D-string tilts due to force balance between the D-string tension and the magnetic force at the endpoint . This system was analyzed in as the solution of the linearly realized BPS equation in the commutative space.
However if we would like to see the tilts directly from the non-commutative viewpoint, it would be a hard task. We can only solve the linearly realized BPS equation in the perturbation expansion with respect to the non-commutativity parameter $`\theta `$. And even if we have solved the BPS equation, we would have to know how to extract the eigenvalues for the brane interpretation of . In our previous works we proposed the non-commutative eigenvalue equation to analyze the asymptotic behavior and confirmed the tilts.
In , it was claimed that the brane interpretation is possible if we transform the results into the commutative viewpoint by the Seiberg-Witten map. Actually they argued that the linearly realized BPS monopole in the non-commutative space is mapped to the non-linearly realized BPS monopole in the commutative space by extending the argument in the instanton case . This relation is persuasive for the following reason. Since we are considering the monopole in the non-commutative space with the property that the field strength and the covariant derivative of the Higgs field vanish at the infinity and this property seems unchanged under the Seiberg-Witten map, it is expected that what relates to the non-commutative monopole by the Seiberg-Witten map should also have the same property. The condition of preserving the combination of supersymmetries which is unbroken at the infinity where the field strength and the covariant derivative vanish is exactly the non-linearly realized BPS equation.
One might worry that the linear BPS equation in the non-commutative space usually derived from the Yang-Mills theory cannot be obtained from the Born-Infeld theory which is of our prime interest in discussing the Seiberg-Witten map. However, it is shown that even in the non-commutative space the linear BPS equation of the Yang-Mills theory reproduces the equation of motion of the Born-Infeld theory by extending earlier discussions for the non-Abelian Born-Infeld theory if we adopt the symmetrized trace prescription for the definition of the determinant. Note also that this kind of transformation is possible owing to another fact that the solution is unchanged even when the derivative corrections to the Born-Infeld action are taken into account .
Moreover it is proposed that the non-linear BPS monopole in the commutative space is related to the linear BPS monopole in the commutative space by the rotation in the target space. In the electric case of , an exact treatment of the soliton was given though there are no discussions on the non-linear BPS equation. In the magnetic case of the discussion on the non-linear BPS equation was restricted to the approximation $`r^22\pi \alpha ^{}(2\pi \alpha ^{})^2B`$.
In this paper we shall extend the works of to explore the non-commutative BPS monopole from the non-linear BPS monopole in the commutative space. First we shall solve the non-linear BPS equation in the commutative space exactly without any approximation. We find that the solution is nothing but the one obtained by rotating the solution of the linear BPS equation in the target space.
After establishing the non-linear BPS monopole in the commutative space, we explicitly write down the first few terms in the expansion of the NS 2-form $`B_{ij}`$. What we find is terms in a mess and at first sight it seems impossible that they are related to the non-commutative monopole by the Seiberg-Witten map. We shall resolve this problem by using the moduli of the open string, namely, the open string metric $`G`$ and the non-commutativity parameter $`\theta `$. This resolution is regarded as an evidence for the claim that the non-linear monopole is transformed into the non-commutative monopole by the Seiberg-Witten map. Finally we map the non-linear monopole into the non-commutative space. We confirm that it satisfies the non-commutative BPS equation up to $`O(\theta ^2)`$.
## 2 Nonlinear BPS equation
In this section, we shall explicitly solve the non-linear BPS equation in the commutative space. First we shall recall the linearly realized supersymmetries $`\delta _\mathrm{L}`$ and non-linearly realized supersymmetries $`\delta _{\mathrm{NL}}`$ of the Born-Infeld action:
$`\delta _\mathrm{L}\lambda `$ $`=`$ $`{\displaystyle \frac{1}{2\pi \alpha ^{}}}M_{mn}^+\sigma ^{mn}\eta ,`$ (1)
$`\delta _\mathrm{L}\overline{\lambda }`$ $`=`$ $`{\displaystyle \frac{1}{2\pi \alpha ^{}}}M_{mn}^{}\sigma ^{mn}\overline{\eta },`$ (2)
$`\delta _{\mathrm{NL}}\lambda `$ $`=`$ $`{\displaystyle \frac{1}{4\pi \alpha ^{}}}\left(1PfM+\sqrt{1TrM^2/2+(PfM)^2}\right)\eta ^{},`$ (3)
$`\delta _{\mathrm{NL}}\overline{\lambda }`$ $`=`$ $`{\displaystyle \frac{1}{4\pi \alpha ^{}}}\left(1+PfM+\sqrt{1TrM^2/2+(PfM)^2}\right)\overline{\eta }^{},`$ (4)
where $`M`$ denotes
$`M=(2\pi \alpha ^{})\left(\begin{array}{cccc}0& _1\mathrm{\Phi }& _2\mathrm{\Phi }& _3\mathrm{\Phi }\\ _1\mathrm{\Phi }& 0& (F_3+B_3)& (F_2+B_2)\\ _2\mathrm{\Phi }& (F_3+B_3)& 0& (F_1+B_1)\\ _3\mathrm{\Phi }& (F_2+B_2)& (F_1+B_1)& 0\end{array}\right),`$ (5)
with the magnetic field $`F_i=ϵ_{ijk}F_{jk}/2`$, a constant NS 2-form background $`B_i=ϵ_{ijk}B_{jk}/2`$ and the Higgs field $`\mathrm{\Phi }`$. Here we turn on only the spatial components of the field strength and the NS 2-form. The matrix $`M`$ has been obtained by regarding the Euclidean time component of the gauge field as the Higgs field $`\mathrm{\Phi }`$ and discarding the time derivatives. We shall set $`2\pi \alpha ^{}=1`$ for simplicity hereafter, however we can restore it on the dimensional ground anytime we like. The non-linear BPS equation is the condition for preserving the linear combination of $`\delta _\mathrm{L}`$ and $`\delta _{\mathrm{NL}}`$ which is unbroken at the infinity where the field strength and the derivative of the Higgs field vanish:
$`{\displaystyle \frac{𝑭+𝑩\mathbf{}\mathrm{\Phi }}{1+(𝑭+𝑩)\mathbf{}\mathrm{\Phi }+\sqrt{1+(𝑭+𝑩)^2+(\mathbf{}\mathrm{\Phi })^2+\left((𝑭+𝑩)\mathbf{}\mathrm{\Phi }\right)^2}}}={\displaystyle \frac{𝑩}{1+\sqrt{1+𝑩^2}}}.`$ (6)
This BPS equation is not so complicated to solve as it looks. The starting point is similar to the case of instanton . First we note that eq. (6) implies $`𝑭\mathbf{}\mathrm{\Phi }`$ is proportional to $`𝑩`$:
$`𝑭\mathbf{}\mathrm{\Phi }=f𝑩,`$ (7)
where $`f`$ is an unknown function. The key point to solve this BPS equation is to rewrite eq. (6) as
$`(f+1)\left(1+\sqrt{1+𝑩^2}\right)1(𝑭+𝑩)\mathbf{}\mathrm{\Phi }=\sqrt{1+(𝑭+𝑩)^2+(\mathbf{}\mathrm{\Phi })^2+\left((𝑭+𝑩)\mathbf{}\mathrm{\Phi }\right)^2}.`$ (8)
Taking the square of this equation (8) and using the relation (7) to eliminate the magnetic field $`𝑭`$ when necessary, we find that eq. (8) is reduced simply to
$`f=(\mathbf{}\mathrm{\Phi })^2+(f+1)𝑩\mathbf{}\mathrm{\Phi }.`$ (9)
Another equation for $`f`$ and $`\mathrm{\Phi }`$ besides (9) is obtained by taking the divergence of the relation (7) and using the Bianchi identity $`\mathbf{}𝑭=0`$,
$`\mathbf{}^2\mathrm{\Phi }=𝑩\mathbf{}f.`$ (10)
Now we have a system of differential equations (9) and (10) for the scalar quantities $`f`$ and $`\mathrm{\Phi }`$. After eliminating $`f`$ we find quite a non-linear equation for $`\mathrm{\Phi }`$:
$`\mathbf{}^2\mathrm{\Phi }\left(1𝑩\mathbf{}\mathrm{\Phi }\right)^2+2𝑩\mathbf{}\mathbf{}\mathrm{\Phi }\mathbf{}\mathrm{\Phi }\left(1𝑩\mathbf{}\mathrm{\Phi }\right)+𝑩\mathbf{}𝑩\mathbf{}\mathrm{\Phi }\left(1+(\mathbf{}\mathrm{\Phi })^2\right)=0.`$ (11)
Hereafter we shall suppose the constant background $`𝑩`$ is in the $`z`$ direction and rewrite the equation (11) in the cylindrical coordinate $`(\rho ,\phi ,z)`$ with $`x=\rho \mathrm{cos}\phi `$ and $`y=\rho \mathrm{sin}\phi `$,
$`(_\rho ^2\mathrm{\Phi }+_\rho \mathrm{\Phi }/\rho +_z^2\mathrm{\Phi })(1B_z\mathrm{\Phi })^2+2B(_\rho _z\mathrm{\Phi }_\rho \mathrm{\Phi }+_z^2\mathrm{\Phi }_z\mathrm{\Phi })(1B_z\mathrm{\Phi })`$
$`+B^2_z^2\mathrm{\Phi }(1+(_\rho \mathrm{\Phi })^2+(_z\mathrm{\Phi })^2)=0.`$ (12)
This differential equation looks impossible to solve. However, we can apply the idea of to find that the solution is exactly the one obtained by rotating the solution of the linear BPS equation in the target space by an angle $`\varphi `$ with $`\mathrm{tan}\varphi =B`$. This idea is convincing for the following reason.<sup>*</sup><sup>*</sup>* We are grateful to T. Hirayama for a valuable discussion on this point. Originally the string theory has the $`SO(1,9)`$ Lorentz symmetry and 32 supersymmetries. Taking the static gauge the Lorentz symmetry is broken into $`SO(1,3)\times SO(6)`$ and half of the supersymmetries are broken. The broken symmetries are realized non-linearly. If we rotate the target space and still take the static gauge by adopting a different worldsheet coordinate, we would find that originally linearly realized symmetries correspond in general to some combinations of the linear and the non-linear ones. Therefore the linear BPS equation and the non-linear BPS equation should be related by a target space rotation.
To see this explicitly we change our variables into those with bars by the target space rotation,
$`\left(\begin{array}{c}\overline{\mathrm{\Phi }}\\ \overline{z}\end{array}\right)=\left(\begin{array}{cc}\mathrm{cos}\varphi & \mathrm{sin}\varphi \\ \mathrm{sin}\varphi & \mathrm{cos}\varphi \end{array}\right)\left(\begin{array}{c}\mathrm{\Phi }\\ z\end{array}\right),`$ (13)
and show that the solution for the variables with bars is the same as that of the linear BPS equation. First we have to rewrite the equation (12) by changing $`\mathrm{\Phi }`$, $`_\rho `$, $`_z`$ into $`\overline{\mathrm{\Phi }}`$, $`\overline{}_\rho /\rho |_{\overline{z}}`$, $`\overline{}_z/\overline{z}|_\rho `$. Note that, though we do not change the coordinate $`\rho `$, $`\overline{}_\rho `$ is different form $`_\rho `$ because the coordinate to be fixed is different between them. The formulas for rewriting $`_z\mathrm{\Phi }`$ and $`_\rho \mathrm{\Phi }`$ into $`\overline{}_z\overline{\mathrm{\Phi }}`$ and $`\overline{}_\rho \overline{\mathrm{\Phi }}`$ read
$`_z\mathrm{\Phi }`$ $`=`$ $`{\displaystyle \frac{\mathrm{cos}\varphi \overline{}_z\overline{\mathrm{\Phi }}\mathrm{sin}\varphi }{\mathrm{cos}\varphi +\mathrm{sin}\varphi \overline{}_z\overline{\mathrm{\Phi }}}},`$ (14)
$`_\rho \mathrm{\Phi }`$ $`=`$ $`{\displaystyle \frac{\overline{}_\rho \overline{\mathrm{\Phi }}}{\mathrm{cos}\varphi +\mathrm{sin}\varphi \overline{}_z\overline{\mathrm{\Phi }}}},`$ (15)
where the first formula is directly obtained from the rotation (13) and the second formula is a consequence of the chain rule formula
$`{\displaystyle \frac{\overline{\mathrm{\Phi }}}{\rho }}|_z={\displaystyle \frac{\rho }{\rho }}|_z{\displaystyle \frac{\overline{\mathrm{\Phi }}}{\rho }}|_{\overline{z}}+{\displaystyle \frac{\overline{z}}{\rho }}|_z{\displaystyle \frac{\overline{\mathrm{\Phi }}}{\overline{z}}}|_\rho ,`$ (16)
and the relations $`\overline{\mathrm{\Phi }}/\rho |_z=\mathrm{cos}\varphi \mathrm{\Phi }/\rho |_z`$ and $`\overline{z}/\rho |_z=\mathrm{sin}\varphi \mathrm{\Phi }/\rho |_z`$. In the same way we also find the similar formulas for higher derivatives,
$`_z^2\mathrm{\Phi }`$ $`=`$ $`\overline{}_z^2\overline{\mathrm{\Phi }}/(\mathrm{cos}\varphi +\mathrm{sin}\varphi \overline{}_z\overline{\mathrm{\Phi }})^3,`$ (17)
$`_z_\rho \mathrm{\Phi }`$ $`=`$ $`\left[\mathrm{cos}\varphi \overline{}_z\overline{}_\rho \overline{\mathrm{\Phi }}+\mathrm{sin}\varphi (\overline{}_z\overline{}_\rho \overline{\mathrm{\Phi }}\overline{}_z\overline{\mathrm{\Phi }}\overline{}_z^2\overline{\mathrm{\Phi }}\overline{}_\rho \overline{\mathrm{\Phi }})\right]/(\mathrm{cos}\varphi +\mathrm{sin}\varphi \overline{}_z\overline{\mathrm{\Phi }})^3,`$ (18)
$`_\rho ^2\mathrm{\Phi }`$ $`=`$ $`[(\mathrm{cos}\varphi )^2\overline{}_\rho ^2\overline{\mathrm{\Phi }}+2\mathrm{cos}\varphi \mathrm{sin}\varphi (\overline{}_z\overline{}_\rho \overline{\mathrm{\Phi }}\overline{}_\rho \overline{\mathrm{\Phi }}+\overline{}_\rho ^2\overline{\mathrm{\Phi }}\overline{}_z\overline{\mathrm{\Phi }})`$ (19)
$`+(\mathrm{sin}\varphi )^2(\overline{}_z^2\overline{\mathrm{\Phi }}(\overline{}_\rho \overline{\mathrm{\Phi }})^22\overline{}_z\overline{}_\rho \overline{\mathrm{\Phi }}\overline{}_z\overline{\mathrm{\Phi }}\overline{}_\rho \overline{\mathrm{\Phi }}+\overline{}_\rho ^2\overline{\mathrm{\Phi }}(\overline{}_z\overline{\mathrm{\Phi }})^2)]/(\mathrm{cos}\varphi +\mathrm{sin}\varphi \overline{}_z\overline{\mathrm{\Phi }})^3.`$
Using these formulas, the terribly non-linear equation (12) now becomes
$`\overline{}_\rho ^2\overline{\mathrm{\Phi }}+\overline{}_\rho \overline{\mathrm{\Phi }}/\rho +\overline{}_z^2\overline{\mathrm{\Phi }}=0,`$ (20)
which is nothing but the three-dimensional laplace equation. The solution to eq. (20) is given by the sum of the Coulomb term and the linear term determined by the boundary condition in the asymptotic region:
$`\overline{\mathrm{\Phi }}={\displaystyle \frac{q}{\sqrt{\rho ^2+\overline{z}^2}}}+B\overline{z}.`$ (21)
Turning back to the variables without bars using the relation (13), our final result for the Higgs field $`\mathrm{\Phi }`$ is given as the solution of the algebraic equation,
$`\left((1+B^2)\rho ^2+z^22Bz\mathrm{\Phi }+B^2\mathrm{\Phi }^2\right)\mathrm{\Phi }^2=q^2,`$ (22)
or its covariant form
$`\left((1+𝑩^2)𝒙^2(𝑩𝒙)^22𝑩𝒙\mathrm{\Phi }+𝑩^2\mathrm{\Phi }^2\right)\mathrm{\Phi }^2=q^2.`$ (23)
The explicit expression of the first few terms in the expansion with respect to $`B`$ is
$`\mathrm{\Phi }={\displaystyle \frac{q}{r}}+{\displaystyle \frac{q^2𝑩𝒙}{r^4}}{\displaystyle \frac{1}{2}}{\displaystyle \frac{q𝑩^2}{r}}+{\displaystyle \frac{1}{2}}{\displaystyle \frac{q(𝑩𝒙)^2}{r^3}}{\displaystyle \frac{1}{2}}{\displaystyle \frac{q^3𝑩^2}{r^5}}+{\displaystyle \frac{5}{2}}{\displaystyle \frac{q^3(𝑩𝒙)^2}{r^7}},`$ (24)
with $`r=\sqrt{\rho ^2+z^2}`$.
Similarly the magnetic field is also obtained from the relations (7) and (9) as
$`𝑭=\mathbf{}\mathrm{\Phi }+{\displaystyle \frac{(\mathbf{}\mathrm{\Phi })^2+𝑩\mathbf{}\mathrm{\Phi }}{1𝑩\mathbf{}\mathrm{\Phi }}}𝑩.`$ (25)
Using our result (23) we can rewrite this expression (25) by eliminating the derivatives of the Higgs field $`\mathbf{}\mathrm{\Phi }`$,
$`𝑭={\displaystyle \frac{(1+𝑩^2)(𝒙+2\mathrm{\Phi }𝑩)\mathrm{\Phi }}{(1+𝑩^2)𝒙^2(𝑩𝒙)^23(𝑩𝒙)\mathrm{\Phi }+2𝑩^2\mathrm{\Phi }^2}}.`$ (26)
## 3 Physical interpretation
In this section, we shall give some comments and physical interpretations to our solution. First, our result is obtained without any approximations and the expansion (24) is consistent with the result obtained in . The behavior of the Higgs field $`\mathrm{\Phi }`$ (22) against the worldsheet coordinate $`(z,\rho )`$ is depicted in Fig. 2 (A). Note that in the right hand side of eq. (8) we do not persist in taking the positive branch of the square root, because it forces us to discard part of the solution given in Fig. 2 (A). The spike-like behavior of the Higgs field represents the D-string attached to the D3-brane in the brane interpretation of . This D-string tilts due to the uniform magnetic field and the tilt angle is exactly the one expected from the force balance . Here we find that the Higgs field is multi-valued due to this tilt (see Fig. 2 (B) which shows the multi-valuedness of $`\mathrm{\Phi }`$ on the $`\rho =0`$ plane). This multi-valuedness is a consequence of the fact that the eq. (22) determining $`\mathrm{\Phi }`$ is a fourth order algebraic equation which in general has four solutions. Another solution not depicted in Fig. 2 (B) is a fake one with $`\mathrm{\Phi }<0`$. This multi-valuedness implies that the Dirac monopole might be ill-defined as a field theoretic soliton in the non-linear BPS equation and probably also in the non-commutative BPS equation via the Seiberg-Witten map. However, the multi-valuedness is inevitable from the string theory viewpoint.
Though we do not know the non-linear BPS equation for the non-Abelian case due to the complexity of the ordering in the determinant, it is expected that the Higgs field related to the non-commutative monopole by the Seiberg-Witten map is that obtained by rotating the solution of the linear BPS equation in the target space. Note that in this case of the ’t Hooft-Polyakov monopole, the problematic multi-valuedness in the Dirac monopole does not necessarily appear. From the behavior near the origin $`r=0`$ of the exact solution in with $`C=\mathrm{\Phi }`$,
$`\mathrm{\Phi }=(Cr/\mathrm{tanh}Cr1)/rC^2r/3,`$ (27)
we can read off that the tangent vector of the deformed D3-brane is $`\stackrel{}{v}=(1,C^2/3+B)`$ and that of the worldsheet parameterization is $`\stackrel{}{w}=(1,B)`$ in the rotated coordinate system depicted in Fig. 3. Therefore the single-valuedness condition is expressed as the positivity of the inner product of these two vectors:
$`\stackrel{}{v}\stackrel{}{w}=1C^2B/3+B^2>0.`$ (28)
This implies that at some value of NS 2-form even the ’t Hooft-Polyakov monopole is not single-valued, which is something we have never experienced in the usual field-theoretical solitons.
Finally, we would like to comment on the small $`B`$ expansion (24). Unlike our experience of the non-commutative monopole in the flat space where the only parameter is $`\theta `$, at $`O(B^2)`$ of (24) we find terms proportional to $`r^1`$ as well as $`r^5`$, which implies the parameter $`2\pi \alpha ^{}`$ also appears. If we expect the present result is transformed to the non-commutative monopole by the Seiberg-Witten map, this kind of double expansion seems impossible.
The resolution to this paradox is given by considering the moduli of the open string metric $`G`$ and the non-commutativity parameter $`\theta `$. When we relate the non-commutative gauge theory to its commutative counterpart, we should also relate the moduli by
$`{\displaystyle \frac{1}{G}}+\theta ={\displaystyle \frac{1}{g+B}}.`$ (29)
Since we set the metric in the commutative space to the flat one $`g_{ij}=\delta _{ij}`$ and turn on only the spatial NS 2-form $`B_i`$, our moduli are
$`G_{ij}=(1+𝑩^2)\delta _{ij}B_iB_j,\theta ^{ij}={\displaystyle \frac{ϵ_{ijk}B_k}{1+𝑩^2}},`$ (30)
with a necessarily non-trivial open string metric. Using these open string moduli we can construct several kinds of scalars:
$`R^2`$ $``$ $`G_{ij}x^ix^j=(1+B^2)\rho ^2+z^2,`$ (31)
$`\theta x`$ $``$ $`\sqrt{G}ϵ_{ijk}\theta ^{jk}x^i=Bz,`$ (32)
$`\theta ^2`$ $``$ $`G^{ij}\sqrt{G}ϵ_{ikl}\theta ^{kl}\sqrt{G}ϵ_{jmn}\theta ^{mn}=B^2.`$ (33)
Note that in the case of a non-trivial metric the $`ϵ`$ tensor should always be accompanied with $`\sqrt{G}`$. In terms of these scalars our result (22) can be rewritten into an expression with only one parameter $`\theta `$:
$`\left(R^2+2\theta x\mathrm{\Phi }+\theta ^2\mathrm{\Phi }^2\right)\mathrm{\Phi }^2=q^2.`$ (34)
Note that from the dimensional ground there can appear no $`2\pi \alpha ^{}`$ in (34). Our observation here shows that we have two viewpoints for the non-linear monopole. One is with the flat space and a NS 2-form and the other is with the non-trivial metric and the non-commutativity parameter. Similarly if we rewrite the magnetic field (26) into the covariant field strength, we will also find an expression without $`2\pi \alpha ^{}`$:
$`F_{ij}={\displaystyle \frac{\sqrt{G}ϵ_{ijk}(x^k2\mathrm{\Phi }\theta ^k)\mathrm{\Phi }}{R^2+3\theta x\mathrm{\Phi }+2\theta ^2\mathrm{\Phi }^2}}.`$ (35)
Similar expression for the gauge field is difficult to find because of the existence of the Dirac string.
Our analysis so far is believed to be related to the non-commutative monopole by the Seiberg-Witten map . In the case of a constant non-trivial metric the BPS equation in the non-commutative space should be given by
$`\widehat{F}_{ij}=\sqrt{G}ϵ_{ijm}G^{mn}\widehat{D}_n\widehat{\mathrm{\Phi }},`$ (36)
as can be seen from the BPS bound arguments . However since the non-trivial metric is constant, we can always orthonormalize it globally by the vielbein:
$`E_\alpha ^iE_\beta ^jG_{ij}=\delta _{\alpha \beta }.`$ (37)
Therefore if we would like to find the non-commutative monopole in the flat space we have to collect all our results of the non-linear BPS equation, rewrite them in the covariant form, make a coordinate transformation $`x^iE_i^\alpha x^i`$ into the flat space, and transform them into the non-commutative space by the Seiberg-Witten map. Our result in the flat space up to $`O(\theta ^2)`$ for the Higgs field is
$`\mathrm{\Phi }={\displaystyle \frac{q}{r}}{\displaystyle \frac{q^2\theta x}{r^4}}{\displaystyle \frac{q^3\theta ^2}{2r^5}}+{\displaystyle \frac{5q^3(\theta x)^2}{2r^7}},`$ (38)
where we have rewritten $`R`$ into $`r`$ because now we are in the flat space. And the gauge field corresponding to the field strength (35) is given by
$`A_i=A_i^0+A_i^1,`$ (39)
with
$`A_1^0={\displaystyle \frac{qy}{r(r+z)}},A_2^0={\displaystyle \frac{qx}{r(r+z)}},A_3^0=0,`$ (40)
$`A_i^1={\displaystyle \frac{q^2ϵ_{ijk}\theta _jx_k}{r^4}}{\displaystyle \frac{5q^3ϵ_{ijk}\theta _jx_k\theta _mx_m}{2r^7}}.`$ (41)
Note that due to the presence of the Dirac string the solution in the zero-th order in $`\theta `$ cannot be written in a spherically symmetric form. We have explicitly transformed this result into the non-commutative space by the Seiberg-Witten map
$`\widehat{A}_i`$ $`=`$ $`A_i{\displaystyle \frac{1}{2}}\theta ^{kl}A_k(_lA_i+F_{li})+{\displaystyle \frac{1}{2}}\theta ^{kl}\theta ^{mn}A_k(_lA_m_nA_i_lF_{mi}A_n+F_{lm}F_{ni}),`$ (42)
$`\widehat{\mathrm{\Phi }}`$ $`=`$ $`\mathrm{\Phi }{\displaystyle \frac{1}{2}}\theta ^{kl}A_k(2_l\mathrm{\Phi })+{\displaystyle \frac{1}{2}}\theta ^{kl}\theta ^{mn}A_k(_lA_m_n\mathrm{\Phi }_l_m\mathrm{\Phi }A_n+F_{lm}_n\mathrm{\Phi }),`$ (43)
and checked that it indeed satisfies the non-commutative BPS equation by using a symbolic manipulation software. However due to the Dirac string the covariant form is not available and the result is too complicated and not suitable to be written here.
## 4 Summary and further directions
In this paper we extended the earlier idea of rotating the system to solve the non-linear BPS equation without any approximation. Since we solved it exactly, the multi-valuedness problem appeared. We also pointed out the open string metric is in general non-trivial and a careful treatment is necessary. Finally we transformed our result into the non-commutative space by the Seiberg-Witten map and confirmed it satisfies the non-commutative BPS equation.
In our exact manipulation we clarified the physical meaning of the Higgs field in the non-linear BPS equation. Hence in the non-Abelian case, even though we do not know the non-linear BPS equation, we expect the solution for the Higgs field related to the non-commutative monopole is that obtained by rotating the solution of the linear BPS equation in the target space. However the meaning of the gauge field is still unclear. So we do not know what to expect for the gauge field. To understand it is an interesting subject.
Acknowledgment
We would like to thank K. Hashimoto, H. Hata and T. Hirayama for valuable discussions and comments and H. Hata for careful reading of the present manuscript. This work is supported in part by Grant-in-Aid for Scientific Research from Ministry of Education, Science, Sports and Culture of Japan (#04633). The author is supported in part by the Japan Society for the Promotion of Science under the Predoctoral Research Program.
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# II. Territory covered by 𝑁 random walkers on stochastic fractals. The percolation aggregate
## I INTRODUCTION
In the preceding paper we studied the problem of evaluating the average number $`S_N(t)`$ of distinct sites (or territory covered) by $`N`$ independent random walkers that, all starting from the same site, diffuse on a deterministic fractal during time $`t`$. An answer to this problem was given in Ref. in terms of an asymptotic series completely analogous to that given in Refs. for Euclidean substrates. In this present paper we turn our attention to the case of stochastic (or disordered) fractal media and again propose the same kind of asymptotic solution, although, in contrast with all the previous cases, we now have to translate the procedure to the topological or chemical language.
Stochastic media are not constructed by the iteration of an invariable rule, such as that used in deterministic fractals, but are rather the result of a random process. Consequently, the concept of generator and rigorous self-similarity is absent and their fractal nature is recognized by the scaling of statistical quantities. Many natural objects share this statistical-fractal structure so that stochastic models seem to be more suitable to represent diffusion in real media. Of particular interest is the percolation model that has been used to characterize many disordered systems . This model is constructed by filling a regular lattice with “occupied” sites with a certain probability. Nearest neighbor occupied sites are supposed to be connected and form a series of clusters. At a certain critical concentration $`p_c`$ an infinite cluster appears, which is called the incipient percolation aggregate or percolation cluster at criticality.
The quantity we are interested in, $`S_N(t)`$, is, for disordered media, the result of a double average: an average over the walks that the $`N`$ random walkers can perform over a given lattice, followed by an average over many (ideally, all) realizations of the random lattice. This situation leads to certain subtleties, absent in deterministic fractals, that require special treatment. In particular, $`S_N(t)`$ can be expressed by
$$S_N(t)=\left\{1\left[\mathrm{\Gamma }_t(𝐫)\right]^N\right\},$$
(1)
where the sum is over all the sites of each fractal lattice, $`\mathrm{\Gamma }_t(𝐫)`$ is the survival probability, i.e., the probability that site $`𝐫`$ has not been visited by time $`t`$ by a single random walker starting from the origin, $`\left\{1\left[\mathrm{\Gamma }_t(𝐫)\right]^N\right\}`$ represents the mean territory explored by the $`N`$ random walkers on a given lattice (the first average) , and $`[\mathrm{}]`$ indicates that the average (the second average) of $`[\mathrm{}]`$ has to be performed over all possible stochastic lattices compatible with the random generation rules. Equation (1) can be rewritten as
$$S_N(t)=\underset{m=0}{\overset{\mathrm{}}{}}\underset{i=1}{\overset{n(m)}{}}\left\{1\left[\mathrm{\Gamma }_t(𝐫_{m,i})\right]^N\right\},$$
(2)
where $`𝐫_{m,i}`$ stands for the $`i`$-th site out of $`n(m)`$ that are separated from the origin by a Euclidean distance between $`m\mathrm{\Delta }rr_m`$ and $`(m+1)\mathrm{\Delta }r`$ with $`\mathrm{\Delta }r`$ small (say, of the order of the lattice spacing). If $`\mathrm{\Gamma }_t(𝐫_{m,i})`$ is almost independent of $`i`$ and the lattice realization, i.e., if the fluctuations in the probability density $`\mathrm{\Gamma }_t(𝐫_{m,i})`$ follow a narrow distribution, then one could approximate $`\mathrm{\Gamma }_t(𝐫_{m,i})\mathrm{\Gamma }_t(𝐫_{m,i})\mathrm{\Gamma }_t(r_m)`$, and therefore estimate $`S_N(t)`$ by
$$S_N(t)=\underset{m=0}{\overset{\mathrm{}}{}}\left\{1\left[\mathrm{\Gamma }_t(r_m)\right]^N\right\}n(m),$$
(3)
where $`n(m)`$ is the average number of fractal sites separated from the origin by a distance bracketed by $`r_m`$ and $`r_m+\mathrm{\Delta }r`$. This is essentially the starting relationship used (implicitly) by Havlin et al. to find that, for large $`N`$,
$$S_N(t)t^{d_s/2}(\mathrm{ln}N)^{d_f/u}$$
(4)
in the non-trivial time regime (or regime II) . Here $`u=d_w/(d_w1)`$, $`d_w`$ is the anomalous diffusion exponent, $`d_f`$ is the fractal dimension of the substrate, and $`d_s=2d_f/d_w`$ is the spectral dimension. However, the hypothesis leading to Eq. (3) is in general false, as we will explicitly show in Sec. II by means of numerical simulations of $`\mathrm{\Gamma }_t(𝐫_{m,i})`$ for the two-dimensional percolation cluster at criticality. Indeed, it is known that the fluctuations of the probability density $`P(r,t)`$ of random walks (also called the propagator or Green’s function), which is a statistical quantity closely related to the survival probability, exhibits a broad logarithmic distribution for some random fractals such as percolation clusters and self-avoiding walks. Bunde et al. have found that the quantity $`P(r,t)^q`$ exhibits multifractal scaling, $`P(r,t)^qP(r,t)^{\tau (q)}`$, where $`\tau (q)q^\gamma `$ and $`\gamma =(d_w^{\mathrm{}}1)/(d_w1)`$ and $`d_w^{\mathrm{}}`$ is the chemical random walk dimension. This behavior is a consequence of the large fluctuations of $`P(r,t)`$ for fixed $`r`$ and $`t`$ from a given aggregate to another. Nevertheless, these authors have also shown that the distribution of the propagator in the chemical $`\mathrm{}`$ space, $`P(\mathrm{},t)`$, is narrow and, consequently, $`P(\mathrm{},t)^qP(\mathrm{},t)^q`$. The chemical distance $`\mathrm{}`$, the length of the shortest path between two sites along lattice bonds, is a more natural measure than the Euclidean distance in disordered systems. It is, for example, the distance used in the calculation of optimum paths in cities.
Let $`\mathrm{}_{m,i}`$ label the $`i`$-th site out of those $`n(m)`$ that are placed at a chemical distance $`\mathrm{}`$ from a given origin with $`\mathrm{}_m\mathrm{}<\mathrm{}_{m+1}`$, $`\mathrm{}_m=m\mathrm{\Delta }\mathrm{}`$ and $`\mathrm{\Delta }\mathrm{}`$ small (say, of the order of the lattice spacing), and let $`\mathrm{\Gamma }_t(\mathrm{}_{m,i})`$ be the survival probability in the chemical space defined as the probability that site $`\mathrm{}_{m,i}`$ has not been visited by time $`t`$ by a single random walker starting from the origin. Then we can rewrite Eq. (2) in the chemical $`\mathrm{}`$ space as
$$S_N(t)=\underset{m=0}{\overset{\mathrm{}}{}}\underset{i=1}{\overset{n(m)}{}}\left\{1\left[\mathrm{\Gamma }_t(\mathrm{}_{m,i})\right]^N\right\}.$$
(5)
One may expect that the distribution of $`\mathrm{\Gamma }_t(\mathrm{}_{m,i})`$ for fixed $`\mathrm{}_m`$ and $`t`$ is as narrow as the distribution of the propagator in the chemical space. (We support this conjecture in Sec. III by means of numerical simulations of $`\mathrm{\Gamma }_t(\mathrm{}_{m,i})`$ in two-dimensional percolation clusters at criticality.) In this case
$$\mathrm{\Gamma }_t(\mathrm{}_{m,i})\mathrm{\Gamma }_t(\mathrm{}_{m,i})\mathrm{\Gamma }_t(\mathrm{}_m)$$
(6)
for all possible lattice realizations so that $`\left[\mathrm{\Gamma }_t(\mathrm{}_{m,i})\right]^N\mathrm{\Gamma }_t(\mathrm{}_{m,i})^N`$ and, therefore, $`S_N(t)`$ can be approximated by
$$S_N(t)=\underset{m=0}{\overset{\mathrm{}}{}}\left\{1\left[\mathrm{\Gamma }_t(\mathrm{}_m)\right]^N\right\}n(m),$$
(7)
where $`n(m)`$ is the average number of fractal sites separated from the origin by a chemical distance with value between $`\mathrm{}_m`$ and $`\mathrm{}_m+\mathrm{\Delta }\mathrm{}`$. From this formula and following the procedure outlined in the preceding paper , in Sec. II we will arrive at an expression for $`S_N(t)`$ for the non-trivial time regime whose leading asymptotic behavior coincides, apart from the value of the prefactor, with the recent proposal, based on a scaling approach, of Dräger and Klafter :
$$S_N(t)t^{d_s/2}(\mathrm{ln}N)^{d_{\mathrm{}}/v}$$
(8)
with $`v=d_w^{\mathrm{}}/(d_w^{\mathrm{}}1)`$. Equation (8) differs from the relationship proposed by Havlin et al. , Eq. (4), for those cases such as that considered in this present paper where $`d_{\text{min}}1`$. Both Havlin et al. and Dräger and Klafter supported their conjectures by means of data collapsing plots of computer simulation results obtained for two- and three-dimensional percolation aggregates, respectively. We here want also draw attention to the risk involved in this method of analysis when the influence of the corrective terms is not properly considered since these terms have a large influence on the final value of $`S_N(t)`$.
The paper is organized as follows. In Sec. II we present the asymptotic evaluation of $`S_N(t)`$ for stochastic fractal lattices as a translation to the chemical space of the procedure implemented in the Euclidean and deterministic fractal cases . A less rigorous but fairly simple and instructive method for obtaining the main asymptotic term and estimating the corrective terms of $`S_N(t)`$ for large $`N`$ is also presented. In Sec. III we report simulation results for the survival probability of a random walker on a two-dimensional incipient percolation aggregate when a trap is placed at a site at a fixed chemical distance $`\mathrm{}`$ or Euclidean distance $`r`$. We find that the distribution is narrow \[broad\] if the traps are located at a fixed chemical \[Euclidean\] distance. The parameters governing (i) the asymptotic behavior of $`\mathrm{\Gamma }_t(\mathrm{})`$ ($`c`$, $`v`$, $`\mu `$ and $`A`$), (ii) how the fractal volume grows ($`V_0^{\mathrm{}}`$ and $`d_{\mathrm{}}`$), and (iii) how fast a single walker diffuses ($`D_{\mathrm{}}`$ and $`d_w^{\mathrm{}}`$) are estimated in this section. In Sec. IV we compare the zeroth- and first-order asymptotic expansion for $`S_N(t)`$ with simulation results obtained for the two-dimensional incipient percolation aggregate. We also criticize the reliability of typical collapsing plots for the determination of the dominant trend of $`S_N(t)`$. We conclude with some remarks in Sec. V.
## II TERRITORY COVERED BY $`N`$ RANDOM WALKERS ON A STOCHASTIC FRACTAL SUBSTRATE
In this section we will translate the results of the previous paper to chemical language. Reasons for this procedure have already been given in Sec. I. We start by replacing Eq. (7) by its continuum approximation
$`S_N(t)`$ $``$ $`{\displaystyle _0^{\mathrm{}}}\left\{1\left[\mathrm{\Gamma }_t(\mathrm{})\right]^N\right\}d_{\mathrm{}}V_0^{\mathrm{}}\mathrm{}^{d_{\mathrm{}}1}𝑑\mathrm{}`$ (9)
$``$ $`V_0^{\mathrm{}}d_{\mathrm{}}(2D_{\mathrm{}})^{d_{\mathrm{}}/2}t^{d_{\mathrm{}}/d_w}{\displaystyle _0^{\mathrm{}}}\left\{1\left[\mathrm{\Gamma }_t(\xi )\right]^N\right\}\xi ^{d_{\mathrm{}}1}𝑑\xi ,`$ (10)
where $`dV(\mathrm{})=V_0^{\mathrm{}}d_{\mathrm{}}\mathrm{}^{d_{\mathrm{}}1}d\mathrm{}`$ is the average number of fractal sites placed at a chemical distance between $`\mathrm{}`$ and $`\mathrm{}+d\mathrm{}`$, and $`\xi \mathrm{}/(\sqrt{2D_{\mathrm{}}}t^{1/d_w^{\mathrm{}}})`$. Here $`D_{\mathrm{}}`$ is the diffusion constant defined by the Einstein relation,
$$L^2=2D_{\mathrm{}}t^{2/d_w^{\mathrm{}}},$$
(11)
where $`L^2\mathrm{}^2`$ is the mean-square chemical distance traveled by a single random walker by time $`t`$ ($`t`$ large). Next, we assume that the asymptotic dependence of $`\mathrm{\Gamma }_t(\mathrm{})`$ for $`\xi 1`$ is given by
$$\mathrm{\Gamma }_t(\mathrm{})1A\xi ^{\mu v}e^{c\xi ^v}\left(1+\underset{n=1}{\overset{\mathrm{}}{}}h_n\xi ^{nv}\right),$$
(12)
with $`v=d_w^{\mathrm{}}/(d_w^{\mathrm{}}1)`$. This functional form holds (at least in its first terms) on Euclidean lattices and agrees with the expression conjectured in Refs. for fractal substrates with $`d_{\text{min}}=1`$ such as the Sierpinski gaskets. The dominant asymptotic behavior of the propagator in chemical space , $`P(l,t)\mathrm{exp}\left(c\xi ^v\right)`$, also coincides with the assumed dominant exponential decay of the mortality function $`1\mathrm{\Gamma }_t(\mathrm{})`$ in Eq. (12). It is known that both the propagator and the mortality function share the same asymptotic behavior for Euclidean lattices and for the Sierpinski lattice and we can expect that this coincidence also is the case for stochastic fractals (we will check this supposition in Sec. III). The rest of the analysis is identical with that carried out for Euclidean and deterministic fractal lattices save for the change of the Euclidean parameters to their chemical analogues. The result for $`S_N(t)`$ is, consequently,
$$S_N(t)\widehat{S}_N(t)\left(1\frac{d_{\mathrm{}}}{v}\underset{n=1}{\overset{\mathrm{}}{}}\mathrm{ln}^nN\underset{m=0}{\overset{n}{}}s_m^{(n)}\mathrm{ln}^m\mathrm{ln}N\right)$$
(13)
with
$$\widehat{S}_N(t)=V_0^{\mathrm{}}(2D_{\mathrm{}})^{d_{\mathrm{}}/2}t^{d_{\mathrm{}}/d_w^{\mathrm{}}}\left(\frac{\mathrm{ln}N}{c}\right)^{d_{\mathrm{}}/v}$$
(14)
and
$`s_0^{(1)}`$ $`=`$ $`\omega `$ (15)
$`s_1^{(1)}`$ $`=`$ $`\mu `$ (16)
$`s_0^{(2)}`$ $`=`$ $`(\beta 1)\left({\displaystyle \frac{\pi ^2}{12}}+{\displaystyle \frac{\omega ^2}{2}}\right)(ch_1\mu \omega )`$ (17)
$`s_1^{(2)}`$ $`=`$ $`\mu ^2+(\beta 1)\mu \omega `$ (18)
$`s_2^{(2)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\beta 1)\mu ^2.`$ (19)
Here $`\omega =\gamma +\mathrm{ln}A+\mu \mathrm{ln}c`$, $`\gamma 0.577215`$ is the Euler constant, and $`\beta =d_{\mathrm{}}/v`$. The dependence on $`t`$ and $`N`$ of the main term of $`S_N(t)`$ as given by Eq. (14), i.e., $`S_N(t)t^{d_s/2}(\mathrm{ln}N)^{d_{\mathrm{}}/v}`$, coincides with the prediction of Ref. .
### A simpler way to estimate the territory covered
We finish this section by showing how to find the full main term of Eq. (13) and even predict the form of the corrective terms by only resorting to extremely simple arguments. The ideas of the following reasoning were already used in Ref. . The crucial point in our argument is that, for a fixed time $`t`$, $`1[\mathrm{\Gamma }_t(\mathrm{})]^N`$ approaches a unit step function $`\mathrm{\Theta }(\mathrm{}\mathrm{}_\times )`$ when $`N\mathrm{}`$, $`\mathrm{}_\times `$ being the step’s width (see Fig. 1). The reason for this behavior is clear: For large $`N`$, $`[\mathrm{\Gamma }_t(\mathrm{})]^N`$ is only non-neglible when $`\mathrm{\Gamma }_t(\mathrm{})`$ is very close to 1. Obviously this occurs when the root-mean square chemical distance $`L(t)`$ traveled by the single random walker by time $`t`$ is small compared with $`\mathrm{}`$, i.e., when $`\xi =\mathrm{}/L(t)`$ is large. This in turn implies that in the evaluation of $`S_N(t)`$ only the behavior of $`\mathrm{\Gamma }_t(\mathrm{})`$ for large $`\xi `$ is relevant. Then, as $`1[\mathrm{\Gamma }_t(\mathrm{})]^N`$ approaches a step function of width $`\mathrm{}_\times `$, the integration of Eq. (9) yields
$$S_N(t)V_0^{\mathrm{}}\mathrm{}_\times ^d_{\mathrm{}},$$
(20)
i.e., the territory covered is just the volume of a chemical hypersphere of radius $`\mathrm{}_\times `$. Defining the width $`\mathrm{}_\times `$ of the step function as the distance at which $`1[\mathrm{\Gamma }_t(\mathrm{})]^N`$ takes the intermediate value $`1/2`$ (any other value between 0 and 1 would also be valid as $`\mathrm{}_\times `$ is not very sensitive to this value when $`N1`$), and assuming that $`1\mathrm{\Gamma }_t(\mathrm{})A\xi ^{\mu v}\mathrm{exp}(c\xi ^v)`$ for large $`\xi `$, we deduce that $`1/2NA\xi _\times ^{\mu v}\mathrm{exp}(c\xi _\times ^v)`$, with $`\xi _\times \mathrm{}_\times /L`$ so that
$$c\xi _\times ^v\mathrm{ln}N\mu v\mathrm{ln}\xi _\times +\mathrm{ln}2A.$$
(21)
The term $`\mathrm{ln}N`$ is dominant on the right-hand side of Eq. (21) for large $`N`$, so that a first-order solution of this equation is
$$c\xi _\times ^v\mathrm{ln}N,$$
(22)
i.e., $`\mathrm{}_\times ^vL^v\mathrm{ln}(N)/c`$. Hence Eq. (20) yields
$$S_N(t)V_0^{\mathrm{}}L^d_{\mathrm{}}\left(\frac{\mathrm{ln}N}{c}\right)^{d_{\mathrm{}}/v}$$
(23)
which is in full agreement with the main term of Eq. (13) when the Einstein relation, Eq. (11), is considered. Inserting the above first-order solution for $`\mathrm{}_\times `$ into the right-hand side of Eq. (21), we get the improved solution $`c\xi _\times ^v\mathrm{ln}N\mu \mathrm{ln}\mathrm{ln}N+\mathrm{ln}Ac^\mu +\mathrm{ln}2`$, so that Eq. (20) becomes
$`S_N(t)`$ $``$ $`V_0^{\mathrm{}}L^d_{\mathrm{}}\left({\displaystyle \frac{\mathrm{ln}N}{c}}\right)^{d_{\mathrm{}}/v}\left[1+{\displaystyle \frac{d_{\mathrm{}}}{v}}{\displaystyle \frac{\mu \mathrm{ln}\mathrm{ln}N+\mathrm{ln}Ac^\mu +\mathrm{ln}2}{\mathrm{ln}N}}\right].`$ (24)
This expression is strikingly close to the first-order approximation of Eq. (13), the only difference being that the term $`\mathrm{ln}2=0.693\mathrm{}`$ in Eq. (24) plays the role of the Euler constant $`\gamma 0.577215`$ in Eq. (13). Finally, notice that this simple method is not limited to disordered media but that it is also valid for estimating $`S_N(t)`$ for the non-disordered substrates (Euclidean and deterministic fractal media) considered in Refs. .
## III SURVIVAL PROBABILITY, FRACTAL VOLUME AND DIFFUSION IN PERCOLATION AGGREGATES
We have carried out simulations for the number of distinct sites visited by $`N`$ independent random walkers on a typical stochastic fractal: the percolation aggregate embedded in two dimensions. In our simulations every random walker makes a jump from a site to one of its nearest neighbors placed at one unit distance in each unit time. The incipient percolation aggregates were constructed by the standard Leath method on a square lattice with side $`400`$. In the Leath method, a seed is placed on the site in the center of this box and the cluster is generated by epidemic spreading to their nearest neighbors with an infecting probability $`p_c=0.5927460`$ corresponding to site percolation in the square lattice. At every step of the generation process a new chemical shell is added to the previous shell by the occupation with probability $`p_c`$ of its empty nearest neighbors. The process continues until we reach a shell whose sites do not infect any of their empty neighbors; in this case the aggregate so generated is rejected. If the cluster grown in this way spans the box from a side to the other, the cluster is accepted as a good representation of a portion of an infinite percolation aggregate. Our simulations were carried out over 2000 aggregates generated in this way.
In order to compare the simulation results for $`S_N(t)`$ with the predictions of our theoretical approach, Eq. (13), we must check that the survival probability or, equivalently, the mortality function, $`h(\mathrm{},t)=1\mathrm{\Gamma }_t(\mathrm{})`$, really behaves in the form conjectured in Eq. (12). Moreover, we must confirm first that, for a given chemical distance $`\mathrm{}`$, the distribution of $`h(\mathrm{},t)`$ over different realizations of the incipient percolation cluster is narrow because our theoretical analysis \[cf. Eq. (7)\] relies on this assumption \[cf. Eq. (6)\]. The numerical evaluation of this quantity as well as the propagator $`P(𝐢,t)`$ (i.e., the probability of finding a single random walker at site $`𝐢`$ at time $`t`$) is performed by the Chapman-Kolmogorov method (also called the exact enumeration method ). Initially the propagator takes the value $`1`$ at the origin site, $`P(\mathrm{𝟎},0)=1`$, and $`0`$ at any other site. This density evolves at every time step by updating its value at every site in the form described by the master equation
$$P(𝐢,t+1)=\frac{1}{b}\underset{\text{neighbor}=1}{\overset{M}{}}P(𝐣,t)+\left(1\frac{M}{b}\right)P(𝐢,t),$$
(25)
where $`b`$ is the maximum coordination number ($`b=4`$ for the two-dimensional percolation aggregate) and $`M`$ is the number of neighbors of site $`𝐢`$ that belong to the aggregate. This evolution equation is valid for the so-called blind “ants” because the random walker “attempts” a jump at time $`t+1`$ to a possible nearest neighbor (selected randomly) of the site it occupied at time $`t`$. If the selected site does not belong to the cluster, the random walker stays at the same site and no jump takes place. The “blindness” of the random walkers is taken into account by the second term of the right-hand side of Eq. (25). The trap is simulated by a special site belonging to the cluster that absorbs all the probability density that enters it without giving back any probability to its neighbors. In the simulation of the mortality function, we located a trap at a chemical distance $`\mathrm{}=30`$ in each of the 2000 percolating clusters. We repeated the experiment for traps located at a fixed Euclidean distance, $`r=80`$. The resulting histogram for $`t=1000`$ is shown in Fig. 2 (to be compared with the histogram of the propagator shown in Fig. $`4`$ of Ref. ). One observes that the distribution corresponding to fixed $`\mathrm{}`$ is very narrow whereas the Euclidean version is broad and exhibits a long tail. The striking contrast between the two histograms in Fig. 2 can be understood qualitatively. A fixed chemical distance $`\mathrm{}`$ between the origin and the trap means that there exists at least a minimum path connecting those sites whose length is $`\mathrm{}`$. Thus, the minimum time taken by a random walker to arrive at the trap site is $`t=\mathrm{}`$ independently of the lattice structure. On the other hand, a fixed Euclidean distance $`r`$ could correspond to many different chemical distances from one cluster to another. This is a consequence of the stochastic lacunarity of the fractal aggregate. A large hole between the origin and the trap implies a large minimum time to travel around the border of that hole and, obviously, a small mortality in comparison with another cluster where no holes hinder the diffusion.
Figure 3 shows the chemical mean-square displacement
$$L^2\mathrm{}^2=\underset{\text{sites}}{}\mathrm{}^2P(\mathrm{},t)$$
(26)
as a function of time. The propagator in the chemical space $`P(\mathrm{},t)`$ is obtained by summing $`P(𝐢,t)`$ over all cluster sites $`𝐢`$ on the chemical shell situated at distance $`\mathrm{}`$ from the origin. The result is compatible with the Einstein relation Eq. (11) with $`2D_{\mathrm{}}=1.20\pm 0.1`$ and $`d_w^{\mathrm{}}=2.40\pm 0.05`$. This value for $`d_w^{\mathrm{}}`$ coincides with that obtained in Ref. and is in agreement with the value reported in Refs. .
In Fig. 4 we plot $`\mathrm{ln}(\mathrm{ln}h(\mathrm{},t))`$ versus $`\widehat{\xi }\mathrm{}/t^{1/d_w^{\mathrm{}}}`$ with $`\mathrm{}=80`$ and, according to the previous discussion, $`d_w^{\mathrm{}}=2.40`$. If the conjecture in Eq. (12) is right, we can take $`h(\mathrm{},t)\mathrm{exp}(\widehat{c}\widehat{\xi }^v)`$ as a first approximation, and hence should observe the linear behavior $`\mathrm{ln}(\mathrm{ln}h(\mathrm{},t))\mathrm{ln}\widehat{c}+v\widehat{\xi }`$ with $`\widehat{\xi }=\sqrt{2D_{\mathrm{}}}\xi `$ and $`\widehat{c}=c/(2D_{\mathrm{}})^{v/2}`$. Certainly the plot seems linear except for a portion in the range $`\widehat{\xi }2.2`$. This is a finite size effect (already analyzed in the case of the two-dimensional Sierpinski gasket in Ref. ) associated with the existence of a minimum arrival time corresponding to a random walker who travels “ballistically” along a chemical path from the origin to the trap, which in turn implies a maximum available value of $`\widehat{\xi }`$ in the simulations (in our simulations this maximum value is $`80/80^{1/d_w^{\mathrm{}}}12.9`$). A reliable interval for numerical fits should exclude this very short time regime. A linear fit in the interval $`1.6\mathrm{ln}\widehat{\xi }2.17`$, corresponding to $`200t1000`$, gives the values $`\widehat{c}=1.2\pm 0.1`$, i.e., $`c=1.3\pm 0.1`$, and $`v=1.6\pm 0.05`$. The dashed line in Fig. 4 corresponds to these values. The good agreement with numerical values in the above interval seems to assure the correctness of the approximation $`h(\mathrm{},t)h_\text{a}(\mathrm{},t)=\mathrm{exp}(\widehat{c}\widehat{\xi }^v)`$ with the values of $`c`$ and $`v`$ given above. However, the solid line in Fig. 4 is a challenge to this interpretation: one sees that the function $`h_\text{b}(\mathrm{},t)=\widehat{A}\widehat{\xi }^{\mu v}\mathrm{exp}(\widehat{c}\widehat{\xi }^v)`$ with $`v=1.70`$, $`\widehat{c}=0.9`$ (i.e. $`c=1.05`$), $`\mu =0.8`$ and $`\widehat{A}=1.1`$ is as good as $`h_\text{a}(\mathrm{},t)`$. Indeed, $`h_\text{b}(\mathrm{},t)`$ is more consistent from a theoretical point of view than $`h_\text{a}(\mathrm{},t)`$ because the expected theoretical value of $`v`$ corresponding to $`d_w^{\mathrm{}}=2.40`$ is $`v=d_w^{\mathrm{}}/(d_w^{\mathrm{}}1)=1.71`$, which is in better agreement with the exponent $`v=1.7`$ of $`h_\text{a}(\mathrm{},t)`$ than with the exponent $`v=1.6`$ of $`h_\text{b}(\mathrm{},t)`$. Finally, it should be noticed that the values $`c=1.05`$, $`v=1.70`$ are also in agreement with the corresponding parameter values of the propagator , thus supporting the guess made in Sec. II \[see below Eq. (12)\] that the dominant exponential term of the propagator and of the mortality function are the same. This leads us to consider that the set of parameters $`\widehat{c}=0.9,v=1.7,\mu =0.8`$ is more reliable than $`\widehat{c}=1.2,v=1.6,\mu =0`$. Obviously, further intensive (and extremely time consuming) computer simulations for the mortality function would be required in order to reliably determine the values of the parameters that appear in Eq. (12) and in the asymptotic corrections of $`S_N(t)`$.
We have also evaluated numerically the fractal volume in terms of the chemical distance $`V(\mathrm{})`$, i.e., the number of lattice sites inside a circumference (in chemical space) of radius $`\mathrm{}`$. The results are shown in Fig. 5. A good fit to $`dV(\mathrm{})=d_{\mathrm{}}V_0^{\mathrm{}}\mathrm{}^{d_{\mathrm{}}1}d\mathrm{}`$ is found with $`V_0^{\mathrm{}}=1.1\pm 0.2`$ and $`d_{\mathrm{}}=1.65\pm 0.05`$. Taking into account that $`d_f=91/48`$, we deduce that $`d_{\text{min}}=d_f/d_{\mathrm{}}=1.15\pm 0.05`$, which agrees with previous estimates .
## IV SIMULATION RESULTS: TERRITORY COVERED BY $`N`$ RANDOM WALKERS ON THE PERCOLATION AGGREGATE
As discussed in Sec. I, the average implicit in the evaluation of $`S_N(t)`$ is double: first, we take an average over experiments performed on the same aggregate and then a second average over different percolation clusters. The $`N`$ random walkers are always placed initially upon the site in the center of the square box. In our simulations we have performed an average over $`100`$ runs per cluster and a second average over $`2000`$ percolation clusters in order to achieve good statistics. The maximum time considered was $`t=1000`$.
According to Eqs. (13) and (14), the quotient $`S_N(t)/(\mathrm{ln}N)^\gamma `$ with $`\gamma =d_{\mathrm{}}/v`$ is only a function of $`t`$. In Fig. 6 the logarithm of that quotient is plotted versus $`\mathrm{ln}t`$ for several values of $`N`$. The data collapse and the slope close to $`0.66`$ seems to support Eq. (8) with $`\gamma =d_{\mathrm{}}/v=0.97`$, which is in agreement with similar recent results for the three-dimensional percolation aggregate . The collapse is, however, slightly poorer when the exponent $`\gamma =d_f/u=1.24`$ ($`d_f=91/48`$ and $`d_w=2.87`$ ) proposed by Havlin et al. \[see Eq. (4)\] is used, as Fig. 6 shows. So, one migth be led to the conclusion that the correct value of $`\gamma `$ as defined above is $`d_{\mathrm{}}/v`$. But in this analysis there was no consideration of the relatively large logarithmic corrections predicted by the asymptotic analysis presented in Section II, so that the reliability of the above conclusion is seriously affected by this omission.
To illustrate this point, let us now carry out the same kind of analysis with the simulation results of $`S_N(t)`$ when the substrate is a three-dimensional Euclidean lattice. For this case it is well known that $`S_N(t)`$ is given by an asymptotic expression with the form of Eq. (13) in which the logarithmic corrective terms are very important even for very large values of $`N`$. Indeed, the main asymptotic term leads to very poor predictions for $`S_N(t)`$, whereas the second-order approximation ($`n=2`$) gives excellent agreement with numerical simulation results. The exponent $`\gamma `$ of the main logarithmic term in $`N`$ and the time exponent $`d_{\mathrm{}}/d_w^{\mathrm{}}=d_s/2`$ are equal to $`3/2`$. We have plotted in Fig. 7 the quotient $`S_N(t)/\mathrm{ln}^\gamma N`$ versus $`\mathrm{ln}t`$ for several values of $`N`$ taking into account that the rigorous value of $`\gamma `$ is $`3/2`$. We see that the collapse is far from being perfect because the logarithmic corrections have been ignored. Nevertheless, an effective (but incorrect!) value of $`\gamma =2.75`$ yields a much better data collapse and a slope close to the theoretical value $`d_s/2=1.5`$. We thus conclude that analysis of data collapse plots based on the form of the main term of quantities such as $`S_N(t)`$ (which typically exhibit large corrective terms) should be performed with caution. The values of the exponents estimated in this way are untrustworthy because the existence of logarithmic corrections to the main term cannot simply be ignored. The value of $`\gamma =2.75`$ obtained before is then only an effective way of including all these corrective terms together but the true expression involves a main term of the form $`(t\mathrm{ln}N)^{3/2}`$ times a series similar to that given in Eq. (13). These considerations should prevent us from drawing hasty conclusions from a simple view of plots such as Figs. 6 and 7.
Finally, in Fig. 8 we show the dependence of $`S_N(t)`$ on $`N`$ and compare simulation results with the zeroth- and first-order asymptotic prediction given by Eq. (13). When the parameter set $`\widehat{c}=0.9,v=1.7,\mu =0.8,A=1`$ (see Sec. III) is used, we get results with a very familiar aspect as they are quite similar (although, perhaps the first-order approximation is too good) to that already found for Euclidean and Sierpinski lattices . This is indeed encouraging. However, when the parameters $`\widehat{c}=1.3`$ and $`v=1.6`$ are used, we obtain a surprising and strikingly accurate zeroth-order approximation. At this point, we again suspect that this last set of parameters are only effective parameters that include the influence of the true logarithmic corrective terms in the range of $`N`$ simulated. Hence, Fig. 8 illustrates again, but from a different perspective, how the omission of important corrective terms could lead to finding effective parameters that, although providing excellent approximations in the (relatively short) range under consideration, are really erroneous.
## V Summary
In this paper, the average fractal territory covered up to time $`t`$ by $`N`$ independent random walkers all starting from the same origin on stochastic fractal lattices is calculated in terms of an asymptotic series expansion, $`_{n=0}^{\mathrm{}}_{m=0}^ns_{nm}(\mathrm{ln}N)^{d_{\mathrm{}}/vn}(\mathrm{ln}\mathrm{ln}N)^m`$ \[see Eq. (13)\], which is formally identical to those obtained for Euclidean and deterministic fractal lattices. Equation (13) is obtained by assuming that (i) the average fractal volume inside a “hypersphere” of chemical radius $`r`$ grows as $`V_0^{\mathrm{}}r^d_{\mathrm{}}`$, (ii) the distribution of the the short-time survival probability of a single random walker in the presence of a trap is narrow, so that Eq. (6) holds, and (iii) this short-time survival probability is asymptotically given by Eq. (12). We performed numerical simulations for the two-dimensional percolation aggregate at criticality which support the validity of the above assumptions. The zeroth- and first-order theoretical asymptotic expression of $`S_N(t)`$ were calculated explicitly and, for the two-dimensional percolation aggregate, they compared reasonably well with numerical simulation results. The agreement is similar to that found for non-disordered media.
In our procedure, the use of the chemical distance turns out to be fundamental because the distribution of the short-time survival probability in the chemical space is so narrow that we can safely replace the power $`N`$ of the mean value of the survival probability by the mean value of the power $`N`$ of the survival probability. (On the contrary, this does not hold at all when Euclidean distances are used.) This allowed us to easily translate the theoretical results previously derived for Euclidean and deterministic fractals to disordered media.
###### Acknowledgements.
This work has been supported by the DGICYT (Spain) through Grant No. PB97-1501 and by the Junta de Extremadura-Fondo Social Europeo through Grant No. IPR99C031.
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# Self force on particle in orbit around a black hole
\[
## Abstract
We study the self force acting on a scalar charge in uniform circular motion around a Schwarzschild black hole. The analysis is based on a direct calculation of the self force via mode decomposition, and on a regularization procedure based on Ori’s mode-sum regularization prescription. We find the four self-force at arbitrary radii and angular velocities (both geodesic and non-geodesic), in particular near the black hole, where general-relativistic effects are strongest, and for fast motion. We find the radial component of the self force to be repulsive or attractive, depending on the orbit.
PACS number(s): 04.25-g, 04.70.-s, 04.70.Bw
\]
The problem of finding the equations of motion for a particle in curved spacetime has become recently extremely important, as the first generation of interferometric gravitational wave detectors will soon be operational, and with the prospects of having a gravitational wave space antenna in the not-very-distant future. The generation of very accurate templates for the waveforms detected from a system of a compact object in orbit around a supermassive black hole is an extremely hard task. It is likely that one would need to have accurate templates for as many as $`5\times 10^5`$ orbits. For such a system, accurate templates are necessary for detection, because the predicted signal-to-noise ratio for LISA is approximately of order $`10`$ for a $`1`$ year integration time. Lack of accurate templates would result in a loss of a factor of roughly the square root of the number of orbits in sensitivity , which would result in signal-to-noise ratio below the detectibility threshold.
In order to generate accurate templates, an important necessary ingredient is the inclusion of radiation reaction in the orbital evolution of the compact object. The radiation-reaction forces need to be calculated locally, i.e., in the neighborhood of the compact object. In the conventional approach, one calculates, at infinity and at the event horizon of the black hole, the fluxes of quantities, which are constants of motion in the absence of radiation reaction. Then, one uses a balance argument to relate these fluxes to the rate of change of a corresponding local quantity of the compact object. This approach generally fails because of the inadditivity of the Carter constant. For very simple cases, e.g. for circular or equatorial orbits around a Kerr black hole, the evolution of the Carter constant is trivial, such that the conventional approach is useful . However, generic orbits around a rotating black hole are neither circular nor equatorial, and consequently a new approach, which is not based on balance arguments, is of great need.
Several approaches have been suggested for the calculation of self forces. Quinn and Wald and Mino, Sasaki, and Tanaka recently proposed general approaches for the calculation of self forces. However, it is not presently clear how to apply these approaches directly for actual computations, the greatest problem being the calculation of the so-called “tail” term of the Green’s function, which arises from the failure of the Huygens principle in curved spacetime. (In addition, there are in general also local, Ricci-curvature coupled and Abraham-Lorentz-Dirac (ALD) type terms , which are much easier to calculate, the former vanishing identically in vacuum.)
Recently, Ori proposed a local, causal approach for the calculation of the self forces , which is based on decomposition of the self force into modes, and on a mode-sum regularization prescription (MSRP). Although MSRP is not fully proven as yet, it has already been shown to be valid for simple cases, such as scalar charges in general orbits in Schwarzschild spacetime, and, in particular, for circular orbits which we consider here. MSRP is likely to be susceptible of generalization also for massive particles in orbit around a Kerr black hole. If robust, MSRP can be of great importance for the generation of accurate templates. We hope that MSRP can be combined with other approaches, which were recently proposed, such as mode decomposition of the self forces which are sourced by just the distant past world line or a normal-neighborhood expansion .
We first describe very briefly the main ideas of MSRP . Then, we apply MSRP for the case of a scalar particle in circular orbit around a Schwarzschild black hole, and calculate the self four-force acting on the particle linearized in its own self field.
The contribution to the physical self force from the tail part of the Green’s function can be decomposed into stationary Teukolsky modes, and then summed over the frequencies $`\omega `$ and the azimuthal numbers $`m`$. The self force equals then the limit $`ϵ0^{}`$ of the sum over all $`\mathrm{}`$ modes, of the difference between the force sourced by the entire world line (the bare force $`{}_{}{}^{\mathrm{bare}}F_{\mu }^{\mathrm{}}`$) and the force sourced by the half-infinite world line to the future of $`ϵ`$, where the particle has proper time $`\tau =0`$, and $`ϵ`$ is an event along the past ($`\tau <0`$) world line. Next, we seek a regularization function $`h_\mu ^{\mathrm{}}`$ which is independent of $`ϵ`$, such that the series $`_{\mathrm{}}({}_{}{}^{\mathrm{bare}}F_{\mu }^{\mathrm{}}h_\mu ^{\mathrm{}})`$ converges. Once such a function is found, the regularized self force is then given by $`{}_{}{}^{\mathrm{ren}}F_{\mu }^{}=_{\mathrm{}}({}_{}{}^{\mathrm{bare}}F_{\mu }^{\mathrm{}}h_\mu ^{\mathrm{}})d_\mu `$, where $`d_\mu `$ is a finite valued function. MSRP then shows, from a local integration of the Green’s function, that the regularization function $`h_\mu ^{\mathrm{}}=a_\mu \mathrm{}+b_\mu +c_\mu \mathrm{}^1`$, and for the case of a scalar charge in circular orbit around a Schwarzschild black hole MSRP yields the values of the functions $`a_\mu ,b_\mu ,c_\mu `$ and $`d_\mu `$ analytically. In particular, it can be shown that for such orbits $`a_\mu =0=c_\mu `$ and $`d_r=0`$, such that in practice the regularization prescription of the radial force is reduced to subtracting $`b_r`$ from each $`\mathrm{}`$ mode of the bare radial force. Note that $`b_\mu `$ is just the limit $`\mathrm{}\mathrm{}`$ of $`{}_{}{}^{\mathrm{bare}}F_{\mu }^{\mathrm{}}`$.
In the following we describe the results obtained from this new approach for the case of a point-like scalar charge in circular orbit around a Schwarzschild black hole. Our results are fully relativistic, i.e., we do not introduce any simplifying assumptions such as far field or slow motion. Because of the fully relativistic nature of this study, our analysis is numerical. However, it is reasonable to expect that analytical solutions will not be available in general, except, possibly, only for very simple cases, such as static configurations . We calculate the contribution to the force which the scalar charge feels, due to its own field, to leading order in the particle’s charge. We use spherical Regge-Wheeler coordinates, for which the Schwarzschild metric is $`ds^2=\left(1\frac{2M}{r}\right)\left(dt^2+dr_{}^{}{}_{}{}^{2}\right)+r^2\left(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2\right)`$, where $`M`$ is the black hole’s mass, and the radial Schwarzschild coordinate $`r(r^{})`$ is given implicitly by $`r^{}=r+2M\mathrm{ln}\left|r/(2M)1\right|`$. The field satisfies the wave equation $`_\mu ^\mu \mathrm{\Phi }(x^\alpha )=4\pi \rho (x^\alpha )`$, where the charge density $`\rho (x^\alpha )=q_{\mathrm{}}^{\mathrm{}}𝑑\tau \delta ^4[x^\alpha x_0^\alpha (\tau )](g)^{1/2}`$. Here, $`q`$ is the charge of the particle, whose world line is $`x_0^\alpha (\tau )`$, $`\tau `$ being its proper time, $`g`$ being the metric determinant, and $`_\mu `$ denotes covariant differentiation. We take the charge to be in circular orbit at $`r^{}=r_0^{}`$, $`\theta =\pi /2`$, and $`d\varphi /dt=\mathrm{\Omega }`$. (We are not restricted to the Keplerian angular velocity $`\mathrm{\Omega }_K`$.) We next decompose the scalar field $`\mathrm{\Phi }`$ into modes according to $`\mathrm{\Phi }=_{\mathrm{}}^{\mathrm{}}𝑑\omega _\mathrm{}me^{i\omega t}\mathrm{\Psi }_\mathrm{}m(r^{})Y^\mathrm{}m(\theta ,\varphi )/r(r^{})`$, such that the equation which governs the field $`\mathrm{\Psi }_\mathrm{}m`$ becomes
$`{\displaystyle \frac{d^2\mathrm{\Psi }_\mathrm{}m}{dr_{}^{}{}_{}{}^{2}}}`$ $`+`$ $`\left\{\omega ^2V_{\mathrm{}}[r(r^{})]\right\}\mathrm{\Psi }_\mathrm{}m=4\pi {\displaystyle \frac{q}{\gamma }}`$ (1)
$`\times `$ $`{\displaystyle \frac{\delta (rr_0)}{r_0}}\delta (\omega m\mathrm{\Omega })Y^\mathrm{}m(\pi /2,\varphi )e^{im\varphi }.`$ (2)
This equation should be solved for each mode $`\mathrm{}m`$ with boundary conditions of ingoing waves at the event horizon $`(r^{}\mathrm{})`$, and outgoing waves at infinity $`(r^{}\mathrm{})`$. The effective potential is given by $`V_{\mathrm{}}(r)=(12M/r)[2M/r^3+\mathrm{}(\mathrm{}+1)/r^2]`$, and $`\gamma =1/\sqrt{12M/rr^2\mathrm{\Omega }^2}`$. The contribution of the $`\mathrm{}m`$ mode to the force is given by $`{}_{}{}^{\mathrm{bare}}F_{\mu }^{\mathrm{}m}=q(\mathrm{\Phi }_{,\mu }^\mathrm{}m+u_\mu u^\nu \mathrm{\Phi }_{,\nu }^\mathrm{}m)`$ (note that for circular orbits $`u^\nu \mathrm{\Phi }_{,\nu }^\mathrm{}m=0`$).
We next find numerically the solutions $`\mathrm{\Psi }_\mathrm{}m^+(\mathrm{\Psi }_\mathrm{}m^{})`$ for the homogeneous equations corresponding to Eq. (2), which satisfy the boundary condition at infinity (the horizon). The components of the force are then given by
$`{}_{}{}^{\mathrm{bare}}F_{r^{}}^{\mathrm{}m}`$ $`=`$ $`2\pi q^2{\displaystyle \frac{\left|Y^\mathrm{}m(\frac{\pi }{2},0)\right|^2}{\gamma r_0^2}}\{\mathrm{Re}[W_\mathrm{}m^1(r_0^{})]\mathrm{Re}[S_\mathrm{}m(r_0^{})]`$ (3)
$`+`$ $`{\displaystyle \frac{2}{r_0}}\left(1{\displaystyle \frac{2M}{r_0}}\right)\mathrm{Re}[W_\mathrm{}m^1(r_0^{})]\mathrm{Re}[T_\mathrm{}m(r_0^{})]`$ (4)
$``$ $`{\displaystyle \frac{2}{r_0}}\left(1{\displaystyle \frac{2M}{r_0}}\right)\mathrm{Im}[W_\mathrm{}m^1(r_0^{})]\mathrm{Im}[T_\mathrm{}m(r_0^{})]`$ (5)
$`+`$ $`\mathrm{Im}[W_\mathrm{}m^1(r_0^{})]\mathrm{Im}[S_\mathrm{}m(r_0^{})]\},`$ (6)
where $`W_\mathrm{}m`$ is the Wronskian determinant of $`\mathrm{\Psi }_\mathrm{}m^{}(r^{})`$ and $`\mathrm{\Psi }_\mathrm{}m^+(r^{})`$, $`T_\mathrm{}m(r^{})=\mathrm{\Psi }_\mathrm{}m^+(r^{})\mathrm{\Psi }_\mathrm{}m^{}(r^{})`$, and $`S_\mathrm{}m(r^{})=\mathrm{\Psi }_\mathrm{}m^+(r^{})\mathrm{\Psi }_{\mathrm{}m,r^{}}^{}(r^{})+\mathrm{\Psi }_\mathrm{}m^{}(r^{})\mathrm{\Psi }_{\mathrm{}m,r^{}}^+(r^{})`$. We find that
$`{}_{}{}^{\mathrm{bare}}F_{t}^{\mathrm{}m}`$ $`=`$ $`4\pi q^2m\mathrm{\Omega }{\displaystyle \frac{\left|Y^\mathrm{}m(\frac{\pi }{2},0)\right|^2}{\gamma r_0^2}}\{\mathrm{Im}[W_\mathrm{}m^1(r_0^{})]`$ (7)
$`\times `$ $`\mathrm{Re}[T_\mathrm{}m(r_0^{})]+\mathrm{Re}[W_\mathrm{}m^1(r_0^{})]\mathrm{Im}[T_\mathrm{}m(r_0^{})]\}.`$ (8)
We also obtain $`{}_{}{}^{\mathrm{bare}}F_{\varphi }^{\mathrm{}m}=\mathrm{\Omega }^1{}_{}{}^{\mathrm{bare}}F_{t}^{\mathrm{}m}`$ and $`{}_{}{}^{\mathrm{bare}}F_{\theta }^{\mathrm{}m}=0`$. It is convenient to define new radial functions $`Z^{\pm \mathrm{}m}(r^{})`$ by $`\mathrm{\Psi }^{\pm \mathrm{}m}(r^{})=e^{\pm i\omega r^{}}Z^{\pm \mathrm{}m}(r^{})`$, which satisfy the homogeneous equations
$$\frac{d^2Z^{\pm \mathrm{}m}}{dr_{}^{}{}_{}{}^{2}}2i\omega \frac{dZ^{\pm \mathrm{}m}}{dr^{}}V_{\mathrm{}}[r^{}(r)]Z^{\pm \mathrm{}m}=0,$$
(9)
with boundary conditions $`Z^{+\mathrm{}m}(r^{}M)=1+a_1^+f_++a_2^+f_+^2+O(f_+^3)`$ and $`Z^\mathrm{}m(r^{}M)=1+a_1^{}f_{}+a_2^{}f_{}^2+O(f_{}^3)`$, where $`f_+=(\omega r)^1`$, $`f_{}=12M/r`$, and $`a_1^+=i\mathrm{}(\mathrm{}+1)/2`$, $`a_2^+=[\mathrm{}(\mathrm{}1)(\mathrm{}+1)(\mathrm{}+2)+4i\omega M]/8`$, $`a_1^{}=[1+\mathrm{}(\mathrm{}+1)]/(14i\omega M)`$, and $`a_2^{}=[\mathrm{}(\mathrm{}+1)(\mathrm{}^2+\mathrm{}+6)+4i\omega M+4]/[4(12i\omega M)(14i\omega M)]`$. We solve Eqs. (9) numerically using both Burlisch-Stoer and fourth-order Runge-Kutta integrations with adaptive step-size controls. Both integrators yield results compatible within the error limits. We place the exterior and interior boundaries at a distance of several mode wavelengths (in $`r^{}`$) from $`r_0^{}`$, and then use successive Richardson extrapolations, with increasing distance to the boundaries, until the extrapolation of the boundaries to $`r^{}\pm \mathrm{}`$ yields an error smaller than a given threshold. Notice that for modes with $`m=0`$ the wavelength is infinite, such that the boundaries cannot be taken far enough from the charge. Indeed, we find that for this case the Richardson extrapolations do not converge. Instead, we can solve for this case analytically. We find that $`f_t^{\mathrm{},m=0}=0`$ and $`f_r^{\mathrm{},m=0}=(2\pi /\gamma )(q/M)^2Y_{}^{\mathrm{}0}{}_{}{}^{2}(\pi /2,0)[2Q_{\mathrm{}}(\rho _0)dP_{\mathrm{}}(\rho _0)/d\rho +1/(1\rho _0^2)]`$, where $`\rho =(rM)/M`$. Here, $`P_{\mathrm{}},Q_{\mathrm{}}`$ are the Legendre functions of the first and second kinds, respectively. Figure 1 shows the functions $`Z_{\mathrm{}=1m=1}^\pm (r^{})`$. Similar qualitative behavior is found also for the other modes. Until the peak of the effective potential barrier the functions $`Z^{\pm \mathrm{}m}`$ vary only slowly, and then start oscillating rapidly.
The temporal component of the bare force is finite. (In the regularization scheme this corresponds to $`b_t=0`$.) MSRP predicts that $`d_t`$ exactly balances the ALD force, such that the full self force is given only by $`F_t=_{\mathrm{}}{}_{}{}^{\mathrm{bare}}F_{t}^{\mathrm{}}`$. We compare our results with their Minkowski spacetime counterparts. In flat spacetime one can solve analytically for each mode, sum over all modes, and find that $`F_t^{\mathrm{Min}}=\frac{1}{3}q^2\mathrm{\Omega }^2r_0^2\gamma _{\mathrm{Min}}^5`$. Here, $`\gamma _{\mathrm{Min}}`$ is the usual flat spacetime Lorentz factor. For the comparison we choose the same values of $`r_0,\mathrm{\Omega }`$ for the curved and flat spacetimes. Figure 2 displays $`F_t`$ as a function of $`r/M`$ for two cases: (2A) Non-geodesic circular orbits, with a fixed angular velocity $`\mathrm{\Omega }`$. (In this case the tangential velocity increases linearly with $`r`$.) When $`r/M`$ is large the value of $`F_t`$ approaches its flat spacetime counterpart. Recall that $`\mathrm{\Omega }=d\varphi /dt`$, where $`t`$ is the time of an observer at infinity. Because of the red-shift effect at small values of $`r`$, orbits with the same value of $`\mathrm{\Omega }`$ have very large proper tangential velocities at small $`r/M`$: a fixed $`\mathrm{\Omega }`$ corresponds to the ultra-relativistic limit when the orbit is close to the black hole. The second case is (2B) geodesic motion, which satisfies Kepler’s law $`\mathrm{\Omega }_K^2r^3=M`$. The innermost (unstable) causal orbit is located at $`r=3M`$. Approaching $`r=3M`$, the motion of the particle approaches the ultra-relativistic limit, which is manifested by the rapid growth of $`F_t`$. At larger radii the value of $`F_t`$ approaches the flat spacetime counterpart. This is detailed in Fig. 2(C).
Next, we study the radial, conservative component of the self force. First, we check the agreement of our numerical results with MSRP. MSRP predicts that $`a_r=0=c_r`$ and that
$`b_r=`$ $``$ $`{\displaystyle \frac{q^2}{2r^2}}{\displaystyle \frac{1}{\gamma \sqrt{g_{tt}}}}[2{}_{2}{}^{}F_{1}^{}({\displaystyle \frac{1}{2}},{\displaystyle \frac{1}{2}};1;{\displaystyle \frac{r^2\mathrm{\Omega }^2}{12M/r}})`$ (10)
$``$ $`{\displaystyle \frac{13M/r}{12M/r}}{}_{2}{}^{}F_{1}^{}({\displaystyle \frac{1}{2}},{\displaystyle \frac{3}{2}};1;{\displaystyle \frac{r^2\mathrm{\Omega }^2}{12M/r}})].`$ (11)
We check the accuracy of our numerical determination of the value of $`b_r`$ by comparison to Eq. (11). This check serves the two-fold purpose of (i) checking the numerical code, and (ii) checking the compatibility of the analytical prediction of MSRP for the regularization function $`h_r^{\mathrm{}}`$ with the numerical determination of the bare force. Figure 3(A) shows the behavior of $`{}_{}{}^{\mathrm{bare}}F_{r}^{\mathrm{}}`$, after summation over $`m`$ and $`\omega `$, as a function of $`\mathrm{}`$, and Fig. 3(B) displays $`{}_{}{}^{\mathrm{bare}}F_{r}^{\mathrm{}}b_r`$ as a function of $`\mathrm{}`$. This difference decreases like $`\mathrm{}^2`$ for large values of $`\mathrm{}`$, which confirms the predictions of MSRP for the values of $`a_r`$, $`b_r`$, and $`c_r`$. (We emphasize that we cannot test numerically the prediction of MSRP that $`d_r=0`$.) Note that the radial ALD force vanishes, such that $`{}_{}{}^{\mathrm{ren}}F_{r}^{}`$ is the full self force. Similar behavior is found also for Keplerian orbits at other values of $`r_0^{}`$, and also for non-geodesic circular motion, with angular velocities both greater or smaller than $`\mathrm{\Omega }_K`$.
The regularized component of the radial self-force $`{}_{}{}^{\mathrm{tail}}F_{r}^{\mathrm{}}`$ is obtained by subtracting $`b_r{}_{}{}^{\mathrm{bare}}F_{r}^{\mathrm{}\mathrm{}}`$ from each mode $`{}_{}{}^{\mathrm{bare}}F_{r}^{\mathrm{}}`$. The total regularized force is then obtained by
$$F_r^{\mathrm{ren}}=\underset{n=0}{\overset{\mathrm{}}{}}{}_{}{}^{\mathrm{tail}}F_{r}^{n}+_r^{\mathrm{}+1},$$
(12)
where the remainder $`_r^{\mathrm{}+1}`$ is given approximately by $`_r^{\mathrm{}+1}\mathrm{}^2({}_{}{}^{\mathrm{bare}}F_{r}^{\mathrm{}}b_r)\psi ^{(1)}(\mathrm{}+1)`$ for sufficiently large values of $`\mathrm{}`$, $`\psi ^{(1)}(x)d^2\mathrm{ln}\mathrm{\Gamma }(x)/dx^2`$ being the trigamma function. Note that for large arguments $`\psi ^{(1)}(x)x^1`$. Figure 4 displays the regularized radial self-force for Keplerian (4A) and non-Keplerian (4B) orbits. The radial self force in the far field limit ($`rM`$) is repulsive, and satisfies $`F_r^{\mathrm{ren}}\alpha q^2M^3r^5`$, where $`\alpha `$ is a dimensionless parameter of order unity. A minimum-$`\chi ^2`$ fit shows the exponent of $`r`$ to equal $`5`$ and found $`\alpha `$ to equal unity, both with $`3\%`$ errors. However, in the strong field, the force law deviates from this simple relation, and grows faster. For the non-Keplerian orbits, we find that in the slow motion limit ($`\mathrm{\Omega }\mathrm{\Omega }_K`$) the radial force is proportional to $`\mathrm{\Omega }^2`$. The exponent of $`\mathrm{\Omega }`$ is found to be $`2`$ with a $`3\%`$ error. Combined with the result for Keplerian orbits, we find that for any circular orbit, in the far field and slow motion limits, the radial force is repulsive, and is given by
$$F_r^{\mathrm{ren}}\alpha q^2(G^3/c^6)M^2\mathrm{\Omega }^2/r^2.$$
(13)
This results explains the vanishing self force in the static limit . For faster motion the $`\mathrm{\Omega }^2`$ law does not hold any more. In fact, for $`\mathrm{\Omega }>\mathrm{\Omega }_K`$ we find that the radial self force varies rapidly, and eventually changes from repulsive to attractive. The radial self force does not cause a net change in the energy of the particle. However, if the orbit has a non-zero eccentricity, this force induces an additional precession of the periastron, which in the slow motion or far field limits is retrograde. This precession has an effect on the frequencies of the emitted radiation.
Our results show that the self force can be calculated for a simple, although non-trivial, problem. MSRP was found to be useful also for other cases, e.g., static scalar and electric charges in Schwarzschild , for scalar and electric charges in circular motion in flat spacetime , and for general radial motion of scalar charges in spherical symmetry . A closely related approach was used also for a mass point in radial free fall in Schwarzschild . We hope that similar methods can be used for more realistic cases, which may be relevant for the orbital evolution of compact objects around black holes.
I thank Scott Hughes and Kip Thorne for valuable discussions. I am indebted to Amos Ori for many stimulating discussions and for letting me use his results before their publication. This research was supported by NSF grants AST-9731698 and PHY-9900776, and by NASA grant NAG5-6840.
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# 1 Introduction
## 1 Introduction
In spite of the fact that the top quark has been discovered already several years ago its interactions are still unknown. Therefore it remains an open question if the top-quark couplings obey the Standard Model (SM) scheme of the electroweak forces or there exists a contribution from physics beyond the SM. In this talk<sup>♯1</sup><sup>♯1</sup>♯1Presented by B. Grzadkowski I will try to use angular and energy distributions of top-quark decay products in the process $`e^+e^{}t\overline{t}\mathrm{}^\pm /\stackrel{()}{b}\mathrm{}`$ in order to estimate how precisely top quark couplings could be determined at future linear collider.
We will parameterize $`t\overline{t}`$ couplings to the photon and the $`Z`$ boson in the following way
$$\mathrm{\Gamma }_{vt\overline{t}}^\mu =\frac{g}{2}\overline{u}(p_t)\left[\gamma ^\mu \{A_v+\delta A_v(B_v+\delta B_v)\gamma _5\}+\frac{(p_tp_{\overline{t}})^\mu }{2m_t}(\delta C_v\delta D_v\gamma _5)\right]v(p_{\overline{t}}),$$
(1)
where $`g`$ denotes the $`SU(2)`$ gauge coupling constant, $`v=\gamma ,Z`$, and
$$A_\gamma =\frac{4}{3}\mathrm{sin}\theta _W,B_\gamma =0,A_Z=\frac{1}{2\mathrm{cos}\theta _W}\left(1\frac{8}{3}\mathrm{sin}^2\theta _W\right),B_Z=\frac{1}{2\mathrm{cos}\theta _W}$$
denote the SM contributions to the vertices. Among the above non-SM form factors, $`\delta A_{\gamma ,Z}`$, $`\delta B_{\gamma ,Z}`$, $`\delta C_{\gamma ,Z}`$ describe $`CP`$-conserving while $`\delta D_{\gamma ,Z}`$ parameterizes $`CP`$-violating interactions.
Similarly, we will adopt the following parameterization of the $`Wtb`$ vertex suitable for the $`t`$ and $`\overline{t}`$ decays:
$`\mathrm{\Gamma }_{Wtb}^\mu ={\displaystyle \frac{g}{\sqrt{2}}}V_{tb}\overline{u}(p_b)\left[\gamma ^\mu (f_1^LP_L+f_1^RP_R){\displaystyle \frac{i\sigma ^{\mu \nu }k_\nu }{M_W}}(f_2^LP_L+f_2^RP_R)\right]u(p_t),`$
$`\overline{\mathrm{\Gamma }}_{Wtb}^\mu ={\displaystyle \frac{g}{\sqrt{2}}}V_{tb}^{}\overline{v}(p_{\overline{t}})\left[\gamma ^\mu (\overline{f}_1^LP_L+\overline{f}_1^RP_R){\displaystyle \frac{i\sigma ^{\mu \nu }k_\nu }{M_W}}(\overline{f}_2^LP_L+\overline{f}_2^RP_R)\right]v(p_{\overline{b}}),`$
(2)
where $`P_{L/R}=(1\gamma _5)/2`$, $`V_{tb}`$ is the $`(tb)`$ element of the Kobayashi-Maskawa matrix and $`k`$ is the momentum of $`W`$.
It will be assumed here that interactions of leptons with gauge bosons are properly described by the SM. Through the calculations all fermions except the top quark will be considered as massless. We will also neglect terms quadratic in non-standard form factors.
## 2 Angular and Energy Distributions
In this section we will present results for $`d^2\sigma /dx_fd\mathrm{cos}\theta _f`$ of the top-quark decay product $`f`$, where $`f`$ could be either $`\mathrm{}^\pm `$ or $`\stackrel{()}{b}`$, $`x_f`$ denotes the normalized energy of $`f`$ and $`\theta _f`$ is the angle between the $`e^{}`$ beam direction and the direction of $`f`$ momentum in the $`e^+e^{}`$ CM frame.
Using the technique developed by Kawasaki, Shirafuji and Tsai and adopting the general formula for the $`t\overline{t}`$ distribution $`d\sigma (s_+,s_{})/d\mathrm{\Omega }_t`$ found by Brzezinski et al., one obtains the following result for the distribution:
$$\frac{d^2\sigma }{dx_fd\mathrm{cos}\theta _f}=\frac{3\pi \beta \alpha _{\text{EM}}^2}{2s}B_f\left[\mathrm{\Theta }_0^f(x_f)+\mathrm{cos}\theta _f\mathrm{\Theta }_1^f(x_f)+\mathrm{cos}^2\theta _f\mathrm{\Theta }_2^f(x_f)\right],$$
(3)
where $`\beta `$ is the top velocity, $`\alpha _{\text{EM}}`$ is the fine structure constant and $`B_f`$ denotes the appropriate branching fraction. The energy dependence is specified by the functions $`\mathrm{\Theta }_i^f(x_f)`$, explicit forms of which could be found in a recent paper by the authors of this contribution. They are parameterized both by production and decay form factors.
The angular distribution $`d\sigma /d\mathrm{cos}\theta _f`$ for $`f`$ could be easy obtained from eq.(3) by the integration over the energy of $`f`$. It turns out that for $`f=\mathrm{}`$ all the non-standard contributions from the decay vertex disappear upon integration over the energy $`x_f`$.
The fact that the angular leptonic distribution is insensitive to corrections to the $`V`$-$`A`$ structure of the decay vertex allows for much more clear tests of the production vertices through a measurement of the distribution, since that way we can avoid a contamination from non-standard structure of the decay vertex. As an illustration, we define a $`CP`$-violating asymmetry which could be constructed using the angular distributions of $`f`$ and $`\overline{f}`$:
$$𝒜_{CP}(\theta _f)=\left[\frac{d\sigma ^+(\theta _f)}{d\mathrm{cos}\theta _f}\frac{d\sigma ^{}(\pi \theta _f)}{d\mathrm{cos}\theta _f}\right]/\left[\frac{d\sigma ^+(\theta _f)}{d\mathrm{cos}\theta _f}+\frac{d\sigma ^{}(\pi \theta _f)}{d\mathrm{cos}\theta _f}\right],$$
(4)
where $`d\sigma ^{+/}`$ is referring to $`f`$ and $`\overline{f}`$ distributions, respectively. The asymmetry for $`f=\mathrm{}`$ is sensitive to $`CP`$ violation originating exclusively from the production mechanism, i.e. $`CP`$-violating form factors $`\delta D_\gamma `$ and $`\delta D_Z`$ while the decay vertex enters with the SM $`CP`$-conserving coupling. For bottom quarks the effect of the modification of the decay vertex is contained in corrections to so called depolarization factors, $`\alpha ^b+\alpha ^{\overline{b}}\mathrm{Re}(f_2^R\overline{f}_2^L)`$ with SM $`CP`$-conserving contribution from the production process. As it is seen from fig.1 for $`\sqrt{s}=1\text{TeV}`$ the asymmetry could be quite large, e.g., reaching for the semileptonic decays $`20\%`$ for $`\mathrm{Re}(\delta D_\gamma )=\mathrm{Re}(\delta D_Z)=0.2`$.
## 3 Optimal Observable Analysis
Using the double angular and energy distributions an expected statistical uncertainty for determination of real parts for all the form factors has been found adopting optimal observables and varying the beam polarizations $`P_e^{}`$ and $`P_{e^+}`$. $`|\mathrm{cos}\theta |0.9`$ has been assumed as a cut for the polar angle. For $`t\overline{t}`$ tagging efficiency in $`\mathrm{}`$ \+ 4 jet channel we adopted 60% and for the integrated luminosity we chose the TESLA design with $`L=500\text{fb}^1`$ at $`\sqrt{s}=500\text{GeV}`$.
Generically we have observed that positive polarization led to smaller statistical errors for the eight form factors in the production vertices. For each form factor we have adjusted the optimal beam polarization such that the statistical error was minimal:
$$\begin{array}{cc}\mathrm{\Delta }[\text{Re}(\delta A_\gamma )]=0.16\hfill & \mathrm{for}P_e^{}=0.7\mathrm{and}P_{e^+}=0.7\hfill \\ \mathrm{\Delta }[\text{Re}(\delta A_Z)]=0.07\hfill & \mathrm{for}P_e^{}=0.5\mathrm{and}P_{e^+}=0.4\hfill \\ \mathrm{\Delta }[\text{Re}(\delta B_\gamma )]=0.09\hfill & \mathrm{for}P_e^{}=0.2\mathrm{and}P_{e^+}=0.2\hfill \\ \mathrm{\Delta }[\text{Re}(\delta B_Z)]=0.27\hfill & \mathrm{for}P_e^{}=0.4\mathrm{and}P_{e^+}=0.4\hfill \\ \mathrm{\Delta }[\text{Re}(\delta C_\gamma )]=0.11\hfill & \mathrm{for}P_e^{}=0.1\mathrm{and}P_{e^+}=0.0\hfill \\ \mathrm{\Delta }[\text{Re}(\delta C_Z)]=1.11\hfill & \mathrm{for}P_e^{}=0.1\mathrm{and}P_{e^+}=0.0\hfill \\ \mathrm{\Delta }[\text{Re}(\delta D_\gamma )]=0.08\hfill & \mathrm{for}P_e^{}=0.2\mathrm{and}P_{e^+}=0.1\hfill \\ \mathrm{\Delta }[\text{Re}(\delta D_Z)]=14.4\hfill & \mathrm{for}P_e^{}=0.2\mathrm{and}P_{e^+}=0.1\hfill \end{array}$$
(5)
As it is seen the precision of $`\delta \{C,D\}_Z`$ measurement would be very poor even for the optimal polarization. In addition, determination of $`\delta D_\gamma `$ would be difficult as well since its error varies rapidly with the polarization. For example, $`\mathrm{\Delta }[\text{Re}(\delta D_\gamma )]`$ becomes 0.86 for $`P_e^{}=0.1/P_{e^+}=0.1`$ and 0.99 for $`P_e^{}=0.3/P_{e^+}=0.1`$. The source of that sensitivity is hidden in the neutral current structure with $`\mathrm{sin}^2\theta _W0.23`$. Indeed, the optimal polarization becomes $`P_e^{}=0.1`$ instead of 0.2 $`(\mathrm{\Delta }[\text{Re}(\delta D_\gamma )]=0.09)`$ for $`\mathrm{sin}^2\theta _W=0.25`$. On the other hand, a good determination (almost independently of the polarization) could be expected for $`f_2^R`$. Indeed, the best precision is
$$\mathrm{\Delta }[\text{Re}(f_2^R)]=0.01\mathrm{for}P_e^{}=0.8\mathrm{and}P_{e^+}=0.8$$
(6)
however, we have $`\mathrm{\Delta }[\text{Re}(f_2^R)]=0.03`$ even for $`P_e^{}=P_{e^+}=0`$.
## 4 Summary and Conclusions
We have presented here the angular and energy distributions for $`\stackrel{()}{f}`$ in the process $`e^+e^{}t\overline{t}\stackrel{()}{f}\mathrm{}`$, where $`f=\mathrm{}`$ or $`b`$ quark. The most general ($`CP`$-violating and $`CP`$-conserving) couplings for $`\gamma t\overline{t}`$, $`Zt\overline{t}`$ and $`Wtb`$ have been assumed. The bottom-quark mass has been neglected and we have kept only terms linear in anomalous couplings.
Test of CP violation has also been discussed, introducing a CP-sensitive asymmetry $`𝒜_{CP}`$ as an example.
Using the double angular and energy distribution of a lepton we have found that at $`\sqrt{s}=500\text{GeV}`$ with the integrated luminosity $`L=500\text{fb}^1`$ the best determined top-quark coupling would be the axial coupling of the Z boson with the error $`\mathrm{\Delta }[\text{Re}(\delta A_Z)]=0.07`$ while the lowest precision is expected for $`\text{Re}(\delta D_Z)`$ with $`\mathrm{\Delta }[\text{Re}(\delta D_Z)]=14.4`$.
## Note added
After this work has been presented, a paper by Rindani appeared where the double angular and energy leptonic distribution has been also found.
## Acknowledgments
BG is grateful to the organizers of the PASCOS99 conference for creating a very warm and inspiring atmosphere during the meeting. This work is supported in part by the State Committee for Scientific Research (Poland) under grant 2 P03B 014 14 and by Maria Skłodowska-Curie Joint Fund II (Poland-USA) under grant MEN/NSF-96-252.
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# Untitled Document
A simple question about a complicated object
The complicated object is the cohomological induction functor (for which the biblical reference is Knapp & Vogan \[KV\]). Here is the simple question. For definiteness sake set $`G:=PU(n,1)`$, $`H:=PU(k,1)`$ with $`kn`$, and let $`_\rho ^n`$ (resp. $`_\rho ^k`$) be the category of Harish-Chandra modules with the generalized infinitesimal character of the trivial module attached to $`G`$ (resp. $`H`$). I’ll surprise nobody by claiming that $`H`$ is a subgroup of $`G`$. What’s far less obvious, but proved by Khoroshkin \[K\], is the existence of a full embedding $`F`$ of $`_\rho ^k`$ into $`_\rho ^n`$ (Fuser \[F,Thm I.4.2\] showed that $`F`$ is even Ext-full, that is compatible with Ext calculus), prompting the question : is there a geometric interpretation of the embedding $`F`$ ? The first candidate for $`F`$ is the (ordinary) induction functor ; but this fails miserably — so let’s break the pseudo-suspense of this introduction by saying that I claim that $`F`$ is (isomorphic to) a certain cohomological induction functor, and conjecture that this phenomenon is general.
1. Statements
Let $`G`$ be a center free connected semisimple Lie group, $`KG`$ a maximal compact subgroup, $`𝔤𝔨`$ the respective complexified Lie algebras. Let’s start by recalling the notion of Harish-Chandra module. Say that a $`𝔤`$-module $`V`$ is $`𝔨`$-finite if it is a sum of finite dimensional sub-$`𝔨`$-modules, and that $`V`$ is an $`(𝔤,𝔨)`$-module if it is $`𝔨`$-finite and $`𝔨`$-semisimple. The category $`=(𝔤,K)`$ of Harish-Chandra modules is the full subcategory of $`𝔤`$-mod whose objects are those $`(𝔤,𝔨)`$-modules of finite length $`V`$ such that for any finite dimensional $`𝔨`$-invariant subspace $`FV`$ the action of $`𝔨`$ on $`F`$ exponentiates to $`K`$. The category $``$ is a $`\mathrm{}`$-category in the sense of Bass \[B\] page 57. Let $`I`$ be the annihilator of the trivial module in the center of $`U(𝔤)`$, let
$$_\rho =_\rho (𝔤,K)$$
be the full sub-$`\mathrm{}`$-category of $``$ whose objects are annihilated by some power of $`I`$, let $``$ be the set of isomorphism classes of simple objects of $`_\rho `$ \[it is a finite set\] ; for each $`i`$ choose a representative $`V_ii`$ and let $`\mathrm{}(i)`$ be the projective dimension of $`V_i`$ \[i.e. the supremum in $`\mathrm{}\{+\mathrm{}\}`$ of the set $`\{n\mathrm{}|Ext^n(V_i,)0\}`$\].
(1) Definition. The $`_\rho `$-ordering is the smallest partial ordering $``$ on $``$ satisfying
$$\begin{array}{c}i,j\\ \\ \mathrm{}(j)=\mathrm{}(i)+1<\mathrm{}\\ \\ Ext^1(V_j,V_i)0\end{array}\}ij.$$
(2) Definition. The sub-$`\mathrm{}`$-category generated by the subset $`𝒥`$ of $`_\rho `$ is the full sub-$`\mathrm{}`$-category $`𝒥__\rho `$ of $`_\rho `$ characterized by the condition that an object $`V`$ of $`_\rho `$ belongs to $`𝒥__\rho `$ iff each simple subquotient of $`V`$ is isomorphic to $`V_j`$ for some $`j𝒥`$.
(3) Definition. Say that a full sub-$`\mathrm{}`$-category $`𝒞`$ of $`_\rho `$ is Ext-full in $`_\rho `$ if for all $`V,W𝒞`$ the natural morphism
$$Ext_𝒞^{}(V,W)Ext__\rho ^{}(V,W)$$
is an isomorphism.
For $`i`$ put $`𝒥_i:=\{j|ji\}`$, let $`\theta `$ be the Cartan involution of $`(𝔤,K)`$, denote by $`d`$ the dimension of $`G/K`$, and consider the following
(4) Property of $`G`$. For each $`i`$ such that $`V_i`$ is unitary the cohomology $`H^{d\mathrm{}(i)}(𝔤,K;V_i)`$ is nonzero and there is a $`\theta `$-stable parabolic subalgebra of $`𝔤`$ with Levi subgroup $`L=L_i`$ (see Vogan \[V2,4.1,4.2\] for definitions) satisfying
(a) the corresponding cohomological induction functor $`F`$ (see \[KV\]) sets up an equivalence
$$_\rho (𝔩,LK)\stackrel{}{}𝒥_i_{_\rho (𝔤,K)};$$
(b) $`F\mathrm{}V_i;`$
(c) $`𝒥_i_{_\rho (𝔤,K)}`$ is Ext-full in $`_\rho (𝔤,K)`$ ;
(d) if $`a`$ is nonzero vector of $`H^{d\mathrm{}(i)}(𝔤,K;V_i)`$ and $`V`$ a simple object of $`_\rho (𝔩,LK)`$, then the map
$$\begin{array}{ccc}H^{}(𝔩,LK;V)& & H^{d\mathrm{}(i)+}(𝔤,K;FV)\\ \\ x& & F(x)a\end{array}$$
\[where $``$ denotes the cup-product\] is an isomorphism \[of $`H^{}(𝔩,LK;\mathrm{})`$-modules\] ;
(e) we have $`2\mathrm{}(i)=d+dimL/(LK)`$.
Note once and for all that (e) follows from (a) by the well known argument which consists in setting $`V:=\mathrm{}`$ and using Poincaré duality.
(6) Conjecture. All center free connected semisimple Lie groups have Property (4).
A partial proof (with explicitly indicated gaps) of the fact that $`PU(n,1)`$, $`P\text{Spin}(n,1)`$ and $`SL(3,\mathrm{})`$ have Property (4) is contained in the expanded version of this text, downloadable from
http://www.iecn.u-nancy.fr/$``$gaillard/Recherche/Ci/ci.html
* * *
\[F\] Fuser A., Autour de la conjecture d’Alexandru, Thèse de l’Univesité Nancy 1 (1997).
\[K\] Khoroshkin S.M., Category of Harish-Chandra modules of the group SU(n,1). Funct. Anal. Appl. 14 (1980) 153-155.
\[KV\] Knapp A. & Vogan D., Cohomological induction and unitary representations, Princeton University Press (1995).
\[V\] Vogan D., Cohomology and group representations, Proc. Symp. Pure Math. 61 (1997) 219-243.
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# The Energy Spectrum of TeV Gamma-Rays from the Crab Nebula as measured by the HEGRA system of imaging air Čerenkov telescopes
## 1 Introduction
The Crab Nebula has been observed and studied over an enormously broad photon energy range embracing the radio, optical, and X-ray bands, as well as, high energy $`\gamma `$-ray region up to hundreds of TeV. The various theoretical scenarios of photon emission are primarily based on the Synchro Compton model (Gould 1965) which combines the synchrotron and inverse Compton (IC) emissions from high-energy electrons, which are accelerated up to $`100`$ TeV and interact with the magnetic field and the low frequency seed photons within the Nebula (de Jager & Harding 1992; Atoyan & Aharonian 1996; de Jager et al. 1996; Hillas et al. 1998). The predicted IC spectrum in the TeV energy domain appears to be very sensitive to the model parameters, such as the value of the magnetic field, the nature of the seed photons, the maximum energy of electrons, etc. Although IC scenarios of the photon emission are widely believed to be appropriate for the Crab Nebula one can not exclude the possible contribution of $`\gamma `$-ray fluxes from $`\pi ^{}`$-decay (see Atoyan & Aharonian 1996, Bednarek & Protheroe 1997). Thus, a measurement of TeV Crab Nebula spectrum sets major constraints on theoretical expectations and precise spectral measurements allow to fix the model parameters.
The imaging air Čerenkov technique was successfully used for observations of the Crab Nebula in TeV $`\gamma `$-rays. Since the time of detection at a 9$`\sigma `$ confidence level by the Whipple group (Weekes et al. 1989) a number of observations of the Crab Nebula have been made at TeV energies (for reviews see Ong 1998; Catanese & Weekes 1999). By now the Crab Nebula is established as the standard candle of steady TeV $`\gamma `$-ray emission. The optical Nebula of the Crab has an angular extension of about 6 arc min (Hester et al. 1995). The standard analysis for the HEGRA system of IACTs provides an angular resolution of about $`0.1^{}`$ (see Figure 1). In the present analysis the Crab Nebula was assumed to be as a $`\gamma `$-ray point source (see Figure 2). The detailed mapping of the TeV $`\gamma `$-rays from the Crab Nebula undertaken by Hofmann (1999) provided an upper limit for the Crab angular extension in TeV $`\gamma `$-rays of $`1.5`$ arcmin. No pulsed emission has been seen so far from the Crab (Gillanders et al. 1997; Burdett et al. 1999; Aharonian et al. 1999a). Vacanti et al (1991) provided a first measurement of the energy spectrum of TeV $`\gamma `$-ray emission from the Crab Nebula using 10 m Whipple telescope data and derived a power-law spectrum index of 2.4 and a differential $`\gamma `$-ray flux at 400 GeV of $`2.510^{10}\mathrm{photons}\mathrm{cm}^2\mathrm{s}^1`$. Note that the absolute calibration of individual imaging air Čerenkov telescopes is quite difficult because there is no “test” beam of TeV $`\gamma `$-rays. The uncertainties in absolute calibration propagate to uncertainties of estimated $`\gamma `$-ray fluxes. Different methods of telescope calibration have been developed recently in order to reduce systematic errors which could influence the estimate of the absolute $`\gamma `$-ray flux and the slope of energy spectrum (see e.g., Frass et al. 1997). The energy spectrum measurements heavily rely on Monte Carlo simulations of the telescope response. Thus in the past the analysis of the same observational data using various methods of telescope calibration revealed, not infrequently, very different estimates of the telescope energy threshold and of the TeV $`\gamma `$-ray fluxes. These estimates may vary by more than a factor of 2. Thus, based on data taken with the prototype HEGRA imaging Čerenkov telescope (CT2), Konopelko et al. (1996) detected the TeV $`\gamma `$-ray signal from the Crab Nebula with a signal-to-noise ratio of 10$`\sigma `$ and derived from the data a power-law index of the energy spectrum of $`2.7\pm 0.1`$ and an integral flux of TeV $`\gamma `$-rays above 1 TeV of $`8(\pm 1)_{\mathrm{Stat}}(\pm 2.4)_{\mathrm{Syst}}10^{12}\mathrm{photons}\mathrm{cm}^2\mathrm{s}^1`$. However, later on, the detailed treatment of the telescope hardware (optical smearing, photon-to-photoelectrons conversion efficiency, etc) allowed a more precise estimate of the TeV $`\gamma `$-ray flux using the single telescope data, as $`1.5(\pm 0.2)_{\mathrm{Stat}}(+1.00.5)_{\mathrm{Syst}}10^{11}\mathrm{photons}\mathrm{cm}^2\mathrm{s}^1`$ above 1 TeV (Petry et al. 1996).
Significant improvements in the telescope hardware as well as in the simulations and the data analysis (Fegan, 1997) recently provided measurements of the Crab Nebula energy spectrum over the energy range from 200 GeV up to 50 TeV by several groups, Whipple (Hillas et al. 1998), CANGAROO (Tanimori et al. 1998), and CAT (see Barrau 1998) exploiting the imaging Čerenkov technique. An additional measurement was provided by a low energy air shower array (Amenomori et al. 1999). However the uncertainties of the $`\gamma `$-ray flux estimates as well as of the energy spectrum slope remain rather large, in particular at the high energy end of the measurements i.e. beyond 10 TeV. For example the Tibet data (Amenomori et al. 1999) show significantly higher (by a factor of 2) $`\gamma `$-ray fluxes compared with the results obtained using the IACTs in the energy range from 3 to $`18`$ TeV, and are in favor of a gradual steepening of the spectral slope at high energies. Thus additional precise measurements of the Crab Nebula energy spectrum in TeV $`\gamma `$-rays are of great importance. The HEGRA system of IACT provides such data. It was primarily designed for detailed spectral measurements in the TeV energy domain utilizing the advantages of stereoscopic observations. Stereo imaging gives several advantages for spectral studies, compared to a single telescope: (i) direct measurement of the shower impact parameter with an accuracy better than 10 m (ii) good energy resolution of 18% (iii) wide dynamic range from 500 GeV to 20 TeV (iv) extended abilities for systematic studies using several images for an individual shower. The detailed systematic studies for the spectrum evaluation technique have been recently made using Mrk 501 1997 observational data (Aharonian et al. 1999b, 1999c). The performance of the system was discussed by Konopelko et al. (1999a). Here we present Crab Nebula data taken with the HEGRA telescope system in 1997/1998 and 1998/1999 observational campaigns. The data were analyzed using a new technique of energy spectrum evaluation for the stereoscopic observations. We also discuss the physics implications of the present results for the modeling of the TeV $`\gamma `$-ray emission from the Crab Nebula.
## 2 Observational data
The Crab Nebula was extensively observed with the HEGRA IACT system in two observing seasons from September 1997 to March 1998 and from October 1998 to April 1999. The observations were made with the stereoscopic system of IACTs which are located on La Palma, Canary Islands (Aharonian et al. 1999c). Each of the telescopes consists of a 8.5 $`\mathrm{m}^2`$ reflector focussing the Čerenkov light onto a photomultiplier tube camera. The 271 photomultipliers in the camera were arranged in a hexagonal matrix covering a field of view with a radius of $`2.15^{}`$. The telescope camera was triggered when the signal in two next neighbors of the 271 photo multiplier tubes exceeded a threshold of 8 photoelectrons, and the system readout started when at least two telescopes were triggered by Čerenkov light from an air shower. The detection rate was 12.6 Hz near the Zenith in December 1997 and dropped down to about 10 Hz in December 1999 for the 4-telescope system due to aging of the PMTs and reduced mirror reflectivity.
The Crab Nebula was observed in a “wobble mode”; i.e., the telescopes were pointed in Declination $`\pm 0.5^{}`$ away from the nominal Crab Nebula position (the sign of the angular shift was altered from one run of 20 min to the next). This is useful for continuous monitoring of the cosmic-ray background because it positions the OFF-source region symmetric to the camera center, and $`1^{}`$ apart from the ON-source region. Observations of the Crab Nebula at zenith angles up to 50 degree were made with 4 telescopes from 1997 September 1 to 1998 March 29, for a total of 82.5 hr of data taken at good weather. Through the fire at the HEGRA site one of the telescopes was damaged and was out of operation for a month in October-November 1997. At that time, Crab Nebula observations were made with only three telescopes in the system, providing an event rate of 10.3 Hz near the zenith. Due to unstable weather and a substantial amount of dust which came from Sahara desert to the island, the average detection rate in 1998 February-March was reduced down to 10.7 Hz in observations near the zenith. At the beginning of 1998/1999 observational period 4 telescopes were operational. Since October 1998 the HEGRA collaboration operates 5 telescopes. However, for technical reasons one telescope (CT2) was out of operation since December 1998 until the end of the first observation campaign (April 1999). During the last observational period (since August 1999) observations at zenith angles less than $`50^{}`$ were taken for about 76.1 hrs. In addition, the observations at large zenith angles (LZA) ($`50^{}<\theta <65^{}`$), were carried out for a total of 24 hrs in order to study the performance of the telescope system at LZA and to extend the measurements of the Crab Nebula energy spectrum beyond 10 TeV. The total exposure times for the three periods are summarized in Table 1.
Only data taken under good weather conditions were used in the analysis. In order to exclude data taken under less than optimal telescope performance conditions the entire database has been checked very carefully as follows: First, each night the compressed protocols of the system performance were transfered to one of the collaboration host institutes where they were scanned by software tools which closely monitor the status of the telescopes’ hardware (single pixel rates, trigger rates of system telescopes, tracking accuracy, etc). This information was accumulated in a corresponding database which was used afterwards for a standard data reduction procedure. The final condensed data file for each particular run contains all the information needed for data analysis. In addition a specific software tool was developed which allows to control a posteriori for each data run (i) the system trigger rate, taking into account the zenith angle dependence (ii) the angular shape of the cosmic ray images, tested by a $`\chi ^2`$–criterion for the deviation of the mean scaled Width distribution for a single run from the corresponding average distribution filled over an extended sample of runs (iii) the flatness of the $`\theta ^2`$–distribution for the isotropic cosmic ray images over the full field of view (iv) the image Size distributions for each individual telescope.
## 3 Analysis
The stereoscopic imaging analysis of the data is based on the geometrical reconstruction of the shower arrival direction and the shower core position in the observation plane, as well as on the joint parameterization of the shape of the Čerenkov light images. The simultaneous registration of several ($``$2) Čerenkov light images from an air shower provides an angular resolution of $`0.^{}1`$ for $`\gamma `$-ray showers. For each individual shower, stereoscopic observations allow to determine the position of the shower axis. Thus, at first only air showers within a certain impact distance $`\mathrm{R}_0`$ from the center of the telescope system were selected. The limiting upper radius of $`\mathrm{R}_0=200\mathrm{m}`$ was used for zenith angles less than 50 degrees, and a significantly larger radius of 400 m for the large zenith angle observations ($`>50`$ degrees). The effective collection area in observations at LZA dramatically increases at high energies, far beyond the limiting radius of 200 m (Konopelko et al. 1999b). For the data taken at zenith angles up to $`50^{}`$ an orientation cut $`\theta ^2<0.05[\mathrm{deg}^2]`$ was applied, where $`\theta ^2`$ is the squared angular distance of the reconstructed source position from the true source position. In addition the data were analyzed using the mean scaled Width parameter, $`<\stackrel{~}{w}>`$. To compensate for the dependence of the image shape on primary shower energy and distance from shower core to the telescope (impact parameter), the standard parameter Width (Fegan, 1997) ($`w^k`$), calculated for each telescope, is scaled according to the Monte Carlo predicted values, for $`\gamma `$-rays $`<w>_{ij}^k`$, taken for the corresponding bin of reconstructed distance from the telescope to the shower core (i) and for the corresponding bin of image size (total number of photoelectrons in the image)(j) (Aharonian et al. 1999b, 1999c). The mean scaled Width parameter is defined for each individual shower as follows
$$<\stackrel{~}{w}>=1/N\underset{k=1}{\overset{N}{}}w^k/<w>_{ij}^k$$
(1)
where $`N`$ is the number of triggered telescopes. This parameter was introduced in order to provide an almost constant $`\gamma `$-ray acceptance over the dynamic energy range of the telescope system. The optimum cut on mean scaled Width is about 1.1, which gives a $`\gamma `$-ray acceptance of $`60`$% at Small Zenith Angle (SZA). However, for a precise determination of $`\gamma `$-ray spectra, a loose cut on mean scaled Width ($`<\stackrel{~}{w}><1.2`$) has been so far used in the data analysis in order to maximize the $`\gamma `$-ray acceptance and to minimize systematic errors related to cut efficiencies. Thus, the second $`\gamma `$-ray selection criterion was $`<\stackrel{~}{w}>1.2`$. This set of cuts was found to be optimal for spectrum studies (Aharonian et al. 1999b, 1999c). These loose analysis cuts provide a Crab Nebula $`\gamma `$-ray rate of 83 $`\gamma `$s/hr at SZA (less that $`25^{}`$) for the 5-IACT system. The corresponding energy threshold of the $`\gamma `$-rays is about 500 GeV. For the LZA data the looser orientation cut of $`\theta ^2<0.1[\mathrm{deg}^2]`$ was used because of the lower accuracy of the arrival direction reconstruction for the $`\gamma `$-ray showers. In order to improve the cosmic ray rejection in observations at LZA an additional parameter, mean scaled Length, $`<\stackrel{~}{l}>`$ was used, which is defined by analogy with $`<\stackrel{~}{w}>`$ (Konopelko et al. 1999b). These two parameters, $`<\stackrel{~}{w}>`$ and $`<\stackrel{~}{l}>`$, can be used for calculating a Mahalanobis distance, MD (Mahalanobis 1963), in two-dimensional space as
$$\mathrm{MD}=((1<\stackrel{~}{\mathrm{w}}>)^2/\sigma _{<\stackrel{~}{\mathrm{w}}>}^2+(1<\stackrel{~}{\mathrm{l}}>)^2/\sigma _{<\stackrel{~}{\mathrm{l}}>}^2)^{1/2}$$
(2)
where $`\sigma _{<\stackrel{~}{w}>}`$ and $`\sigma _{<\stackrel{~}{l}>}`$ are the standard deviations for the corresponding distributions of $`<\stackrel{~}{w}>`$ and $`<\stackrel{~}{l}>`$. The optimum value of the MD cut for LZA is found to be 1.5. Note that this analysis improves the enhancement factor by $`30`$% (it gives $`50`$% acceptance of $`\gamma `$-rays) in observations at LZA, whereas it gives only marginal improvement for the data taken at SZA. The Crab Nebula $`\gamma `$-ray rate in observations at LZA ($`60^{}`$) is about 16 $`\gamma `$’s/hr with a corresponding energy threshold of $``$5 TeV. Note that SZA observations give a $`\gamma `$-ray rate at high energies (above 3 TeV) of $`8\gamma `$s/hr. A summary of the data is shown in Table 2.
The observations of the Crab Nebula have been made during 6 periods which differ in the system configuration, reflectivity of the mirrors, light reflection by the pixel funnels and camera protecting plate etc. All that affects the hardware event rate of the telescope system, $`\mathrm{R}_{\mathrm{exp}}`$. These changes of system performance were implemented in the Monte Carlo simulations. Assuming the standard chemical composition of the primary cosmic rays (Wiebel, 1994) the calculated detection rates, $`\mathrm{R}_{\mathrm{MC}}`$, were adjusted to the measured rates (see Table 3).
The collection areas, as a function of energy and zenith angle, for $`\gamma `$-ray showers have been inferred from Monte Carlo simulations (Konopelko et al. 1999a). The rms error of the energy determination is $`\mathrm{\Delta }\mathrm{E}/\mathrm{E}0.18`$. The Monte Carlo studies show that for a good energy resolution of 18% this approach does not distort the initial spectrum shape. The collection area for $`\gamma `$-rays rises very quickly in the energy range near the energy threshold of the telescope system, which is 500 GeV, whereas it is almost constant at the energies above $`3`$ TeV. Even slight variations of the trigger threshold could lead to noticeable systematic changes in the predicted spectral behavior in the energy range of $`0.51`$ TeV. This effect leads to a noticeable probability for “sub-threshold” triggers. In addition, the trigger level for different camera pixels is slightly different even after very accurate adjustment of the high voltage using the calibration laser runs. Measurements of the trigger setting for a number of camera pixels revealed variation in the trigger threshold of order $`10`$%. These variations were implemented into the simulations in order to estimate the corresponding systematic error of the energy spectrum at energies below 1 TeV. The fine tuning of the Monte Carlo simulations with respect to the IACT system data provided measurements of the flux of the cosmic ray protons in the energy range from 1.3 to 10 TeV (Aharonian et al. 1999d). The proton fluxes as measured by the HEGRA IACT system are perfectly consistent with the results of a bulk of satellite experiments held in this energy range.
The procedure for the evaluation of the energy spectrum using the stereoscopic observations was discussed in detail by Aharonian et al. (1997); Hofmann (1997), and, more recently in Aharonian et al. (1999c). In the stereoscopic observations the impact distance of the shower axis to a system telescope can be measured with an accuracy $`10`$ m. The energy E of a $`\gamma `$-ray shower is defined by interpolation over the “size” parameter $`\mathrm{S}`$ (total number of photoelectrons in Čerenkov light image) at a fixed impact distance $`R`$, as $`\mathrm{E}=\mathrm{f}_{\mathrm{MC}}(\mathrm{S},\mathrm{R},\theta )`$, where $`\theta `$ is the zenith angle and $`\mathrm{f}_{\mathrm{MC}}`$ is a function obtained from Monte Carlo simulations. Note that the Monte Carlo simulations used here include the sampling of detector response in great detail (Hemberger 1998). The energy distribution for the ON- and OFF-source events, after the orientation and shape image cuts, were histogrammed over the energy range from 500 GeV to 30 TeV with 8 bins per decade. The $`\gamma `$-ray energy spectrum was obtained by subtracting ON- and OFF-histograms and dividing the resulting energy distribution by the corresponding collection area and the $`\gamma `$-ray acceptance. In the present Crab Nebula analysis the energy spectrum measurements were extended up to large zenith angles ($`65^{}`$). The data were processed independently for each of the four zenith angle bins: $`(0^{}25^{}),(25^{}40^{}),(40^{}50^{}),(50^{}65^{})`$. The corresponding effective collection areas as well as the cut efficiencies were calculated as a function of the zenith angle. First, the energy spectra were derived for all zenith angle bins independently. Note that the spectra evaluated at different zenith angles are in a good agreement. For the final energy spectrum the different zenith angle bins were joined according to the prescription:
$`dJ_\gamma ^i/dE={\displaystyle \underset{j=1}{\overset{4}{}}}w_j(dJ_\gamma ^i/dE)_j\mathrm{\Theta }(E^iE_{th}^j),`$
$`w_j=t_j/t_0,i=1,n;`$ (3)
where $`dJ_\gamma ^i/dE`$, $`(dJ_\gamma ^i)_j/dE`$ are the differential energy spectra at energy $`E^i`$ as measured over all zenith angle ranges, and for the particular zenith angle bin (j), respectively. $`E_{th}^j`$ is an estimated energy threshold for the zenith angle bin $`j`$, $`t_j`$ is the observation time for the $`j`$-bin on the zenith angle, and $`t_0`$ is the total observation time. The first three zenith angle bins were joined using the time dependent weights $`\mathrm{w}_\mathrm{j}=\mathrm{t}_\mathrm{j}/\mathrm{t}_0`$. Finally the spectrum measured in the zenith angle range of $`(0^{}50^{})`$ was combined with the spectrum derived from large zenith angle data, $`(50^{}65^{})`$ using the weights based on the estimated statistical errors for both spectra. Such procedure takes into account the advantageous $`\gamma `$-ray rate in LZA observations at high energies.
The statistics of the $`\gamma `$-rays from the Crab Nebula provides a measurement of the energy spectrum up to a few tens of TeV. However, detection of Čerenkov light images with extremely large amplitudes - several thousands of ph.e. - is complicated by the nonlinearity in the PMT response as well as by the saturation in the 8 bit Flash-ADC readout. Measurements of the photomultiplier response under high light loads over the extended sample of the EMI 9073 PMTs gave a calibration function which was used to correct the image amplitudes. The readout of the HEGRA IACT is based on the sampling of Čerenkov light time impulse by the 16 FADC bins of $`8`$ ns each (Hess et al. 1998). The time pulses from the air showers with a full width at half maximum of a few ns were widened using an electronic scheme in order to fit into several FADC bins for the accurate measurement of the time profile. The smoothing of the FADC signal was unfolded back to the impulse, which almost always fits 2 FADC bins. The calibrated amplitude, summed over two FADC bins, is used as a measure of the pixel signal. For the high energy air showers the FADC signals run into saturation and the simple unfolding procedure fails. For such pulses the initial amplitude is reconstructed using the additional calibration function obtained by simultaneous measurements of light flashes with FADCs and a 14 bit ADC. This procedure drastically extends the dynamic range of the FADC readout.
To avoid the saturation problem one might only use images detected from air showers at large impact distances from the telescope system (e.g. beyond 150 m). The size of these images is very small even for high energy events because of the low Čerenkov light density far off the shower axis. However these images are very often truncated by the camera edge and do not allow a proper reconstruction of the shower impact point and of the shower energy. This effect becomes less important in observations at LZA because of the high shower maximum height (the images shrink to the camera center). In the present analysis the maximum impact distance of the shower core from the center of the system was extended up to 400 m for observations at LZA. Observations at LZAs permit measurements of the energy spectrum far beyond 10 TeV. The images of $`\gamma `$-ray air showers observed at LZAs have small Size and are not influenced by the saturation effect.
## 4 Results
We have observed the Crab Nebula extensively in two observational seasons with the HEGRA IACT system. The HEGRA system of 5 IACTs currently has a sensitivity which allows the detection of a $``$5$`\sigma `$ signal from the Crab Nebula within 1 hr of observation time (see Figure 1). The integral $`\gamma `$-ray fluxes measured during the different observational periods are consistent within the estimated statistical and systematic error (see Table 3). The differential energy spectrum of the Crab Nebula has been derived from the HEGRA data for two observational campaigns using recently developed advanced techniques for the measurements of the spectrum using stereoscopic data taken at small and large zenith angles. The Crab Nebula differential energy spectrum derived from SZA data matches quite well the spectrum derived at LZA (see Figure 3). The $`\gamma `$-ray rate measured at energies above 10 TeV in observations at LZAs exceeds the corresponding rate measured at SZA by a factor of 3. The LZA data are not affected by saturation effects. At 3.7$`\sigma `$ confidence level 27 $`\gamma `$-ray events from the Crab Nebula were detected in the highest energy bin from 17.8 to 23.7 TeV (see Figure 4). One may expect that a number of $`\gamma `$-ray events in the highest energy bin are spilled over from the lower energies. However, given the good energy resolution of 20% and the power law energy spectrum, such effect is very small and is compensated almost by the backwards influx of the $`\gamma `$-rays from the energies above (see e.g., Aharonian et al. 1995). Finally, the simulations show that spilling over of low energy $`\gamma `$-rays does not influence the resulting fluxes measured at the upper end of the power law spectrum.
The analysis for the different system configurations as well as for different trigger threshold values gives a differential energy spectrum of the Crab Nebula measured at zenith angles up to $`60^{}`$
$$\mathrm{dJ}_\gamma /\mathrm{dE}=(2.79\pm 0.02\pm 0.5)10^7(\frac{\mathrm{E}}{1\mathrm{TeV}})^{2.59\pm 0.03\pm 0.05}\mathrm{ph}\mathrm{m}^2\mathrm{s}^1\mathrm{TeV}^1$$
(4)
The statistical and systematic errors are also given. The final Crab Nebula spectrum as measured by the HEGRA collaborations is shown in Figure 4. The measured $`\gamma `$-ray fluxes are given in Table 4. The Crab Nebula energy spectrum is best fitted by a pure power law in the energy range 1-20 TeV. It does not exclude a possible slight steepening of the energy spectral usually predicted by inverse Compton modeling of TeV $`\gamma `$-ray emission. A fit with a logarithmic steepening of the power law spectrum gives the following result
$$\mathrm{dJ}_\gamma /\mathrm{dE}=(2.67\pm 0.01\pm 0.5)10^7(\frac{\mathrm{E}}{1\mathrm{TeV}})^{2.47\pm 0.1\pm 0.05(0.11\pm 0.10)\mathrm{log}(\mathrm{E})}\mathrm{ph}\mathrm{m}^2\mathrm{s}^1\mathrm{TeV}^1$$
(5)
Such a fit indicates the slight flattening of the spectrum at low energies as predicted by the IC calculations. However, the change of the energy spectrum slope is within the current statistical and systematic errors, and the data for the overall energy range are consistent with a simple power law fit in 0.5-20 TeV. The HEGRA Crab Nebula data match well the recent 20 TeV data published by the CANGAROO group (Tanimori et al. 1998), and are consistent with a flat power law index of $`2.5`$ beyond 20 TeV. The compilation of the world data is given in Figure 5. All data are consistent within statistical and systematic errors, except possibly for the Tibet data which show relatively higher fluxes.
## 5 Astrophysics implications
The TeV energy spectrum of the Crab Nebula as measured by the HEGRA system of IACTs is consistent with the expectations for the TeV $`\gamma `$-ray emission from pulsar-driven Nebulae (plerions). According to this scenario the ultra relativistic electrons, accelerated in the pulsar wind shock, produce TeV $`\gamma `$-rays through the IC scattering with soft photons within the Nebula. The predicted fluxes of TeV $`\gamma `$-rays rely on the spatial distribution of the magnetic field within the Nebula as determined by the parameter $`\sigma `$ (ratio of energy density of magnetic field to the particle energy density) and/or by the average magnetic field $`<\mathrm{B}>`$ in the optical nebula. The HEGRA data are shown in Figure 6 together with predicted spectra using two SSC models of TeV $`\gamma `$-ray emission (de Jager et al. 1996; Atoyan & Aharonian 1996). One may conclude that both models fit the HEGRA data rather well.
According to the calculations of the TeV $`\gamma `$-ray emission by de Jager et al. (1996) the $`\gamma `$-ray flux from the Crab Nebula at TeV energies constrains the choice of the parameter $`\sigma `$. The IC spectrum computed by de Jager et al. (1996), assuming for the parameter $`\sigma `$ a value $`0.003`$ gives a good fit to the HEGRA data. This value of the parameter $`\sigma `$ corresponds to the best-fitting magneto-hydrodynamic (MHD) solution of electron propagation in the Crab Nebula as found by Kennel & Coroniti (1984).
In another approach, assuming the spatial distribution of magnetic field in the Crab Nebula, one can determine the average magnetic field $`<\mathrm{B}>`$ in the optical nebula (Gould 1965). According to the calculations of Atoyan & Aharonian (1996), made within the framework of the MHD model of Kennel & Coroniti (1984), the average magnetic field $`<\mathrm{B}>`$ is determined by the TeV $`\gamma `$-ray flux as $`<\mathrm{B}>\mathrm{J}_\gamma ^{0.5}10^5\mathrm{G}`$ where $`\mathrm{J}_\gamma ,\mathrm{ph}\mathrm{cm}^2\mathrm{s}^1\mathrm{TeV}^1`$ is a differential flux at 1 TeV. Thus the HEGRA spectrum gives an average magnetic field strength $`<\mathrm{B}>(1.7\pm 0.3)10^4\mathrm{G}`$. This value is consistent with the estimate derived by Hillas et al. (1998) from the Crab Nebula date taken with the 10 m Whipple telescope.
The IC energy spectrum of $`\gamma `$ rays from the Crab Nebula, measured in the energy range from 1 TeV to 10 TeV, is likely to be a power law $`\mathrm{dJ}_\gamma /\mathrm{dE}\mathrm{E}^\alpha =\mathrm{E}^{2.6}`$. At the same time, for the energy range above 10 TeV calculations predict gradual steepening with $`\alpha `$2.7 and 2.9 at 10 and 30 TeV, respectively. That is due to both the energy loss of the ultra high energy electrons by fast synchrotron cooling and the Klein-Nishina effect in the cross-section of the inverse Compton scattering. Atoyan & Aharonian (1996) and Bednarek & Protheroe (1997) have shown that $`\pi ^0`$-decay $`\gamma `$-ray fluxes, due to the relativistic protons accelerated in the Crab Nebula, may noticeably contribute at energies above 10 TeV. However the HEGRA Crab Nebula data expanded up to 20 TeV are still consistent with the pure IC spectrum. To assess the contribution of $`\pi ^0`$-produced $`\gamma `$-rays from the Crab Nebula measurements above 30 TeV are needed. Note that the LZA technique could help to perform such observations.
The predicted IC $`\gamma `$-ray spectrum of the Crab Nebula is rather flat in the energy range below $`1`$ TeV. It could be well approximated by $`\mathrm{dJ}_\gamma /\mathrm{dE}\mathrm{E}^{2.0}`$ at 100 GeV. Detection of a gradual flattening in this energy range will prove the SSC scenario of TeV $`\gamma `$-ray emission. However the low energy points at the HEGRA Crab Nebula spectra ($`\mathrm{E}_\gamma <1\mathrm{TeV}`$) are strongly affected by possible systematic errors ($`50`$%) and do not allow such a conclusion. Future observations of the Crab Nebula with the forthcoming low threshold high sensitivity Čerenkov detectors (see Catanese & Weekes 1999) will offer precise measurements in these energy range.
## Acknowledgments
The support of the German ministry for Research and technology BMBF and of the Spanish Research Council CYCIT is gratefully acknowledged. We thank the Instituto de Astrophysica de Canarias for the use of the site and for supplying excellent working conditions at La Palma. We gratefully acknowledge the technical support staff of the Heidelberg, Kiel, Munich, and Yerevan Institutes.
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# MATRIX REPRESENTATIONS OF OCTONIONS AND THEIR APPLICATIONS
## Abstract
As is well-known, the real quaternion division algebra $``$ is algebraically isomorphic to a 4-by-4 real matrix algebra. But the real division octonion algebra $`𝕆`$ can not be algebraically isomorphic to any matrix algebras over the real number field $``$, because $`𝕆`$ is a non-associative algebra over $``$. However since $`𝕆`$ is an extension of $``$ by the Cayley-Dickson process and is also finite-dimensional, some pseudo real matrix representations of octonions can still be introduced through real matrix representations of quaternions. In this paper we give a complete investigation to real matrix representations of octonions, and consider their various applications to octonions as well as matrices of octonions.
AMS Mathematics Subject Classification: 15A33; 15A06; 15A24; 17A35
Key Words: quaternions, octonions, matrix representations, linear equations, similarity, eigenvalues, Cayley-Hamilton theorem
1. Introduction
Let $`𝕆`$ be the octonion algebra over the real number field $``$. Then it is well known by the Cayley-Dickson process that any $`a𝕆`$ can be written as
$$a=a^{}+a^{\prime \prime }e,$$
$`(1.1)`$
where $`a^{},a^{\prime \prime }=\{a=a_0+a_1i+a_2j+a_3k|i^2=j^2=k^2=1,ijk=1,a_0`$$`a_3\}`$, the real quaternion division algebra. The addition and multiplication for any $`a=a^{}+a^{\prime \prime }e,b=b^{}+b^{\prime \prime }e𝕆`$ are defined by
$$a+b=(a^{}+a^{\prime \prime }e)+(b^{}+b^{\prime \prime }e)=(a^{}+b^{})+(a^{\prime \prime }+b^{\prime \prime })e,$$
$`(1.2)`$
and
$$ab=(a^{}+a^{\prime \prime }e)(b^{}+b^{\prime \prime }e)=(a^{}b^{}\overline{b^{\prime \prime }}a^{\prime \prime })+(b^{\prime \prime }a^{}+a^{\prime \prime }\overline{b^{}})e,$$
$`(1.3)`$
where $`\overline{b^{}},\overline{b^{\prime \prime }}`$ denote the conjugates of the quaternions $`b^{}`$ and $`b^{\prime \prime }`$. In that case, $`𝕆`$ is an eight-dimensional non-associative but alternative division algebra over its center field $``$, and the canonical basis of $`𝕆`$ is
$$1,e_1=i,e_2=j,e_3=k,e_4=e,e_5=ie,e_6=je,e_7=ke.$$
$`(1.4)`$
The multiplication rules for the basis of $`𝕆`$ are listed in the following matrix
$$E_8^TE_8=\left[\begin{array}{cccccccc}1& \hfill e_1& \hfill e_2& \hfill e_3& \hfill e_4& \hfill e_5& \hfill e_6& \hfill e_7\\ e_1& \hfill 1& \hfill e_3& \hfill e_2& \hfill e_5& \hfill e_4& \hfill e_7& \hfill e_6\\ e_2& \hfill e_3& \hfill 1& \hfill e_1& \hfill e_6& \hfill e_7& \hfill e_4& \hfill e_5\\ e_3& \hfill e_2& \hfill e_1& \hfill 1& \hfill e_7& \hfill e_6& \hfill e_5& \hfill e_4\\ e_4& \hfill e_5& \hfill e_6& \hfill e_7& \hfill 1& \hfill e_1& \hfill e_2& \hfill e_3\\ e_5& \hfill e_4& \hfill e_7& \hfill e_6& \hfill e_1& \hfill 1& \hfill e_3& \hfill e_2\\ e_6& \hfill e_7& \hfill e_4& \hfill e_5& \hfill e_2& \hfill e_3& \hfill 1& \hfill e_1\\ e_7& \hfill e_6& \hfill e_5& \hfill e_4& \hfill e_3& \hfill e_2& \hfill e_1& \hfill 1\end{array}\right],$$
$`(1.5)`$
where $`E_8=[\mathrm{\hspace{0.17em}1},e_1,\mathrm{},e_7]`$. Under Eq.(1.4) all elements of $`𝕆`$ take the form
$$a=a_0+a_1e_1+\mathrm{}+a_7e_7,$$
$`(1.6)`$
where $`a_0`$$`a_7`$, which can also simply be written as $`a=\mathrm{Re}a+\mathrm{Im}a,`$ where $`\mathrm{Re}a=a_0.`$ The conjugate of $`a`$ is defined to be
$$\overline{a}=\overline{a^{}}a^{\prime \prime }e=\mathrm{Re}a\mathrm{Im}a.$$
$`(1.7)`$
This operation satisfies
$$\overline{\overline{a}}=a,\overline{a+b}=\overline{a}+\overline{b},\overline{ab}=\overline{b}\overline{a}$$
$`(1.8)`$
for all $`a,b𝕆.`$ The norm of $`a`$ is defined to be $`|a|:=\sqrt{a\overline{a}}=\sqrt{\overline{a}a}=\sqrt{a_0^2+a_1^2+\mathrm{}+a_7^2}.`$ Although $`𝕆`$ is nonassociative, it is still an alternative, flexible, quadratic, composition and division algebra over $`,`$ that is, for all $`a,b𝕆`$, the following equalities hold:
$$\begin{array}{c}\text{ }a(ab)=a^2b,(ba)a=ba^2,(ab)a=a(ba):=aba,(1.9)\text{ }\hfill \\ \text{ }a^1=\frac{\overline{a}}{|a|^2},(1.10)\text{ }\hfill \\ \text{ }a^22(\mathrm{Re}a)a+|a|^2=0,(\mathrm{Im}a)^2=|\mathrm{Im}a|^2,(1.11)\text{ }\hfill \\ \text{ }|ab|=|a||b|.(1.12)\text{ }\hfill \end{array}$$
As is well known, any finite-dimensional associative algebra over an arbitrary field $`𝔽`$ is algebraically isomorphic to a subalgebra of a total matrix algebra over the field. In other words, any element in a finite-dimensional associative algebra over $`𝔽`$ has a faithful matrix representation over the field. For the real quaternion algebra $``$, it is well known that through the bijective map
$$\varphi :a=a_0+a_1i+a_2j+a_3k\varphi (a)=\left[\begin{array}{cccc}\hfill a_0& \hfill a_1& \hfill a_2& \hfill a_3\\ \hfill a_1& \hfill a_0& \hfill a_3& \hfill a_2\\ \hfill a_2& \hfill a_3& \hfill a_0& \hfill a_1\\ \hfill a_3& \hfill a_2& \hfill a_1& \hfill a_0\end{array}\right],$$
$`(1.13)`$
$``$ is algebraically isomorphic to the matrix algebra
$$=\left\{\left[\begin{array}{cccc}\hfill a_0& \hfill a_1& \hfill a_2& \hfill a_3\\ \hfill a_1& \hfill a_0& \hfill a_3& \hfill a_2\\ \hfill a_2& \hfill a_3& \hfill a_0& \hfill a_1\\ \hfill a_3& \hfill a_2& \hfill a_1& \hfill a_0\end{array}\right]\right|a_0,a_1,a_2,a_3\},$$
$`(1.14)`$
and $`\varphi (a)`$ is a faithful real matrix representation of $`a`$. Our consideration for matrix representations of octonions are based on Eqs.(1.1)—(1.3) and the result in Eq.(1.13).
We next present some basic results related to matrix representations of quaternions, which will be serve as a tool for our examination in the sequel.
Lemma 1.1. Let $`a=a_0+a_1i+a_2j+a_3k`$ be given, where $`a_0`$$`a_3`$. Then the diagonal matrix $`\mathrm{diag}(a,a,a,a)`$ satisfies the following unitary similarity factorization equality
$$Q\left[\begin{array}{cccc}a& & & \\ & a& & \\ & & a& \\ & & & a\end{array}\right]Q^{}=\left[\begin{array}{cccc}\hfill a_0& \hfill a_1& \hfill a_2& \hfill a_3\\ \hfill a_1& \hfill a_0& \hfill a_3& \hfill a_2\\ \hfill a_2& \hfill a_3& \hfill a_0& \hfill a_1\\ \hfill a_3& \hfill a_2& \hfill a_1& \hfill a_0\end{array}\right]^{4\times 4},$$
$`(1.15)`$
where the matrix $`Q`$ has the following independent expression
$$Q=Q^{}=\frac{1}{2}\left[\begin{array}{cccc}\hfill 1& \hfill i& \hfill j& \hfill k\\ \hfill i& \hfill 1& \hfill k& \hfill j\\ \hfill j& \hfill k& \hfill 1& \hfill i\\ \hfill k& \hfill j& \hfill i& \hfill 1\end{array}\right],$$
$`(1.16)`$
which is a unitary matrix over $``$.
Lemma 1.2. Let $`a,b`$, and $`\lambda `$. Then
(a) $`a=b\varphi (a)=\varphi (b).`$
(b) $`\varphi (a+b)=\varphi (a)+\varphi (b),\varphi (ab)=\varphi (a)\varphi (b),\varphi (\lambda a)=\lambda \varphi (a),\varphi (1)=I_4.`$
(c) $`a=\frac{1}{4}E_4\varphi (a)E_4^{},`$ where $`E_4:=[\mathrm{\hspace{0.17em}1},i,j,k]`$ and $`E_4^{}:=[\mathrm{\hspace{0.17em}1},i,j,k]^T.`$
(d) $`\varphi (\overline{a})=\varphi ^T(a).`$
(e) $`\varphi (a^1)=\varphi ^1(a),`$ if $`a0.`$
(f) $`\mathrm{det}[\varphi (a)]=|a|^4.`$
We can also introduce from Eq.(1.13) another real matrix representation of $`a`$ as follows
$$\tau (a):=K\varphi ^T(a)K=\left[\begin{array}{cccc}\hfill a_0& \hfill a_1& \hfill a_2& \hfill a_3\\ \hfill a_1& \hfill a_0& \hfill a_3& \hfill a_2\\ \hfill a_2& \hfill a_3& \hfill a_0& \hfill a_1\\ \hfill a_3& \hfill a_2& \hfill a_1& \hfill a_0\end{array}\right],$$
$`(1.17)`$
where $`K=\mathrm{diag}(\mathrm{\hspace{0.17em}1},1,1,1)`$. Some basic operation properties on $`\tau (a)`$ are
$$\tau (a+b)=\tau (a)+\tau (b),\tau (ab)=\tau (b)\tau (a),\tau (\overline{a})=\tau ^T(a),$$
$`(1.18)`$
$$det[\varphi (a)]=|a|^4,\varphi (a^1)=\varphi ^1(a)\mathrm{if}a0.$$
$`(1.19)`$
Combining the two real matrix representations of quaternions with their real vector representations, we have the following important result.
Lemma 1.3.Let $`x=x_0+x_1i+x_2j+x_3k,`$ and denote $`\stackrel{}{x}=[x_0,x_1,x_2,x_3]^T,`$ called the vector representation of $`x`$. Then for all $`a,b,x`$, we have
$$\stackrel{}{ax}=\varphi (a)\stackrel{}{x},\stackrel{}{xb}=\tau (b)\stackrel{}{x},\stackrel{}{axb}=\varphi (a)\tau (b)\stackrel{}{x}=\tau (b)\varphi (a)\stackrel{}{x},$$
$`(1.20)`$
and the equality
$$\varphi (a)\tau (b)=\tau (b)\varphi (a)$$
$`(1.21)`$
always holds.
Proof. Observe that
$$\stackrel{}{x}=\varphi (x)\alpha _4^T,\stackrel{}{x}=\tau (x)\alpha _4^T,\alpha _4=[\mathrm{\hspace{0.17em}1},0,0,0].$$
We find by Lemma 1.1 and Eq.(1.2) that
$$\stackrel{}{ax}=\varphi (ax)\alpha _4^T=\varphi (a)\varphi (x)\alpha _4^T=\varphi (a)\stackrel{}{x},\stackrel{}{xb}=\tau (xb)\alpha _4^T=\tau (b)\tau (x)\alpha _4^T=\tau (b)\stackrel{}{x},$$
and
$$\stackrel{}{axb}=\stackrel{}{a(xb)}=\varphi (a)\stackrel{}{(xb)}=\varphi (a)\tau (b)\stackrel{}{x},\stackrel{}{axb}=\stackrel{}{(ax)b}=\tau (b)\stackrel{}{(ax)}=\tau (b)\varphi (a)\stackrel{}{x}.$$
These four equalities are exactly the results in Eqs.(1.20) and (1.21). $`\mathrm{}`$
Lemma 1.4. Let $`a,b,x𝕆`$ be given. Then
(a) $`\mathrm{Re}(ab)=\mathrm{Re}(ba),\mathrm{Re}((ax)b)=\mathrm{Re}(a(xb)).`$
(b) $`(aba)x=a(b(ax)),x(aba)=((xa)b)a.`$
(c) $`(ab)(xa)=a(bx)a,(bx)(ab)=b(xa)b.`$
(d) $`(a,b,x)=(a,x,b)=(x,a,b),`$ where $`(a,b,x)=(ab)xa(bx).`$
2. The real matrix representations of octonions
Based on the results on the real matrix representation of quaternions, we now can introduce real matrix representation of octonions.
Definition 2.1. Let $`a=a^{}+a^{\prime \prime }e𝕆,`$ where $`a^{}=a_0+a_1i+a_2j+a_3k,a^{\prime \prime }=a_4+a_5i+a_6j+a_7k.`$ Then the $`8\times 8`$ real matrix
$$\omega (a):=\left[\begin{array}{cc}\varphi (a^{})& \tau (a^{\prime \prime })K_4\\ \varphi (a^{\prime \prime })K_4& \tau (a^{})\end{array}\right],$$
$`(2.1)`$
is called the left matrix representation of $`a`$ over $`,`$ where $`K_4=\mathrm{diag}(1,1,1,1)`$. Written in an explicit form,
$$\omega (a)=\left[\begin{array}{cccccccc}a_0& \hfill a_1& \hfill a_2& \hfill a_3& \hfill a_4& \hfill a_5& \hfill a_6& \hfill a_7\\ a_1& \hfill a_0& \hfill a_3& \hfill a_2& \hfill a_5& \hfill a_4& \hfill a_7& \hfill a_6\\ a_2& \hfill a_3& \hfill a_0& \hfill a_1& \hfill a_6& \hfill a_7& \hfill a_4& \hfill a_5\\ a_3& \hfill a_2& \hfill a_1& \hfill a_0& \hfill a_7& \hfill a_6& \hfill a_5& \hfill a_4\\ a_4& \hfill a_5& \hfill a_6& \hfill a_7& \hfill a_0& \hfill a_1& \hfill a_2& \hfill a_3\\ a_5& \hfill a_4& \hfill a_7& \hfill a_6& \hfill a_1& \hfill a_0& \hfill a_3& \hfill a_2\\ a_6& \hfill a_7& \hfill a_4& \hfill a_5& \hfill a_2& \hfill a_3& \hfill a_0& \hfill a_1\\ a_7& \hfill a_6& \hfill a_5& \hfill a_4& \hfill a_3& \hfill a_2& \hfill a_1& \hfill a_0\end{array}\right],$$
$`(2.2)`$
Theorem 2.1.Let $`x=x_0+x_1e_1+\mathrm{}+x_7e_7𝕆`$, and denote $`\stackrel{}{x}=[x_0,x_1,\mathrm{},x_7]^T,`$ called the vector representation of $`x`$. Then
$$\stackrel{}{ax}=\omega (a)\stackrel{}{x}$$
$`(2.3)`$
holds for $`a,x𝕆`$.
Proof. Write $`a,x𝕆`$ as $`a=a^{}+a^{\prime \prime }e,x=x^{}+x^{\prime \prime }e,`$ where $`a^{},a^{\prime \prime },x^{},x^{\prime \prime }`$. We know by Eq.(1.3) that $`ax=(a^{}x^{}\overline{x^{\prime \prime }}a^{\prime \prime })+(x^{\prime \prime }a^{}+a^{\prime \prime }\overline{x^{}})e.`$ Thus it follows by Eq.(1.20) that
$`\stackrel{}{ax}=\left[\begin{array}{c}\stackrel{}{a^{}x^{}\overline{x^{\prime \prime }}a^{\prime \prime }}\\ \stackrel{}{x^{\prime \prime }a^{}+a^{\prime \prime }\overline{x^{}}}\end{array}\right]`$ $`=`$ $`\left[\begin{array}{c}\stackrel{}{a^{}x^{}}\stackrel{}{\overline{x^{\prime \prime }}a^{\prime \prime }}\\ \stackrel{}{x^{\prime \prime }a^{}}+\stackrel{}{a^{\prime \prime }\overline{x^{}}}\end{array}\right]`$
$`=`$ $`\left[\begin{array}{c}\varphi (a^{})\stackrel{}{x^{}}\tau (a^{\prime \prime })K_4\stackrel{}{x^{\prime \prime }}\\ \tau (a^{})\stackrel{}{x^{\prime \prime }}+\varphi (a^{\prime \prime })K_4\stackrel{}{x^{}}\end{array}\right]`$
$`=`$ $`\left[\begin{array}{cc}\varphi (a^{})& \tau (a^{\prime \prime })K_4\\ \varphi (a^{\prime \prime })K_4& \tau (a^{})\end{array}\right]\left[\begin{array}{c}\stackrel{}{x^{}}\\ \stackrel{}{x^{\prime \prime }}\end{array}\right],`$
as required for Eq.(2.3). $`\mathrm{}`$
Theorem 2.2.Let $`a𝕆`$ be given. Then
$$aE_8=E_8\omega (a),andE_8^{}a=\omega (a)E_8^{},$$
$`(2.4)`$
where $`E_8:=[\mathrm{\hspace{0.17em}1},e_1,\mathrm{},e_7],`$ and $`E_8^{}:=[\mathrm{\hspace{0.17em}1},e_1,\mathrm{},e_7]^T.`$
Proof. Follows from a direct verification. $`\mathrm{}`$
We can also introduce from Eq.(2.1) another matrix representation for an octonion as follows.
Definition 2.2. Let $`a=a^{}+a^{\prime \prime }e=a_0+a_1e_1+\mathrm{}+a_7e_7𝕆`$ be given, where $`a^{},a^{\prime \prime }`$. Then we call the $`8\times 8`$ real matrix
$$\nu (a):=K_8\omega ^T(a)K_8=\left[\begin{array}{cc}\tau (a^{})& \varphi (\overline{a^{\prime \prime }})\\ \varphi (a^{\prime \prime })& \tau (\overline{a^{}})\end{array}\right],$$
$`(2.5)`$
the right matrix representation of $`a`$, where $`K_8=\mathrm{diag}(K_4,I_4),`$ an orthogonal matrix. Written in an explicit form,
$$\nu (a)=\left[\begin{array}{cccccccc}a_0& \hfill a_1& \hfill a_2& \hfill a_3& \hfill a_4& \hfill a_5& \hfill a_6& \hfill a_7\\ a_1& \hfill a_0& \hfill a_3& \hfill a_2& \hfill a_5& \hfill a_4& \hfill a_7& \hfill a_6\\ a_2& \hfill a_3& \hfill a_0& \hfill a_1& \hfill a_6& \hfill a_7& \hfill a_4& \hfill a_5\\ a_3& \hfill a_2& \hfill a_1& \hfill a_0& \hfill a_7& \hfill a_6& \hfill a_5& \hfill a_4\\ a_4& \hfill a_5& \hfill a_6& \hfill a_7& \hfill a_0& \hfill a_1& \hfill a_2& \hfill a_3\\ a_5& \hfill a_4& \hfill a_7& \hfill a_6& \hfill a_1& \hfill a_0& \hfill a_3& \hfill a_2\\ a_6& \hfill a_7& \hfill a_4& \hfill a_5& \hfill a_2& \hfill a_3& \hfill a_0& \hfill a_1\\ a_7& \hfill a_6& \hfill a_5& \hfill a_4& \hfill a_3& \hfill a_2& \hfill a_1& \hfill a_0\end{array}\right].$$
$`(2.6)`$
Theorem 2.3.Let $`a,x𝕆`$ be given. Then
$$\stackrel{}{xa}=\nu (a)\stackrel{}{x}$$
$`(2.7)`$
holds.
Proof. Write $`a,x𝕆`$ as $`a=a^{}+a^{\prime \prime }e,x=x^{}+x^{\prime \prime }e,`$ where $`a^{},a^{\prime \prime },x^{},x^{\prime \prime }`$. we know by (1.3) that $`xa=(x^{}a^{}\overline{a^{\prime \prime }}x^{\prime \prime })+(a^{\prime \prime }x^{}+x^{\prime \prime }\overline{a^{}})e.`$ Thus we find by Eq.(1.20) that
$`\stackrel{}{xa}=\left[\begin{array}{c}\stackrel{}{x^{}a^{}\overline{a^{\prime \prime }}x^{\prime \prime }}\\ \stackrel{}{a^{\prime \prime }x^{}+x^{\prime \prime }\overline{a^{}}}\end{array}\right]`$ $`=`$ $`\left[\begin{array}{c}\stackrel{}{x^{}a^{}}\stackrel{}{\overline{a^{\prime \prime }}x^{\prime \prime }}\\ \stackrel{}{a^{\prime \prime }x^{}}+\stackrel{}{x^{\prime \prime }\overline{a^{}}}\end{array}\right]=\left[\begin{array}{c}\tau (a^{})\stackrel{}{x^{}}\varphi (\overline{a^{\prime \prime }})\stackrel{}{x^{\prime \prime }}\\ \varphi (a^{\prime \prime })\stackrel{}{x^{}}+\tau (\overline{a^{}})\stackrel{}{x^{\prime \prime }}\end{array}\right]=\left[\begin{array}{cc}\tau (a^{})& \varphi (\overline{a^{\prime \prime }})\\ \varphi (a^{\prime \prime })& \tau (\overline{a^{}})\end{array}\right]\left[\begin{array}{c}\stackrel{}{x^{}}\\ \stackrel{}{x^{\prime \prime }}\end{array}\right],`$
as required for Eq.(2.7). $`\mathrm{}`$
Theorem 2.4.Let $`a𝕆`$ be given. Then
$$aF_8=F_8\nu ^T(a),andF_8^{}a=\nu ^T(a)F_8^{},$$
$`(2.8)`$
where $`F_8:=[\mathrm{\hspace{0.17em}1},e_1,\mathrm{},e_7]`$ and $`F_8^{}:=[\mathrm{\hspace{0.17em}1},e_1,\mathrm{},e_7]^T.`$
Proof. Follows from a direct verification. $`\mathrm{}`$
Observe from Eqs.(2.1) and (2.5) that the two real matrix representations of an octonion $`a=a^{}+a^{\prime \prime }e`$ are in fact constructed by the real matrix representations of two quaternions $`a^{}`$ and $`a^{\prime \prime }`$. Hence the operation properties for the two matrix representations of octonions can easily be established through the results in Lemmas 1.2 and 1.3.
Theorem 2.5.Let $`a,b𝕆,\lambda `$ be given. Then
(a) $`a=b\omega (a)=\omega (b).`$
(b) $`\omega (a+b)=\omega (a)+\omega (b),\omega (\lambda a)=\lambda \omega (a),\omega (1)=I_8.`$
(c) $`\omega (\overline{a})=\omega ^T(a).`$
Proof. Follows from a direct verification. $`\mathrm{}`$
Theorem 2.6.Let $`a,b𝕆,\lambda `$ be given. Then
(a) $`a=b\nu (a)=\nu (b).`$
(b) $`\nu (a+b)=\nu (a)+\nu (b),\nu (\lambda a)=\lambda \nu (a),\nu (1)=I_8.`$
(c) $`\nu (\overline{a})=\nu ^T(a).`$
Proof. Follows from a direct verification. $`\mathrm{}`$
Theorem 2.7.Let $`a𝕆`$ be given. Then
$$a=\frac{1}{8}E_8\omega (a)E_8^{},anda=\frac{1}{8}F_8\nu ^T(a)F_8^{},$$
$`(2.9)`$
where $`E_8,E_8^{},F_8`$ and $`F_8^{}`$ are as in Eqs.(2.4) and (2.8).
Proof. Note that $`\omega (a)`$ and $`\nu (a)`$ are real matrices. Thus we get from Eqs.(2.4) and (2.8) that
$$E_8(E_8^{}a)=E_8[\omega (a)E_8^{}]=E_8\omega (a)E_8^{},andF(F_8^{}a)=F_8[\nu ^T(a)F_8^{}]=F_8\nu ^T(a)F_8^{}.$$
On the other hand, note that $`𝕆`$ is alternative. It follows that
$$E_8(E_8^{}a)=ae_1(e_1a)\mathrm{}e_7(e_7a)=ae_1^2a\mathrm{}e_7^2a=8a,$$
and
$$F_8(F_8^{}a)=ae_1(e_1a)\mathrm{}e_7(e_7a)=ae_1^2a\mathrm{}e_7^2a=8a.$$
Thus we have Eq.(2.9). $`\mathrm{}`$
Theorem 2.8.Let $`a𝕆`$ be given. Then
$$\mathrm{det}[\omega (a)]=\mathrm{det}[\nu (a)]=|a|^8.$$
$`(2.10)`$
Proof. Write $`a=a^{}+a^{\prime \prime }e`$. Then we easily find by Eqs.(1.21) and (2.5) that
$`\mathrm{det}[\omega (a)]=\mathrm{det}[\nu (a)]=\left|\begin{array}{cc}\tau (a^{})& \varphi (\overline{a^{\prime \prime }})\\ \varphi (a^{\prime \prime })& \tau (\overline{a^{}})\end{array}\right|`$ $`=`$ $`\mathrm{det}[\tau (a^{})\tau (\overline{a^{}})+\varphi (a^{\prime \prime })\varphi (\overline{a^{\prime \prime }})]`$
$`=`$ $`\mathrm{det}[\tau (\overline{a^{}}a^{})+\varphi (a^{\prime \prime }\overline{a^{\prime \prime }})]`$
$`=`$ $`\mathrm{det}[|a^{}|^2I_4+|a^{\prime \prime }|^2I_4]`$
$`=`$ $`(|a^{}|^2+|a^{\prime \prime }|^2)^4=|a|^8,`$
as required for Eq.(2.10). $`\mathrm{}`$
Theorem 2.9.Let $`a𝕆`$ be given. Then the two matrix representations of $`a`$ satisfy the following three identities
$$\omega (a^2)=\omega ^2(a),\nu (a^2)=\nu ^2(a),\omega (a)\nu (a)=\nu (a)\omega (a).$$
$`(2.11)`$
Proof. Applying Eqs.(2.3) and (2.7) to the both sides of the three identities in Eq.(1.9) leads to
$$\omega ^2(a)\stackrel{}{b}=\omega (a^2)\stackrel{}{b},\nu ^2(a)\stackrel{}{b}=\nu (a^2)\stackrel{}{b},\omega (a)\nu (a)\stackrel{}{b}=\nu (a)\omega (a)\stackrel{}{b}.$$
Note that $`\stackrel{}{b}`$ is an arbitrary $`8\times 1`$ real vector when $`b`$ runs over $`𝕆`$. Thus Eq.(2.11) follows. $`\mathrm{}`$
Theorem 2.10.Let $`a𝕆`$ be given with $`a0`$. Then
$$\omega (a^1)=\omega ^1(a),and\nu (a^1)=\nu ^1(a).$$
$`(2.12)`$
Proof. Note from Eqs.(1.10) and (1.11) that
$$a^1=\frac{\overline{a}}{|a|^2}=\frac{1}{|a|^2}[\mathrm{\hspace{0.17em}2}(\mathrm{Re}a)a]$$
and
$$a^22\mathrm{R}\mathrm{e}a+|a|^2=0.$$
Applying Theorems 2.5 and 2.6, as well as the first two equalities in Eq.(2.11) to the both sides of the above two equalities, we obtain
$$\omega (a^1)=\frac{1}{|a|^2}[\mathrm{\hspace{0.17em}2}(\mathrm{Re}a)I_8\omega (a)],\nu (a^1)=\frac{1}{|a|^2}[\mathrm{\hspace{0.17em}2}(\mathrm{Re}a)I_8\nu (a)]$$
and
$$\omega ^2(a)2(\mathrm{Re}a)\omega (a)+|a|^2I_8=0,\nu ^2(a)2(\mathrm{Re}a)\nu (a)+|a|^2I_8=0.$$
Contrasting them yields Eq.(2.12). $`\mathrm{}`$
Because $`𝕆`$ is non-associative, the operation properties $`\omega (ab)=\omega (a)\omega (b)`$ and $`\nu (ab)=\nu (b)\nu (a)`$ do not hold in general, otherwise $`𝕆`$ will be algebraically isomorphic to or algebraically anti-isomorphic to an associative matrix algebra over $``$, this is impossible. Nevertheless, some other kinds of identities on the two real matrix representations of octonions can still be established from the identities in Lemma 1.4(a)—(d).
Theorem 2.11.Let $`a,b𝕆`$ be given. Then their matrix representations satisfy the following two identities
$$\omega (aba)=\omega (a)\omega (b)\omega (a),and\nu (aba)=\nu (a)\nu (b)\nu (a).$$
$`(2.13)`$
Proof. Follows from applying Eqs.(2.3) and (2.7) to the Moufang identities in Lemma 1.4(b) . $`\mathrm{}`$
Theorem 2.12.Let $`a,b𝕆`$ be given. Then their matrix representations satisfy the following identities
$$\begin{array}{c}\text{ }\omega (ab)+\omega (ba)=\omega (a)\omega (b)+\omega (b)\omega (a),(2.14)\text{ }\hfill \\ \text{ }\nu (ab)+\nu (ba)=\nu (a)\nu (b)+\nu (b)\nu (a),(2.15)\text{ }\hfill \\ \text{ }\omega (ab)+\nu (ab)=\omega (a)\omega (b)+\nu (b)\nu (a),(2.16)\text{ }\hfill \\ \text{ }\omega (a)\nu (b)+\omega (b)\nu (a)=\nu (a)\omega (b)+\nu (b)\omega (a),(2.17)\text{ }\hfill \\ \text{ }\omega (ab)=\omega (a)\omega (b)+\omega (a)\nu (b)\nu (b)\omega (a),(2.18)\text{ }\hfill \\ \text{ }\nu (ab)=\nu (b)\nu (a)+\omega (b)\nu (a)\nu (a)\omega (b).(2.19)\text{ }\hfill \end{array}$$
Proof. The identities in Lemma 1.4(d) can clearly be written as the following six identities
$$(ab)xa(bx)=(ba)x+b(ax),(xa)bx(ab)=(xb)a+x(ba),$$
$$(ab)xa(bx)=(bx)a+b(xa),(ab)xa(bx)=(xa)b+x(ab),$$
$$(ab)xa(bx)=(ax)b+a(xb),(xa)bx(ab)=(ax)b+a(xb).$$
Applying Eqs.(2.3) and (2.7) to the both sides of the above identities, we obtain
$$\begin{array}{c}\text{ }[\omega (ab)\omega (a)\omega (b)]\stackrel{}{x}=[\omega (ba)+\omega (b)\omega (a)]\stackrel{}{x},\text{ }\hfill \\ \text{ }[\nu (b)\nu (a)\nu (ab)]\stackrel{}{x}=[\nu (a)\nu (b)+\nu (ba)]\stackrel{}{x},\text{ }\hfill \\ \text{ }[\nu (b)\omega (a)\omega (a)\nu (b)]\stackrel{}{x}=[\nu (a)\omega (b)+\omega (b)\nu (a)]\stackrel{}{x},\text{ }\hfill \\ \text{ }[\omega (ab)\omega (a)\omega (b)]\stackrel{}{x}=[\nu (b)\nu (a)+\nu (ab)]\stackrel{}{x},\text{ }\hfill \\ \text{ }[\omega (ab)\omega (a)\omega (b)]\stackrel{}{x}=[\nu (b)\omega (a)+\omega (a)\nu (b)]\stackrel{}{x},\text{ }\hfill \\ \text{ }[\nu (b)\nu (a)\nu (ab)]\stackrel{}{x}=[\nu (b)\omega (a)+\omega (a)\nu (b)]\stackrel{}{x}.\text{ }\hfill \end{array}$$
Notice that $`\stackrel{}{x}`$ is an arbitrary real $`8\times 1`$ real matrix when $`x`$ runs over $`𝕆`$. Therefore Eqs.(2.14)—(2.19) follow. $`\mathrm{}`$
Theorem 2.13.Let $`a,b𝕆`$ be given with $`a0,b0.`$ Then their matrix representations satisfy the following two identities
$$\omega (ab)=\nu (a)[\omega (a)\omega (b)]\nu ^1(a),and\nu (ab)=\omega (b)[\nu (b)\nu (a)]\omega ^1(b).$$
$`(2.20)`$
which imply that
$$\omega (ab)\omega (a)\omega (b),and\nu (ab)\nu (b)\nu (a).$$
$`(2.21)`$
Proof. Applying Eqs.(2.3) and (2.7) to the both sides of the two identities in Lemma 1.4(c), we obtain
$$\omega (ab)\nu (a)\stackrel{}{x}=\nu (a)\omega (a)\omega (b)\stackrel{}{x},and\nu (ab)\omega (b)\stackrel{}{x}=\omega (b)\nu (b)\nu (a)\stackrel{}{x},$$
which are obviously equvalent to Eq.(2.20). $`\mathrm{}`$
Note from Eqs.(2.3) and (2.7) that any linear equation of the form $`axxb=c`$ over $`𝕆`$ can equivqlently be written as $`[\omega (a)\nu (b)]\stackrel{}{x}=\stackrel{}{a}`$, which is a linear equation over $``$. Thus it is necessary to consider the operation properties of the matrix $`\omega (a)\nu (b)`$, especially the determinant of $`\omega (a)\nu (b)`$ for any $`a,b𝕆`$. Here we only list the expression of the determinant of $`\omega (a)\nu (b)`$. Its proof is quite tedious and is, therefore, omitted here.
Theorem 2.14.Let $`a,b𝕆`$ be given and define $`\delta (a,b):=\omega (a)\nu (b).`$ Then
$$\begin{array}{c}\text{ }\mathrm{det}[\delta (a,b)]=|a\overline{b}|^4[s^2+(|\mathrm{Im}a||\mathrm{Im}b|)^2][s^2+(|\mathrm{Im}a|+|\mathrm{Im}b|)^2](2.22)\text{ }\hfill \\ \text{ }\mathrm{det}[\delta (a,b)]=(s^2+|\mathrm{Im}a+\mathrm{Im}b|^2)^2[s^4+2s^2(|\mathrm{Im}a|^2+|\mathrm{Im}b|^2)+(|\mathrm{Im}a|^2|\mathrm{Im}b|^2)^2],(2.23)\text{ }\hfill \end{array}$$
where $`s=\mathrm{Re}a\mathrm{Re}b.`$ The characteristic polynomial of $`\delta (a,b)`$ is
$$\begin{array}{c}\text{ }|\lambda I_8\delta (a,b)|\text{ }\hfill \\ \text{ }=[(\lambda s)^2+|\mathrm{Im}a+\mathrm{Im}b|^2]^2[(\lambda s)^2+(|\mathrm{Im}a||\mathrm{Im}b|)^2][(\lambda s)^2+(|\mathrm{Im}a|+|\mathrm{Im}b|)^2].(2.24)\text{ }\hfill \end{array}$$
In particular, if $`\mathrm{Re}a=\mathrm{Re}b`$ and $`|\mathrm{Im}a|=|\mathrm{Im}b|,`$ but $`a\overline{b},`$ then
$$\mathrm{rank}\delta (a,b)=6.$$
$`(2.25)`$
Theorem 2.15.Let $`a,b𝕆`$ be given with $`a0`$ and $`b0.`$ Then $`\delta (a,b)=\omega (a)\nu (b)`$ is a real normal matrix over $`,`$ that is, $`\delta (a,b)\delta ^T(a,b)=\delta ^T(a,b)\delta (a,b)`$.
Proof. Follows from
$`\delta (a,b)+\delta ^T(a,b)`$ $`=`$ $`\omega (a)\nu (b)+\omega ^T(a)\nu ^T(b)`$
$`=`$ $`\omega (a)\nu (b)+\omega (\overline{a})\nu (\overline{b})`$
$`=`$ $`\omega (a+\overline{a})\nu (b+\overline{b})=2(\mathrm{Re}a\mathrm{Re}b)I_8.\mathrm{}`$
Theorem 2.16.Let $`a𝕆`$ be given with $`a`$. Then
$$\delta ^3(a,a)=4|\mathrm{Im}a|^2\delta (a,a),$$
$`(2.26)`$
and $`\delta (a,a)`$ has a generalized inverse as follows
$$\delta ^{}(a,a)=\frac{1}{4|\mathrm{Im}a|^2}\delta (a,a).$$
$`(2.27)`$
Proof. Observe that $`\delta (a,a)=\omega (a)\nu (a)=\omega (\mathrm{Im}a)\nu (\mathrm{Im}a)`$ and $`(\mathrm{Im}a)^2=|\mathrm{Im}a|^2`$. Thus we find that
$`\delta ^2(a,a)`$ $`=`$ $`[\omega (\mathrm{Im}a)\nu (\mathrm{Im}a)]^2`$
$`=`$ $`[\omega ^2(\mathrm{Im}a)2\omega (\mathrm{Im}a)\nu (\mathrm{Im}a)+\nu ^2(\mathrm{Im}a)]`$
$`=`$ $`[\omega ((\mathrm{Im}a)^2)2\omega (\mathrm{Im}a)\nu (\mathrm{Im}a)+\nu ((\mathrm{Im}a)^2)]`$
$`=`$ $`2[|\mathrm{Im}a|^2I_8+\omega (\mathrm{Im}a)\nu (\mathrm{Im}a)],`$
and
$`\delta ^3(a,a)`$ $`=`$ $`2[|\mathrm{Im}a|^2I_8+\omega (\mathrm{Im}a)\nu (\mathrm{Im}a)][\omega (\mathrm{Im}a)\nu (\mathrm{Im}a)]`$
$`=`$ $`4|\mathrm{Im}a|^2[\omega (\mathrm{Im}a)\nu (\mathrm{Im}a)]=4|\mathrm{Im},a|^2\delta (a,a),`$
as required for Eq.(2.26). $`\mathrm{}`$
3. Some linear equations over $`𝕆`$
The matrix expressions of octonions and their properties introduced in Section 2 enable us to easily deal with various problems related to octonions. One of the most fundamental topics on octonions is concerning solutions of various linear equations over $`𝕆`$. In this section, we shall give a complete discussion for this problem. Our first result is concerning the linear equation $`ax=xb`$, which was examined by the author in .
Thoerem 3.1. Let $`a=a_0+a_1e_1+\mathrm{}+a_7e_7,b=b_0+b_1e_1+\mathrm{}+b_7e_7𝕆`$ be given. Then the linear equation $`ax=xb`$ has a nonzero solution if and only if
$$\mathrm{Re}a=\mathrm{Re}band|\mathrm{Im}a|=|\mathrm{Im}b|.$$
$`(3.1)`$
(a) In that case, if $`b\overline{a},`$ i. e., $`\mathrm{Im}a+\mathrm{Im}b0,`$ then the general solution of $`ax=xb`$ can be expressed as
$$x=(\mathrm{Im}a)p+p(\mathrm{Im}b),$$
$`(3.1)`$
where $`p𝒜(a,b)`$, the subalgebra generated by $`a`$ and $`b`$, is arbitrary. or equivalently
$$x=\lambda _1(\mathrm{Im}a+\mathrm{Im}b)+\lambda _2[|\mathrm{Im}a||\mathrm{Im}b|(\mathrm{Im}a)(\mathrm{Im}b)],$$
$`(3.2)`$
where $`\lambda _1,\lambda _2`$ are arbitrary.
(b) If $`b=\overline{a},`$ then the general solution of $`ax=xb`$ is
$$x=x_1e_1+x_2e_2+\mathrm{}+x_7e_7,$$
$`(3.3)`$
where $`x_1`$$`x_7`$ satisfy $`a_1x_1+a_2x_2+\mathrm{}+a_7x_7=0.`$
The correctness of this result can be directly verified by substitution.
Based on the equation $`ax=xb`$, we can define the similarity of two octonions. Two octonions are said to be similar if there is a nonzero $`p𝕆`$ such that $`a=pbp^1`$, which is written as $`ab`$. Theorem 3.1 shows that two octonions are similar if and only if $`\mathrm{Re}a=\mathrm{Re}b`$ and $`|\mathrm{Im}a|=|\mathrm{Im}b|`$. Thus the similarity defined here is also an equivalence relation on octonions. In addition, we have the following.
Theorem 3.2.Let $`a,b𝕆`$ be given with $`b\overline{a}`$. Then
$$ab\omega (a)\omega (b).$$
$`(3.4)`$
Proof. Suppose first that $`ab`$. Then it follows by Eq.(1.11) that
$$a^22(\mathrm{Re}a)a=|a|^2=|b|^2=b^22(\mathrm{Re}b)b.$$
Applying Theorem 2.5(a) and Eq.(2.11) to the both sides of the above equality and we get
$$\omega ^2(a)2(\mathrm{Re}a)\omega (a)=\omega ^2(b)2(\mathrm{Re}b)\omega (b).$$
Thus
$$\omega ^2(a)+\omega (a)\omega (b)2(\mathrm{Re}a)\omega (a)=\omega ^2(b)+\omega (a)\omega (b)2(\mathrm{Re}b)\omega (b),$$
which is equivalent to
$$\omega (a)[\omega (a)+\omega (b)2(\mathrm{Re}a)I_8]=[\omega (a)+\omega (b)2(\mathrm{Re}b)I_8]\omega (b),$$
or simply
$$\omega (a)\omega (\mathrm{Im}a+\mathrm{Im}b)=\omega (\mathrm{Im}a+\mathrm{Im}b)\omega (b).$$
Note that $`\mathrm{Im}a+\mathrm{Im}b0`$. Thus $`\omega (\mathrm{Im}a+\mathrm{Im}b)`$ is invertible. The above equality shows that $`\omega (a)\omega (b)`$. Conversely, if $`\omega (a)\omega (b)`$, then trace $`\omega (a)=`$ trace $`\omega (b)`$ and $`|\omega (a)|=|\omega (b)|`$, which are equivalent to Eq.(3.1). $`\mathrm{}`$
Next we consider some nonhomogeneous linear equations over $`𝕆`$.
Theorem 3.3.Let $`a,b𝕆`$ be given with $`a`$. Then the linear equation $`axxa=b`$ has a solution in $`𝕆`$ if and only if The equality $`ab=b\overline{a}`$ holds. In this case, the general solution of $`axxa=b`$ is
$$x=\frac{1}{4|\mathrm{Im}a|^2}(baab)+p\frac{1}{|\mathrm{Im}a|^2}(\mathrm{Im}a)p(\mathrm{Im}a),$$
$`(3.5)`$
where $`p𝕆`$ is arbitrary.
Proof. According to Eqs.(2.3) and (2.7), the equation $`axxa=b`$ can equivalently be written as
$$[\omega (a)\nu (a)]\stackrel{}{x}=\delta (a,a)\stackrel{}{x}=\stackrel{}{b}.$$
$`(3.6)`$
This equation is solvable if and only if $`\delta (a,a)\delta ^{}(a,a)\stackrel{}{b}=\stackrel{}{b}.`$ In that case, the general solution of Eq.(3.6) can be expressed as
$$\stackrel{}{x}=\delta ^{}(a,a)\stackrel{}{c}+2[I_8\delta ^{}(a,a)\delta (a,a)]\stackrel{}{p},$$
where $`\stackrel{}{p}`$ is an arbitrary real vector. Substituting
$$\delta ^{}(a,a)=\frac{1}{4|\mathrm{Im}a|^2}\delta (a,a),and\delta ^2(a,a)=2[|\mathrm{Im}a|^2+\omega (\mathrm{Im}a)\nu (\mathrm{Im}a)]$$
in the above two equalities and then returning them to octonion forms by Eqs.(2.3) and (2.7) produce the equality in Part (b) and Eq.(3.5). $`\mathrm{}`$
Theorem 3.4.Let $`a=a_0+a_1e_1+\mathrm{}+a_7e_7,b=b_0+b_1e_1+\mathrm{}+b_7e_7𝕆`$ be given with $`a.`$ Then the equation
$$axx\overline{a}=b$$
$`(3.7)`$
has a solution if and only if there exist $`\lambda _0,\lambda _1`$ such that
$$b=\lambda _0+\lambda _1a,$$
$`(3.8)`$
in which case, the general solution of Eq.(3.7) is
$$x=\frac{\lambda _1}{2}+x_1e_1+\mathrm{}+x_7e_7,$$
$`(3.9)`$
where $`x_1`$$`x_7`$ satisfy
$$a_1x_1+\mathrm{}+a_7x_7=\frac{1}{2}\mathrm{Re}b.$$
$`(3.10)`$
Proof. According to Eqs.(2.3) and (2.7), the equation (3.7) is equivalent to
$$[\omega (a)\nu (\overline{a})]\stackrel{}{x}=\delta (a,\overline{a})\stackrel{}{x}=\stackrel{}{b},$$
$`(3.11)`$
namely
$$\left[\begin{array}{cccc}0& 2a_1& \mathrm{}& 2a_7\\ 2a_1& 0& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 2a_7& 0& \mathrm{}& 0\end{array}\right]\left[\begin{array}{c}x_0\\ x_1\\ \mathrm{}\\ x_7\end{array}\right]=\left[\begin{array}{c}b_0\\ b_1\\ \mathrm{}\\ b_7\end{array}\right].$$
Obviously, this equation is solvable if and only if there is a $`\lambda _1`$ such that
$$b_1=\lambda _1a_1,b_2=\lambda _1a_2,\mathrm{},b_7=\lambda _1a_7,$$
i. e., $`\mathrm{Im}\mathrm{b}=\lambda _1\mathrm{Im}\mathrm{a}`$, which is equivalent to Eq.(3.8). In that case, the solution to $`x_0`$ is $`x_0=\frac{\lambda _1}{2}`$, and $`x_1`$$`x_7`$ are determined by Eq.(3.9). $`\mathrm{}`$
Next we consider the linear equation
$$axxb=c$$
$`(3.12)`$
under the condition $`ab`$. Clearly Eq.(3.12) is equivalent to
$$[\omega (a)\nu (b)]\stackrel{}{x}=\delta (a,b)\stackrel{}{x}=\stackrel{}{c}.$$
$`(3.13)`$
Under $`ab`$, we know by Theorem 3.3 that $`ax=xb`$ has a nonzero solution. Hence $`\delta (a,b)`$ is singular under $`ab`$. In that case, Eq.(3.12) is solvable if and only if
$$\delta (a,b)\delta ^{}(a,b)\stackrel{}{c}=\stackrel{}{c},$$
$`(3.14)`$
and the general solution of Eq.(3.13) is
$$\stackrel{}{x}=\delta ^{}(a,b)\stackrel{}{c}+2[I_8\delta ^{}(a,b)\delta (a,b)]\stackrel{}{p},$$
$`(3.15)`$
where $`\stackrel{}{p}`$ is an arbitrary real vector. If $`a`$ is not not similar to $`b`$. Clearly Eq.(3.13) has a unique solution
$$\stackrel{}{x}=\delta ^1(a,b)\stackrel{}{c}$$
$`(3.16)`$
Eqs.(3.15) and (3.16) show that the solvability and solution of the octonion equation (3.12) can be completely determined by its real adjoint linear system of equations (3.13). Through the characteristic polynomial (2.24), one can also retern Eqs.(3.15) and (3.16) to octonion forms. But their expressions are quite tedious in form, and are omitted here.
Another instinctive linear equation over $`𝕆`$ is
$$a(xb)(ax)b=c,$$
$`(3.17)`$
which is also equivalent to
$$(ab)xa(bx)=c,$$
$`(3.18)`$
as well as
$$x(ab)(xa)b=c,$$
$`(3.19)`$
because $`(ab)xa(bx)=(ab)xa(bx)=x(ab)(xa)b`$ hold for all $`a,b,x𝕆`$. Now applying Eqs.(2.3) and (2.7) to the both sides of Eq.(3.17), we obtain an equivalent equation
$$[\omega (a)\nu (b)\nu (b)\omega (a)]\stackrel{}{x}=\stackrel{}{c}.$$
$`(3.20)`$
Here we set $`\mu (a,b)=\omega (a)\nu (b)\nu (b)\omega (a)`$. Then it is easy to see that Eq.(3.20) is solvable if and only if
$$\mu (a,b)\mu ^{}(a,b)\stackrel{}{c}=\stackrel{}{c},$$
where $`\mu ^{}(a,b)`$ is a generalized inverse of $`\mu (a,b)`$. In that case, the general solution of Eq.(3.20) is
$$\stackrel{}{x}=\mu ^{}(a,b)\stackrel{}{c}+[I_8\mu ^{}(a,b)\mu (a,b)]\stackrel{}{p},$$
$`(3.21)`$
where $`\stackrel{}{p}`$ is an arbitrary real vector. Numerical computation for Eq.(3.21) can reveal some interesting facts on Eq.(3.17). The reader can try to find them.
Theoretically speaking, any kind of two-sided linear equations or systems of linear equations over $`𝕆`$ can be equivalently transformed into systems of linear equations over $``$ by the two equalities in Eqs.(2.3) and (2.7). Thus the problems related to linear equations over $`𝕆`$ now have a complete resolution.
4. Real adjoint matrices of octonion matrices
In this section, we consider how to extend the work in Sections 2 and 3 to octonion matrices and use them to deal with various octonion matrix problems. Since octonion algebra is non-associative, the matrix operations in $`𝕆`$ is much different from what we are familiar with in an associative algebra. Even the simplest matrix multiplication rule $`A^2A=AA^2`$ does not hold over $`𝕆`$, that is to say, multiplication of matrices over $`𝕆`$ is completely not associative. Thus nearly all the known results and methods on matrices over associative algebras can hardly be extended to matrices over $`𝕆`$. In that case, a unique method available to deal with matrices over $`𝕆`$ is to establish real matrix representations of octonion matrices, and then to transform matrix problems over $`𝕆`$ to various equivalent real matrix problems.
Based on the two matrix representations of octonions shown in Eqs.(2.2) and (2.6), we now introduce two adjoints for a octonion matrix as follows.
Definition 4.1. Let $`A=(a_{st})𝕆^{m\times n}`$ be given . Then the left adjoint matrix of $`A`$ is defined to be
$$\omega (A)=[\omega (a_{st})]=\left[\begin{array}{ccc}\omega (a_{11})& \mathrm{}& \omega (a_{1n})\\ \mathrm{}& & \mathrm{}\\ \omega (a_{m1})& \mathrm{}& \omega (a_{mn})\end{array}\right]^{8m\times 8n},$$
$`(4.1)`$
the right adjoint matrix of $`A`$ is defined to be
$$\nu (A)=[\nu (a_{ts})]=\left[\begin{array}{ccc}\nu (a_{11})& \mathrm{}& \nu (a_{m1})\\ \mathrm{}& & \mathrm{}\\ \nu (a_{1n})& \mathrm{}& \nu (a_{mn})\end{array}\right]^{8n\times 8m},$$
$`(4.2)`$
and the adjoint vector of $`A`$ is defined to be
$$\mathrm{vec}A:=[\stackrel{}{a_{11}}^T,\mathrm{},\stackrel{}{a_{m1}}^T,\stackrel{}{a_{12}}^T,\mathrm{},\stackrel{}{a_{m2}}^T,\mathrm{},\stackrel{}{a_{1n}}^T,\mathrm{},\stackrel{}{a_{mn}}^T]^T.$$
$`(4.3)`$
Definition 4.2. Let $`A=(A_{st})_{m\times n}`$ and $`B=(B_{st})_{p\times q}`$ are two block matrices over $``$, where $`A_{st},B_{st}^{8\times 8}`$. Then the left and right block Kronecker products of $`A`$ and $`B`$, denoted respectively by $`A\widehat{}B`$ and $`A\stackrel{~}{}B`$, are defined to be
$$A\widehat{}B=\left[\begin{array}{ccc}A_{11}_LB& \mathrm{}& A_{1n}_LB\\ \mathrm{}& \mathrm{}& \mathrm{}\\ A_{m1}_LB& \mathrm{}& A_{mn}_LB\end{array}\right]^{8mp\times 8nq},$$
$`(4.4)`$
and
$$A\stackrel{~}{}B=\left[\begin{array}{ccc}A_RB_{11}& \mathrm{}& A_RB_{1q}\\ \mathrm{}& \mathrm{}& \mathrm{}\\ A_RB_{p1}& \mathrm{}& A_RB_{pq}\end{array}\right]^{8mp\times 8nq},$$
$`(4.5)`$
where
$$A_{st}_LB=\left[\begin{array}{ccc}A_{st}B_{11}& \mathrm{}& A_{st}B_{1q}\\ \mathrm{}& \mathrm{}& \mathrm{}\\ A_{st}B_{p1}& \mathrm{}& A_{st}B_{pq}\end{array}\right]^{8p\times 8q},$$
$`(4.6)`$
$$A_RB_{st}=\left[\begin{array}{ccc}A_{11}B_{st}& \mathrm{}& A_{1n}B_{st}\\ \mathrm{}& \mathrm{}& \mathrm{}\\ A_{m1}B_{st}& \mathrm{}& A_{mn}B_{st}\end{array}\right]^{8m\times 8n}.$$
$`(4.7)`$
Noticing the equality (2.5), we see the two adjoint matrices $`\omega (A)`$ and $`\nu (A)`$ of an octonion matrix $`A`$ satisfy the following equality
$$\nu (A)=K_{8n}\omega ^T(A)K_{8m},$$
$`(4.8)`$
where
$$K_{8t}=\mathrm{diag}(K_8,\mathrm{},K_8),K_8=\mathrm{diag}(\mathrm{\hspace{0.17em}1},1,\mathrm{},1),t=m,n.$$
$`(4.9)`$
It is easy to see from Eqs.(4.4) and (4.5) that the two kinds of block Kronecker products are actually constructed by replacing all elements in the standard Kronecker product of matrices with $`8\times 8`$ matrices. Hence the operation properties on these two kinds of products are much similar to those on the standard Kronecker product of matrices. We do not intend to list them here.
We next present some operation properties on the two real matrix representations of octonion matrices.
Theorem 4.1.Let $`A,B𝕆^{m\times n},\lambda `$ be given. Then
(a) $`A=B\omega (A)=\omega (B)\nu (A)=\nu (B),`$ i. e., $`\omega `$ and $`\nu `$ are 1-1.
(b) $`\omega (A+B)=\omega (A)+\omega (B),and\nu (A+B)=\nu (A)+\nu (B).`$
(c) $`\omega (\lambda A)=\lambda \omega (A),and\nu (\lambda A)=\lambda \nu (A).`$
(d) $`\omega (I_m)=I_{8m},and\nu (I_m)=I_{8m}.`$
(e) $`\omega (A^{})=\omega ^T(A),and\nu (A^{})=\nu ^T(A),`$ where $`A^{}=(\overline{a_{ts}})`$ is the conjugate transpose of $`A`$.
Theorem 4.2.Let $`A𝕆^{m\times n}`$ be given. Then
$$A=\frac{1}{8}E_{8m}\omega (A)E_{8n}^T,$$
$`(4.10)`$
where
$$E_{8t}=\mathrm{diag}(E_8,\mathrm{},E_8),andE_8=\mathrm{diag}(\mathrm{\hspace{0.17em}1},e_1,\mathrm{},e_7),t=m,n.$$
Proof. Follows directly from Corollary 2.7. $`\mathrm{}`$
Since the multiplication of matrices over $`𝕆`$ is completely not associative, no identities on products of octonions matrices can be established over $`𝕆`$ in general. Consequently, no identities on products of the two kinds of real matrix representations of octonion matrices can be established. In spit of this, we can still apply Eqs.(4.1) and (4.2) to deal with various problems related to octonion matrices. Next are some results on the relationship of $`\omega ()`$, $`\nu ()`$ and vec$`()`$ for matrices over $`𝕆.`$
Lemma 4.3.Let $`A𝕆^{n\times 1},B𝕆^{1\times n}`$ and $`x𝕆`$ be given. Then
$$\mathrm{vec}(Ax)=\omega (A)\stackrel{}{x}and\mathrm{vec}(xB)=\nu (B^T)\stackrel{}{x}.$$
$`(4.11)`$
Proof. Let $`A=[a_1,\mathrm{},a_n]^T`$ and $`B=[b_1,\mathrm{},b_n]^T`$. Then by Eqs.(2.3), (2.7) and Eqs.(4.1)—(4.3) we find
$$\mathrm{vec}(Ax)=\left[\begin{array}{c}\stackrel{}{a_1x}\\ \mathrm{}\\ \stackrel{}{a_nx}\end{array}\right]=\left[\begin{array}{c}\omega (a_1)\stackrel{}{x}\\ \mathrm{}\\ \omega (a_n)\stackrel{}{x}\end{array}\right]=\left[\begin{array}{c}\omega (a_1)\\ \mathrm{}\\ \omega (a_n)\end{array}\right]\stackrel{}{x}=\omega (A)\stackrel{}{x},$$
and
$$\mathrm{vec}(xB)=\left[\begin{array}{c}\stackrel{}{xb_1}\\ \mathrm{}\\ \stackrel{}{xb_n}\end{array}\right]=\left[\begin{array}{c}\nu (b_1)\stackrel{}{x}\\ \mathrm{}\\ \nu (b_n)\stackrel{}{x}\end{array}\right]=\left[\begin{array}{c}\nu (b_1)\\ \mathrm{}\\ \nu (b_n)\end{array}\right]\stackrel{}{x}=\nu (B^T)\stackrel{}{x}.\mathrm{}$$
Lemma 4.4.Let $`A𝕆^{m\times n},X𝕆^{n\times 1}`$ and $`a𝕆`$ be given. Then
$$\mathrm{vec}(AX)=\omega (A)\mathrm{vec}Xand\mathrm{vec}(Xa)=[\nu (a)\widehat{}I_{8n}]\mathrm{vec}X=\nu (a)\widehat{}\mathrm{vec}X.$$
$`(4.12)`$
Proof. Let $`A=[A_1,\mathrm{},A_n]`$ and $`X=[x_1,\mathrm{},x_n]^T`$. Then by Eq.(4.11) we find
$`\mathrm{vec}(AX)`$ $`=`$ $`\mathrm{vec}(A_1x_1+\mathrm{}+A_nx_n)`$
$`=`$ $`\mathrm{vec}(A_1x_1)+\mathrm{}+\mathrm{vec}(A_nx_n)`$
$`=`$ $`\omega (A_1)\mathrm{vec}x_1+\mathrm{}+\omega (A_n)\mathrm{vec}x_n`$
$`=`$ $`[\omega (A_1),\mathrm{},\omega (A_n)]\left[\begin{array}{c}\mathrm{vec}x_1\\ \mathrm{}\\ \mathrm{vec}x_n\end{array}\right]=\omega (A)\mathrm{vec}X,`$
as required for the first equality in (4.12). On the other hand,
$$\mathrm{vec}(Xa)=\left[\begin{array}{c}\stackrel{}{x_1a}\\ \mathrm{}\\ \stackrel{}{x_na}\end{array}\right]=\left[\begin{array}{c}\nu (a)\stackrel{}{x_1}\\ \mathrm{}\\ \nu (a)\stackrel{}{x_n}\end{array}\right]=[\nu (a)\widehat{}I_{8n}]\mathrm{vec}X=\nu (a)\widehat{}\mathrm{vec}X,$$
as required for the second equality in (4.12). $`\mathrm{}`$
Lemma 4.5.Let $`B𝕆^{p\times 1}`$ and $`X𝕆^{n\times p}`$ be given. Then
$$\mathrm{vec}(XB)=[\nu (B^T)\widehat{}I_{8n}]\mathrm{vec}X.$$
$`(4.13)`$
Proof. Let $`X=[X_1,\mathrm{},X_p]`$ and $`B=[b_1,\mathrm{},b_p]^T`$. Then it follows from the second equality in (4.12) that
$`\mathrm{vec}(XB)`$ $`=`$ $`\mathrm{vec}(X_1b_1+\mathrm{}+X_pb_p)`$
$`=`$ $`\mathrm{vec}(X_1b_1)+\mathrm{}+\mathrm{vec}(X_pb_p)`$
$`=`$ $`(\nu (b_1)\widehat{}I_{8n})\mathrm{vec}X_1+\mathrm{}+(\nu (b_p)\widehat{}I_{8n})\mathrm{vec}X_p`$
$`=`$ $`([\nu (b_1),\mathrm{},\nu (b_p)]\widehat{}I_{8n})\left[\begin{array}{c}\mathrm{vec}X_1\\ \mathrm{}\\ \mathrm{vec}X_p\end{array}\right]=[\nu (B^T)\widehat{}I_{8n}]\mathrm{vec}X,`$
as required for Eq.(4.13). $`\mathrm{}`$
Based on the above several lemmas, we can find the following three general results.
Theorem 4.6.Let $`A=(a_{st})𝕆^{m\times n}`$ and $`X𝕆^{n\times p}`$ be given. Then
$$\mathrm{vec}(AX)=[I_{8p}\widehat{}\omega (A)]\mathrm{vec}X.$$
$`(4.14)`$
Proof. Let $`X=[X_1,\mathrm{},X_p]`$. Then we find by Eq.(4.12) that
$`\mathrm{vec}(AX)`$ $`=`$ $`\mathrm{vec}[AX_1,\mathrm{},AX_p]`$
$`=`$ $`[\mathrm{vec}(AX_1),\mathrm{},\mathrm{vec}(AX_p)]`$
$`=`$ $`[\omega (A)\mathrm{vec}X_1,\mathrm{},\omega (A)\mathrm{vec}X_p]`$
$`=`$ $`\mathrm{diag}(\omega (A),\mathrm{},\omega (A))[\mathrm{vec}X_1,\mathrm{},\mathrm{vec}X_p]=[I_{8p}\widehat{}\omega (A)]\mathrm{vec}X,`$
establishing Eq.(4.14). $`\mathrm{}`$
Theorem 4.7.Let $`B=(b_{st})𝕆^{p\times q}`$ and $`X𝕆^{n\times p}`$ be given. Then
$$\mathrm{vec}(XB)=[\nu (B^T)\widehat{}I_{8n}]\mathrm{vec}X.$$
$`(4.15)`$
Proof. Let $`B=[B_1,\mathrm{},B_q]`$. Then we find by Eq.(4.13) that
$`\mathrm{vec}(XB)`$ $`=`$ $`\mathrm{vec}[XB_1,\mathrm{},XB_q]`$
$`=`$ $`\left[\begin{array}{c}\mathrm{vec}XB_1\\ \mathrm{}\\ \mathrm{vec}XB_q\end{array}\right]`$
$`=`$ $`\left[\begin{array}{c}\left[\nu (B_1^T)\widehat{}I_{8n}\right]\mathrm{vec}X\\ \mathrm{}\\ \left[\nu (B_q^T)\widehat{}I_{8n}\right]\mathrm{vec}X\end{array}\right]=\left[\begin{array}{c}\left[\nu (B_1^T)\widehat{}I_{8n}\right]\\ \mathrm{}\\ \left[\nu (B_q^T)\widehat{}I_{8n}\right]\end{array}\right]\mathrm{vec}X=[\nu (B^T)\widehat{}I_{8n}]\mathrm{vec}X,`$
as rerquired for Eq.(4.15). $`\mathrm{}`$
Theorem 4.8.Let $`A=(a_{st})𝕆^{m\times n},B=(b_{st})𝕆^{p\times q},`$ and $`X𝕆^{n\times p}`$ be given. Then
$$\mathrm{vec}[(AX)B]=[\nu (B^T)\widehat{}\omega (A)]\mathrm{vec}X,and\mathrm{vec}[A(XB)]=[\omega (A)\stackrel{~}{}\nu (B^T)]\mathrm{vec}X.$$
$`(4.16)`$
Proof. According to Eqs.(4.14) and (4.15), we find that
$`\mathrm{vec}[(AX)B]`$ $`=`$ $`[\nu (B^T)\widehat{}I_{8m}]\mathrm{vec}(AX)`$
$`=`$ $`[\nu (B^T)\widehat{}I_{8m}][I_{8p}\widehat{}\omega (A)]\mathrm{vec}X=[\nu (B^T)\widehat{}\omega (A)]\mathrm{vec}X,`$
and
$`\mathrm{vec}[A(XB)]`$ $`=`$ $`[I_{8p}\widehat{}\omega (A)]\mathrm{vec}(XB)`$
$`=`$ $`[I_{8p}\widehat{}\omega (A)][\nu (B^T)\widehat{}I_{8n}]\mathrm{vec}X=[\omega (A)\stackrel{~}{}\nu (B^T)]\mathrm{vec}X,`$
as required for Eq.(4.16). $`\mathrm{}`$
Theorem 4.9.Let $`A=(a_{st})𝕆^{n\times n},X=(b_{st})𝕆^{n\times p},Y𝕆^{q\times n}`$ be given, and denote
$$A^{(k|}X=A(A\mathrm{}(AX)\mathrm{})),andYA^{|k)}=((\mathrm{}(YA)\mathrm{})A)A.$$
Then
$$\mathrm{vec}(A^{(k|}X)=[I_{8p}\widehat{}\omega ^k(A)]\mathrm{vec}X,and\mathrm{vec}(YA^{|k)})=[\nu ^k(A^T)\widehat{}I_{8q}]\mathrm{vec}Y.$$
$`(4.17)`$
Just as the standard Kronecker products for matrices over any field, the three formulas in Eqs.(4.14)—(4.16) can directly be used for transforming any linear matrix equations over $`𝕆`$ into an ordinary linear system of equation over $``$. For example,
$$\begin{array}{c}\text{ }AX=B[I\widehat{}\omega (A)]\mathrm{vec}X=\mathrm{vec}B,\text{ }\hfill \\ \text{ }XA=B[\nu (A^T)\widehat{}I]\mathrm{vec}X=\mathrm{vec}B,\text{ }\hfill \\ \text{ }A(BX)=C[I\widehat{}\omega (A)\omega (B)]\mathrm{vec}X=\mathrm{vec}C,\text{ }\hfill \\ \text{ }(XA)B=C[\nu (B^T)\nu (A^T)\widehat{}I]\mathrm{vec}X=\mathrm{vec}C,\text{ }\hfill \\ \text{ }(AX)B=C[\nu (B^T)\widehat{}\omega (A)]\mathrm{vec}X=\mathrm{vec}C,\text{ }\hfill \\ \text{ }A(XB)=C[\omega (A)\stackrel{~}{}\nu (B)]\mathrm{vec}X=\mathrm{vec}C,\text{ }\hfill \\ \text{ }AXXB=C[I\widehat{}\omega (A)\nu (B)\widehat{}I]\mathrm{vec}X=\mathrm{vec}C,\text{ }\hfill \\ \text{ }(AX)AA(XA)=B[\nu (A^T)\widehat{}\omega (A)\omega (A)\stackrel{~}{}\nu (A^T)]\mathrm{vec}X=\mathrm{vec}B.\text{ }\hfill \end{array}$$
Theoreticlly speaking, various problems related to linear matrix equations over the octonion algebra now have a complete resolution.
Below are several simple results related to solutions of linear matrix equations over $`𝕆`$.
Definition 4.3. Let $`A𝕆^{n\times n}`$ be given. If its left adjoint matrix $`\omega (A)`$ is invertible, then $`A`$ is said to be completely invertible.
Theorem 4.10.Let $`A=(a_{st})𝕆^{m\times m}`$ and $`B=(b_{st})𝕆^{m\times n}`$ be given. If $`A`$ is completely invertible, then the matrix equation
$$AX=B,$$
$`(4.18)`$
has a unique solution over $`𝕆`$. In that case, if the real characteristic polynomial of $`\omega (A)`$ is
$$p(\lambda )=\lambda ^t+r_{t1}\lambda ^{t1}+\mathrm{}+r_1\lambda +r_0,$$
$`(4.19)`$
where $`r_0`$ is the determinant of $`\omega (A),`$ then the unique solution of Eq.(4.18) can be expressed as
$$X=\frac{1}{r_0}[A^{(t1|}B+r_{t1}(A^{(t2|}B)+\mathrm{}+r_3A(AB)+r_2AB+r_1B].$$
$`(4.20)`$
Proof. According to Eq.(4.14), the matrix equation (4.18) is equivalent to
$$[I_{8n}\widehat{}\omega (A)]\mathrm{vec}X=\mathrm{vec}B.$$
$`(4.21)`$
Because $`\omega (A)`$ is invertible, $`I_{8m}\widehat{}\omega (A)`$ is also invertible. Hence the solution of Eq.(4.25) is unique and this solution is
$$\mathrm{vec}X=[I_{8n}\widehat{}\omega (A)]^1\mathrm{vec}B=[I_{8m}\widehat{}\omega ^1(A)]\mathrm{vec}B.$$
Observe that
$$\omega ^t(A)+r_{t1}\omega ^{t1}(A)+\mathrm{}+r_1\omega (A)+r_0I_{8m}=0$$
holds. We then have
$$\omega ^1(A)=\frac{1}{r_0}\left[\omega ^{t1}(A)+r_{t1}\omega ^{t2}(A)+\mathrm{}+r_2\omega (A)+r_1I_{8m}\right].$$
Thus
$$I_{8n}\widehat{}\omega ^1(A)=\frac{1}{r_0}[I_{8n}\widehat{}\omega ^{t1}(A)+r_{t1}(I_{8n}\widehat{}\omega ^{t2}(A))+\mathrm{}+r_2(I_{8n}\widehat{}\omega (A))+(r_1I_{8n}\widehat{}I_{8m})],$$
and
$`\mathrm{vec}X=[I_{8n}\widehat{}\omega (A)]^1\mathrm{vec}B`$ $`=`$ $`{\displaystyle \frac{1}{r_0}}[(I_{8n}\widehat{}\omega ^{t1}(A))\mathrm{vec}B+r_{t1}(I_{8n}\widehat{}\omega ^{t2}(A))\mathrm{vec}B`$
$`+\mathrm{}+r_2(I_{8n}\widehat{}\omega (A))\mathrm{vec}B+r_1(I_{8n}\widehat{}I_{8m})\mathrm{vec}B].`$
Retuning it to octonion matrix expression by Eq.(4.17), we obtain Eq.(4.24). $`\mathrm{}`$
Similarly we have the following.
Theorem 4.11.Let $`A=(a_{st})𝕆^{m\times m}`$ and $`B=(b_{st})𝕆^{n\times m}`$ be given. If $`A`$ is completely invertible, then the matrix equation $`XA=B`$ has a unique solution over $`𝕆`$. In that case, if the real characteristic polynomial of $`\omega (A)`$ is
$$p(\lambda )=\lambda ^t+r_{t1}\lambda ^{t1}+\mathrm{}+r_1\lambda +r_0,$$
$`(4.22)`$
then the unique solution of $`XB=A`$ can be expressed as
$$X=\frac{1}{r_0}[BA^{|t1)}+r_{t1}(BA^{|t2)})+\mathrm{}+r_3(BA)A+r_2BA+r_1B].$$
$`(4.23)`$
For simplicity, the two solutions in Eqs.(4.20) and (4.23) can also be written as
$$X=L_A^1B,X=BR_A^1,$$
$`(4.24)`$
where $`L_A^1`$ and $`R_A^1`$ are, respectively, called the left and the right inverse operators of the completely invertible octonion matrix $`A`$. Some properties on these two inverse operators are listed below.
Theorem 4.12.Let $`A𝕆^{m\times m}`$ be an completely invertible matrix, $`B𝕆^{m\times n}`$ and $`C𝕆^{n\times m}`$ be given. Then
$$A(L_A^1B)=B,A(L_A^1I_m)=I_m,$$
$`(4.25)`$
$$L_A^1(AB)=B,L_A^1A=I_m,$$
$`(4.26)`$
$$(CR_A^1)A=C,(I_mR_A^1)A=I_m,$$
$`(4.27)`$
$$(CA)R_A^1=C,AR_A^1=I_m.$$
$`(4.28)`$
Proof. Follows from Theorems 4.10 and 4.11. $`\mathrm{}`$
We can also consider the inverses of octonion matrices in the usual sense. Let $`A𝕆^{m\times m}`$ be given. If there are $`X,Y𝕆^{m\times m}`$ such that $`XA=I_m`$ and $`AY=I_m`$, then $`X`$ and $`Y`$ are, respectively, called the left inverse and the right inverse of $`A`$, and denoted by $`A_L^1:=X`$ and $`A_R^1:=Y`$. From Theorems 4.10 and 4.11, we know that a square matrix of order $`m`$ over $`𝕆`$ has a left inverse if and only if the equation $`[\nu (A^T)\widehat{}I_{8m}]\mathrm{vec}X=\mathrm{vec}I_m`$ is solvable, and $`A`$ has a right inverse if and only if the equation $`[I_{8m}\widehat{}\omega (A)]\mathrm{vec}Y=\mathrm{vec}I_m`$ is solvable. These two facts imply that the left and the right inverses of a square matrix may not be unique, even both of them exist. As two special cases, we have the following.
Theorem 4.13.Let $`A𝕆^{m\times m}`$ be given. Then the left and the right inverses of $`A`$ are unique if and only if $`A`$ is completely invertible. In that case, if the real characteristic polynomial of $`\omega (A)`$ is
$$p(\lambda )=\lambda ^t+r_{t1}\lambda ^{t1}+\mathrm{}+r_1\lambda +r_0,$$
then the unique left and the unique right inverses $`A`$ can be expressed as
$$A_L^1=\frac{1}{r_0}[A^{(t1|}+r_{t1}A^{(t2|}+\mathrm{}+r_3A(A^2)+r_2A^2+r_1I_m],$$
and
$$A_R^1=\frac{1}{r_0}[A^{|t1)}+r_{t1}A^{|t2)}+\mathrm{}+r_3(A^2)A+r_2A^2+r_1I_m],$$
where $`A^{(s|}:=A(A(\mathrm{}(AA)\mathrm{}))`$ and $`A^{|s)}:=((\mathrm{}(AA)\mathrm{})A)A.`$
Proof. Follows directly from Theorems 4.10 and 4.11. $`\mathrm{}`$
Based on Theorems 4.10 and 4.12, as well as Eqs.(4.25)—(4.28), we can also derive the following two simple results.
Corollary 4.14.If $`A𝕆^{m\times m}`$ is completely invertible, and $`AB_1=AC_1`$ and $`B_2A=C_2A,`$ then $`B_1=C_1`$ and $`B_2=C_2.`$ In other words, the left and the right cancellation rules hold for completely invertible matrices.
Corollary 4.15.Suppose that $`A𝕆^{m\times m},B𝕆^{n\times n}`$ are completely invertible and $`C𝕆^{m\times n}`$. Then
(a) The matrix equation $`A(XB)=C`$ has a unique solution $`X=(L_A^1C)R_B^1.`$
(b) The matrix equation $`(AX)B=C`$ has a unique solution $`X=L_A^1(CR_B^1),`$
where $`L_A^1`$ and $`R_B^1`$ are the left and the right inverse operators of $`A`$ and $`B`$ respectively.
Our next result is concerned with the extension of the Cayley-Hamilton theorem to octonion matrices, which could be regarded as one of most successful applications of matrix representations of octonions.
Theorem 4.16.Let $`A𝕆^{m\times m}`$ be given and suppose that the real characteristic polynomial of $`\omega (A)`$ is
$$p(\lambda )=\lambda ^t+r_{t1}\lambda ^{t1}+\mathrm{}+r_1\lambda +r_0.$$
Then $`A`$ satisfies the following two identities
$$A^{(t|}+r_{t1}A^{(t1|}+\mathrm{}+r_3A(AA)+r_2A^2+r_1A+r_0I_m=0,$$
$`(4.29)`$
$$A^{|t)}+r_{t1}A^{|t1)}+\mathrm{}+r_3(AA)A+r_2A^2+r_1A+r_0I_m=0.$$
$`(4.30)`$
Proof. Observe that $`p[\omega (A)]=0`$. It follows that
$$[I_{8m}\widehat{}p[\omega (A)]]\mathrm{vec}I_m=0.$$
$`(4.31)`$
On the other hand, it is east to see by Eq.(4.17) that
$$\mathrm{vec}A^{(s|}=\mathrm{vec}(A^{(s|}I_m)=[I_{8m}\widehat{}\omega ^s(A)]\mathrm{vec}I_m,s=1,\mathrm{\hspace{0.17em}2},\mathrm{}.$$
Thus we find that
$$\begin{array}{c}\text{ }[I_{8m}\widehat{}p(\omega (A))]\mathrm{vec}I_m\text{ }\hfill \\ \text{ }=[I_{8m}\widehat{}\omega ^t(A)+r_{t1}(I_{8m}\widehat{}\omega ^{t1}(A))+\mathrm{}+r_1(I_{8m}\widehat{}\omega (A))+r_0(I_{8m}\widehat{}I_{8m})]\mathrm{vec}I_m\text{ }\hfill \\ \text{ }=(I_{8m}\widehat{}\omega ^t(A))\mathrm{vec}I_m+r_{t1}(I_{8m}\widehat{}\omega ^{t1}(A))\mathrm{vec}I_m+\mathrm{}+r_1(I_{8m}\widehat{}\omega (A))\mathrm{vec}I_m+r_0(I_{8m}\widehat{}I_{8m})\mathrm{vec}I_m\text{ }\hfill \\ \text{ }=\mathrm{vec}A^{(t|}+r_{t1}\mathrm{vec}A^{(t1|}+\mathrm{}+r_1\mathrm{vec}A+r_0\mathrm{vec}I_m\text{ }\hfill \\ \text{ }=\mathrm{vec}[A^{(t|}+r_{t1}A^{(t1|}+\mathrm{}+r_1A+r_0I_m].\text{ }\hfill \end{array}$$
The combination of this equality with Eq.(4.31) results in Eq.(4.29). The identity in Eq.(3.30) can be established similarly. $`\mathrm{}`$
Finally we present a result on real eigenvalues of Hermitian octonion matrices.
Theorem 4.17.Suppose that $`A𝕆^{m\times m}`$ is Hermitian, that is, $`A^{}=A`$. Then $`A`$ and its real adjoint $`\omega (A)`$ have identical real eigenvalues.
Proof. Since $`A=A^{}`$, we know by Theorem 4.1(e) that $`\omega (A)=\omega (A^{})=\omega ^T(A),`$ that is, $`\omega (A)`$ is a real symmetric matrix. In that case, all eigenvalues of $`\omega (A)`$ are real. Now suppose that
$$\omega (A)X=X\lambda ,$$
$`(4.32)`$
where $`\lambda `$ and $`X^{8m\times 1}`$. Then there is unique $`Y𝕆^{m\times 1}`$ such that $`\mathrm{vec}Y=X`$. In that case, it is easy to find by Theorem 4.1(a) and Eq.(4.12) that
$$\omega (A)X=X\lambda \omega (A)\mathrm{vec}Y=\mathrm{vec}Y\lambda \mathrm{vec}(AY)=\mathrm{vec}(Y\lambda )AY=Y\lambda ,$$
$`(4.33)`$
which implies that $`\lambda `$ is a real eigenvalue of $`A`$, and $`Y`$ is a eigenvector of $`A`$ corresponding to this $`\lambda `$. Conversely suppose that $`AY=Y\lambda ,`$ where $`\lambda `$, $`Y𝕆^{m\times 1}`$. Then taking vec operation on its both sides according to Eq.(4.12) yields
$$\omega (A)\mathrm{vec}Y=\mathrm{vec}Y\lambda .$$
This implies that $`\lambda `$ is also a real eigenvalue of $`\omega (A)`$ and $`\mathrm{vec}Y`$ is a real eigenvector of $`\omega (A)`$ associcated with this $`\lambda `$. $`\mathrm{}`$
The above result clearly shows that real eigenvalues and the corresponding eigenvectors of a Hermitian octonion matrix $`A`$ can all be determined by its real adjoint $`\omega (A)`$. Since $`\omega (A)`$ is a real symmetric $`8m\times 8m`$ matrix, it has $`8m`$ eigenvalues and $`8m`$ corresponding orthogonal eigenvectors.
Now a fundamental problem would naturally be asked: how many different real eigenvalues can a Hermitian octonion matrix $`A`$ have at most? For a $`2\times 2`$ Hermitian octonion matrix $`A=\left[\begin{array}{cc}a& b\\ \overline{b}& c\end{array}\right]`$, where $`a,c`$, its real adjoint is
$$\omega (A)=\left[\begin{array}{cc}aI_8& \omega (b)\\ \omega ^T(b)& cI_8\end{array}\right].$$
Clearly the characteristic polynomial of $`\omega (A)`$ is
$$\mathrm{det}(\lambda I_{16}\omega (A))=[(\lambda a)(\lambda c)|b|^2]^8.$$
This shows that $`\omega (A)`$, and correspondingly $`A`$, has 2 eigenvalues, each of which has a multiplicity 8.
The eigenvalue problem for $`3\times 3`$ Hermitian octonion matrices was recently examined by Dray and Manogue and Okubo . They showed by algebraic methods that every $`3\times 3`$ Hermitian octonion matrix has 24 real eigenvalues which are divided into 6 groups, each of them has multiplicity 4. Now according to Theorem 4.17, the real eigenvalues of any $`3\times 3`$ Hermitian octonion matrix
$$A=\left[\begin{array}{ccc}a_{11}& a_{12}& a_{13}\\ \overline{a}_{12}& a_{22}& a_{23}\\ \overline{a}_{13}& \overline{a}_{23}& a_{33}\end{array}\right],a_{11},a_{22},a_{33},$$
can be completely determined by its real adjoint
$$\omega (A)=\left[\begin{array}{ccc}\omega (a_{11})& \omega (a_{12})& \omega (a_{13})\\ \omega ^T(a_{12})& \omega (a_{22})& \omega (a_{23})\\ \omega ^T(a_{13})& \omega ^T(a_{23})& \omega (a_{33})\end{array}\right].$$
Obviously this matrix has 24 real eigenvalues and the $`24`$ corresponding real orthogonal eigenvectors. Numerical computation shows that these 24 eigenvalues are divided into 6 groups, each of them has multiplicity 4, which is consistent with the fact revealed in and . Moreover the 24 real orthogonal eigenvectors can also be converted to octonion expressions by (4.33).
Furthermore, numerical computation reveals an interesting fact that the 32 real eigenvalues any $`4\times 4`$ Hermitian octonion matrix are divided into 16 groups, each of them has multiplicity 2; the 40 real eigenvalues of any $`5\times 5`$ Hermitian octonion matrix are divided into 20 groups, each of them has multiplicity 2.
In general, we guess that for any $`m\times m`$ Hermitian octonion matrix with $`m>3`$, its $`8m`$ real eigenvalues can be divided into $`4m`$ groups, each of them has multiplicity 2.
As a subsequent work of Thereom 4.17, one might naturally ask hwo to establish a possible factorization for a Hermitian octonion matrix using its real eigenvalues and corresponding octonion orthogonal eigenvectors, speak more precisely, for an $`m\times m`$ Hermitian octonion matrix $`A`$, how construct a complete invertible octonion matrix $`P`$ (unitary?) and a real diagonal matrix $`D`$ such that $`A=PDP^1`$ using its $`8m`$ real eigenvalues and $`8m`$ corresponding octonion orthogonal eigenvectors. However, this problem seems quite curious, because the number of different real eigenvalues of an Hermitian octonion matrix is more than its order. This problem is also quite challenging, because various traditional methods in associative matrix theory are not applicable to this non-associative case.
As pointed out in , Hermitian octonion matrices can also have non-real right eigenvalues. Theoretically speaking, the non-real eigenvalue problem of Hermitian octonion matrices may also be converted to a problem related to real representations of octonion matrices. In fact, suppose that $`AX=X\lambda ,`$ where $`\lambda 𝕆`$ and $`X𝕆^{m\times 1}`$. Then according to Eq.(4.12), it is equivalent to
$$\omega (A)\mathrm{vec}X=\nu (\lambda )\widehat{}\mathrm{vec}X,$$
or alternatively
$$[\omega (A)\mathrm{diag}(\nu (\lambda ),\mathrm{},\nu (\lambda )]\mathrm{vec}X=0.$$
How to find $`\nu (\lambda )`$ satisfying the equation remains to further study.
Conclusions. In this paper, we have introduced two pseudo real matrix representations for octonions. Based on them we have made a complete investigation to their operation properties and have considered their various applications to octonions and matrices of octonions. However our work could only be regarded as a first step in the research of octonion matrix analysis and its applications. Numerous problems related to matrices of octonions remain to further examine, such as:
* How to determine eigenvalues and eigenvectors of a square octonion matrix, not necessarily Hermitian, and what is the relationship of eigenvalues and eigenvectors of a octonion matrix and its real adjoint matrices?
* Besides Eq.(4.29) and (4.30), how to establish some other identities for octonion matrices through their adjoint matrices?
* How to establish similarity theory for octonion matrices, and how to determine the relationship between the similarity of octonions matrices and the similarity of their adjoint matrices?
* How to consider various possible decompositions of octonion matrices, such as, LU decomposition, singular value decomposition and Schur decomposition?
* How to characterize various particular octonion matrices, such as, idempotent matrices, nipoltent matrices, involutary matrices, unitary matrices, normal matrices, and so on?
* How to define generalized inverses of octonion matrices when they are not completely invertible?
and so on. As mentioned in the beginning of the section, matrix multiplication for octonion matrices is completely not associative. In that case, any further research to problems related matrices of octonions is extremely difficult, but is also quite challenging. Any advance in solving the problems mentioned above could lead to remarkable new development in the real octonion algebra and its applications in mathematical physics.
Finally we should point out that the results obtained in the paper can use to establish pseudo matrix representations for real sedenions, as well as, in general, for elements in any $`2^n`$-dimensional real Cayley-Dickson algebras.
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# Magnetic Properties of HTSC with Weak Interlayer Coupling.
## I Introduction.
Scaling properties of layered high temperature superconductors (HTSC) around the mean–filed transition line $`H_{c2}(T)`$ has been under intensive experimental and theoretical study during the last decade. The layered YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-δ</sub> compounds exhibit 3D properties. All thermodynamic and transport quantities obey the 3D scaling law as a function of the scaling variable $`B(TT_{c2}(H))/(TH)^{2/3}`$, when the applied magnetic field is perpendicular to the layers plane. In the case of Bi<sub>2</sub>Sr<sub>2</sub>Ca<sub>2</sub>Cu<sub>3</sub>O<sub>10</sub> experimental results show 2D scaling behaviour with respect to the 2D scaling variable $`A(TT_{c2}(H))/(TH)^{1/2}`$. The fact that the phenomenon is observed in the vicinity of the mean–field upper critical field indicates that the problem can be studied in the framework of the Ginzburg–Landau (GL) theory with the order parameter projected onto the space of the lowest Landau levels (LLL). First calculations of the free energy scaling function were done perturbatively. The non-perturbative approach for pure 2D systems was developed by Tesanović et. al.. It was shown that around the mean–field transition line only the fluctuations of the total amplitude of the order parameter play important role in the superconductor thermodynamics. The remaining part, fluctuations in the position of the vortices enters through the dimensionless Abrikosov geometric factor $`\beta _A`$. In the most realistic cases this quantity weakly depends on the vortex configuration and can be taken as a constant from the very beginning. This approach was generalized for the layered 2D and 3D systems by Tesanović and Andreev. They showed that the system dimensionality depends on the strength of the interlayer coupling constant. When coupling is absent, the system becomes pure two dimensional. In the limit of strong interlayer coupling, when the superconducting correlation length $`\xi _c`$ along the $`c`$-axis direction becomes much larger then the effective interlayer separation $`s`$ the 3D description becomes appropriate. In the intermediate case $`\xi _c<s`$ the system becomes quasi two–dimensional for which scaling is impossible.
In this Paper we study quasi two dimensional layered superconductor in the region of the critical fluctuations. In this region we use GL-LLL description, which allows us to build a solvable model. The interlayer coupling is assumed to be small and therefore can be taken into account perturbatively. Within this approach we calculated the sample magnetization to the second order in the interlayer coupling constant. If this constant is extremely small one can consider only the first order correction. We show that because of the nearest neighboring coupling this correction leads to shift of the critical temperature, preserving the form of usual 2D scaling function of the magnetization. Only the next order correction destroys scaling and leads to disappearance of the magnetization crossing point. The effective interlayer coupling constant turns out to be proportional to the square root of the inverse applied field, which results in effective suppression of the interlayer interaction in high field region and, therefore, leads to recovering of the 2D scaling law.
## II The Model.
Consider a layered type II superconductor in the region of critical fluctuations around its mean-field transition line $`T_{c2}(H)`$ (or $`H_{c2}(T)`$). The applied magnetic field $`𝐇`$ is assumed to be normal to the layers plane: $`𝐇\widehat{c}`$. Then the superconductor thermodynamics at the temperature $`T`$ can be described by the following partition function:
$$𝒵𝒟[\mathrm{\Psi }]𝒟[𝐀]\mathrm{exp}\left\{\frac{[\mathrm{\Psi },𝐀]}{k_BT}\right\},$$
(1)
where $`[\mathrm{\Psi },𝐀]`$ is the GL free energy functional of the layered system with a nearest neighboring Josephson coupling between pancake vortices given by
$`[\mathrm{\Psi },𝐀]=s{\displaystyle \underset{n}{}}{\displaystyle }d^2𝐫\{\alpha _0|\mathrm{\Psi }_n|^2+{\displaystyle \frac{\beta }{2}}|\mathrm{\Psi }_n|^4+`$ (2)
$`\gamma _{ab}|_n\mathrm{\Psi }_n|^2+\gamma _c|\mathrm{\Psi }_n\mathrm{\Psi }_{n+1}|^2+{\displaystyle \frac{(𝐇𝐁_n)^2}{8\pi }}\}`$ (3)
and $`k_B`$ is the Boltzmann constant. The quantity $`s`$ is an effective interlayer spacing, $`\alpha _0=a(TT_{c0})`$ and $`\beta =const`$ are the first and the second GL coefficients correspondingly. The third GL coefficient $`\gamma `$ is assumed to be anisotropic, where the quantities $`\gamma _{ab}`$ and $`\gamma _c`$ define its value in the layer plane and $`\widehat{c}`$ direction correspondingly. In what follows, we refer to the quantity $`\gamma _c`$ as an interlayer coupling constant. The quantity $`\mathrm{\Psi }_n(𝐫)`$ is the order parameter of the $`n`$th layer and $`𝐁_n(𝐫)\widehat{c}`$ is the magnetic induction induced in the $`n`$th layer of the superconductor. Two-dimensional gauge invariant gradient $`_n`$ is defined as:
$$_n=i\mathrm{}\frac{}{𝐫}\frac{2e}{c}𝐀_n(𝐫),$$
where $`𝐀_n(𝐫)=\times 𝐁_n(𝐫)`$. The sample magnetization is given by
$$𝐌=\frac{1}{4\pi N_L}\underset{n}{}\frac{𝒟[\mathrm{\Psi },𝐀](𝐁_n𝐇)\mathrm{exp}\left\{{\displaystyle \frac{}{k_BT}}\right\}}{𝒟[\mathrm{\Psi },𝐀]\mathrm{exp}\left\{{\displaystyle \frac{}{k_BT}}\right\}},$$
(4)
where $`N_L`$ is the number of layers. In the weak coupling regime, it is convenient to introduce renormalized critical temperature:
$$T_c=T_{c0}2\frac{\gamma _c}{a}T_{c0}\delta T_c.$$
(5)
This modifies the expression (3) for the free energy, which now reads as
$`[\mathrm{\Psi },𝐀]=s{\displaystyle \underset{n}{}}{\displaystyle }d^2𝐫\{\alpha |\mathrm{\Psi }_n|^2+{\displaystyle \frac{\beta }{2}}|\mathrm{\Psi }_n|^4+`$ (6)
$`\gamma _{ab}|_n\mathrm{\Psi }_n|^2\gamma _c\mathrm{\Psi }_n^{}(\mathrm{\Psi }_{n1}+\mathrm{\Psi }_{n+1})+{\displaystyle \frac{(𝐇𝐁_n)^2}{8\pi }}\},`$ (7)
where $`\alpha =a(TT_c)`$. Thus, in order to solve the problem one has to calculate extremely complicated integrals over the order parameter $`\mathrm{\Psi }`$ and the vector potential $`𝐀`$ appearing in the expressions (1) and (4) for the partition function and the magnetization. However, there are number of simplifications, which can be done.
In the limit of large values of the GL parameter $`\kappa `$ ($`\kappa 100`$ for the most HTSC) one can neglect fluctuations of the magnetic induction $`𝐁_n`$ . Then we minimize the last expression (7) for the GL free energy with respect to the vector potential $`𝐀_n`$ as it is done in the case of conventional superconductor. This leads to the set of decoupled (with respect to the layer index) GL equations for the vector potential $`𝐀_n`$ and the order parameter $`\mathrm{\Psi }_n(𝐫)`$. In general, the order parameter can be expanded over the electron eigenfunctions of the Landau levels in the applied magnetic field $`H`$. However, close to the mean-field transition line $`H_{c2}(T)=\alpha (T)c/(2\mathrm{}e\gamma _{ac})`$ one can restrict oneself to the zeroth Landau level only. This approximation can be used, at least, if $`H>1/3H_{c2}(T)`$, until the first Landau level becomes important. Recently, it was shown that the lowest Landau level (LLL) approximation works good even if $`HH_{c2}(T)`$ for $`\kappa 1`$. In the LLL approximation the equations, described above, can be solved analytically . After substituting these solutions into the expression (7) for the free energy, we finally obtain:
$`=sS{\displaystyle \underset{n}{}}\{\alpha (1{\displaystyle \frac{H}{H_{c2}}})\overline{|\mathrm{\Psi }_n|^2}+{\displaystyle \frac{\beta }{2}}\overline{|\mathrm{\Psi }_n|^4}`$ (8)
$`\gamma _c\overline{\mathrm{\Psi }_n^{}(\mathrm{\Psi }_{n1}+\mathrm{\Psi }_{n+1})}\},`$ (9)
where the bar means averaging over the layer area $`S`$. In this case only the integrals over the order parameter $`\mathrm{\Psi }`$ remain in the expression for the partition function (1):
$$𝒵𝒟[\mathrm{\Psi }]\mathrm{exp}\left[\frac{[\mathrm{\Psi }]}{k_BT}\right],$$
where the quantity $`[\mathrm{\Psi }]`$ is given now by (9). In the LLL approximation the sample magnetization (4) can be calculated using the following formula:
$$M=\frac{\beta H_{c2}}{8\pi \alpha \kappa ^2}N_L^1\underset{n}{}\frac{𝒟[\mathrm{\Psi }]\overline{|\mathrm{\Psi }_n|^2}\mathrm{exp}\left({\displaystyle \frac{}{k_BT}}\right)}{𝒟[\mathrm{\Psi }]\mathrm{exp}\left({\displaystyle \frac{}{k_BT}}\right)}.$$
(10)
In order to proceed further, we define the Abrikosov geometric factor for each layer separately: $`\beta _A(n)\overline{|\mathrm{\Psi }_n|^4}/\left(\overline{|\mathrm{\Psi }_n|^2}\right)^2`$, and its average $`\beta _A=N_L^1\beta _A(n)`$. Following the refs. we assume that the quantity $`\beta _A(n)`$ only slightly depends on the actual vortex configuration and therefore we put $`\beta _A(1)=\beta _A(2)=\mathrm{}=\beta _A`$. Then we replace the quantities $`\overline{|\mathrm{\Psi }_n|^4}`$ by $`\beta _A\left(\overline{|\mathrm{\Psi }_n|^2}\right)^2`$ in the expression (9) for the free energy. With this replacement the model becomes exactly solvable:
$`𝒵{\displaystyle }𝒟[\mathrm{\Delta }]\mathrm{exp}\{N_v{\displaystyle \underset{n}{}}[x\overline{|\mathrm{\Delta }_n|^2}+{\displaystyle \frac{1}{4}}\left(\overline{|\mathrm{\Delta }_n|^2}\right)^2`$ (11)
$`\mu \overline{\mathrm{\Delta }_n^{}(\mathrm{\Delta }_{n+1}+\mathrm{\Delta }_{n1})}]\},`$ (12)
where $`N_v=\mathrm{\Phi }/\mathrm{\Phi }_0`$ is the number of vortices. The standard 2D scaling variable $`x`$ is given by
$$x=A\frac{TT_{c2}(H)}{\sqrt{TH}},$$
where $`A=\sqrt{s\mathrm{\Phi }_0/(16\pi \kappa ^2\beta _Ak_B)}H_{c2}^{}`$ and $`H_{c2}^{}=dH_{c2}(T)/dT|_{T=T_c}`$. The dimensionless coupling constant $`\mu `$ is
$$\mu =A\frac{\delta T_c}{2\sqrt{TH}}$$
(13)
and the dimensionless order parameter $`\mathrm{\Delta }_n`$ reads as
$$|\mathrm{\Delta }_n|^2=A\frac{2\beta _A\beta }{a\sqrt{TH}}|\mathrm{\Psi }_n|^2.$$
With these new variables the sample magnetization (10) can be expressed as
$$\frac{M}{\sqrt{HT}}=\frac{k_BA}{s\mathrm{\Phi }H_{c2}^{}}N_L^1\frac{d\mathrm{ln}𝒵}{dx}.$$
(14)
## III Calculation of the Partition Function.
In order to compute magnetization of the superconductor we, first, have to calculate the partition function (12). This involves evaluation of the integrals over the order parameter $`\mathrm{\Delta }`$. The main difficulty in such calculation is the quartic term appearing in the exponent of the right hand side of the formula (12). In order to decouple it, we introduce a set of additional integration variables $`\{\lambda _n\}`$:
$`𝒵{\displaystyle }𝒟[\mathrm{\Delta }]{\displaystyle _\mathrm{\Gamma }}𝒟[\lambda ]\mathrm{exp}\{N_v{\displaystyle \underset{n}{}}[\lambda _n^2+`$ (15)
$`(x+i\lambda _n)\overline{|\mathrm{\Delta }_n|^2}\mu \overline{\mathrm{\Delta }_n^{}(\mathrm{\Delta }_{n+1}+\mathrm{\Delta }_{n1})}]\},`$ (16)
where
$$_\mathrm{\Gamma }𝒟[\lambda ]=\underset{n}{}_{\mathrm{\Gamma }_n}𝑑\lambda _n.$$
The contours $`\mathrm{\Gamma }_n`$ are parallel to the real axis, standing on some distance from it, in order to insure convergence of the integrals over the order parameter $`\mathrm{\Delta }`$. The formula (16) is the result of use of the simplified version of the Hubbard-Stratonovich transformation, usually applied in the field theory. As it was explained above, the order parameter $`\mathrm{\Delta }_n`$ is the linear combination of the electron eigenfunctions of the lowest Landau level:
$$\mathrm{\Delta }_n(𝐫)=\underset{k=0}{\overset{N_v}{}}C_{nk}L_k(𝐫),$$
and
$$L_k(𝐫)=\frac{1}{\sqrt{k!}}\left(\frac{r}{l}\right)^k\mathrm{exp}\left\{ik\vartheta \frac{r^2}{2l^2}\right\},$$
where $`l`$ is the magnetic length corresponding to the charge $`2e`$. Then the meaning of the integration over the order parameter becomes clear:
$$𝒟[\mathrm{\Delta }]\underset{n,k}{}𝑑C_{nk}^{}𝑑C_{nk}.$$
The integrals over these expansion coefficients in (16) are of the generalized gaussian type and can be evaluated analytically. As a result, we obtain:
$$𝒵_\mathrm{\Gamma }𝒟[\lambda ]\mathrm{exp}\left\{N_v\right\}.$$
(17)
The action $``$ is given by
$$=\text{tr}\left[\widehat{\lambda }^2+\mathrm{ln}\left(x\widehat{\text{I}}+i\widehat{\lambda }\mu \widehat{\gamma }\right)\right],$$
where $`\widehat{\lambda }`$ is diagonal matrix, consisting from the elements $`\lambda _n`$. The matrix $`\widehat{\gamma }`$ is the real symmetric matrix, arising as a result of the interlayer coupling:
$$\gamma _{mn}=\delta _{m,n+1}+\delta _{m,n1}.$$
(18)
The integrals over the expansion coefficients $`C_{nk}`$ converge, if along the integration contours $`\mathrm{\Gamma }_n`$ the following inequality is satisfied:
$$\text{Re}(f_n)>0,$$
(19)
where $`f_n`$ are eigenvalues of the complex symmetric matrix $`x\widehat{\text{I}}+i\widehat{\lambda }\mu \widehat{\gamma }`$. In the thermodynamic limit $`N_v\mathrm{}`$ the integrals over $`\lambda _n`$ in the partition function (17) can be evaluated within saddle point approximation. It will be shown below that due to particular properties of the saddle-point manifold the condition (19) can be satisfied in the weak-coupling regime in rather large range of the scaling variable $`x`$.
The saddle point equation is found from the condition $`\delta =0`$ and reads as
$$\text{tr}\left\{\left[2\widehat{\lambda }+i\left(x\widehat{\text{I}}+i\widehat{\lambda }\mu \widehat{\gamma }\right)^1\right]\delta \widehat{\lambda }\right\}=0.$$
(20)
The position of the saddle point depends now on the temperature $`T`$ and the applied field $`H`$ via two parameters $`x(H,T)`$ and $`\mu (H,T)`$. In this case, 2D scaling of the magnetization (24) becomes impossible, as it was predicted in the ref. . If the coupling is weak $`2\mu <|x+i\lambda |`$ (this inequality is similar to Tesanović-Andreev criterion for quasi 2D systems), the second term in the left hand side of the equation (20) can be expanded in the powers of $`\mu `$. Then the saddle point equation (20) can be rewritten as follows:
$`2\lambda _n+{\displaystyle \frac{i}{x+i\lambda _n}}+`$ (21)
$`{\displaystyle \frac{i\mu ^2}{(x+i\lambda _n)^2}}{\displaystyle \frac{2x+i(\lambda _{n1}+\lambda _{n+1})}{(x+i\lambda _{n1})(x+i\lambda _{n+1})}}+o(\mu ^4)=0.`$ (22)
The terms of the order of $`\mu `$ and $`\mu ^3`$ drop out in the right hand side of the last equation, since $`\text{tr}(\widehat{\gamma }^{2n+1})=0`$. The structure of the saddle point equation is such that, at least, up to the second order in $`\mu `$ the saddle point solution for $`\widehat{\lambda }`$ is proportional to the unit matrix, namely
$$\lambda _1=\lambda _2=\mathrm{}\lambda .$$
Then the equation (22) can be rewritten in the following simple form:
$$2\lambda +\frac{i}{x+i\lambda }+\frac{2i\mu ^2}{(x+i\lambda )^3}+o(\mu ^4)=0.$$
(23)
As soon as the saddle point solution is found, the sample magnetization (14) can be calculated using the following simple relation:
$$\frac{M}{\sqrt{HT}}=\frac{2Ak_B}{s\mathrm{\Phi }_0H_{c2}^{}}i\lambda .$$
(24)
From the last equation we conclude that the physically meaningful solution for $`\lambda `$ lies on the imaginary axis. Further, using the Rayleigh-Ritz theorem , one can show that both matrices $`2\widehat{\text{I}}\widehat{\gamma }`$ and $`2\widehat{\text{I}}+\widehat{\gamma }`$ are positive definite. Then the eigenvalues of the matrix $`\widehat{\gamma }`$ belong to the interval $`[2,2]`$. In this case, the conditions (19) for convergence of the integrals over expansion coefficients $`C_{kn}`$ in the partition function can be written as:
$$2\mu <|x+i\lambda |,$$
(25)
which is automatically satisfied in the weak coupling regime.
## IV Magnetic Properties.
In the previous section we derived the saddle point equation (23) for the layered superconductor under assumption that the interlayer coupling constant $`\mu `$ is small. The sample magnetization $`M(H,T)`$ is proportional to the saddle point solution and is given by the formula (24).
In the ”zeroth” approximation only one of the two saddle point solutions of the eq. (23) can be reached by allowed deformation of the integration contour:
$$\lambda _0(x)=\frac{i}{2}(x\sqrt{x^2+2}).$$
Then the magnetization preserves the 2D scaling law and is given by the same formula as in the noninteracting case :
$$\frac{M_0}{\sqrt{HT}}=\frac{Ak_B}{s\mathrm{\Phi }_0H_{c2}^{}}(x\sqrt{x^2+2}).$$
(26)
Actually, this ”zeroth” order result includes the first order correction in $`\mu `$ by means of the critical temperature shift (5). The second order correction to the saddle point solution is calculated as a small perturbation:
$$\lambda (x,\mu )=\lambda _0(x)\left(1+4\mu ^2\frac{i\lambda _0(x)}{\sqrt{x^2+2}}\right).$$
(27)
Then the sample magnetization is given by
$$M=M_0\left(12\mu ^2\frac{x\sqrt{x^2+2}}{\sqrt{x^2+2}}\right),$$
(28)
where $`M_0(H,T)`$ is the magnetization of the ”decoupled” sample given by the equation (26). Using the last formula, we plotted the quantity $`M/\sqrt{HT}`$ as a function of 2D scaling variable $`x`$ (see the figure 1) for five different values of the applied field $`H`$ between $`10`$ and $`50kOe`$. The interlayer coupling constant is chosen to be small, such that $`\delta T_c=1K`$(see the eq. 5 for definition of $`\delta T_c`$). We used $`T_c=111K`$, $`\kappa =100`$, $`H_{c2}^{}=40kOeK^1`$ and $`s=2nm`$, which are of the order of the typical experimental parameters . The whole range of the scaling variable $`x`$ in the fig. 1 satisfies the applicability condition of the theory (25). As it was expected, the 2D scaling is destroyed, due to the dimensionless coupling constant $`\mu (T,H)`$ appearing in the right hand side of the expression (28) for the magnetization.
The plot of the magnetization as a function of temperature can be found in the fig. 2, in which we used the same set of the phenomenological parameters, as for the fig. 1. Like in the previous case, the whole temperature range in this figure satisfies the applicability condition (25). Form this figure follows that even account of the small interlayer coupling destroys the magnetization crossing point, which is the property of the scaling form (26) (see ref. ). Actually, the crossing point is noticeably destroyed only in the low–field region. The higher the applied field, the better the crossing point is pronounced. This is the consequence of the specific form of the effective coupling constant $`\mu `$ (see eq. (13)). The interaction correction to the magnetization in the formula (28) is proportional to the square of this constant and therefore is proportional to $`H^1`$. Then, at the large values of the applied field interlayer interaction is effectively suppressed. Indeed, one can treat the magnetization data plotted in the figure 2 as if they were obtained from the experiment. Then, we try to fit them to 2D scaling form variating phenomenological parameters $`T_c`$ and $`H_{c2}^{}`$, as it is usually done in experiments. The fitting results are given in figure 3. The best fit is obtained for $`T_c111.5K`$. Like in the experiment, the scaling results are insensitive to the value of the phenomenological parameter $`\mu _0H_{c2}^{}`$ in relatively large range of its values around $`4TK^1`$. It can be observed from the figure 3 that the scaling is satisfactory good for fields larger than $`30kOe`$ and is destroyed in the weak-field region. It must be stressed that the mentioned above fit is made essentially by hand and without any error estimation. However, it can serve as demonstration of effective suppression of interlayer interaction in the high field region.
## V Summary.
In this paper we considered HTSC assuming weak interlayer nearest neighboring coupling. In the region of the critical fluctuations, close to the mean–field transition line $`H_{c2}(T)`$ the order parameter can be taken as a linear combination of the electron eigenfunctions of the lowest Landau level. This approximation together with the assumption that the Abrikosov geometric factor only weakly depends on the actual vortex configuration allows to reduce the problem to the simpler one, namely to the saddle point equation (23). This equation can be solved perturbatively. We calculated the magnetization of the superconducting sample to the second order in the effective coupling constant $`\mu (T,H)`$. If the coupling is sufficiently weak, one can consider the first order correction only. It turns out that this correction does not modify the scaling properties of the sample, leading to the trivial renormalization of the critical temperature (see eq. (5)). Account of the second order correction leads to violation of the 2D scaling and destroys the magnetization crossing point in the low field region. At sufficiently high values of the applied field the interlayer interaction is effectively suppressed. This leads to recovering of 2D scaling with a well pronounced crossing point.
## Acknowledgments
The author would like to thank Sergey A. Gredeskul for helpful discussions.
The author gratefully acknowledges the MINERVA foundation for the financial support.
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# On the possibility of quantum computation based on photon exchange interactions
## I Introduction
The fact that photons almost do not interact with each other limits our ability to build photonic logical gates for two and more qubits. Nonlinear effects whereby one light beam influences another require large numbers of photons or else photon confinement in high-$`Q`$ cavities . Therefore a gate in which one photon-qubit would influence another is difficult to construct. This difficulty has motivated Franson et al. to search for a new effect that would enable the creation of a simple gate operating at the two-photon level.
Their basic proposal for the photonic computer is introduced in . Each qubit is represented by a photon which can travel along two alternate paths in an interferometer. Single qubit manipulation can be easily performed by setting proper phases in the relevant interferometer. The reading of the computed result is done by photodetection. The initial state should be prepared as a multi-mode product of single-photon Fock states. The authors of propose using quantum non-demolition photon number measurements to select such states from a larger set of coherent states.
The essential and most difficult part of the proposed scheme is an efficient two-bit quantum gate. Physically, the phase which one photon picks up in an interferometer should be determined by the path which the other photon chooses in another interferometer. For this to occur both interferometers should share one common branch in a nonlinear medium, such that when both photons travel along this branch an extra phase equal to $`\pi `$ would arise. This would require a giant cross-Kerr effect with negligible photon loss. The main advantage of such a scheme would be experimental simplicity, as compared to cavity-based schemes .
Here we examine the main assumptions and results of the schemes detailed in Ref. . We concentrate on the physical principle for the logic gates, leaving out the questions of the preparation and detection stages. In Sec. II an overview of these schemes is outlined. In Sec. III we present our detailed evaluation. Our conclusions are given in Sec. IV.
## II Overview of non-local photon exchange schemes
The underlying model of can be described as follows. Two light pulses simultaneously enter a medium of $`N`$ off-resonant atoms. The two pulses contain $`n_{1,2}`$ photons, respectively. Applying perturbation theory we can find the energy eigenvalues of the joint system of atoms and photons (within the dipole and rotating-wave approximation). Some of the fourth-order perturbation terms contain products of the photon numbers $`n_1n_2`$. Such terms can be characterized by Feynman diagrams in which a photon from pulse 1 is virtually absorbed by atom A and re-emitted by atom B, whereas a photon from pulse 2 is absorbed by atom B and re-emitted by atom A. The number of such terms is proportional to the number of various atomic pairs, i.e., $`N(N1)/2`$. If such terms contributed to the total energy, the physical result would be a nonlinear refractive index of the medium with very interesting properties: (i) the non-linear term would be proportional to the number of atoms squared $`N^2`$; (ii) the index of refraction caused by pulse 1 would be proportional to the photon number in pulse 2 and vice versa. However, after summing them up, all the terms containing $`n_1n_2`$ exactly cancel each other. We may ask whether such non-linear terms are just a mathematical artifact of perturbation theory or whether they correspond to some real physical situation. If the latter is true, the question is how to suppress some of the terms so that the remaining terms contribute to an experimentally observable non-linear effect.
Franson’s first suggestion for suppressing some of the $`n_1n_2`$ terms was to take advantage of collisional line broadening . The model used $`N`$ two-level atoms and the resulting non-linear part of the total atom-photon energy was claimed to be
$`\mathrm{\Delta }E{\displaystyle \frac{2M^4N^2n_1n_2f_R}{\delta ^3}}{\displaystyle \frac{w^2}{(\delta _1\delta _2)^2}},`$ (1)
where $`M`$ is the transition matrix element, $`f_R`$ ($`<1`$) is a factor taking into account decoherence due to a possible which-way information about the position of the photon absorption and re-emission, $`\delta _{1,2}`$ $``$ $`\delta `$ are the detunings of the modes 1 and 2 from the atomic resonance, and $`w`$ is the collisional line-width. For an efficient non-linear coupling between single-photon pulses one should find a sufficiently dense medium ($`N`$ large) and a broad collisional line-width, so that $`w^2/(\delta _1\delta _2)^2`$ $`1`$, which would yield a considerable non-linear energy shift.
The next suggestion was to manipulate the atomic resonance frequencies . The authors considered three-level atomic media, where strong laser pulses coupled to one of the atomic transitions would manipulate the resonance frequency of another transition, e.g., by Stark shifts. Thus, the medium would be turned on and off resonance with the incident photons, which would be absorbed and re-emitted in a controlled way.
In Ref. it is assumed that the photonic states would be coupled to collective $`N`$-atom excitations (Dicke states). The coupling frequency is then proportional to $`\sqrt{N}`$ and is assumed to be larger than any decay and decoherence rates. In both schemes with external driving pulses the authors claim that the photon of one kind would exhibit Rabi oscillations whose frequency depends on the presence or absence of the photon of the other kind. A proper choice of the external strong laser pulses is then claimed to effectively induce the required non-linear coupling of two single-photon pulses.
## III Detailed examination of the schemes
### A Collisional scheme
In the derivation of the results of there are several unjustified assumptions:
(1) The model of assumes two weak off-resonant light pulses propagating in a medium of two-level atoms. It is supposed that the effective Hilbert space describing the system is spanned by quantum states with different numbers of photons in the two relevant optical modes (1 and 2) and with different excited and de-excited atoms. All other optical modes are ignored: processes where photons can be re-emitted to modes other than 1 and 2 are disregarded. This is an arbitrary assumption: in open space photons are re-emitted into a continuum of modes, so that the main effect would be scattering. This would, of course, invalidate the potential application of the proposed effect as a quantum gate.
(2) Even if we go along with the model where only two optical modes are present, we cannot accept the main result Eq. (1). Its derivation is based on replacing in the fourth-order perturbation expansion the energy levels $`ϵ_m`$ which are influenced by collisions by the complex values $`ϵ_m`$ $``$ $`iw`$. However, in doing so, we always obtain zero for the non-linear $`n_1n_2`$ terms. The only way to obtain Eq. (1) is to assume that the states with no excited atoms ($`n_1\pm 1`$ photons in mode 1 and $`n_21`$ in mode 2) suffer from collisional decoherence and their energy levels should be modified by adding the $`iw`$ term. Of course, this assumption is not physically justified.
(3) Finally, even under the unlikely assumption that states with all atoms in the ground level decohere by collisions, the non-linear term (1) would be accompanied by an imaginary part
$`\mathrm{\Delta }E^{}=2if_RM^4N^2n_1n_2{\displaystyle \frac{w}{\delta ^2(\delta _1\delta _2)^2}},`$ (2)
whose magnitude is larger by a factor of $`\delta /w`$ than the real part. Hence, decay would always dominate any such non-linear phase shift and render the effect unobservable. We note that in Ref. the authors mention that there are difficulties with the collisional scheme.
### B Laser-induced nonlinear phase shift in “ladder” systems
In the scheme of one applies strong laser pulses which induce AC Stark shifts and thereby change the detuning of the near-resonant atomic transition from the relevant single-photon-carrying modes. Again, this model assumes that after photons 1 and 2 are absorbed by the atoms (not only virtually but also really, when the atoms are on resonance), they can only be re-emitted into the useful modes 1 and 2. Thus, it is assumed in that in open space, a single optical photon on resonance with the atomic medium can perform a Rabi oscillation without being scattered to the continuum of other modes. In other words, a resonant atom which absorbs the only photon from the traveling field would re-emit it to exactly the same (now empty) mode. In Ref. , no mechanism has been presented that would justify the assumed mode selectivity.
On the other hand, if the mode selectivity is guaranteed (e.g., by a cavity or by the Dicke cooperativity mechanism described in , see below), then for off resonant photons a weak non-linear phase shift may occur, whose magnitude is of the same order as their absorption probability. We think that a two-photon interference experiment revealing such a phase shift would be an interesting (though demanding) challenge. We feel, however, that such a phase shift would be too small to be useful for quantum logic gates.
### C The Dicke cooperative mechanism in Raman transitions
A mechanism supporting mode selectivity and the elimination of the mode continuum is presented in the e-print . It is argued there that a field mode of the wavevector $`𝐤`$ is effectively only coupled to a particular superposition of atomic excited states (Dicke state ), namely
$`|p(𝐤)={\displaystyle \frac{1}{\sqrt{N}}}{\displaystyle \underset{j}{}}\mathrm{exp}(i𝐤𝐫_j)|e_j,`$ (3)
where $`|e_j`$ determines the state with $`j`$th atom excited and all the other atoms being in the ground state; the summation runs over all $`N`$ atoms. The coupling between the field mode and the corresponding Dicke state is $`\sqrt{N}M`$, $`M`$ being the coupling between the field and a single atom. In an original extension of the standard formalism in the same cooperative enhancement is shown to apply to Raman transitions: The Dicke collective state will then be excited for the ground state by the operator
$`\widehat{R}_+(\stackrel{}{k}\stackrel{}{K})={\displaystyle \underset{j=1}{\overset{N}{}}}\widehat{R}_+^{(j)}e^{(\stackrel{}{k}\stackrel{}{K})\stackrel{}{r}_j}`$ (4)
where $`\stackrel{}{k}`$ and $`\stackrel{}{K}`$ are the wavevectors of the incident photon and the external laser field, respectively, and $`\widehat{R}_+^{(j)}`$ is the $`j`$th atom raising operator.
The spectacular feature of the Dicke formalism is that, in the absence of decoherence or losses, any distribution of atomic positions is guaranteed to have a state of maximal cooperation, such that the corresponding transition rate (Rabi frequency) is enhanced by $`\sqrt{N}`$ compared to that of a single atom. For Raman transitions we estimate that the cooperatively enhanced Rabi frequency is (in SI units)
$`\mathrm{\Omega }_{\mathrm{Raman}}^{(\mathrm{coop})}\sqrt{\rho \omega /(ϵ_0\mathrm{})}\mu \mathrm{\Omega }_0/\mathrm{\Delta }`$ (5)
where $`\rho `$ is the atomic density, $`\omega `$ the transition frequency, $`ϵ_0`$ the vacuum permittivity, $`\mu `$ the transition dipole moment, $`\mathrm{\Delta }`$ the detuning from the single-photon resonance and $`\mathrm{\Omega }_0=\mu E/\mathrm{}`$ is the strong-field Rabi frequency.
The crux of the effect is that during the relevant time, energy can only be exchanged between two states - a photon in a single field mode and a single-excitation atomic collective state (3), thus exhibiting Rabi oscillations. The decay and decoherence rates leading to scattering into other modes could be presumably slower, due to the $`\sqrt{N}`$ proportionality of the Rabi frequency. Because of linearity and the weak dependence of the Rabi frequency on the photonic frequency $`\omega `$, the same would be valid for photonic wavepackets. Under these assumptions, any sufficiently dense medium of resonant (or Raman-resonant) atoms could be transparent for single-photon light pulses: the photon would travel inside the medium “dressed” by the atomic excitations.
#### Evaluation
The cooperative Dicke effect discussed in is essentially the well-known excitonic enhancement of absorption and emission in crystals except that Raman transitions are discussed in instead of the standard direct transitions. Thus far, single-photon Rabi oscillations associated with cooperative (excitonic) enhancement have only been observed in semiconducting cavities , where they have the character of the single-mode Tavis-Cummings cooperative effect known for atoms in high-$`Q`$ cavities . By contrast, single-photon cooperative Rabi oscillations in open space of mode continuum suggested in have never been observed. The reason is that decoherence usually prevails, i.e., is faster than the achievable Rabi oscillation. This can be clarified using the estimate (5) in the two possible regimes:
(a) In the high-density regime, $`\rho /k^31`$, corresponding to interatomic distances $`r_{ij}\rho ^{1/3}`$ much smaller than the photon wavelength, the dominant source of decoherence detrimental to cooperation are resonant dipole-dipole interactions whose rate is $`\mathrm{\Omega }_{\mathrm{dip}}\gamma /(kr_{ij})^3`$, $`\gamma `$ being the radiative linewidth (for direct transitions), i.e., $`\mathrm{\Omega }_{\mathrm{dip}}`$ scales as $`\rho /k^3`$. For both direct transitions and Raman transitions Eq. (5) typically yields a lower rate than the dipole-dipole rate. Only for spatially symmetric atomic arrangements cooperative effects prevail over the dipole-dipole dephasing . For example, for interatomic distances of 10 nm $`\rho 10^{18}`$ cm<sup>-3</sup>, $`\mathrm{\Omega }_{\mathrm{Raman}}^{(\mathrm{coop})}10^{12}`$ s$`{}_{}{}^{1}\mathrm{\Omega }_{\mathrm{dip}}10^6\gamma `$.
(b) In the low-density regime $`\rho /k^31`$, the dipole-dipole rate is less than $`\gamma `$ and does not have to hamper cooperation. However, other sources of dephasing set $`T_2`$ (the decoherence time during which the off-diagonal density matrix elements go to zero) to be shorter than the cooperative Rabi period (typically $`\mathrm{\Omega }_{\mathrm{Raman}}^{(\mathrm{coop})}10^6`$ s<sup>-1</sup>) both in thermal gases ($`T_210^{10}`$ s) and in semiconductors ($`T_210^{12}`$ s).
However, the most important, ingenious assumption in concerns the initial condition for the atom-field system: the atoms are suddenly switched on resonance in the vicinity of the photon wavepacket, which is already within the medium. This is achieved by an appropriate geometry in which the photonic wavepacket overlaps with the strong rapidly-switched laser pulse, propagating perpendicularly to the single photon. Thus, only a photon being initially within the resonant medium can perform Rabi oscillations with the corresponding Dicke state. This is in contrast with the usual situation when a photon arrives at the medium which already has been resonant: the photon is then reflected or absorbed at the medium boundary, but cannot enter inside. We think that if this intriguing effect is realizable, its experimental observation would be very interesting. The following experiment could be planned: Photons would be sent into a transparent medium one by one. When inside the medium, a strong perpendicularly propagating pulse would bring the medium into resonance with the photons. The probability of detecting the photons as they exit the medium would be a periodic function of the duration of the strong pulse, thus demonstrating the single-photon Rabi oscillations.
Let us mention that several mechanisms of mapping the quantum state of light onto a collective state of atoms in the presence of a strong field have been suggested recently . The challenge of such mechanisms would be to allow coherent excitation and de-excitation of collective states with single photons. A detailed comparison of these intriguing approaches would be of great interest.
### D Two-photon entanglement and conditional phase shifts
Notwithstanding the chances of realizing the cooperative effect discussed above, the question is whether this effect can be used to produce the required conditional phase shift of the photonic states. In the authors argue that a proper sequence of external laser pulses would accomplish this task. Concrete suggestions for the pulse sequences were given in (five-pulse sequence), and in (three-pulse sequence). In the following we discuss these two suggestions separately.
#### 1 Five-pulse sequence
The first pulse brings the medium into resonance with photon 1 causing a $`\pi `$ Rabi transition: if photon 1 was initially present, it is absorbed creating a single-excitation Dicke state. The second pulse brings the medium into resonance with photon 2. The Rabi frequency now depends on whether the medium is excited or not. The authors of assume that the second pulse causes a $`2\pi `$ transition if there is no initial excitation (photon 1 is absent) or a $`\sqrt{2}2\pi `$ transition if there is initial medium excitation (photon 1 is present). In the latter case it is assumed that a superposition of two states is produced: a state with two atomic excitations and no photon, and a state with a single atomic excitation and a single photon in mode 2. The third pulse is used to produce a phase-shift in this superposition, and the remaining two pulses reverse the evolution of the first two pulses, to a state with the initial number of photons.
Our objection is that during the second pulse, also a state with two photons in mode 2 can be produced (as can be seen from the Hamiltonian, Eq. (31) of ): stimulated photon emission would occur with the same rate as photon absorption. Of course, the presence of such a two-photon state would invalidate the function of a quantum gate, where qubits are represented by single-photon states.
#### 2 Three-pulse sequence
The above flaw is removed in the scheme of : during the second pulse in the case of initially one atomic excitation and one photon in mode 2, the system oscillates between this state and the superposition of the state with two atomic excitations and no photons and the state with two photons and no atomic excitations. The frequency of this oscillation is twice as large as the Rabi frequency in the absence of the initial atomic excitation. Thus, if there was no photon in mode 1, the photon 2 exhibits a $`2\pi `$ Rabi transition during the second pulse, whereas if there was a photon in mode 1, the system exhibits a $`4\pi `$ Rabi transition. A $`2\pi `$ Rabi transition returns the original state with additional sign $`1`$, whereas a $`4\pi `$ Rabi transition simply reproduces the original state. The authors of claim that this difference in sign produces the required conditional phase shift.
Let us investigate in more detail the evolution of the states, denoting the basis of the logical gate as $`|0,0`$, $`|0,1`$, $`|1,0`$, and $`|1,1`$. Here, e.g., $`|1,0`$ means that photon 1 goes through the medium whereas photon 2 does not, etc. The evolution is then as follows. (a) State $`|0,0`$: there are no Rabi transitions, therefore $`|0,0`$ $``$ $`|0,0`$. (b) State $`|0,1`$: no change during the first pulse (photon 1 is absent), a $`2\pi `$ Rabi transition during the second pulse and no change during the third pulse, therefore $`|0,1`$ $``$ $`|0,1`$. (c) State $`|1,0`$: a $`\pi `$ transition during the first pulse, a $`2\pi `$ transition during the second pulse and a $`\pi `$ transition during the third pulse, therefore $`|1,0`$ $``$ $`|1,0`$. (d) State $`|1,1`$: a $`\pi `$ transition during the first pulse, a $`4\pi `$ transition during the second pulse and a $`\pi `$ transition during the third pulse, therefore $`|1,1`$ $``$ $`|1,1`$. Thus, the transformation matrix is
$`U=\left(\begin{array}{cccc}\hfill 1& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 1& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 1& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill 1\end{array}\right),`$ (10)
which is not the required conditional phase shift. This transformation does not entangle the two photons and is thus not suitable for building quantum logical gates.
## IV Conclusion
In the works we have not found a convincing proof that the suggested mechanisms could produce the conditional phase shift required for a quantum gate. On the other hand, we cannot disprove such a mechanism altogether, i.e., claim that it is principally impossible. Be it as it may, we find the idea of coupling the photonic state to an atomic Dicke state by fast switching, which would result in single-photon Rabi oscillations, very interesting and ingenious, even though the prospects for its realization are presently unclear.
Note added in proof: A proof has been given that exchange interactions between photons in a large ensemble of atoms cannot generate entanglement: if the process generates two photons in distinguishable modes, then the two-mode state can be factorized.
## Acknowledgments
We thank A. Ben-Reuven, I. Cirac, M. Fleischhauer, Ph. Grangier, A. Kofman, A. Kozhekin, M. Lukin, J. Peřina, S. Scheel, E. Schmidt, Y. Silberberg and D.-G. Welsch for stimulating discussions. This work was supported by ISF and DFG grants.
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# Cosmological Reionization
## 1 Introduction
At epochs corresponding to $`z1000`$ the intergalactic medium (IGM) is expected to recombine and remain neutral until sources of radiation and heat develop that are capable of reionizing it. The detection of transmitted flux shortward of the Ly$`\alpha `$wavelength in the spectra of sources at $`z5`$ implies that the hydrogen component of this IGM was ionized at even higher redshifts. There is some evidence that the double reionization of helium may have occurred later, but this is still controversial. It appears then that substantial sources of ultraviolet photons and mechanical energy were already present when the Universe was less than 7% of its current age, perhaps quasars and/or young star-forming galaxies: an episode of pre-galactic star formation may provide a possible explanation for the widespread existence of heavy elements (like carbon, oxygen, and silicon) in the IGM, while the integrated radiation emitted from quasars is likely responsible for the reionization of the intergalactic helium. Establishing the epoch of reionization and reheating is crucial for determining its impact on several key cosmological issues, from the role reionization plays in allowing protogalactic objects to cool and make stars, to determining the small-scale structure in the temperature fluctuations of the cosmic microwave background. Conversely, probing the reionization epoch may provide a means for constraining competing models for the formation of cosmic structures, and of detecting the onset of the first generation of stars, galaxies, and black holes in the Universe.
## 2 The transition from a neutral to an ionized Universe
Popular cosmological models predict that most of the intergalactic hydrogen was reionized by the first generation of stars or accreting black holes at $`z=715`$. One should note, however, that while numerical N-body$`+`$hydrodynamical simulations have convincingly shown that the IGM is expected to fragment into structures at early times in cold dark matter (CDM) cosmogonies (e.g. Cen et al. 1994; Zhang, Anninos, & Norman 1995; Hernquist et al. 1996), the same simulations are much less able to predict the efficiency with which the first gravitationally collapsed objects lit up the Universe at the end of the ‘dark age’ (Rees, this volume).
### 2.1 Photo- versus collisional ionization
The scenario that has received the most theoretical studies is one where hydrogen is photoionized by the UV radiation emitted either by quasars or by stars with masses $`\mathrm{}>10\mathrm{M}_{}`$, rather than ionized by collisions with electrons heated up by, e.g. supernova-driven winds from early pregalactic (‘Pop III’) objects. In the former case a high degree of ionization requires about $`13.6\times (1+t/\overline{t}_{\mathrm{rec}})`$eV per hydrogen atom, where $`\overline{t}_{\mathrm{rec}}`$ is the volume-averaged hydrogen recombination timescale, $`t/\overline{t}_{\mathrm{rec}}`$ being much greater than unity already at $`z10`$ according to the numerical simulations of Gnedin & Ostriker (1997), and Gnedin (2000). Collisional ionization to a neutral fraction of only few parts in $`10^5`$ requires a comparable energy input, i. e. an IGM temperature close to $`10^5`$K or about $`25`$eV per atom.
Massive stars will deposit both radiative and mechanical energy into the interstellar medium of Pop III objects. A complex network of ‘feedback’ mechanisms is likely at work in these systems, as the gas in shallow potential is more easily blown away thereby quenching further star formation (Mac Low & Ferrara 1999), and the blastwaves produced by supernova explosions reheat the surrounding intergalactic gas and enrich it with newly formed heavy elements and dust. It is therefore difficult to establish whether an early input of mechanical energy will actually play a major role in determining the thermal and ionization state of the IGM on large scales (Tegmark, Silk, & Evrard 1993). What can be easily shown is that, during the evolution of a a ‘typical’ stellar population, more energy is lost in ultraviolet radiation than in mechanical form. This is because in nuclear burning from zero to solar metallicity ($`Z_{}=0.02`$), the energy radiated per baryon is $`0.02\times 0.007\times m_\mathrm{H}c^2`$; about one third of it goes into H-ionizing photons. The same massive stars that dominate the UV light also explode as supernovae (SNe), returning most of the metals to the interstellar medium and injecting about $`10^{51}`$ergs per event in kinetic energy. For a Salpeter initial mass function (IMF), one has about one SN every $`150\mathrm{M}_{}`$ of baryons that forms stars. The mass fraction in mechanical energy is then approximately $`4\times 10^6`$, ten times lower than the fraction released in photons above 1 ryd.
The relative importance of photoionization versus shock ionization will depend, however, on the efficiency with which radiation and mechanical energy actually escape into the IGM. Consider, for example, the case of an early generation of halos with circular speed $`v_c=50\mathrm{km}\mathrm{s}^1`$, corresponding in top-hat spherical collapse to a virial temperature $`T_v=0.5\mu m_pv_c^2/k10^{5.3}`$K and halo mass $`M=0.1v_c^3/GH10^9[(1+z)/10]^{3/2}h^1\mathrm{M}_{}`$.<sup>1</sup><sup>1</sup>1This assumes an Einstein-de Sitter (EdS) Universe with $`H_0=100h\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$. In these systems rapid cooling by atomic hydrogen can take place and a significant fraction, $`f\mathrm{\Omega }_B`$, of their total mass may be converted into stars over a dynamical timescale (here $`\mathrm{\Omega }_B`$ is the baryon density parameter). For $`f=0.05`$, $`\mathrm{\Omega }_Bh^2=0.02`$, and $`h=0.5`$, the explosive output of $`50,000`$ SNe would inject an energy $`E_010^{55.7}`$ergs. The hot gas will escape its host, shock the IGM, and eventually form a cosmological blast wave. If the explosion occurs at cosmic time $`t=4\times 10^8`$yr, corresponding in the adopted cosmology (EdS with $`h=0.5`$) to $`z=9`$, at time $`\mathrm{\Delta }t=0.2t`$ after the event the proper radius of the (adiabatic) shock is given by the standard Sedov-Taylor self-similar solution,
$$R_s\left(\frac{12\pi GE_0}{\mathrm{\Omega }_B}\right)^{1/5}t^{2/5}\mathrm{\Delta }t^{2/5}23\mathrm{kpc}.$$
$`(1)`$
At this instant the shock velocity relative to the Hubble flow is
$$v_s2R_s/5\mathrm{\Delta }t110\mathrm{km}\mathrm{s}^1,$$
$`(2)`$
still much higher than the escape velocity from the halo. The gas temperature just behind the shock front is $`T_s=3\mu m_pv_s^2/16k4\times 10^5`$K, more than enough to efficiently ionize all the incoming hydrogen. At these redshifts, it is the onset of Compton cooling off cosmic microwave background photons that ends the adiabatic stage of blast wave propagation. According to the Press-Schechter formalism, the comoving abundance of collapsed dark halos with mass $`M=10^9h^1M_{}`$ at $`z=9`$ is $`dn/d\mathrm{ln}M5h^3`$Mpc<sup>-3</sup>, corresponding to a mean proper distance between neighboring halos of $`40h^1`$kpc, and to a total mass density parameter of order $`0.02`$. With the assumed star formation efficiency, only a small fraction, about one percent, of the stars seen today would have to be formed at these early epochs. Still, our simple analysis shows that the blast waves from such a population of pregalactic objects could overlap with large enough velocities to initially drive the intergalalactic medium to a significantly higher adiabat, $`T\mathrm{}>10^5`$K, than expected from photoionization, and pollute the entire IGM with metal-enriched material. A lower density of sources – which would therefore have to originate from higher amplitude peaks – would suffice if the typical efficiency of star formation were larger than assumed here.
Quasar-driven blast waves ($`E_0\mathrm{}>10^{60}`$ergs) are instead quite inefficient at ionizing the IGM, since much of the initial explosion energy is lost into the collisionless component (Voit 1996). They would also be too rare to fill the IGM without violating the COBE limit on the $`y`$-distortion of the microwave background.
### 2.2 Cosmological $`\mathrm{II}`$regions
In the following sections we will focus our attention to the photoionization of the IGM, i.e. we will assume that UV photons from an early generation of stars and/or quasars are the main source of energy for the reionization and reheating of the Universe, and that star formation and quasar activity occurs in collapsed galaxy halos. The process then begins as individual sources start to generate expanding $`\mathrm{II}`$regions in the surrounding IGM; throughout an $`\mathrm{II}`$region, H is ionized and He is either singly or doubly ionized. As more and more sources of ultraviolet radiation switch on, the ionized volume grows in size while the neutral phase shrinks. Reionization is completed when the $`\mathrm{II}`$regions overlap, and every point in the intergalactic space gets exposed for the first time to a nearly uniform Lyman-continuum (Lyc) background.
When an isolated point source of ionizing radiation turns on, an ionization (I) front separating the $`\mathrm{II}`$and $`\mathrm{I}`$regions propagates into the neutral gas, and the proper volume $`V_I`$ of the ionized zone grows according to the equation
$$\frac{dV_I}{dt}3HV_I=\frac{\dot{N}_{\mathrm{ion}}}{\overline{n}_\mathrm{H}}\frac{V_I}{\overline{t}_{\mathrm{rec}}},$$
$`(3)`$
(Shapiro & Giroux 1987), where $`\dot{N}_{\mathrm{ion}}`$ is the number of ionizing photons emitted by the central source per unit time, $`\overline{n}_\mathrm{H}(0)=1.7\times 10^7`$ $`(\mathrm{\Omega }_Bh^2/0.02)`$ cm<sup>-3</sup> is today’s mean hydrogen density, and all other symbols have their usual meaning. Most photons travel freely in the ionized gas, and are absorbed in a transition layer. In the case of stellar sources the I-front is quite sharp, and the degree of ionization changes on a short distance of the order of the mean free path for an ionizing photon. When $`\overline{t}_{\mathrm{rec}}t`$, the growth of the $`\mathrm{II}`$region is slowed down by recombinations in the highly inhomogeneous IGM, and its evolution can be decoupled from the Hubble expansion. Just like in the static case, the ionized bubble then fills its time-varying Strömgren sphere after a few recombination timescales,
$$V_I=\frac{\dot{N}_{\mathrm{ion}}\overline{t}_{\mathrm{rec}}}{\overline{n}_\mathrm{H}}(1e^{t/\overline{t}_{\mathrm{rec}}}).$$
$`(4)`$
While the volume that is ionized depends on the luminosity of the central source, the time it takes to produce an ionization-bounded region is only a function of $`\overline{t}_{\mathrm{rec}}`$.
In the presence of a population of ionizing sources, the transition from a neutral IGM to one that is almost fully ionized can be statistically described by the evolution with redshift of the volume filling factor (or porosity) $`Q`$ of $`\mathrm{II}`$, He $`\mathrm{II}`$, and He $`\mathrm{III}`$regions. The radiation emitted by spatially clustered stellar-like and quasar-like sources – the number densities and luminosities of which may change rapidly as a function of redshift – coupled with absorption processes in a medium that becomes more and more clumpy owing to the non-linear collapse of structures (Figure 1), all determine the complex topology of neutral and ionized zones in the Universe (Gnedin 2000; Ciardi et al. 2000; Abel, Norman, & Madau 1999). When $`Q1`$ and the radiation sources are randomly distributed, the ionized regions are spatially isolated, every UV photon is absorbed somewhere in the IGM, and the UV radiation field is highly inhomogeneous. As $`Q`$ grows, the crossing of ionization fronts becomes more and more common, until percolation occurs at $`Q=1`$.
Since the mean free path of Lyc radiation is always much smaller than the horizon (this is also true after ‘overlapping’ because of the residual $`\mathrm{I}`$still present in the Ly$`\alpha `$forest clouds and the Lyman-limit systems), the filling factor of cosmological $`\mathrm{II}`$regions is equal at any given time $`t`$ to the total number of ionizing photons emitted per hydrogen atom by all radiation sources present at earlier epochs, $`_0^t\dot{n}_{\mathrm{ion}}𝑑t^{}/\overline{n}_\mathrm{H}`$, minus the total number of radiative recombinations per atom, $`_0^tQ𝑑t^{}/\overline{t}_{\mathrm{rec}}`$. This statement reflects the simple fact that every ultraviolet photon that is emitted is either absorbed by a newly ionized hydrogen atom or by a recombining one. Differentiating one gets
$$\frac{dQ}{dt}=\frac{\dot{n}_{\mathrm{ion}}}{\overline{n}_\mathrm{H}}\frac{Q}{\overline{t}_{\mathrm{rec}}}$$
$`(5)`$
(Madau, Haardt, & Rees 1999). It is this differential equation – and its equivalent for expanding helium zones – that statistically describes the transition from a neutral Universe to a fully ionized one independently, for a given UV photon emissivity per unit cosmological volume $`\dot{n}_{\mathrm{ion}}`$, of the complex and possibly short-lived emission histories of individual radiation sources, e.g. on whether their comoving space density is constant or varies with cosmic time. Initially, when the filling factor is $`1`$, recombinations can be neglected and the ionized volume increases at a rate fixed solely by the ratio $`\dot{n}_{\mathrm{ion}}/\overline{n}_\mathrm{H}`$. As time goes on and more and more Lyc photons are emitted, radiative recombinations become important and slow down the growth of the ionized volume, until $`Q`$ reaches unity, the recombination term saturates, and reionization is finally completed (except for the high density regions far from any source which are only gradually eaten away, Miralda-Escudé, Haehnelt, & Rees 2000). In the limit of a fast recombining IGM ($`\overline{t}_{\mathrm{rec}}t`$), one can neglect the derivative on the left-hand side of equation (5) and derive
$$Q\mathrm{}<\frac{\dot{n}_{\mathrm{ion}}}{\overline{n}_\mathrm{H}}\overline{t}_{\mathrm{rec}},$$
$`(6)`$
i.e. the volume filling factor of ionized bubbles must be less (or equal) to the number of Lyc photons emitted per hydrogen atom in one recombination time. In other words, because of radiative recombinations, only a fraction $`\overline{t}_{\mathrm{rec}}/t1`$ of the photons emitted above 1 ryd is actually used to ionize new IGM material. The Universe is completely reionized when
$$\dot{n}_{\mathrm{ion}}\overline{t}_{\mathrm{rec}}\mathrm{}>\overline{n}_\mathrm{H},$$
$`(7)`$
i.e. when emission rate of ultraviolet photons exceeds the mean rate of recombinations.
### 2.3 A clumpy Universe
The simplest way to treat reionization in a inhomogeneous medium is in terms of a clumping factor that increases the effective gas recombination rate. In this case the volume-averaged recombination time is
$$\overline{t}_{\mathrm{rec}}=[(1+2\chi )\overline{n}_p\alpha _BC]^1=0.06\mathrm{Gyr}\left(\frac{\mathrm{\Omega }_Bh^2}{0.02}\right)^1\left(\frac{1+z}{10}\right)^3\frac{\overline{n}_\mathrm{H}}{\overline{n}_p}C_{10}^1,$$
$`(8)`$
where $`\alpha _B`$ is the recombination coefficient to the excited states of hydrogen (at an assumed gas temperature of $`10^4`$K), $`\chi `$ the helium to hydrogen abundance ratio, and the factor $`Cn_p^2/\overline{n}_p^2>1`$ takes into account the degree of clumpiness of photoionized regions (hereafter $`C_{10}C/10`$). If ionized gas with density $`n_p`$ filled uniformly a fraction $`1/C`$ of the available volume, the rest being empty space, the mean square density would be $`n_p^2=n_p^2/C=\overline{n}_p^2C`$. More in general, if $`f_m`$ is the fraction of baryonic mass in photoionized gas at an overdensity $`\delta `$ relative to the mean, and the remaining (underdense) medium is distributed uniformly, then the fractional volume occupied by the denser component is
$$f_v=f_m/\delta ,$$
$`(9)`$
the density of the diffuse component is
$$\overline{n}_p\frac{1f_m}{1f_v},$$
$`(10)`$
and the recombination rate is larger than that of a homogeneous Universe by the factor
$$C=f_m\delta +\frac{(1f_m)^2}{1f_v}$$
$`(11)`$
(e.g. Chiu & Ostriker 1999; Valageas & Silk 1999). It is difficult to estimate the clumping factor accurately. According to hydrodynamics simulations of structure formation in the IGM (within the framework of CDM-dominated cosmologies), Ly$`\alpha `$forest clouds with moderate overdensities, $`5\mathrm{}<\delta \mathrm{}<10`$, occupy a fraction of the available volume which is too small for them to dominate the clumping at high redshifts (e.g. Zhang et al. 1998; Theuns et al. 1998). In hierarchical clustering models, it is the virialized gas (with $`\delta 180`$ if one ignores the slope of the density profile) in dark matter halos with temperatures $`\mathrm{}<10^4`$K (masses $`M\mathrm{}<10^7h^1M_{}`$) which will plausibly boost the recombination rate by large factors as soon as the collapsed mass fraction exceeds 0.5%. Halos or halo cores which are dense and thick enough to be self-shielded from UV radiation will stay neutral and will not contribute to the recombination rate. This is also true of gas in more massive halos, which will be virialized to higher temperatures and ionized by collisions with thermal electrons. With a large comoving space density at $`z=9`$ of $`dn/d\mathrm{ln}M1000h^3`$Mpc<sup>-3</sup>, corresponding to a mean proper distance of only $`6h^1`$kpc, and to a mass fraction of $`0.04`$, halos with $`T_v10^4`$K will contribute significantly, $`f_m\delta 7`$, to the clumping. Recent calculations by Benson et al. (2000), which instead include all halos with $`T_v>10^4`$K and adopt an isothermal density profile with a flat core, give $`C30`$ already at $`z=9`$. Because of finite resolution effects, numerical simulations may underestimate clumping: in those of Gnedin & Ostriker (1997), for example, $`C`$ rises above unity at $`z\mathrm{}<20`$, and grows to $`C10`$ (40) at $`z9`$ (5).
It is important to note that the use of the volume-averaged clumping factor in the recombination timescale is only justified when the size of the $`\mathrm{II}`$regions is much larger compared to the scale of the clumping, so that the effect of many halos within the ionized volume can be averaged over. This will be a good approximation either at late epochs, when the $`\mathrm{II}`$zones have had time to grow (or when overlapping ionized regions from an ensemble of sources are able to proper sample the small-scale density fluctuations), or at earlier epochs if the ionized bubbles are produced by more luminous sources like quasars or the stars within halos collapsing from high-$`\sigma `$ peaks. As mentioned above, the mean free path between halos having $`T_v10^4`$K is $`\lambda 6h^1`$kpc at $`z=9`$, but their virial radius is only $`r_v0.4h^1`$kpc. It is only on scales greater than $`\lambda ^3/r_v^22h^1`$Mpc that the clumping can then be averaged over, and the covering factor of halos within the Strömgren sphere exceeds unity.
## 3 Sources of UV photons
### 3.1 Quasars
In recent years, several optical surveys (Warren, Hewett, & Osmer 1994; Schmidt, Schneider, & Gunn 1995; Kennefick, Djorgovski, & de Carvalho 1995) have consistently provided evidence for a turnover in the QSO counts. The space density of radio-loud quasars also appears to decrease strongly for $`z>3`$ (Shaver et al. 1996), suggesting that the turnover is indeed real and not an effect on optically-selected QSOs induced by dust along the line of sight. The density of optically bright and flat-spectrum radio-loud quasars has a relatively flat maximum at $`1.8\mathrm{}<z\mathrm{}<2.8`$, and declines gradually at higher redshifts (Figure 2).
The QSO emission rate of hydrogen Lyc photons per unit comoving volume, $`\dot{𝒩}_Q`$, is also shown in Figure 2. The procedure adopted to derive this quantity implies a large correction for incompleteness at high-$`z`$. With a fit to the quasar luminosity function (LF) which goes as $`\varphi (L)L^\beta `$, with $`\beta =1.64`$ at the faint end (Pei 1995), the contribution to the emissivity converges rather slowly, as $`L^{0.36}`$. At $`z=4`$, for example, the blue magnitude at the break of the LF is $`M_{}25.4`$, comparable or slightly fainter than the limit of current high-$`z`$ QSO surveys. While a large fraction, about 90% at $`z=4`$ and even higher at earlier epochs, of the ionizing emissivity shown in the figure is therefore produced by quasars that have not been actually observed, and are assumed to be present based on an extrapolation from lower redshifts, it is also fair to ask whether an excess of low-luminosity QSOs, relative to the best-fit LF, could actually boost the estimated Lyc emissivity at early epochs. The interest in models where the quasar LF significantly steepens with lookback time, and therefore predict many more QSOs at faint magnitudes than the extrapolation of Pei’s (1995) fitting functions, stems from recent claims of a strong linear correlation between bulge and observed black hole masses (Magorrian et al. 1998), linked to the steep mass function of dark matter haloes predicted by hierarchical cosmogonies (e.g. Haehnelt, Natarajan, & Rees 1998; Haiman & Loeb 1998). As discussed by Haiman, Madau, & Loeb (1999), the space density of low-luminosity quasars at high-$`z`$ is constrained by the observed lack of red, unresolved faint objects in the Hubble Deep Field (HDF). Down to a 50% completeness limit of $`V_{\mathrm{AB}}=29.6`$ ($`I_{\mathrm{AB}}=28.6`$), no $`z>4`$ quasar candidates have actually been found by Conti et al. (1999): by contrast, about 10 objects would be predicted by a QSO evolution model characterized by a steep LF with slope $`\beta =2`$ and a comoving space density that remains constant above $`z=2.5`$ instead of dropping (Figure 3), and choosen to boost the emission rate of ultraviolet photons at $`z5`$ by a factor of 5. A large population of faint AGNs at high-$`z`$ would still be consistent with the data if, at these faint magnitude levels and high image resolution, the host galaxies of active nuclei could actually be resolved by the Hubble Space Telescope.
### 3.2 Star-forming galaxies
Galaxies with ongoing star-formation are another obvious source of Lyc photons. The recent progress in our understanding of faint galaxy data made possible by the identification of star-forming galaxies at $`2\mathrm{}<z\mathrm{}<4`$ in ground-based surveys and in the HDF has provided new clues to the long-standing issue of whether galaxies at high redshifts can provide a significant contribution to the ionizing background flux. Since the rest-frame UV continuum at 1500 Å (redshifted into the visible band for a source at $`z3`$) is dominated by the same short-lived, massive stars which are responsible for the emission of photons shortward of the Lyman edge, the needed conversion factor, about one Lyc photon every 10 photons at 1500 Å, is fairly insensitive to the assumed IMF and is independent of the galaxy history for $`t10^{7.3}`$ yr.
Composite ultraviolet luminosity functions of Lyman-break galaxies (LBG) at $`z3`$ and $`z4`$ have been recently derived by Steidel et al. (1999). They are based on a large catalog of spectroscopically and photometrically selected galaxies from the ground-based and HDF samples, and span about a factor of 40 in luminosity from the faint to the bright end. Integrating these LF over all luminosities $`L>0.1L^{}`$, and using the conversion $`L(1500)/L(912)6`$ valid for a Salpeter mass function and constant star formation rate, we derive for the comoving emissivities at 1 ryd the values of $`9\pm 2\times 10^{25}h\mathrm{ergs}\mathrm{s}^1\mathrm{Hz}^1\mathrm{Mpc}^3`$ at $`z3`$, and $`7\pm 2\times 10^{25}h\mathrm{ergs}\mathrm{s}^1\mathrm{Hz}^1\mathrm{Mpc}^3`$ at $`z4`$, about 4 times higher than the estimated quasar contribution at $`z=3`$. These numbers do not include any correction for local $`\mathrm{I}`$absorption (since the color excess $`E_{9121500}`$ is expected to be small, dust exinction can probably be neglected in correcting from observed rest-frame far-UV to the Lyman edge). The data points plotted in Figure 2 assumes a value of $`f_{\mathrm{esc}}=0.5`$ for the unknown fraction of Lyc photons which escapes the dense sites of star formation (not included in our clumping factor) into the halos and the intergalalactic space. Note that, at $`z=3`$, Lyman-break galaxies radiate more ionizing photons than QSOs for $`f_{\mathrm{esc}}\mathrm{}>25\%`$.
## 4 Implications
### 4.1 First light
We have seen in the previous sections that, in the approximation the clumping can be averaged over, only the photons emitted within one recombination timescale can actually be used to ionize new material. As $`\overline{t}_{\mathrm{rec}}t`$ at high redshifts, it is possible to compute using equation (7) a critical value for the photon emission rate per unit cosmological comoving volume at a given epoch, $`\dot{𝒩}_c`$, independently of the (unknown) previous emission history of the Universe: only rates above this value will provide enough UV photons to keep the IGM ionized at that epoch. Equation (7) can then be rewritten as
$$\dot{𝒩}_c(z)=\frac{\overline{n}_\mathrm{H}(0)}{\overline{t}_{\mathrm{rec}}(z)}=(10^{51.4}\mathrm{s}^1\mathrm{Mpc}^3)C_{10}\left(\frac{1+z}{10}\right)^3\left(\frac{\mathrm{\Omega }_Bh^2}{0.02}\right)^2.$$
$`(12)`$
The uncertainty on this value is difficult to estimate, as it depends on the clumping factor and the nucleosynthesis constrained baryon density. It is interesting to convert this rate into a ‘minimum’ star formation rate per unit (comoving) volume, $`\dot{\rho }_{}`$ (for $`\mathrm{\Omega }_Bh^2=0.02`$):
$$\dot{\rho }_{}=\frac{\dot{𝒩}_c\times 10^{53.1}}{f_{\mathrm{esc}}}(0.12\mathrm{M}_{}\mathrm{yr}^1\mathrm{Mpc}^3)\left(\frac{0.5}{f_{\mathrm{esc}}}\right)C_{10}\left(\frac{1+z}{10}\right)^3.$$
$`(13)`$
(The conversion factor can be understood by noting that, for each 1 $`M_{}`$ of stars formed, 8% goes into massive stars with $`M>20M_{}`$ that dominate the Lyc luminosity of a stellar population. At the end of the C-burning phase, roughly half of the initial mass is converted into helium and carbon, with a mass fraction released as radiation of 0.007. About 25% of the energy radiated away goes into Lyc photons of mean energy 20 eV. For each 1 $`M_{}`$ of stars formed every year, we then expect $`0.08\times 0.5\times 0.007\times 0.25\times M_{}c^2/20\mathrm{eV}\mathrm{yr}10^{53}\mathrm{phot}\mathrm{s}^1`$ to be emitted shortward of 1 ryd.)
Taken at face value, equations (12) and (13) have perhaps a surprising implication. In a inhomogeneous Universe, early reionization at $`z9`$ requires an ionizing emissivity which is comparable or larger than that radiated by QSOs at the peak of their activity, $`z3`$. In a similar manner, photoionization by massive stars can only play a role if the star formation density at this epoch were significantly larger than the value directly ‘observed’ (i.e. uncorrected for dust reddening) at $`z=2`$ (Madau, Pozzetti, & Dickinson 1998).
### 4.2 Delayed He $`\mathrm{II}`$reionization
Because of its higher ionization potential and the steep spectra of UV radiation sources, the most abundant (by a factor $`100`$) absorbing ion in the post-reionization Universe is not $`\mathrm{I}`$but He $`\mathrm{II}`$. The importance of intergalactic helium in the context of this study stems from the possibility of detecting the effect of ‘incomplete’ He $`\mathrm{II}`$reionization in the spectra of $`z3`$ quasars as, depending on the clumpiness of the IGM (Madau & Meiksin 1994), the photoionization of singly ionized helium may be delayed until much later than for $`\mathrm{I}`$.
Since $`\mathrm{I}`$and He $`\mathrm{I}`$do not absorb a significant fraction of $`h\nu >54.4`$ eV photons, the problem of He $`\mathrm{II}`$reionization can be decoupled from that of other ionizations, and the equivalent of equation (5) for expanding He $`\mathrm{III}`$regions becomes
$$\frac{dQ}{dt}=\frac{\dot{n}_{\mathrm{ion4}}}{\overline{n}_{\mathrm{He}}}\frac{Q}{\overline{t}_{\mathrm{HeIII}}},$$
$`(14)`$
where $`\dot{n}_{\mathrm{ion4}}`$ now includes only photons above 4 ryd, and $`\overline{t}_{\mathrm{HeIII}}`$ is $`6.5`$ times shorter than the hydrogen recombination timescale if ionized hydrogen and doubly ionized helium have similar clumping factors.<sup>2</sup><sup>2</sup>2This last assumption appears, however, rather dubious: the reason is that self-shielding of He $`\mathrm{II}`$Lyc radiation occurs at much lower hydrogen columns than self-shielding of photons at $`1`$ryd (by about a factor of $`S/2`$, where the spectral ‘softness’ $`S`$ is the the ratio of the radiation flux at the hydrogen Lyman edge to the flux at 4 ryd), and self-shielded gas will remain neutral and not add to the recombination rate. Ionized hydrogen may then be more clumpy than doubly ionized helium. It is interesting to note that, if the intrinsic photon spectrum of ionizing sources has slope $`\dot{n}(\nu )\nu ^{2.8}`$, the first terms on the right-hand side of equations (5) and (14) are actually equal, and a significant delay between the complete overlapping of $`\mathrm{II}`$and He $`\mathrm{III}`$regions can only arise if recombinations are important. This effect is illustrated in Figure 4, where the expected evolution of the He $`\mathrm{III}`$filling factor (obtained by numerical integration of eq. 14) is plotted for a QSO-photoionization model with a source decline at high redshifts: He $`\mathrm{II}`$reionization is never completed before $`z=3`$ in models with $`C\mathrm{}>10`$. A significant contribution to the UV background at 4 ryd from massive stars, which could push the helium reionization epoch to higher redshifts, has been traditionally ruled out on the basis that the ratio between the number of He $`\mathrm{II}`$and $`\mathrm{I}`$Lyc photons emitted from low-metallicity starbursts is only about two percent (Leitherer & Heckman 1995), five times smaller than in typical QSO spectra. It has been recently pointed out by Tumlinson & Shull (2000), however, that metal-free stars exhibit higher effective temperatures and dramatically harder stellar spectra, particularly in the He $`\mathrm{II}`$continuum. This enhanced He-ionizing capabilities of Pop III stars could have interesting implications for reionization.
To date, various studies of the HeII Ly$`\alpha `$forest in the spectra of distant QSOs (Hogan, Anderson, & Rugers 1997; Reimers et al. 1997; Heap et al. 2000) have revealed patchy absorption with low He $`\mathrm{II}`$opacity ‘voids’ alternating several Mpc sized regions with vanishing flux. These observations suggest that helium absorption does not increase smoothly with lookback time, but rather in the abrupt manner expected in the final stages of inhomogeneous reionization by quasar sources. Radiative transfer effects during He $`\mathrm{II}`$reionization could affect the thermal history of the IGM (Abel & Haehnelt 1999, Efstathiou, this volume). Here it is important to remark that, while delayed He $`\mathrm{II}`$reionization in a clumpy Universe appears to be naturally linked to the observed decline in the space density of quasars beyond $`z3`$, the complete overlapping of He $`\mathrm{III}`$regions occurs instead much earlier ($`z5`$) in models that predict many more faint QSOs at high redshifts (Haiman & Loeb 1998).
I would like to thank my collaborators, T. Abel, F. Haardt, Z. Haiman, and M. Rees, for many useful discussions on the topics discussed here.
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# 1 𝐹(𝜃)=𝑠𝑖𝑛𝜃
I. Introduction
General covariance and the notion of absolute parallelism of vector fields in space-time are the two major geometrical concepts of teleparallel theories of gravity. Møller was probably the first to make use of the concept of absolute parallelism to put forward an alternative geometrical framework for general relativity that could satisfactorily address the problem of the definition of the gravitational energy.
The idea of absolute parallelism can be established by considering a space-time vector field $`V^\mu (x)`$ and a set of orthonormal tetrad fields $`e_\mu ^a(x)`$. At the space-time point $`x^\lambda `$ the tetrad components of the vector field are given by $`V^a(x)=e_\mu ^a(x)V^\mu (x)`$, and at $`x^\lambda +dx^\lambda `$ by $`V^a(x+dx)=e_\mu ^a(x+dx)V^\mu (x+dx)=V^a(x)+DV^a(x)`$, where $`DV^a(x)=e_\mu ^a(_\lambda V^\mu )dx^\lambda `$. The covariant derivative $``$ is constructed out of the connection
$$\mathrm{\Gamma }_{\mu \nu }^\lambda =e^{a\lambda }_\mu e_{a\nu }.$$
$`(1)`$
The vector field $`V^\mu (x)`$ is said to be autoparallel if its tetrad components at distant points coincide. Thus $`V^\mu (x)`$ is autoparallel if $`_\lambda V^\mu `$ vanishes. Therefore connection (1) defines a condition for absolute parallelism, or teleparallelism, in space-time. Such connection only makes sense if the tetrad field transforms under the global SO(3,1) group.
A gravity theory based on (1) obviously depart from the Riemannian geometry because the curvature tensor constructed out of it vanishes identically. Inspite of this fact, there does exist a theory based on (1) that describes the dynamics of the gravitational field in agreement with Einstein’s general relativity. Møller called such theory the “tetrad theory of gravity”, but for a long time it has been known as the teleparallel equivalent of general relativity (TEGR). The Lagrangian density for this alternative description of general relativity is constructed by means of a quadratic combination of the torsion tensor $`T_{\mu \nu }^a=_\mu e_\nu ^a_\nu e_\mu ^a`$, which is related to the antisymmetric part of connection (1).
It is possible to establish a theory for the gravitational field directly from (1). However, in order to make contact with more recent analysis, it is instructive to approach the TEGR firstly considering it with a local SO(3,1) symmetry in the Lagrangian context. Eventually we will return to the geometrical framework determined by (1).
Although not extensively investigated in the literature, the TEGR in Lagrangian form has been considered as a viable formulation of the gravitational dynamics inspite of troubles that may spoil the initial value problem. Such problems take place if the Lagrangian density of the TEGR is formulated with a local SO(3,1) symmetry, in which case the torsion tensor is defined by $`T_{\mu \nu }^a=_\mu e_\nu ^a_\nu e_\mu ^a+\omega _\mu ^a{}_{b}{}^{}e_{\nu }^{b}\omega _\nu ^a{}_{b}{}^{}e_{\mu }^{b}`$. The spin connection $`\omega _{\mu ab}`$ is totally independent of $`e_\mu ^a`$ and satisfies the condition of vanishing curvature: $`R_{b\mu \nu }^a(\omega )0`$.
In order to solve problems with respect to the initial value problem, the Hamiltonian formulation of the TEGR with the local SO(3,1) symmetry was considered. By working out the constraint algebra one ultimately concludes that in order to formulate consistently the theory in terms of first class constraints it is mandatory to break the local SO(3) symmetry of the action in Hamiltonian form. Therefore a consistent Hamiltonian formulation of the theory displays invariance under the global SO(3) group that acts on triads restricted to three-dimensional spacelike hypersurfaces.
The major feature of the Hamiltonian formulation of the TEGR is that the integral form of the Hamiltonian constraint equation can be written as an energy equation of the type $`C=HE=0`$. The reason for this follows from the fact that the Hamiltonian constraint contains a scalar density in the form of a total divergence, whose integral over the whole three-dimensional hypersurface yields the ADM energy. As we will argue, the integral of this scalar density over finite volumes of the three-dimensional hypersurface yields a natural definition for the gravitational field energy. Therefore in the TEGR one can make definite statements about the localizability of the gravitational energy, inspite of claims according to which the latter is not localizable. In fact the very concept of a black hole lends support to the idea that the gravitational energy is localizable, since there is no process by means of which the gravitational mass inside a black hole can be made to vanish. A remarkable application of the energy expression of the TEGR has been made in the evaluation of the irreducible mass of a rotating black hole.
The TEGR exhibits two specific properties: the emergence of a possible definition for the gravitational field energy and the global SO(3) symmetry of the theory. We believe that these two features are intimately related. Since the symmetry of the theory is global, triads related by a local SO(3) transformation are inequivalent and a priori we have no means to select the one that actually describes the spacelike hypersurface. We conjecture that the requirement of a minimum gravitational energy for a given space volume is one condition that singles out uniquely the correct set of triads. In this paper we investigate this conjecture in the realistic context of Bondi’s radiating metric. This conjecture was already put forward in a previous investigation of Bondi’s energy in the framework of the TEGR. An additional, essential requirement for a consistent expression of the gravitational energy is the boundary conditions on the triads. It was noted that the ADM energy is obtained from the energy expression of the TEGR if the asymptotic behaviour of the triads at spacelike infinity is given by
$$e_{(i)j}\eta _{ij}+\frac{1}{2}h_{ij}(\frac{1}{r}),$$
$`(2)`$
irrespective of any symmetry of the tensor $`h_{ij}`$.
A second possible, independent condition for assigning a set of triads to a given three-dimensional metric tensor amounts to requiring a symmetric tensor $`h_{ij}=h_{ji}`$ in the asymptotic expansion of $`e_{(i)j}`$. We will prove that this symmetry condition uniquely associates $`e_{(i)j}`$ to the metric tensor of the spacelike hypersurface of Bondi’s metric. The unique character of such triads strongly supports this second conjecture.
The two conjectures above are not mutually excluding. In this paper we argue that the set of triads that yield the minimum gravitational energy within a space volume containing the radiating source is the one whose asymptotic behaviour is determined by the symmetry condition $`h_{ij}=h_{ji}`$. We will show that this fact is indeed verified by analyzing several configurations for the triads. Unfortunately we have not found it possible to prove on general grounds that the “symmetrized” triad yields the minimum energy.
In order to obtain numerical values for the gravitational energy for a finite three-dimensional volume of Bondi’s space-time we need an explicit expression of the news function. However the existing expressions in the literature are not suitable for our purposes. In particular, the expression of the news function given by Hobill (to be presented ahead) is very intricate to the extent of not allowing a computer evaluation of numerical values of the gravitational energy. Therefore we have obtained an original and simpler expression for the news function that: (i) satisfies all necessary regularity conditions related to the axial symmetry of the system, (ii) makes the initial and final states (space-times) described by Bondi’s metric to be nonradiative and (iii) yields an expression for the energy density that can be numerically integrated.
In Section II we describe the TEGR in Lagrangian and Hamiltonian formulations, show the emergence of the definition of the gravitional energy and further discuss the troubles that arise in the initial value problem of the theory if it is constructed with a local SO(3,1) symmetry. In section III we present Bondi’s radiating metric and three expressions for triads restricted to the three-dimensional spacelike hypersurface. In this section we also prove that the symmetry condition $`h_{ij}=h_{ji}`$ uniquely associates a set of triads (whose asymptotic behaviour is given by (2)) with the metric tensor for the spacelike section of asymptotically flat space-times. The news function and the related mass aspect that will be needed for the calculations of the gravitational energy are obtained in section IV. In section V we carry out several calculations that lead to the main conclusion regarding the selection of triads.
Notation: spacetime indices $`\mu ,\nu ,\mathrm{}`$ and local Lorentz indices $`a,b,\mathrm{}`$ run from 0 to 3. In the 3+1 decomposition latin indices from the middle of the alphabet indicate space indices according to $`\mu =0,i,a=(0),(i)`$. The tetrad field $`e_\mu ^a`$ and the spin connection $`\omega _{\mu ab}`$ yield the usual definitions of the torsion and curvature tensors: $`R_{b\mu \nu }^a=_\mu \omega _\nu ^a{}_{b}{}^{}+\omega _\mu ^a{}_{c}{}^{}\omega _{\nu }^{c}{}_{b}{}^{}\mathrm{}`$, $`T_{\mu \nu }^a=_\mu e_\nu ^a+\omega _\mu ^a{}_{b}{}^{}e_{\nu }^{b}\mathrm{}`$. The flat space-time metric is fixed by $`\eta _{(0)(0)}=1`$.
II. The Lagrangian and Hamiltonian formulations of the TEGR
The Lagrangian density of the TEGR in empty space-time, displaying a local SO(3,1) symmetry, is given by
$$L(e,\omega ,\lambda )=ke(\frac{1}{4}T^{abc}T_{abc}+\frac{1}{2}T^{abc}T_{bac}T^aT_a)+e\lambda ^{ab\mu \nu }R_{ab\mu \nu }(\omega ).$$
$`(3)`$
where $`k=\frac{1}{16\pi G}`$, $`G`$ is the gravitational constant; $`e=det(e_\mu ^a)`$, $`\lambda ^{ab\mu \nu }`$ are Lagrange multipliers and $`T_a`$ is the trace of the torsion tensor defined by $`T_a=T_{ba}^b`$. The tetrad field $`e_{a\mu }`$ and the spin connection $`\omega _{\mu ab}`$ are completely independent field variables. The latter is enforced to satisfy the condition of zero curvature. Therefore this Lagrangian formulation is in no way similar to the usual Palatini formulation, in which the spin connection is related to the tetrad field via field equations. Later on we will introduce the tensor $`\mathrm{\Sigma }_{abc}`$ defined by
$$\frac{1}{4}T^{abc}T_{abc}+\frac{1}{2}T^{abc}T_{bac}T^aT_aT^{abc}\mathrm{\Sigma }_{abc}.$$
The equivalence of the TEGR with Einstein’s general relativity is based on the identity
$$eR(e,\omega )=eR(e)+e(\frac{1}{4}T^{abc}T_{abc}+\frac{1}{2}T^{abc}T_{bac}T^aT_a)\mathrm{\hspace{0.17em}2}_\mu (eT^\mu ),$$
which is obtained by just substituting the arbitrary spin connection $`\omega _{\mu ab}=^o\omega _{\mu ab}(e)+K_{\mu ab}`$ in the scalar curvature tensor $`R(e,\omega )`$ in the left hand side; $`{}_{}{}^{o}\omega _{\mu ab}^{}(e)`$ is the Levi-Civita connection and $`K_{\mu ab}=\frac{1}{2}e_a^\lambda e_b^\nu (T_{\lambda \mu \nu }+T_{\nu \lambda \mu }T_{\mu \nu \lambda })`$ is the contorsion tensor. The vanishing of $`R_{b\mu \nu }^a(\omega )`$, which is one of the field equations derived from (3), implies the equivalence of the scalar curvature $`R(e)`$, constructed out of $`e_\mu ^a`$ only, and the quadratic combination of the torsion tensor. It also ensures that the field equation arising from the variation of $`L`$ with respect to $`e_\mu ^a`$ is strictly equivalent to Einstein’s equations in tetrad form. Let $`\frac{\delta L}{\delta e^{a\mu }}=0`$ denote the field equations satisfied by $`e^{a\mu }`$. It can be shown by explicit calculations that
$$\frac{\delta L}{\delta e^{a\mu }}=\frac{1}{2}e\{R_{a\mu }(e)\frac{1}{2}e_{a\mu }R(e)\}.$$
$`(4)`$
We refer the reader to Ref. for additional details.
Throughout this section we will be interested in asymptoticaly flat space-times. The Hamiltonian formulation of the TEGR can be successfully implemented if we fix the gauge $`\omega _{0ab}=0`$ from the outset, since in this case the constraints (to be shown below) constitute a first class set. The condition $`\omega _{0ab}=0`$ is achieved by breaking the local Lorentz symmetry of (3). We still make use of the residual time independent gauge symmetry to fix the usual time gauge condition $`e_{(k)}^0=e_{(0)i}=\mathrm{\hspace{0.17em}0}`$. Because of $`\omega _{0ab}=0`$, $`H`$ does not depend on $`P^{kab}`$, the momentum canonically conjugated to $`\omega _{kab}`$. Therefore arbitrary variations of $`L=p\dot{q}H`$ with respect to $`P^{kab}`$ yields $`\dot{\omega }_{kab}=0`$. Thus in view of $`\omega _{0ab}=0`$, $`\omega _{kab}`$ drops out from our considerations.
As a consequence of the above gauge fixing the canonical action integral obtained from (3) becomes
$$A_{TL}=d^4x\{\mathrm{\Pi }^{(j)k}\dot{e}_{(j)k}H\},$$
$`(5)`$
$$H=NC+N^iC_i+\mathrm{\Sigma }_{mn}\mathrm{\Pi }^{mn}+\frac{1}{8\pi G}_k(NeT^k)+_k(\mathrm{\Pi }^{jk}N_j).$$
$`(6)`$
$`N`$ and $`N^i`$ are the lapse and shift functions, and $`\mathrm{\Sigma }_{mn}=\mathrm{\Sigma }_{nm}`$ are Lagrange multipliers. The constraints are defined by
$$C=_j(2keT^j)ke\mathrm{\Sigma }^{kij}T_{kij}\frac{1}{4ke}(\mathrm{\Pi }^{ij}\mathrm{\Pi }_{ji}\frac{1}{2}\mathrm{\Pi }^2),$$
$`(7a)`$
$$C_k=e_{(j)k}_i\mathrm{\Pi }^{(j)i}\mathrm{\Pi }^{(j)i}T_{(j)ik},$$
$`(7b)`$
where $`e=det(e_{(j)k})`$, $`T^i=g^{ik}e^{(j)l}T_{(j)lk}`$, $`T_{(i)jk}=_je_{(i)k}_ke_{(i)j}`$, and
$$\mathrm{\Sigma }^{ijk}=\frac{1}{4}(T^{ijk}+T^{jik}T^{kij})+\frac{1}{2}(\eta ^{ik}T^j\eta ^{ij}T^k).$$
We remark that (5) and (6) are now invariant under global SO(3) and general coordinate transformations. Therefore the torsion tensor restricted to the three-dimensional spacelike hypersurface is ultimately related to the antisymmetric part of the spatial components of (1). Had we dispensed with the connection $`\omega _{\mu ab}`$ from the outset we would have arrived at precisely the canonical formulation determined by (5), (6) and (7). Such connection has been considered in previous investigations of teleparallel theories, but in fact it is eventually unnecessary for the establishment of the theory.
If we assume the asymptotic behaviour
$$e_{(j)k}\eta _{jk}+\frac{1}{2}h_{jk}(\frac{1}{r})$$
$`(2)`$
for $`r\mathrm{}`$, then in view of the relation
$$\frac{1}{8\pi G}d^3x_j(eT^j)=\frac{1}{16\pi G}_S𝑑S_k(_ih_{ik}_kh_{ii})E_{ADM}$$
$`(8a)`$
where the surface integral is evaluated for $`r\mathrm{}`$, the integral form of the Hamiltonian constraint $`C=0`$ may be rewritten as
$$d^3x\left\{ke\mathrm{\Sigma }^{kij}T_{kij}+\frac{1}{4ke}(\mathrm{\Pi }^{ij}\mathrm{\Pi }_{ji}\frac{1}{2}\mathrm{\Pi }^2)\right\}=E_{ADM}.$$
$`(8b)`$
The integration is over the whole three-dimensional space. Given that $`_j(eT^j)`$ is a scalar density, from (8a,b) we define the gravitational energy density enclosed by a volume V of the space as
$$E=\frac{1}{8\pi G}_Vd^3x_j(eT^j).$$
$`(9)`$
It must be noted that $`E`$ depends only on the triads $`e_{(k)i}`$ restricted to a three-dimensional spacelike hypersurface; the inverse quantities $`e^{(k)i}`$ can be written in terms of $`e_{(k)i}`$. From the right hand side of equation (4) we observe that the dynamics of the triads does not depend on $`\omega _{\mu ab}`$. Therefore $`E`$ given above does not depend on the fixation of any gauge for $`\omega _{\mu ab}`$. The reference space which defines the zero of gravitational energy has been defined in ref.. We briefly remark that the differences between (9) and Møller’s expression for the gravitational energy have been thoroughly discussed in ref..
We make now the important assumption that the general form of the canonical structure of the TEGR is the same for any class of space-times, irrespective of the peculiarities of the latter (for the de Sitter space, for example, there is an additional term in the Hamiltonian constraint $`C`$). Therefore we assert that the integral form of the Hamiltonian constraint equation can be written as $`C=HE=0`$ for any space-time, and that (9) represents the gravitational energy for arbitrary space-times with any topology.
We recall finally that Müller-Hoissen and Nitsch and Kopczyński have shown that in general the theory defined by (3) faces difficulties with respect to the Cauchy problem. They have shown that in general six components of the torsion tensor are not determined from the evolution of the initial data. On the other hand, the constraints of the theory constitute a first class set provided we fix the six quantities $`\omega _{0ab}=0`$ before varying the action. This condition is mandatory and does not merely represent one particular gauge fixing of the theory. Since the fixing of $`\omega _{0ab}`$ yields a well defined theory with first class constraints, we cannot assert that the field configurations of the latter are gauge equivalent to configurations whose time evolution is not precisely determined. The requirement of local SO(3,1) symmetry plus the addition of $`\lambda ^{ab\mu \nu }R_{ab\mu \nu }(\omega )`$ in (3) has the ultimate effect of discarding the connection.
Constant rotations constitute a basic feature of the teleparallel geometry. According to Møller, in the framework of the abolute parallelism tetrad fields, together with the boundary conditions, uniquely determine a tetrad lattice, apart from an arbitrary constant rotation of the tetrads in the lattice.
III. Bondi’s radiating metric and the associated triads.
Bondi’s metric describes the asymptotic form of a radiating solution of Einstein’s equations. It is not an exact solution; it holds only in the asymptotic region. In terms of radiation coordinates $`(u,r,\theta ,\varphi )`$, where $`u`$ is the retarded time and $`r`$ is the luminosity distance, Bondi’s metric is written as
$$ds^2=\left(\frac{V}{r}e^{2\beta }U^2r^2e^{2\gamma }\right)du^22e^{2\beta }dudr2Ur^2e^{2\gamma }dud\theta $$
$$+r^2\left(e^{2\gamma }d\theta ^2+e^{2\gamma }sin^2\theta d\varphi ^2\right).$$
$`(10)`$
This metric tensor displays axial symmetry and reflection invariance. By requiring $`u=constant`$, (10) describes null hypersurfaces. Each null radial (light) ray is labelled by particular values of $`u,\theta `$ and $`\varphi `$. At spacelike infinity $`u`$ takes the standard form $`u=tr`$. The four quantities appearing in (10), $`V,U,\beta `$ and $`\gamma `$ are functions of $`u,r`$ and $`\theta `$. A more general form of this metric has been given by Sachs, who showed that the most general metric tensor describing asymptotically flat gravitational waves depends on six functions of the coordinates.
The functions in (10) have the following asymptotic behaviour:
$$\beta =\frac{c^2}{4r^2}+\mathrm{}$$
$$\gamma =\frac{c}{r}+O(\frac{1}{r^3})+\mathrm{}$$
$$\frac{V}{r}=\mathrm{\hspace{0.33em}1}\frac{2M}{r}\frac{1}{r^2}\left[\frac{d}{\theta }+dcot\theta \left(\frac{c}{\theta }\right)^24c\left(\frac{c}{\theta }\right)cot\theta \frac{1}{2}c^2\left(1+8cot^2\theta \right)\right]+\mathrm{}$$
$$U=\frac{1}{r^2}\left(\frac{c}{\theta }+2ccot\theta \right)+\frac{1}{r^3}\left(2d+3c\frac{c}{\theta }+4c^2cot\theta \right)+\mathrm{}$$
where $`M=M(u,\theta )`$ and $`d=d(u,\theta )`$ are the mass aspect and the dipole aspect, respectively. From the function $`c(u,\theta )`$ we define the news function $`\frac{c(u,\theta )}{u}`$.
The functions $`U,V,\beta `$ and $`\gamma `$ must satisfy regularity conditions along the $`z`$ axis ($`\theta =0,\pi `$). We must require
$$V,\beta ,\frac{U}{sin\theta },\frac{\gamma }{sin^2\theta }$$
to be regular functions of $`cos\theta `$ for $`\theta =0,\pi `$. The regularity conditions will be necessary for the construction of the news function, in section IV.
The application of (9) to Bondi’s metric requires transforming it to spherical coordinates ($`t,r,\theta ,\varphi `$) for which $`t=constant`$ defines a space-like hypersurface. Therefore we carry out a coordinate transformation such that the new timelike coordinate is given by $`t=u+r`$. We arrive at
$$ds^2=\left(\frac{V}{r}e^{2\beta }U^2r^2e^{2\gamma }\right)dt^22Ur^2e^{2\gamma }dtd\theta $$
$$+2\left[e^{2\beta }\left(\frac{V}{r}1\right)U^2r^2e^{2\gamma }\right]drdt$$
$$+\left[e^{2\beta }\left(2\frac{V}{r}\right)+U^2r^2e^{2\gamma }\right]dr^2+2Ur^2e^{2\gamma }drd\theta +r^2\left(e^{2\gamma }d\theta ^2+e^{2\gamma }sin^2\theta d\varphi ^2\right).$$
$`(11)`$
Therefore the metric restricted to a three-dimensional spacelike hypersurface is given by
$$ds^2=\left[e^{2\beta }\left(2\frac{V}{r}\right)+U^2r^2e^{2\gamma }\right]dr^2+2Ur^2e^{2\gamma }drd\theta $$
$$+r^2\left(e^{2\gamma }d\theta ^2+e^{2\gamma }sin^2\theta d\varphi ^2\right).$$
$`(12)`$
We recall that Goldberg and Papapetrou have already considered Bondi’s metric in cartesian coordinates.
The crucial point of the present investigation is the determination, in the framework of the TEGR, of the correct set of triads that lead to (12). If the metric tensor has only diagonal components, such as the metric tensor restricted to the three-dimensional spacelike section of Kerr’s space-time, then the simplest construction has proven to be the correct one. However, for metric tensors that contain off-diagonal terms, the determination of the unique triad is by no means a trivial procedure. In fact there is an infinity of triads that satisfy the boundary conditions determined by (2) and lead to (12). In Ref. two sets of triads that comply with (2) are presented. They are given by
$$e_{(k)i}=\left(\begin{array}{ccc}Asin\theta cos\varphi +Bcos\theta cos\varphi & rCcos\theta cos\varphi & rDsin\theta sin\varphi \\ Asin\theta sin\varphi +Bcos\theta sin\varphi & rCcos\theta sin\varphi & rDsin\theta cos\varphi \\ Acos\theta Bsin\theta & rCsin\theta & 0\end{array}\right),$$
$`(13)`$
where
$$A=e^\beta \sqrt{2\frac{V}{r}},$$
$`(14a)`$
$$B=rUe^\gamma ,$$
$`(14b)`$
$$C=e^\gamma ,$$
$`(14c)`$
$$D=e^\gamma ,$$
$`(14d)`$
and
$$e_{(k)i}=\left(\begin{array}{ccc}A^{}sin\theta cos\varphi & rB^{}cos\theta cos\varphi +rC^{}sin\theta cos\varphi & rD^{}sin\theta sin\varphi \\ A^{}sin\theta sin\varphi & rB^{}cos\theta sin\varphi +rC^{}sin\theta sin\varphi & rD^{}sin\theta cos\varphi \\ A^{}cos\theta & rB^{}sin\theta +rC^{}cos\theta & 0\end{array}\right),$$
$`(15)`$
where
$$A^{}=\left[e^{2\beta }\left(2\frac{V}{r}\right)+U^2r^2e^{2\gamma }\right]^{\frac{1}{2}},$$
$`(16a)`$
$$B^{}=\frac{1}{A^{}}e^{\beta +\gamma }\sqrt{2\frac{V}{r}},$$
$`(16b)`$
$$C^{}=\frac{1}{A^{}}Ure^{2\gamma },$$
$`(16c)`$
$$D^{}=e^\gamma .$$
$`(16d)`$
It is easy to see that both (13) and (15) yield the metric tensor (12) through the relation $`e_{(i)j}e_{(i)k}=g_{jk}`$. They are related by a local SO(3) transformation.
We have presented (13) and (15) because they are the simplest constructions that satisfy two basic requirements: (i) the triads must have the asymptotic behaviour given by (2); (ii) by making the physical parameters of the metric vanish we must have $`T_{(k)ij}=0`$ everywhere. In the present case if we make $`M=d=c=0`$ both (13) and (15) acquire the form
$$e_{(k)i}=\left(\begin{array}{ccc}sin\theta cos\varphi & rcos\theta cos\varphi & rsin\theta sin\varphi \\ sin\theta sin\varphi & rcos\theta sin\varphi & rsin\theta cos\varphi \\ cos\theta & rsin\theta & 0\end{array}\right).$$
$`(17)`$
In cartesian coordinates the expression above can be reduced to the diagonal form $`e_{(k)i}(x,y,z)=\delta _{ik}`$. The requirement (ii) above is essentialy equivalent to the establishment of reference space triads, as discussed in . A proper definition of gravitational energy requires the notion of reference space triads, which in the present case are given by (17). Note that by a suitable choice of a local SO(3) rotation we can make the flat space triads (17) satisfy the requirement (i), but not (ii).
We proceed now to obtain the triads whose aymptotic expansion is given by (2) with the symmetry condition $`h_{ij}=h_{ji}`$. The procedure is the following. We consider, for instance, triads (13) and transform it to cartesian coordiates. Then we perform a local, asymptotic transformation
$$\stackrel{~}{e}_{(k)i}(t,x,y,z)=\mathrm{\Lambda }_{(k)}^{(j)}e_{(j)i}(t,x,y,z),$$
$`(18)`$
where $`\mathrm{\Lambda }_{(k)}^{(j)}`$ satisfies
$$\mathrm{\Lambda }_{(k)}^{(j)}\delta _{(k)}^{(j)}+\omega _{(k)}^{(j)},$$
$`(19a)`$
$$\omega _{(j)(k)}=\omega _{(k)(j)},$$
$`(19b)`$
and $`\omega _{(j)(k)}O(\frac{1}{r})`$ for $`r\mathrm{}`$ ($`r=\sqrt{x^2+y^2+z^2}`$). Transformation (19) preserves the asymptotic behaviour of the triads. By substituting (19) in (18) we find
$$\stackrel{~}{e}_{(k)i}=e_{(k)i}+\omega _{(k)i},$$
from what follows
$$\stackrel{~}{h}_{ki}=h_{ki}+2\omega _{ki},$$
$`(20)`$
where $`h_{ki}`$ is given by the asymptotic expansion of (13) in cartesian coordinates. By requiring
$$\stackrel{~}{h}_{ki}=\stackrel{~}{h}_{ik},$$
$`(21)`$
and making use of (19b) we find that
$$\omega _{ki}=\frac{1}{4}(h_{ki}h_{ik}).$$
$`(22)`$
Substituting now (22) in (20) we arrive at
$$\stackrel{~}{h}_{ki}=\frac{1}{2}(h_{ki}+h_{ik})h_{(ki)}.$$
$`(23)`$
Thus we ultimately obtain the symmetrized triads in cartesian coordinates:
$$\stackrel{~}{e}_{(k)i}\eta _{ki}+\frac{1}{2}\stackrel{~}{h}_{ki}.$$
$`(24)`$
It must be noted that we arrive at a symmetrized triad irrespective of the triads we consider initially. The only requirement is that the unrotated triad must satisfy the asymptotic behaviour (2). Had we considered (15) we would arrive at the same result. In particular, from (22) we observe that if the initial triad is already symmetrized, then no local, asymptotic rotation is necessary.
By transforming (24) into spherical coordinates $`t,r,\theta ,\varphi `$ (in which case the triads are no longer symmetrized), we finally arrive at
$$\stackrel{~}{e}_{(1)1}(1+\frac{M}{r})sin\theta cos\varphi \frac{f}{2r}cos\theta cos\varphi ,$$
$$\stackrel{~}{e}_{(2)1}(1+\frac{M}{r})sin\theta sin\varphi \frac{f}{2r}cos\theta sin\varphi ,$$
$$\stackrel{~}{e}_{(3)1}(1+\frac{M}{r})cos\theta +\frac{f}{2r}sin\theta ,$$
$$\stackrel{~}{e}_{(1)2}r(1+\frac{c}{r})cos\theta cos\varphi \frac{f}{2}sin\theta cos\varphi ,$$
$$\stackrel{~}{e}_{(2)2}r(1+\frac{c}{r})cos\theta sin\varphi \frac{f}{2}sin\theta sin\varphi ,$$
$$\stackrel{~}{e}_{(3)2}r(1+\frac{c}{r})sin\theta \frac{f}{2}cos\theta ,$$
$$\stackrel{~}{e}_{(1)3}r(1\frac{c}{r})sin\theta sin\varphi ,$$
$$\stackrel{~}{e}_{(2)3}r(1\frac{c}{r})sin\theta cos\varphi ,$$
$$\stackrel{~}{e}_{(3)3}0,$$
$`(25)`$
where $`f`$ is given by
$$f=\frac{c}{\theta }+2ccotg\theta .$$
We observe that by making $`M=c=0`$ triads (25) reduce to (17). We also note that (25), as well as (13) and (15), only make sense in the asymptotic region where Bondi’s metric is valid.
In view of (23) and (24) it is now easy to prove the uniqueness of the symmetrized triad. Suppose that the asymptotic behaviour of the metric tensor is given by
$$g_{ij}\eta _{ij}+h_{ij}^{}(\frac{1}{r}).$$
$`(26)`$
On the other hand by making use of (2) it follows from the relation $`g_{ij}=e_i^{(k)}e_{(k)j}`$ and from (23) that
$$g_{ij}\eta _{ij}+\frac{1}{2}(h_{ij}+h_{ji})=\eta _{ij}+h_{(ij)}.$$
$`(27)`$
Since $`h_{ij}^{}`$ is unique, by comparing (26) and (27) we are led to conclude that there exists a unique symmetrized triad associated to the spatial section of an asymptotically flat metric tensor.
IV. Construction of the news function and of the mass aspect
In order to establish explicit expressions for the functions $`c(u,\theta )`$ and $`M(u,\theta )`$, the suplementary field equations $`R_{00}=R_{02}=0`$ for Bondi’s metric are considered. In simplified form they are given by
$$\frac{M}{u}=\left(\frac{c}{u}\right)^2+\frac{1}{2}\frac{}{u}\left[\frac{^2c}{\theta ^2}+3\frac{c}{\theta }cotg\theta 2c\right],$$
$`(28)`$
$$3\frac{d}{u}=\frac{M}{\theta }+3c\frac{^2c}{\theta u}+4c\frac{c}{u}cotg\theta +\frac{c}{\theta }\frac{c}{u}.$$
$`(29)`$
From (28) we observe that if $`\frac{c}{u}=0`$, the mass aspect $`M`$ does not depend on $`u`$.
For a family of null hypersurfaces Bondi’s mass is defined by
$$m(u)=\frac{1}{2}_0^\pi M(u,\theta )sin\theta 𝑑\theta .$$
$`(30)`$
It represents the mass of the system at the retarded time $`u`$. Multiplying both sides of (28) by $`sin\theta `$, integrating in $`\theta `$ and making use of the regularity conditions on the $`z`$ axis stated in section III, we arrive at
$$\frac{dm}{du}=\frac{1}{2}_0^\pi \left(\frac{c}{u}\right)^2sin\theta 𝑑\theta ,$$
$`(31)`$
which expresses the loss of mass. This equation was first obtained by Bondi et. al.. Therefore if the news function is nonvanishing, the mass of the system decreases in time. We remark that not only $`c(u,\theta )`$, but also the news function must be a regular function on the $`z`$ axis.
The function $`c(u,\theta )`$ determines not only the loss of mass, but in fact it determines the whole structure of Bondi’s metric, since via (28) and (29) it also determines the functions $`M`$, $`d`$, and all other functions that arise in the process of integration of the field equations. Therefore its construction deserves special attention.
The news function proposed by Bondi et. al. is given by
$$\frac{c}{u}=\underset{n=0}{\overset{\mathrm{}}{}}f_n(u)P_n(\mu ),$$
where $`\mu =cos\theta `$. $`P_n(\mu )`$ are the Legendre polynomials and $`f_n(u)`$ are functions to be determined. These functions become more and more intricate for increasing $`n`$, and eventually the above expression cannot be manipulated analytically.
Bonnor and Rotenberg, Papapetrou and Hallidy and Janis have attempted at establishing an expression for $`c(u,\theta )`$, but none of these proposals worked out satisfactorily, either because of the complexity of the structure of $`c(u,\theta )`$, or because the loss of mass does not occur for for a finite number of terms in the expansion in $`n`$. Bonnor and Rotenberg’s expression describes a radiative period between an initial nonradiative, static state and a final nonradiative, nonstatic state. The other approaches mentioned above describe a radiative period between static initial and final states.
Hobill has provided an expression for $`c(u,\theta )`$ that describes a radiative period between initial and final static states, and that leads to a loss of mass that is exactly equal to the total variation of the mass aspect during the period. Again the expression is not simple. It is given by
$$c(u,\mu )=\frac{[m_bf(\mu )e^{\eta u}+m_a(1\mu ^2)](1+\eta \mu ^2)(1\mu ^2)}{f^2(\mu )e^{2\eta u}+(n1)\mu ^2n\mu ^4+1},$$
$`(32)`$
where $`n`$ and $`\eta `$ are constants and $`m_a`$ is identified as one fourth of the total mass loss. It relates to the other constants according to
$$\frac{18n}{n+1}=\eta m_a,$$
$$m_b=\pm \sqrt{\frac{2}{n}}m_a.$$
$`f(\mu )`$ is an arbitrary function that must be everywhere regular and positive definite in the interval $`1\mu 1`$. We note that $`c`$ vanishes over the symmetry axis $`(\mu =\pm 1)`$. The news function obtained from (32) reads
$$\frac{c}{u}=\frac{\eta f(\mu )e^{\eta u}(1\mu ^2)}{[e^{2\eta u}f^2(\mu )(1+n\mu ^2)^1+(1\mu ^2)]^2}\{\frac{m_bf^2(\mu )e^{2\eta u}}{1+n\mu ^2}$$
$$\frac{2m_a(1\mu ^2)f(\mu )e^{\eta u}}{1+n\mu ^2}+m_b(1\mu ^2)\}.$$
$`(33)`$
Expressions (32) and (33) lead, via integration of equation (28), to an extremely complicated expression for the mass aspect:
$$M(u,\mu )=M(\mathrm{})4m_a$$
$$+\left[\frac{f^2(\mu )e^{2\eta u}}{1+n\mu ^2}+(1\mu ^2)\right]^3\left[4m_a(1\mu ^2)^3+\underset{l=1}{\overset{5}{}}H_l(\mu )e^{l\eta u}\right].$$
$`(34)`$
The intricate expressions for $`H_l(\mu )`$ are given in the Appendix of Ref. . The problem with (32), (33) and (34) is that they are too complicated to yield numerical values for integrals (to be considered in the next section) that contain these expressions, even by means of computer calculations.
Therefore we attempted at obtaining a simpler expression for $`c(u,\mu )`$. We damanded two conditions on $`c(u,\mu )`$. First, it must satisfy the regularity conditions on the $`z`$ axis, which guarantees that the system is permanently isolated and leads to a well defined loss of mass. Second, that the initial and final states are nonradiative. However, as we will show, our expression leads to a nonstatic final state. It is given by
$$c(u,\mu )=\frac{ae^{nu}(1\mu ^2)F(\mu )}{e^{2nu}+1},$$
$`(35)`$
where $`n`$ and $`a`$ are constants ($`n^1`$ and $`a`$ have dimension of length) and $`F(\mu )`$ is a function that must be choosen such that $`c(1\mu ^2)^1=c(sin\theta )^2`$ is a regular function. Note that (35) vanishes for $`\mu \pm 1`$. The news function associated to (35) is given by
$$\frac{c}{u}=\frac{na(1\mu ^2)F(\mu )e^{nu}}{e^{2nu}+1}\left[1\frac{2e^{2nu}}{e^{2nu}+1}\right].$$
$`(36)`$
In the limit $`u\pm \mathrm{}`$ we have
$$\underset{u\pm \mathrm{}}{lim}c=\underset{u\pm \mathrm{}}{lim}\frac{c}{u}=0,$$
$`(37)`$
for any nonvanishing value of $`n`$. The property above is a necessary but not sufficient condition for having a static final state. Equation (37) does not determine whether the mass aspect $`M`$ depends on $`\theta `$, and so from (29) there is the possibility that $`d`$ depends on $`u`$. The initial state is assumed to be of the Schwarzschild type.
Integrating equation (28) from $`\mathrm{}`$ to $`u`$, considering $`c(u,\mu )`$ given by (35), we obtain
$$M(u,\mu )=M(\mathrm{})n^2a^2(1\mu ^2)^2F^2(\mu )_{\mathrm{}}^u\left[\frac{e^{nu}}{e^{2nu}+1}\frac{2e^{3nu}}{(e^{2nu}+1)^2}\right]^2𝑑u$$
$$+\frac{1}{2}\frac{ae^{nu}}{(e^{2nu}+1)}\left[(1\mu ^2)^2F^{\prime \prime }8\mu (1\mu ^2)F^{}+4(3\mu ^21)F(\mu )\right]+\lambda (\mu ),$$
$`(38)`$
where $`F^{}=\frac{dF}{d\mu }`$, $`F^{\prime \prime }=\frac{d^2F}{d\mu ^2}`$ and $`\lambda (\mu )`$ is an integrating function that depends only on $`\mu `$.
The integral in (38) depends on the sign of $`n`$. For $`n<0`$ we have
$$M(u,\mu )=M(\mathrm{})+\frac{1}{6}na^2(1\mu ^2)^2F^2(\mu )\left[\frac{3e^{4nu}+1}{(e^{2nu}+1)^3}\right]$$
$$+\frac{1}{2}\frac{ae^{nu}}{(e^{2nu}+1)}\left[(1\mu ^2)^2F^{\prime \prime }8\mu (1\mu ^2)F^{}+4(3\mu ^21)F(\mu )\right]+\lambda (\mu ),$$
$`(39)`$
and for $`n>0`$,
$$M(u,\mu )=M(\mathrm{})+na^2(1\mu ^2)^2F^2(\mu )\left[\frac{1}{6}\frac{3e^{4nu}+1}{(e^{2nu}+1)^3}+\frac{23}{48}\right]$$
$$+\frac{1}{2}\frac{ae^{nu}}{(e^{2nu}+1)}\left[(1\mu ^2)^2F^{\prime \prime }8\mu (1\mu ^2)F^{}+4(3\mu ^21)F(\mu )\right]+\lambda (\mu ).$$
$`(40)`$
By requiring the inital state to be static, expressions (39) and (40) must satisfy
$$\underset{u\mathrm{}}{lim}M=M(\mathrm{})M_0.$$
$`(41)`$
Applying the condition above to (39) we find that the function $`\lambda (\mu )`$ vanishes, and therefore
$$M(u,\mu )=M_0+\frac{1}{6}na^2(1\mu ^2)^2F^2(\mu )\left[\frac{3e^{4nu}+1}{(e^{2nu}+1)^3}\right]$$
$$+\frac{1}{2}\frac{ae^{nu}}{(e^{2nu}+1)}\left[(1\mu ^2)^2F^{\prime \prime }8\mu (1\mu ^2)F^{}+4(3\mu ^21)F(\mu )\right],$$
$`(42)`$
for $`n<0`$. Similarly, applying (41) to expression (40) we find that the integration function $`\lambda (\mu )`$ is given by
$$\lambda (\mu )=\left(\frac{1}{6}+\frac{23}{48}\right)na^2(1\mu ^2)F^2(\mu ).$$
Substituting it back in (40) for $`n>0`$ we arrive at
$$M(u,\mu )=M_0+\frac{1}{6}na^2(1\mu ^2)^2F^2(\mu )\left[\frac{3e^{4nu}+1}{(e^{2nu}+1)^3}1\right]$$
$$+\frac{1}{2}\frac{ae^{nu}}{(e^{2nu}+1)}\left[(1\mu ^2)^2F^{\prime \prime }8\mu (1\mu ^2)F^{}+4(3\mu ^21)F(\mu )\right].$$
$`(43)`$
Let us now check the limiting value of (42) and (43) for $`u\mathrm{}`$. Considering first (42) we obtain
$$\underset{u\mathrm{}}{lim}M=M_0+\frac{1}{6}na^2(1\mu ^2)^2F^2(\mu ),$$
where $`n<0`$. It follows from the expression above that the total variation of the mass aspect $`\mathrm{\Delta }M_T`$ associated with (42) is given by
$$\mathrm{\Delta }M_T=\frac{1}{6}na^2(1\mu ^2)^2F^2(\mu ).$$
$`(44)`$
The mass aspect $`M(u,\mu )`$ given by (43), for which $`n>0`$, leads to a similar expression for $`\mathrm{\Delta }M_T`$:
$$\mathrm{\Delta }M_T=\frac{1}{6}na^2(1\mu ^2)^2F^2(\mu ).$$
$`(45)`$
Hobill obtained a result analogous to (44). However his expression, in the limit $`u\mathrm{}`$, does not depend on $`\mu `$. Consequently the final state in his approach is static. In the present case, the mass aspect will still depend on $`\mu `$ in the limit $`u\mathrm{}`$. In view of (29) this fact implies a final nonradiative, nonstatic state, which according to Hobill may describe a even more realistic situation. In any event expressions (35) and (36) describe an isolated physical system whose initial and final states are nonradiative, and whose loss of mass is precisely determined.
V. Evaluation of the gravitational energy
In this section we obtain numerical values for the gravitational energy given by expression (9). The relevance of such expression is that we can apply it to finite spacelike volumes. In the present case we will calculate the gravitational energy inside a large but finite surface of constant radius $`r_0`$, centered at the radiating source. Expression (9) will be evaluated as a surface integral in the asymptotic region where the components of Bondi’s metric are precisely determined. It reduces to just one integral given by
$$E=\frac{1}{8\pi }_S𝑑\theta 𝑑\varphi eT^1,$$
$`(46)`$
where $`S`$ is a surface of fixed radius $`r_0`$, assumed to be large as compared with the dimension of the source, and the determinant $`e`$ is given by $`e=r^2Asin\theta `$.
We will consider triads (13), (15) and (25), for which the functions $`c(u,\theta )`$ and $`M(u,\theta )`$ are given by (35) and (42), respectively. Thus we take the mass aspect determined by the condition $`n<0`$, only. Moreover we will consider four possibilities for the function $`F(\mu )`$. Therefore we are effectively analyzing twelve distinct sets of triads.
We will follow here the steps of section V of . The energy expression that results from (13) and (15) have already been evaluated. They are given by
$$E_1=\frac{r_0}{4}_0^\pi 𝑑\theta \left\{sin\theta \left[e^\gamma +e^\gamma \frac{2}{A}\right]+\frac{1}{A}\frac{}{\theta }(Ursin\theta )\right\},$$
$`(47)`$
and
$$E_2=\frac{r_0}{4}_0^\pi d\theta \frac{1}{A}\{sin\theta [e^\gamma A^{}+e^{2\gamma }A^{}B^{}2+e^{2\gamma }\frac{A^{}}{\theta }C^{}BC^{}Be^\gamma \frac{\gamma }{\theta }]$$
$$Bcos\theta [B^{}e^\gamma ]\},$$
$`(48)`$
together with definitions (14) and (16). $`E_1`$ and $`E_2`$ correspond to (13) and (15), respectively. Unfortunately the sign of the first term in the expansion of the function $`U(u,r,\theta )`$ of the metric tensor (10), in Ref. , is changed. Therefore expressions $`E_1`$ and $`E_2`$ given in the latter reference must be corrected. Minor modifications (such as the modification of some numerical coefficients) are necessary. The correct expressions of $`E_1`$ and $`E_2`$ in terms of $`c(u,\theta )`$ and $`M(u,\theta )`$ are given by
$$E_1=\frac{1}{2}_0^\pi 𝑑\theta sin\theta M\frac{1}{4r_0}_0^\pi 𝑑\theta sin\theta \left[\left(\frac{c}{\theta }\right)^2+4c\left(\frac{c}{\theta }\right)cot\theta +4c^2cot^2\theta \right]$$
$$+\frac{1}{4r_0}_0^\pi 𝑑\theta M\frac{}{\theta }\left[sin\theta \left(\frac{c}{\theta }+2ccot\theta \right)\right],$$
$`(49)`$
$$E_2=\frac{1}{2}_0^\pi d\theta Msin\theta \frac{1}{4r_0}_0^\pi d\theta sin\theta [M^2+\frac{1}{2}\left(\frac{c}{\theta }\right)^2+4c\left(\frac{c}{\theta }\right)cot\theta $$
$$+6c^2cot^2\theta +\left(\frac{M}{\theta }\right)(\frac{c}{\theta }+2ccot\theta )]+\frac{1}{4r_0}_0^\pi d\theta cos\theta [2c\left(\frac{c}{\theta }\right)+4c^2cot\theta ].$$
$`(50)`$
It is clear that the total gravitational energy given by both (49) and (50), in the limit $`r_0\mathrm{}`$, which corresponds to the limit $`u\mathrm{}`$, yield the same value, i.e., the total mass:
$$\underset{r\mathrm{}}{lim}E_1=\underset{r\mathrm{}}{lim}E_2=M(\mathrm{})M_0.$$
$`(51)`$
In fact it is proven in that the total gravitational energies calculated out of triads related by a local SO(3) transformation, and that have the asymptotic behaviour given by (2), are the same.
We consider next triads given by (25). The components of the torsion tensor are given by
$$\stackrel{~}{T}_{(1)12}=\left(\frac{c}{r}\frac{M}{r}+r_1(\frac{c}{r})+\frac{1}{2r}_2f\right)cos\theta cos\varphi \left(\frac{1}{2}_1f+\frac{1}{2r}f+\frac{1}{r}_2M\right)sin\theta cos\varphi ,$$
$$\stackrel{~}{T}_{(1)13}=\left(\frac{c}{r}+\frac{M}{r}+r_1(\frac{c}{r})\right)sin\theta sin\varphi \frac{1}{2r}fcos\theta sin\varphi ,$$
$$\stackrel{~}{T}_{(1)23}=2ccos\theta sin\varphi +(_2c\frac{1}{2}f)sin\theta sin\varphi ,$$
$$\stackrel{~}{T}_{(2)12}=\left(\frac{c}{r}\frac{M}{r}+r_1(\frac{c}{r})+\frac{1}{2r}_2f\right)cos\theta sin\varphi \left(\frac{1}{2}_1f+\frac{1}{2}\frac{f}{r}+\frac{1}{r}_2M\right)sin\theta sin\varphi ,$$
$$\stackrel{~}{T}_{(2)13}=\left(\frac{c}{r}+\frac{M}{r}+r_1(\frac{c}{r})\right)sin\theta cos\varphi +\frac{1}{2r}fcos\theta cos\varphi ,$$
$$\stackrel{~}{T}_{(2)23}=2ccos\theta cos\varphi (_2c\frac{1}{2}f)sin\theta cos\varphi ,$$
$$\stackrel{~}{T}_{(3)12}=\left(\frac{M}{r}\frac{c}{r}r_1(\frac{c}{r})\frac{1}{2r}_2f\right)sin\theta \left(\frac{1}{2}_1f+\frac{1}{2r}f+\frac{1}{r}_2M\right)cos\theta ,$$
$$\stackrel{~}{T}_{(3)13}=\stackrel{~}{T}_{(3)23}=0,$$
$`(52)`$
where $`_1`$ and $`_2`$ denote partial derivatives with respect to $`r`$ and $`\theta `$, respectively.
The calculation of (46) out of the components above does not pose any particular problem, except that the calculation is very long. Denoting by $`\stackrel{~}{E}`$ the gravitational energy that follows from (25) and (52), we have
$$\stackrel{~}{E}=E_1\frac{1}{2r}[_0^\pi d\theta M^2sin\theta \frac{1}{4}_0^\pi d\theta (\frac{c}{\theta }+2ccotg\theta )\frac{c}{\theta }sin\theta $$
$$\frac{1}{2}_0^\pi 𝑑\theta (\frac{c}{\theta }+2ccotg\theta )ccos\theta $$
$$+\frac{1}{8}_0^\pi d\theta (\frac{c}{\theta }+2ccotg\theta )^2sin\theta (3cos^2\theta 1)].$$
$`(53)`$
In the calculation above we have made use of the regularity condition $`c(u,\theta )=0`$ for $`\theta =0,\pi `$. As expected, (53) also satisfies (51).
The comparison of (49), (50) and (53) is crucial for the selection of triads. We are now in a position of obtaining numerical values for these expressions. With this purpose we make use of data based on the work of Saenz and Shapiro for assigning values to $`M_0`$ and $`u`$. In the latter reference it is discussed a model for stellar colapse of white-dwarves, with or without axial symmetry, rotating or nonrotating. In this model gravitational waves are produced during a burst of $`10^3s`$. Thus $`u`$ may be taken to vary from $`10^{10}s`$ to $`10^4s`$. Based on the work of Saenz and Shapiro we consider a white-dwarf whose total mass is $`M_0=1.4M_{}`$, where $`M_{}`$ is the solar mass. It is believed that the maximum value for a white-dwarf is $`1.4M_{}`$. We will use geometrical unities (G=c=1) for the evaluation of the energies. The value of $`M_{}`$ in these unities is $`M_{}=1.47664\times 10^5cm`$. The distance $`r`$ will be taken to vary from $`10^{18}cm`$ to $`10^{13}cm`$. Finally, the parameters $`a`$ and $`n`$ assume the values $`100`$ and $`0.5`$, respectively, in proper unities.
We have used the MAPLE V computer package to obtain $`E_1`$, $`E_2`$ and $`\stackrel{~}{E}`$. For the function $`F(\mu )`$, which we now denote $`F(\theta )`$, we take $`sin\theta `$, $`sin^2\theta `$, $`cos\theta `$ and $`cos^2\theta `$. The resulting values are listed in tables 1-4 for the particular value $`u=10^4s`$. The numerical values were obtained with a precision of 20 digits.
We have verified that for all these functions, for all distances considered, $`\stackrel{~}{E}`$ is always the minimum energy:
$$\stackrel{~}{E}<E_2<E_1.$$
Note that since $`E_1`$, $`E_2`$ and $`\stackrel{~}{E}`$ are calculated at constant $`u=tr`$, the radial dependence of these expression is given just by the $`\frac{1}{r_0}`$ coefficient in (49), (50) and (53). Thus for constant $`u`$ the intricacy of these expressions reside in the angular dependence. We have further verified numerically that the result above is also obtained for any value of $`u`$ between $`10^{10}s`$ and $`10^4s`$. The constants $`a`$ and $`n`$ must be chosen such that the absolute values of $`\mathrm{\Delta }M_T`$ given by expressions (44) and (45) are not greater than $`M_0`$, otherwise the total gravitational energy will be negative. In any case, at present we do not know enough about axially-symmetric, nonrotating isolated sources in order to provide realistic values for these constants.
VI. Discussion
The result of the previous section constitutes a strong indication that the set of triads with asymptotic behaviour given by equation (2), and that satisfies the symmetry condition $`h_{ij}=h_{ji}`$, yields the minimum value for the gravitational energy that is computed from expression (9). In addition to the fact that the symmetrized triads are unique (namely, there does not exist a second set of triads that satisfy (2) and the symmetry condition), the present analysis indicates that the correct description of the gravitational field in terms of orthonormal triads, in the realm of the TEGR, is given by (24). The results described above amount to an interesting interplay between the energy properties of the space-time and its tetrad description.
As long as we are interested in the dynamics of the gravitational field only, as described by the metric tensor, it is irrelevant which configuration of triads we adopt. However, if we consider the dynamics of spinor fields, such as the Dirac field, then the correct choice of triads is crucial. Recall that the theory defined by (5) and (6) was established under the asumption of the time gauge condition. The latter, together with (25), establishes the complete set of tetrad fields. And finally, if the detection of the emission of gravitational energy carried by gravitational waves is experimentally feasible, then the whole scheme developed here will play a relevant role.
Bondi’s radiating metric is valid only in the asymptotic region. Therefore we can do no better than determining the asymptotic behaviour of the triads. However in the more general case where the metric tensor is valid everywhere, except for singularities, we still have a prescription for assigning a unique set of triads to a given metric tensor restricted to the spacelike section. Such prescription is due to Møller, who called it “supplementary conditions”. Although he established these conditions by still requiring de Donder relations for the metric tensor, we can dispense with the latter relations and ascribe generality to his proposal. Møller suggested as supplementary conditions for the space-time tetrad field the weak field condition
$$e_{a\mu }\eta _{a\mu }+\frac{1}{2}h_{a\mu },$$
$`(54)`$
with $`h_{a\mu }`$ satisfying the symmetry condition $`h_{a\mu }=h_{\mu a}`$. It differs from (2), which is restricted to spacelike sections only, in that the symmetry condition must be verified everywhere. Thus (54) is stronger than (2). Nevertheless we can require the set of triads to satisfy (54) in the general case. Every metric tensor for the spacelike section of a space-time admits a unique set of triads that satisfy a relation similar to (54), which is no longer a boundary condition and therefore it can be applied to space-times with arbitrary topology. We finally mention that the set of triads presented in ref. , in the analysis of the irreducible mass of a rotating black hole, satisfies a relation similar to (54), but in the three-dimensional spacelike hypersurface.
Acknowledgements
K. H. C. B. is supported by CAPES, Brazil.
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# Glassy Dynamics in a Frustrated Spin System: Role of Defects
## 1 Introduction
In recent years, there has been considerable interest in spin models without quenched-in disorder which, however, exhibit glassy dynamics. The model discussed in this paper belongs to this category. Frustration is a key concept in theories of structural glasses such as the curved-space pictures of metallic glass. The two concepts taken together, glassy dynamics in non-random spin systems and frustration as a key to glassy behavior, suggest that frustrated spin systems could play a role in understanding the nature of structural glasses.
The compressible, triangular-lattice, Ising antiferromagnet is an interesting example of a system where the frustration is removed by elastic distortions. The phase transition from a disordered, paramagnetic phase to an ordered striped phase has been studied in detail, and in the last couple of years, interest has shifted to studying the dynamics of the supercooled state. Monte Carlo simulations of the model have shown that following an instantaneous quench from an initial temperature $`T_i>>T_F`$ (the first-order transition temperature), to a final temperature $`T_f<T_F`$, the dynamics of the system changes character as $`T_f`$ is lowered below a temperature $`T^{}`$. The observed changes are consistent with a glass transition occurring at $`T^{}`$. As $`T^{}`$ is approached from above, the characteristic relaxation time increases dramatically, and the spin autocorrelation function develops a plateau at intermediate times (Fig. 1). Analysis of the fluctuation metric and the energy metric shows that the system becomes non-ergodic below $`T^{}`$. Although, the dynamics in this non-ergodic phase is reminiscent of “coarsening”, the relaxation towards equilibrium is unusual and is characterized by long-time periods of quiescent behavior interrupted by rare jump events. The distribution of jump-intervals obeys a power law which suggests that the model belongs to the category of weak-ergodicity breaking.
In this paper, the approach to the “glass transition” at $`T^{}`$ is analyzed through a study of the spin-spin autocorrelation functions and the behavior of defects which arise naturally when the spin system is mapped onto a solid-on-solid model (SOS). In the simulations, $`T^{}`$ is operationally defined as the temperature below which the system fails to reach equilibrium within the simulation time.
## 2 Model
The compressible Ising antiferromagnet is described by a nearest-neighbor spin-spin interaction which depends linearly on the distance between the spins. The lattice distortions are described by a homogeneous, three component strain field which describes the change in the three nearest-neighbor bond lengths. In the model studied in this paper, no fluctuations of these strain fields are allowed. This implies that the strain fields are being treated within a mean-field approximation but the spins are being treated exactly. It has been shown that, in this approximation, there is a first order transition to a striped phase in which there are rows of aligned spins alternating in sign. The ferromagnetic bonds are elongated and the anti-ferromagnetic bonds are shortened. It has also been shown that, in this mean-field model, there is an instability to this lattice distortion at a temperature below the first-order transition temperature.
The strain fields which remove the degeneracy in the compressible model play a role that is similar to the anisotropy of interactions or a staggered field which is conjugate to one of the degenerate ground-states. The crucial difference between the models is that the strain fields are annealed variables whereas the anisotropy and the staggered field are quenched variables. The dynamics that we observe in the compressible model as the glass transition is approached can, however, be rationalized on the basis of known results for the models with the quenched variables.
Monte Carlo simulations were used to study the dynamics following instantaneous quenches from a high-temperature disordered phase to a range of temperatures below the ordering transition (which was strongly first-order) at $`T_F`$. Standard spin-exchange dynamics was extended to include moves which attempt global changes of the shape and size of the box. These global changes were attempted after a complete sweep of all the spins in the lattice. This dynamics was adopted because the homogeneous (global) strain fields are expected to respond only to extensive changes of the nearest neighbor spin correlations. There is an “effective” spin model that the compressible Ising antiferromagnet can be mapped onto. This model involves long-range four-spin interactions. This mapping shows the connection between the current model and the p-spin models and the Bernasconi model. The dynamical model, as defined in this work, however, does not correspond to the dynamics of this four-spin model but instead keeps the strain field and the spins explicitly.
## 3 Spin-Spin autocorrelation functions and non-linear dynamic susceptibility
As mentioned in the introduction, the system is in equilibrium above $`T^{}`$ but falls out of equilibrium at $`T^{}`$. The behavior of the spin-spin autocorrelation function, as $`T^{}`$ is approached from above, is shown in Fig. 1. These correlation functions show clear non-exponential relaxations. The best stretched exponential fits ($`\mathrm{exp}(t/\tau )^\beta `$) are also shown in Fig. 1. The stretching parameter $`\beta 0.3`$ with a small dependence on temperature. There are indications of a plateau developing at long times, however, data over much longer times and larger systems will be needed to bear this out. In addition to the correlation functions, Fig. 1 also shows dynamic susceptibility associated with the fluctuations of the spin-spin correlation function. This susceptibility, associated with the time-dependent overlap, shows the same features similar to that observed in a Lennard-Jones binary mixture. These results suggest that the ergodicity-breaking transition at $`T^{}`$ is akin to a structural glass transition even though the system being studied is a spin system. The hope then is that the study of this simple spin model could provide some insight into the structural glass transition. It has already been shown that a trapping model, based on the observations in this spin system, can provide a qualitatively correct description of the observed frequency-dependent susceptibility in structural glasses. The SOS mapping provides an elegant way of analyzing the structures which develop in this spin model as the glass transition is approached. In the following, this mapping will be used to probe the nature of the glass transition.
## 4 Defects and Strings
The ground-state of the triangular-lattice Ising antiferromagnet is a critical state and can be mapped onto the rough phase of an SOS model. In this representation, a line is drawn between two spins which are connected by an anti-ferromagnetic bond. At zero temperature, this defines a tiling of the plane by three different types of rhombi (blue, red and green in Fig. 2). At finite temperatures, defects appear which correspond to elementary triangles with three spins of the same sign and correspond to screw dislocations in the SOS surface. Fig. 2 shows this representation for a configuration generated in the compressible Ising model after a quench above $`T^{}`$. A convenient way of representing the SOS surface is by strings running in the “vertical” direction. Choosing one of the nearest-neighbor directions in the triangular lattice as “horizontal”, these strings are defined by drawing lines connecting the middles of horizontal edges of the rhombi. In Fig. 2, the blue and red rhombi have horizontal edges but the green do not. These strings can end at the defects which change the number of strings by two. The ground-state of the pure triangular lattice antiferromagnet corresponds to a rough SOS surface with no average tilt and is characterized by the number of strings $`N_s=2L/3`$ where $`L`$ is the linear dimension of a finite lattice with periodic boundary conditions. The strings and defects (dislocations) provide an useful way of visualizing the transition at $`T^{}`$. To understand what could be happening at $`T^{}`$, we need to briefly discuss the properties of the compressible Ising antiferromagnet and its ground-state.
The ground-state of the compressible model is the “striped phase”. In the SOS picture this corresponds to a surface with a tilt and no “vertical” strings. There are, obviously, three ways of defining vertical strings in the triangular lattice (corresponding to the three nearest neighbor directions) and the striped phase is three-fold degenerate. These states have only one variety of the rhombi.
In terms of the SOS picture, the first-order transition at $`T_F`$ corresponds to a discontinuous change in the number of strings. Strings can end only at the defects (two strings end at a defect) and therefore the kinetics of this transition is dominated by the interplay between defects and strings. Similarly, one expects that the dynamics of the supercooled phase is dictated by the dynamics of defects and strings. The time evolution of the number of strings is shown in Fig. 3 for a temperature above $`T^{}`$ and a temperature below $`T^{}`$. Above $`T^{}`$, but close to it, the string density fluctuates around $`N_s=2L/3`$ but below $`T^{}`$ the string density is evolving towards zero in a step-wise fashion. The temperature $`T^{}`$ seems to correspond to an instability of the system towards the disappearance of strings. It however does not seem to correspond to a simple spinodal since the system stays in states with a fixed number of strings for very long times before hopping to a state with a fewer number of strings. This dynamics suggests a free-energy surface in this string-number space which has multiple minima with a distribution of barrier heights and a bias, below $`T^{}`$, towards smaller number of strings. This is exactly the type of picture that was used to construct the simple Langevin model of the frequency-dependent response in glasses. The order parameter in that model corresponds to the average number of strings in the SOS model.
Before analyzing the dynamics of the defects and strings, it is useful to look at some static quantities above $`T^{}`$ (where the system is in equilibrium and static averages can be defined. Fig. 4 shows the histogram characterizing the probability distribution of the average number of strings in a particular vertical direction. As the temperature is lowered, the distribution develops a pronounced tail and secondary maxima. At the same time, the distribution of defects changes from a nearly Gaussian distribution to a nearly exponential one with a peak close to zero. The combination of these two features gives a clue to the origin of the slow dynamics. The instability at $`T^{}`$ implies a tendency of the system to have fewer strings than that characterizing the ideal supercooled state with $`N_s=2L/3`$. To achieve this it needs to create fluctuations which have a large amplitude and last for long times. These string fluctuations can be created only via correlated defect events. The type of defect events that change the number of strings by a large number are the appearance of more than one pair of defects (defects always get created in pairs in the dynamics that is being used in the simulations). These events are extremely rare and can lead to long correlation times.
The correlation between defect and string histories is shown in Fig. 5. To change the average number of strings in the system by two or more, a pair of strings running the length of the system has to be created or destroyed. The history shown in Fig. 5 indicates that this is not accomplished by one pair of defects running through the system and creating or destroying strings in a “zipper” type action. Instead, the strings are created or destroyed through the creation of more than one defect pair in succession. As the temperature approaches $`T^{}`$, these events get rarer and the structures with significantly more (or less) strings than the ideal supercooled state get frozen in for longer and longer times. The glass transition, in this frustrated spin system, seems to be pinned by an underlying instability (true instability in mean-field) towards a deformation of the lattice or the spontaneous disappearance of strings. The dynamics approaching this transition is dictated by rare events which involve correlated defects.
## 5 Conclusion
In this preliminary study of correlation between defects and glassy dynamics in a frustrated spin model, the observations suggest that defects are crucial to understanding the glass transition. The dynamics approaching the transition is anomalous because of correlations between defects and strings. The defects (dislocations) are local and owe their origin to the frustration in the system. The strings run the length of the system and define the global characteristics of the SOS surface and hence the ordering in the compressible model. In the mean-field model, there is a true instability towards the disappearance of these strings. In a non-mean-field model one expects to see a pseudo instability. To understand the slowing down of the dynamics as the temperature approaches $`T^{}`$, the instability temperature, it is useful to analyze the nature of this instability a little further. The results from the models with the quenched frustration-removing fields prove to be useful in understanding the instability at $`T^{}`$.
The ground-states of the triangular lattice antiferromagnet can be classified into sectors characterized by the string density. In the zero-defect sector, the free energy of the quenched models can be calculated as a function of the string-density and the field strength. This free energy is minimized by a particular value of the string density for a given value of the staggered field or the strength of the anisotropy. The analysis of Chen and Kardar can be repeated using the language of the strings (at least within the zero-defect sector) and indeed one finds that the “effective” free energy expressed as a function of the strain fields exhibits an instability as the temperature is lowered. The curvature of the effective free energy goes to zero at the instability temperature $`T^{}`$. Above this temperature there is a well-defined free-energy valley centered at zero strain field and a string density of $`2/3`$. As this state loses its stability the system becomes free to explore regions with finite strain. At each realization of the strain field there is, however, a well-defined minimum at a definite string -density different from $`2/3`$. The instability temperature marks the point where the system starts to see the multivalleyed nature of the free-energy surface. The exploration of this phase space can take place only through activated processes which take the system from one string sector to another. The process is activated since the number of strings can change only through the creation of defects. This picture provides a qualitative understanding of the dynamics observed in our simulations.
In continuing studies of this system, we are exploring the spatial correlations of defects. The spatio-temporal correlations which develop in this system as $`T^{}`$ is approached would be crucial in understanding the detailed nature of the long-lived structures and their similarity to the ones observed in structural glass.
This work has been supported in part by the DOE grant DE-FG02-ER45495.
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# 1 Introduction
## 1 Introduction
Some time ago it was proposed by Beneke, Buchalla and Dunietz to study mixing induced CP violation in $`B`$ mesons by using asymmetries of inclusive decay rates. This method although theoretically clean suffers from an experimental difficulty since it supposes a completely inclusive measurement. In order to solve this problem, we propose to study CP asymmetries of one-particle inclusive decays. In this paper, we shall consider CP asymmetries in the $`B_s^0`$ system. This can also be done in the $`B_d`$ system . We shall study time integrated and time dependent CP asymmetries of one-particle inclusive $`B_s^0D_sX`$ decays. The technique involved is nevertheless theoretically not as clean as the one used in for inclusive decays. Nevertheless, the measurement of CP asymmetries of one-particle inclusive decay widths would be clean.
We shall not neglect the decay width difference in the $`B_s^0`$ system, since it is sizable although recent calculations show that the decay width difference in that system could be smaller than previously expected. A method was recently proposed to calculate the decay rates appearing in these asymmetries. The predicted decay rates are compatible with current experimental knowledge and we can therefore be confident that the results we obtain for the CP asymmetries will be meaningful.
We shall concentrate on asymmetries involving one-particle inclusive decays of a $`B_s^0`$ meson. We start by introducing our notations, in section 3 we shall develop the formalism needed for time dependent CP asymmetries in one-particle CP asymmetries. We then present and discuss our results and conclude.
## 2 Definitions
In this section we shall introduce our notations. We basically use notations similar to those introduced in . The proper time evolution of an initial pure $`B_s^0`$ or $`\overline{B}_s^0`$ reads
$`B_{sphys}^0(t)`$ $`=`$ $`g_+(t)B_s^0{\displaystyle \frac{q}{p}}g_{}(t)\overline{B}_s^0`$ (1)
$`\overline{B}_{sphys}^0(t)`$ $`=`$ $`{\displaystyle \frac{p}{q}}g_{}(t)B_s^0+g_+(t)\overline{B}_s^0,`$
where the time dependent functions
$`g_+(t)`$ $`=`$ $`e^{iMt\frac{1}{2}\mathrm{\Gamma }t}\left[\mathrm{cosh}{\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{4}}\mathrm{cos}{\displaystyle \frac{\mathrm{\Delta }Mt}{2}}+i\mathrm{sinh}{\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{4}}\mathrm{sin}{\displaystyle \frac{\mathrm{\Delta }Mt}{2}}\right]`$ (2)
and
$`g_{}(t)`$ $`=`$ $`e^{iMt\frac{1}{2}\mathrm{\Gamma }t}\left[\mathrm{sinh}{\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{4}}\mathrm{cos}{\displaystyle \frac{\mathrm{\Delta }Mt}{2}}+i\mathrm{cosh}{\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{4}}\mathrm{sin}{\displaystyle \frac{\mathrm{\Delta }Mt}{2}}\right]`$ (3)
describe the particle anti-particle mixing. The mass difference $`\mathrm{\Delta }M`$ and the width difference $`\mathrm{\Delta }\mathrm{\Gamma }`$ between the neutral $`B_s^0`$ mesons are given by
$`\mathrm{\Delta }M`$ $`=`$ $`M_HM_L`$ (4)
$`\mathrm{\Delta }\mathrm{\Gamma }`$ $`=`$ $`\mathrm{\Gamma }_H\mathrm{\Gamma }_L.`$
The off diagonal term of the mass-mixing matrix is given by $`M_{12}+i\mathrm{\Gamma }_{12}`$. We define the quantity
$`{\displaystyle \frac{q}{p}}={\displaystyle \frac{\mathrm{\Delta }Mi/2\mathrm{\Delta }\mathrm{\Gamma }}{2\left(M_{12}i/2\mathrm{\Gamma }_{12}\right)}}={\displaystyle \frac{M_{12}^{}}{|M_{12}|}}\left(1{\displaystyle \frac{1}{2}}a\right),a=\mathrm{Im}{\displaystyle \frac{\mathrm{\Gamma }_{12}}{M_{12}}}.`$ (5)
The second expression for $`q/p`$ in equation (5) is valid to first order in the small quantity $`\mathrm{\Gamma }_{12}/M_{12}=𝒪(m_b^2/m_t^2)`$.
## 3 CP asymmetries in one-particle inclusive decays
In this section, we shall consider CP asymmetries of rates of one-particle inclusive $`B_s^0D_sX`$ decays. The time dependent CP asymmetry is defined by
$`𝒜_{CP}(t)={\displaystyle \frac{\mathrm{\Gamma }(B_s^0(t)D_sX)\mathrm{\Gamma }(\overline{B}_s^0(t)\overline{D}_s\overline{X})}{\mathrm{\Gamma }(B_s^0(t)D_sX)+\mathrm{\Gamma }(\overline{B}_s^0(t)\overline{D}_s\overline{X})}}.`$ (6)
In the following, we shall neglect all effects due to direct CP violation. The time dependent decay widths therefore read
$`\mathrm{\Gamma }(B_s^0(t)D_sX)=|g_+(t)|^2\mathrm{\Gamma }(B_s^0D_sX)+|{\displaystyle \frac{q}{p}}g_{}(t)|^2\mathrm{\Gamma }(B_s^0\overline{D}_sX)`$ (7)
$`2\mathrm{R}\mathrm{e}\left(g_+^{}(t){\displaystyle \frac{q}{p}}g_{}(t)T_{D_s}^{B_s^0\overline{B}_s^0}\right)`$
and
$`\mathrm{\Gamma }(\overline{B}_s^0(t)\overline{D}_s\overline{X})=|g_+(t)|^2\mathrm{\Gamma }(B_s^0D_sX)+|{\displaystyle \frac{p}{q}}g_{}(t)|^2\mathrm{\Gamma }(B_s^0\overline{D}_sX)`$ (8)
$`2\mathrm{R}\mathrm{e}\left(g_+^{}(t){\displaystyle \frac{p}{q}}g_{}(t)T_{\overline{D}_s}^{\overline{B}_s^0B_s^0}\right).`$
The decay widths appearing in these formulas were calculated in . It remains to compute the transition matrix elements of $`\mathrm{\Delta }B_s^0=2`$ given by
$`T_{D_s}^{B_s^0\overline{B}_s^0}`$ $`=`$ $`{\displaystyle \frac{1}{2m_{B_s^0}}}{\displaystyle d^4x𝑑\varphi _{D_s}\underset{X}{}(2\pi )^4\delta ^4(P_{B_s^0}P_{D_s}P_X)}`$
$`B_s^0|H_{eff}(x)|D_sXD_sX|H_{eff}^{}(0)|\overline{B}_s^0`$
and
$`T_{\overline{D}_s}^{\overline{B}_s^0B_{}^{0}{}_{s}{}^{}}`$ $`=`$ $`{\displaystyle \frac{1}{2m_{B_s^0}}}{\displaystyle d^4x𝑑\varphi _{\overline{D}_s}\underset{\overline{X}}{}(2\pi )^4\delta ^4(P_{B_s^0}P_{\overline{D}_s}P_{\overline{X}})}`$
$`\overline{B}_s^0|H_{eff}(x)|\overline{D}_s\overline{X}\overline{D}_s\overline{X}|H_{eff}^{}(0)|B_s^0,`$
where $`d\varphi _{D_s}`$ is the phase space of the $`D_s`$ meson. The part of the effective Hamiltonian which is relevant for this work is given by
$`H_{eff}`$ $`=`$ $`{\displaystyle \frac{G_F}{\sqrt{2}}}V_{cb}^{}V_{cs}\left(C_1(\overline{b}c)_{VA}(\overline{c}s)_{VA}+C_2(\overline{b}T^ac)_{VA}(\overline{c}T^as)_{VA}\right)`$
$`+{\displaystyle \frac{G_F}{\sqrt{2}}}V_{cb}^{}V_{us}\left(C_1(\overline{b}c)_{VA}(\overline{u}s)_{VA}+C_2(\overline{b}T^ac)_{VA}(\overline{u}T^as)_{VA}\right)`$
$`+{\displaystyle \frac{G_F}{\sqrt{2}}}V_{ub}^{}V_{cs}\left(C_1(\overline{b}u)_{VA}(\overline{c}s)_{VA}+C_2(\overline{b}T^au)_{VA}(\overline{c}T^as)_{VA}\right)`$
where $`(\overline{q}_1q_2)_{VA}`$ stands for $`(\overline{q}_1\gamma ^\mu (1\gamma _5)q_2)`$, the $`T^a`$ matrices are the $`SU(3)_C`$ Gell-Mann matrices and $`C_1`$ and $`C_2`$ are the Wilson coefficients. The matrix elements $`T_{D_s}^{B_s^0\overline{B}_s^0}`$ and $`T_{\overline{D}_s}^{\overline{B}_s^0B_{}^{0}{}_{s}{}^{}}`$ can be parametrized in the same way as the wrong charm decay in . Applying the Fierz transformation we can rewrite the operators into the following form $`(\overline{b}s)_{VA}(\overline{c}c)_{VA}`$, $`(\overline{b}s)_{VA}(\overline{u}c)_{VA}`$ and $`(\overline{b}s)_{VA}(\overline{c}u)_{VA}`$. Using the large $`N_C`$ limit, we can factorize $`T_{\overline{D}_s}^{\overline{B}_s^0B_s^0}`$
$`T_{\overline{D}_s}^{B_s^0\overline{B}_s^0}(M^2)`$ $`=`$ $`{\displaystyle \underset{q}{}}{\displaystyle \frac{1}{2m_{B_s^0}}}{\displaystyle \frac{G_F^2}{2}}|C_1|^2V_{cb}^{}V_{qs}V_{qb}^{}V_{cs}`$
$`{\displaystyle d^4QK^{\mu \nu }(p_{B_s^0},M,Q)𝑑\varphi _{D_s}P_{\mu \nu }^{(q)}(Q)}`$
with
$`K_{\mu \nu }(p_{B_s^0},M,Q)`$ $`=`$ $`{\displaystyle \underset{X}{}}(2\pi )^4\delta ^4(p_{B_s^0}p_XQ)`$
$`\overline{B}_s^0(p_{B_s^0})|(\overline{b}\gamma _\mu (1\gamma _5)s)|XX|(\overline{b}\gamma _\nu (1\gamma _5)s)|B_s^0(p_{B_s^0})`$
and
$`P_{\mu \nu }^{(q)}(Q)`$ $`=`$ $`{\displaystyle \underset{X^{}}{}}(2\pi )^4\delta ^4(Qp_{D_s}p_X^{})`$
$`0|(\overline{c}\gamma _\mu (1\gamma _5)q)|\overline{D}_s(p_{\overline{D}_s})X^{}\overline{D}_s(p_{\overline{D}_s})X^{}|(\overline{q}\gamma _\nu (1\gamma _5)c)|0.`$
We get a similar expression for $`T_{D_s}^{B_s^0\overline{B}_s^0}(M^2)`$. The nature of the quark $`q`$ depends on the operator under consideration. The tensor $`K_{\mu \nu }(p_{B_s^0},M,Q)`$ is fully inclusive and can be parametrized using the decay constant $`f_{B_s^0}`$ of the $`B_s^0`$ meson and $`P_{\mu \nu }^{(q)}(Q)`$ is very similar to the expression we had encountered in the calculation of the wrong charm decay width . It involves a projection on a state containing a $`D_s`$ meson. A priori we do not know how to calculate that kind of matrix element, a solution is to contract the spinor indices as in the parton calculation and to multiply this tensor by a channel dependent form factor $`f_q`$ which has to be fitted. We obtain
$$P_{\mu \nu }^{(q)}(p_D,Q)=2\pi \delta ((Qp_D)^2m_q^2)\mathrm{Tr}\{\begin{array}{c}/\hfill \\ p\hfill \end{array}_D\gamma {}_{\mu }{}^{}(\begin{array}{c}/\hfill \\ Q\hfill \end{array}\begin{array}{c}/\hfill \\ p\hfill \end{array}_D)\gamma {}_{\nu }{}^{}\}f_q.$$
(15)
We can reproduce the inclusive case by setting $`f_q=1`$.
### 3.1 CP asymmetry $`\mathrm{\Gamma }(B_s^0(t)D_s^+X)`$ vs. $`\mathrm{\Gamma }(\overline{B}_s^0(t)D_s^{}\overline{X})`$
The CP asymmetry reads
$`𝒜(t)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }(B_s^0(t)D_s^+X)\mathrm{\Gamma }(\overline{B}_s^0(t)D_s^{}\overline{X})}{\mathrm{\Gamma }(B_s^0(t)D_s^+X)+\mathrm{\Gamma }(\overline{B}_s^0(t)D_s^{}\overline{X})}}.`$ (16)
We have two quark transitions contributing to $`T_{D_s^+}^{B_s^0\overline{B}_s^0}`$, namely $`bc\overline{c}s`$ interfering with itself and $`bu\overline{c}s`$ interfering with $`bc\overline{u}s`$. These matrix elements can be modeled in the same way as it was proposed in . It corresponds to a rescaling of the parton calculation. The complication due to isospin symmetry which led to the introduction of strong phases in is not present in the $`B_s^0`$ system. We can use the same parameterization for the $`\mathrm{\Delta }B_s^0=2`$ as the one we have used for the wrong charm decay widths . The parton calculation was performed in . We obtain
$`T_{D_s^+}^{B_s^0\overline{B}_s^0}`$ $`=`$ $`{\displaystyle \frac{G_F^2m_b^2}{24\pi M_{B_s^0}}}f_{B_s^0}^2M_{B_s^0}^2V_{cb}V_{us}^{}V_{ub}V_{cs}^{}(1z)^2`$
$`\left(C_1^2(1z)B+C_1^2(1+2z){\displaystyle \frac{M_{B_s^0}^2}{m_b^2}}2B_S^0\right)𝒢`$
$`{\displaystyle \frac{G_F^2m_b^2f_{B_s^0}^2}{24\pi M_{B_s^0}}}M_{B_s^0}^2(V_{cs}^{}V_{cb})^2\sqrt{14z}`$
$`\left[(14z)C_1^2B+(1+2z)C_1^2{\displaystyle \frac{M_{B_s^0}^2}{m_b^2}}2B_S\right],`$
in the leading order of the short distance expansion and where $``$ and $`𝒢`$ are decay channel dependent non-perturbative form factors. We shall define $``$ and $`𝒢`$ later. The parameters $`B`$ and $`B_S`$ are the bag factors and $`z`$ is defined by
$`z={\displaystyle \frac{m_c^2}{m_b^2}}.`$ (18)
We have neglected the penguin operators. We also have two operators contributing to $`T_{D_s^{}}^{\overline{B}_s^0B_{}^{0}{}_{s}{}^{}}`$, as in the previous case $`bc\overline{c}s`$ interfering with itself and $`bc\overline{u}s`$ interfering with $`bu\overline{c}s`$. In the leading order of the short distance expansion we then have
$`T_{D_s^{}}^{\overline{B}_s^0B_s^0}`$ $`=`$ $`{\displaystyle \frac{G_F^2m_b^2}{24\pi M_{B_s^0}}}f_{B_s^0}^2M_{B_s^0}^2V_{cb}^{}V_{us}V_{ub}^{}V_{cs}(1z)^2`$
$`\left(C_1^2(1z)B+C_1^2(1+2z){\displaystyle \frac{M_{B_s^0}^2}{m_b^2}}2B_S\right)𝒢`$
$`{\displaystyle \frac{G_F^2m_b^2f_{B_s^0}^2}{24\pi M_{B_s^0}}}M_{B_s^0}^2(V_{cb}^{}V_{cs})^2\sqrt{14z}`$
$`\left[(14z)C_1^2B+(1+2z)C_1^2{\displaystyle \frac{M_{B_s^0}^2}{m_b^2}}2B_S\right],`$
where $``$ and $`𝒢`$ are decay channel dependent non-perturbative form factors. The function $``$ was given in . We had
$`^{B_s^0D_s^+}=^{B_s^0D_s^{}}=4f,`$ (20)
where $`f=0.121`$ was fitted in . The value of this parameter could be slightly different in the $`B_s^0`$ case and will eventually have to be extracted from experiment. It can be extracted easily from the decay channel $`B_s^0D_s^+X`$. For the decays involving a $`D_s^{}`$ we have
$`^{B_s^0D_s^+}=^{B_s^0D_s^{}}=3f.`$ (21)
We shall make the very same assumptions to model $`𝒢`$ as we made to model $``$. To do so we have to consider the contribution of $`bu\overline{c}s`$ to the decay rates of the form $`\mathrm{\Gamma }(B_s^0D_sX)`$ which we shall denote by $`\mathrm{\Gamma }(B_s^0D_sX)^{bu\overline{c}s}`$. Using the same model as the one we used for the wrong charm case in , we have
$`\mathrm{\Gamma }(B_s^0D_s^+X)^{bu\overline{c}s}=3\mathrm{\Gamma }(B_s^0D_s^+X)_{dir}^{bu\overline{c}s},`$ (22)
where $`\mathrm{\Gamma }(B_s^0D_s^+X)_{dir}^{bu\overline{c}s}=\mathrm{\Gamma }_{dir}^{bu\overline{c}s}`$ is the direct contribution to this decay over the current $`bu\overline{c}s`$. If we take into account the contribution from the decay over a $`D_s^+`$, we obtain
$`\mathrm{\Gamma }(B_s^0D_s^+X)^{bu\overline{c}s}`$ $`=`$ $`4\mathrm{\Gamma }_{dir}^{bu\overline{c}s}`$ (23)
since a $`D_s^+`$ always decays into a $`D_s^+`$. We assume that the direct contribution $`\mathrm{\Gamma }_{dir}^{bu\overline{c}s}`$ to $`\mathrm{\Gamma }(B_s^0D_s^+X)^{bu\overline{c}s}`$ is equal to the direct contribution to $`\mathrm{\Gamma }(B^0D^+X)^{bu\overline{c}s}`$ which is certainly the case in the limit of the heavy quark symmetry. We assume that every $`c`$ quark eventually hadronizes into a $`D`$ meson. Spin counting and isospin symmetry in the decays $`BD^{}X`$ and $`BD_{dir}X`$ through the transition $`bu\overline{c}s`$ allow us to deduce easily that $`\mathrm{\Gamma }_{dir}`$ is equal to $`1/8`$ times the result of the parton calculation. We then obtain the following non-perturbative form factors
$`𝒢^{B_s^0D_s^+}=𝒢^{B_s^0D_s^{}}=1/2`$ (24)
and
$`𝒢^{B_s^0D_s^+}=𝒢^{B_s^0D_s^{}}=3/8.`$ (25)
In the $`B_s^0`$ system, we can use the following approximations
$`a=0\mathrm{and}{\displaystyle \frac{M_{12}^{}}{|M_{12}|}}={\displaystyle \frac{M_{12}}{|M_{12}|}}=1.`$ (26)
These are very good approximations, theory predicts $`a<10^3`$, we then neglect the weak phase in $`\frac{M_{12}^{}}{|M_{12}|}`$. We see that the only weak phases appearing are $`e^{i\gamma }`$ in equation (3.1) and $`e^{i\gamma }`$ in equation (3.1). We can therefore rewrite $`T_{D_s}^{B_s^0\overline{B}_s^0}`$ and $`T_{\overline{D}_s}^{\overline{B}^0B^0}`$ as
$`T_{D_s^+}^{B_s^0\overline{B}_s^0}=n_1+e^{i\gamma }n_2`$ (27)
and
$`T_{D_s^{}}^{\overline{B}_s^0B_s^0}=n_1+e^{i\gamma }n_2,`$ (28)
where $`n_1`$ corresponds to the contribution from $`bc\overline{c}s`$ interfering with it-self and $`e^{i\gamma }n_2`$ corresponds to the contribution from $`bc\overline{u}s`$ interfering with $`bu\overline{c}s`$. Inserting the formulas for the time dependent decay rates in the definition for the CP asymmetry, we obtain
$`𝒜_{CP}(t)={\displaystyle \frac{2n_2\mathrm{sin}\left(\mathrm{\Delta }Mt\right)\mathrm{sin}(\gamma )}{M_1(t)2\mathrm{\Gamma }(B_s^0D_s^+X)+M_2(t)2\mathrm{\Gamma }(B_s^0D_s^{}X)+M_3(t)}},`$ (29)
where $`M_1(t),M_2(t)`$ and $`M_3(t)`$ are given by
$`M_1(t)=\mathrm{cos}^2\left({\displaystyle \frac{\mathrm{\Delta }Mt}{2}}\right)+\mathrm{sinh}^2\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{4}}\right),`$ (30)
$`M_2(t)=\mathrm{sin}^2\left({\displaystyle \frac{\mathrm{\Delta }Mt}{2}}\right)+\mathrm{sinh}^2\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{4}}\right)`$ (31)
and
$`M_3(t)=(2n_1+2n_2)\mathrm{sinh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right).`$ (32)
The time integrated CP asymmetry is the given by
$`𝒜_{CP}={\displaystyle \frac{2n_2\frac{x}{1+x^2}\mathrm{sin}(\gamma )}{N_1(x,y)\mathrm{\Gamma }(B_s^0D_s^+X)+N_2(x,y)\mathrm{\Gamma }(B_s^0D_s^{}X)+N_3(x,y)}},`$ (33)
where $`N_1(x,y)`$, $`N_2(x,y)`$ and $`N_3(x,y)`$ are given by
$`N_1(x,y)={\displaystyle \frac{2+x^2}{1+x^2}}+{\displaystyle \frac{y^2}{4y^2}},`$ (34)
$`N_2(x,y)={\displaystyle \frac{x^2}{(1+x^2)}}+{\displaystyle \frac{y^2}{(4y^2)}}`$ (35)
and
$`N_3(x,y)={\displaystyle \frac{2y}{4y^2}}\left(2n_1+2n_2\mathrm{cos}(\gamma )\right),`$ (36)
where we have introduced the parameters $`x=\mathrm{\Delta }M/\mathrm{\Gamma }`$ and $`y=\mathrm{\Delta }\mathrm{\Gamma }/\mathrm{\Gamma }`$ which can be measured. At the present time, there is only a lower bound for the parameter $`x`$ in the $`B_s^0`$ system namely $`|x|>14`$ . For numerical calculations, we shall use $`|x|=20`$ . It is also not clear what the actual value of $`y`$ is, for numerical calculations, we shall use the value computed in , $`|y|=0.054`$. We set $`B_S=B=1`$, $`m_b=4.8\mathrm{GeV}`$, $`m_{B_s^0}=5.3693\mathrm{GeV}`$ , $`m_c=1.4\mathrm{GeV}`$, $`m_{D_s}=1.9685\mathrm{GeV}`$, $`f_{B_s^0}=210\mathrm{MeV}`$ and $`C_1=1`$. We obtain numerically
$`n_1=4.931\%,`$ (37)
$`n_2=\mathrm{9.526\; 10}^2\%,`$
where $`n_1`$ and $`n_2`$ are normalized to the decay width of the $`B_s^0`$ meson, using $`\tau _{B_s^0}=\mathrm{1.54\; 10}^{12}\mathrm{s}`$. The decay width needed were computed in , we had
$`\mathrm{\Gamma }(B_s^0D_s^+X)=3.3\%`$ (38)
$`\mathrm{\Gamma }(B_s^0D_s^{}X)=64.9\%.`$
The time integrated CP asymmetry is then given by
$`𝒜_{CP}=\mathrm{1.4\; 10}^4\mathrm{sin}(\gamma ),`$ (39)
neglecting the term proportional to $`\mathrm{cos}(\gamma )`$ in the denominator.
### 3.2 CP asymmetry $`\mathrm{\Gamma }(B_s^0(t)D_s^{}X)`$ vs. $`\mathrm{\Gamma }(\overline{B}_s^0(t)D_s^+\overline{X})`$
In this case the CP asymmetry reads
$`𝒜(t)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }(B_s^0(t)D_s^{}X)\mathrm{\Gamma }(\overline{B}_s^0(t)D_s^+\overline{X})}{\mathrm{\Gamma }(B_s^0(t)D_s^{}X)+\mathrm{\Gamma }(\overline{B}_s^0(t)D_s^+\overline{X})}}.`$ (40)
We need to parameterize $`T_{D_s^{}}^{B_s^0\overline{B}_s^0}`$ and $`T_{D_s^+}^{\overline{B}_s^0B_s^0}`$. They are given by
$`T_{D_s^{}}^{B_s^0\overline{B}_s^0}`$ $`=`$ $`\left(T_{D_s^{}}^{\overline{B}_s^0B_s^0}\right)^{}=n_1+e^{i\gamma }n_2`$ (41)
$`T_{D_s^+}^{\overline{B}_s^0B_s^0}`$ $`=`$ $`\left(T_{D_s^+}^{B_s^0\overline{B}_s^0}\right)^{}=n_1+e^{i\gamma }n_2.`$
The form factors appearing in the transition matrix element were given previously. The time dependent CP asymmetry then reads
$`𝒜_{CP}(t)={\displaystyle \frac{2n_2\mathrm{sin}\left(\mathrm{\Delta }Mt\right)\mathrm{sin}(\gamma )}{M_1(t)2\mathrm{\Gamma }(B_s^0D_s^{}X)+M_2(t)2\mathrm{\Gamma }(B_s^0D_s^+X)+M_3(t)}},`$ (42)
where $`M_1(t)`$, $`M_2(t)`$ and $`M_3(t)`$ were defined in equations (30), (31) and (32). The time integrated CP asymmetry is then given by
$`𝒜_{CP}={\displaystyle \frac{2n_2\frac{x}{1+x^2}\mathrm{sin}(\gamma )}{N_1(x,y)\mathrm{\Gamma }(B_s^0D_s^{}X)+N_2(x,y)\mathrm{\Gamma }(B_s^0D_s^+X)+N_3(x,y)}},`$
where $`N_1(x,y)`$, $`N_2(x,y)`$ and $`N_3(x,y)`$ are given in equations (34), (35) and (36). The time integrated CP asymmetry is numerically given by
$`𝒜_{CP}=\mathrm{1.4\; 10}^4\mathrm{sin}(\gamma ),`$ (44)
neglecting the term proportional to $`\mathrm{cos}(\gamma )`$ in the denominator.
### 3.3 CP asymmetry $`\mathrm{\Gamma }(B_s^0(t)D_s^{}X)`$ vs. $`\mathrm{\Gamma }(\overline{B}_s^0(t)\overline{D}_s^{}\overline{X})`$
We can without any difficulty deduce from the preceding sections the formulas for the CP asymmetries in the one-particle inclusive decays of the form $`\mathrm{\Gamma }(B_s^0(t)D_s^{}X)`$ versus $`\mathrm{\Gamma }(\overline{B}_s^0(t)\overline{D}_s^{}\overline{X})`$. We have
$`n_1=3.698\%`$ (45)
$`n_2=\mathrm{7.145\; 10}^2\%,`$
which are normalized to the decay width of the $`B_s^0`$ meson. The decay widths needed were computed in , we had
$`\mathrm{\Gamma }(B_s^0D_s^+X)=2.5\%`$ (46)
$`\mathrm{\Gamma }(B_s^0D_s^{}X)=49.6\%.`$
The next CP asymmetry we shall consider is defined by
$`𝒜(t)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }(B_s^0(t)D_s^{}X)\mathrm{\Gamma }(\overline{B}_s^0(t)D_s^+\overline{X})}{\mathrm{\Gamma }(B_s^0(t)D_s^{}X)+\mathrm{\Gamma }(\overline{B}_s^0(t)D_s^+\overline{X})}}.`$ (47)
The time integrated CP asymmetry is then given by
$`𝒜_{CP}=\mathrm{1.4\; 10}^4\mathrm{sin}(\gamma ),`$ (48)
neglecting the term proportional to $`\mathrm{cos}(\gamma )`$ in the denominator. For the asymmetry defined by
$`𝒜(t)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }(B_s^0(t)D_s^{}X)\mathrm{\Gamma }(\overline{B}_s^0(t)D_s^+\overline{X})}{\mathrm{\Gamma }(B_s^0(t)D_s^{}X)+\mathrm{\Gamma }(\overline{B}_s^0(t)D_s^+\overline{X})}},`$ (49)
we obtain
$`𝒜_{CP}=\mathrm{1.4\; 10}^4\mathrm{sin}(\gamma ),`$ (50)
neglecting the term proportional to $`\mathrm{cos}(\gamma )`$ in the denominator.
## 4 Discussion of the results
We see that in principle we can extract information on $`\mathrm{sin}(\gamma )`$ from the asymmetries calculated in the previous section. But they are small, nevertheless the decay widths involved are sizable. It has been proposed to extract information on $`\mathrm{sin}(\gamma )`$ from exclusive decays (see e.g. ) but the decay widths involved are very small, typically of the order $`10^4`$ and one would have to deal with strong phases which would make the extraction of $`\mathrm{sin}(\gamma )`$ even more difficult. It could therefore be worth to try to extract $`\mathrm{sin}(\gamma )`$ from one-particle inclusive decays in the $`B_s^0`$ system. If the present method is chosen to extract $`\mathrm{sin}(\gamma )`$, it would be interesting to test its precision. This could be done by comparing the results obtained for $`\mathrm{sin}(2\beta )`$ in one-particle inclusive CP asymmetries in the $`B_d`$ system with some more conventional extraction technique like the “gold-plated” $`BJ/\mathrm{\Psi }K_S`$, although one-particle inclusive CP asymmetries in the $`B_d`$ system are theoretically not as clean as the ones in the $`B_s^0`$ system due to the presence of strong phases .
We still have some large uncertainties, some of them due to the method. The corrections to the decay widths could be fairly large, in the worth case of the order of $`30\%`$. But remember that the decay widths calculated in are compatible with current experimental knowledge. On the other hand, we have large experimental uncertainties in the values of $`x`$, $`y`$ and $`f_{B_s^0}`$.
Time dependent CP asymmetries could also allow to extract $`\mathrm{sin}(\gamma )`$, but it is not yet clear if the oscillations can be resolved. Predictions depend on the decay width difference in the $`B_s^0`$ system and it has not yet been possible to measure this quantity.
A way to improve the magnitude of the CP asymmetries would be to do an anti-lepton tagging. The decay widths in the denominator are partially responsible for the low magnitude of the asymmetries, this is particularly true for the CP asymmetries in the $`B_d^0`$ system considered in . Thus, doing an anti-lepton tagging, would slightly improve their magnitude. We would then consider asymmetries of the type
$`𝒜(t)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }(B_s^0(t)D_sX)_{\mathrm{NL}}\mathrm{\Gamma }(\overline{B}_s^0(t)\overline{D}_s\overline{X})_{\mathrm{NL}}}{\mathrm{\Gamma }(B_s^0(t)D_sX)_{\mathrm{NL}}+\mathrm{\Gamma }(\overline{B}_s^0(t)\overline{D}_s\overline{X})_{\mathrm{NL}}}},`$ (51)
where $`\mathrm{NL}`$ stands for non-leptonic. In the $`B_s^0`$ system, we would have time integrated CP asymmetries of the order $`\mathrm{2\; 10}^4\mathrm{sin}(\gamma )`$. The factor $`x=20`$ is responsible for the low magnitude of the asymmetries. In the $`B_d`$ system, this effect would obviously be larger.
## 5 Conclusion
We have discussed CP asymmetries in one-particle inclusive $`B_s^0D_sX`$ decays. The asymmetries are small but would allow to extract $`\mathrm{sin}(\gamma )`$ which is known to be difficult. So any new method is probably welcome. It has the advantage, in comparison to CP asymmetries of exclusive decays, to have some large decay widths and in comparison to CP asymmetries of inclusive decays, of being experimentally clean.
## Acknowledgements
The author is grateful to Z.Z. Xing for long discussions on $`B_s^0`$ physics and CP violation in that system and for his critical reading of this manuscript. He would also like to thank A. Leike for his useful comments.
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# Untitled Document
SPIN-2000/06
hep-th/0003004
The Holographic Principle
Opening Lecture
Gerard ’t Hooft
Institute for Theoretical Physics
University of Utrecht, Princetonplein 5
3584 CC Utrecht, the Netherlands
and
Spinoza Institute
Postbox 80.195
3508 TD Utrecht, the Netherlands
e-mail: g.thooft@phys.uu.nl
internet: http://www.phys.uu.nl/~thooft/
Abstract
After a pedagogical overview of the present status of High-Energy Physics, some problems concerning physics at the Planck scale are formulated, and an introduction is given to a notion that became known as “the holographic principle” in Planck scale physics, which is arrived at by studying quantum mechanical features of black holes.
1. Introduction. To open an International School at which many important issues of modern elementary particle physics will be discussed, it seems appropriate to start with a bird’s eye view of the recent developments in the field. Another motivation to do so is that fundamental physics today appears to have reached a new stage, at which some reflection is needed over the past, in order to explain our present standpoints and views, and to justify the kinds of questions that we think we have to ask today, in order to enable us to proceed in our field.
Before the ’70s, we had the following picture of the fundamental forces. First, there was Quantum Electrodynamics (QED), a very successful scheme to describe (nearly) all electric and magnetic features of our beloved particles. It was understood how to perform impressively accurate calculations by using perturbation expansions with respect to $`\alpha =e^2/4\pi \mathrm{}c1/137`$, a small parameter$`^\text{1}`$. Although it was understood how to renormalise the apparent divergences of the theory, it was still something of a mystery why this procedure worked at all, and indeed, sometimes (for instance when electromagnetic mass differences were to be calculated) it did not appear to work.
As for the weak forces, we only had an ‘effective’ expression for the interaction$`^\text{2}`$, that however was not renormalisable, which made it impossible to calculate any of the higher order radiative corrections.
The strong force was in an even worse state. It was usually treated as a ‘black box’, for which only the symmetry pattern was well established. Several simplistic but quite instructive models could be written down (the Gell-Mann-Lévy model$`^\text{3}`$, the dual resonance model$`^\text{4}`$), but they were mutually incompatible, and since the coupling strength was large, perturbation expansions, even if you could renormalise, appeared to be meaningless.
Then, the revolution of the ’70s came. The discovery that non-Abelian gauge theories are renormalisable$`^\text{5}`$ enabled us to make a very important step. We could now ask the question: “What is the most general perturbatively renormalisable quantum field theory?”$`^\text{6}`$
The answer turned out to be that our theory must consist of three kinds of basic particles, to be represented by fundamental fields. They are distinguished by the value $`S`$ of the intrinsic spin:
$`S=1`$: These particles cannot be described unless you have a (Abelian or non-Abelian) gauge theory. The gauge group may be any local, compact Lie group $`G`$, for instance $`G=SU(3)SU(2)U(1)\mathrm{}`$.
$`S=\frac{1}{2}`$: These particles must be described by Dirac fields, $`\psi ^L`$ and $`\psi ^R`$, of which $`\psi ^L`$ transforms as a $`\mathrm{𝟐}\times \mathrm{𝟏}`$ representation of the algebra $`SU(2)^{\mathrm{Left}}SU(2)^{\mathrm{Right}}`$ of the Lorentz group (in Euclidean notation) and $`\psi ^R`$ transforms as a $`\mathrm{𝟏}\times \mathrm{𝟐}`$. Each of these fields may be in any kind of finite-dimensional representation of the gauge group $`G`$, but there is a very important restriction: an anomaly may arise in the contributions of triangle diagrams to the matrix elements of axial currents. It is not allowed to have currents with anomalies in them, coupled to gauge fields. If $`\psi ^L`$ and $`\psi ^R`$ are in different representations of the gauge group $`G`$, one must require that their contributions to the axial anomalies nevertheless cancel out$`^\text{7}`$. In the current version of the Standard Model, this is indeed the case.
$`S=0`$: Any set of scalar fields $`\varphi `$ may be present, in any finite representation of $`G`$. Its self-interactions must be polynomial of degree four, and its interaction with the fermions must be through gauge-invariant Yukawa terms of the form $`\overline{\psi }^L\varphi \psi ^R+`$h.c. Via the Higgs mechanism$`^\text{8}`$, these fields may produce masses for vector particles and Dirac particles.$`^\text{9}`$The present Standard Model obeys all of these requirements — and more: the renormalization group equations tell us that the coupling strengths vary as the energies increase, in such a way that perturbation expansions can be applied up to extremely high energies An uncertain factor here is the Higgs self-coupling, since the Higgs mass is still unknown.$`^{\text{10}}`$$`p_(`$. Yet many questions are still not answered:
— What determined Nature’s choice for the gauge group $`G`$, the fermionic and scalar representations of $`G`$, the number of leptonic and quark generations, and the values of the coupling strengths, in particular the details of the Kobayashi-Maskawa matrix?$`^{\text{11}}`$
— How do we explain the large hierarchy of scales in Nature? The disparity between the Planck scale, the electro-weak scale, and the scale(s) of the neutrino masses is just one of several questions of this sort.
— What is the role of supersymmetry, and how is supersymmetry broken?$`^{\text{12}}`$
— How do we couple the gravitational force, and why is the cosmological constant as small as it is, or if it vanishes altogether, again, why? There is no known symmetry that ‘protects’ the cosmological constant against renormalisation effects.Different answers to some or all of these questions are presently being investigated. Judging from past experiences, it must be of extreme importance to ask the right questions.
Are there any further useful results to be expected from experiments? Three classes of experimental avenues have not yet been completed, and may give us great improvements in our understanding, although all of these are becoming more and more difficult, demanding increasing skills of the experimenters:
$`a`$) At increasing, higher energies, the following is to be expected, and I think will be done:
— The Higgs is there to be discovered.
— Supersymmetry partners of all presently known particles may be detected, hopefully some time soon. At first sight, the fact that supersymmetric patterns were discovered in nuclear physics$`^{\text{13}}`$ has little to do with the question of supersymmetry among elementary particles, but it may indicate that, as the spectrum of particles is getting more and more complex, some supersymmetric patterns might easily arise, even if there is no ‘fundamental’ reason for their existence.
— Other new structures may also be found at higher energies. The most pleasant surprises will be the unexpected ones, which may open up new fields. It is generally believed that the present model will break down beyond a TeV or so, and this energy level will be within reach in a decade or so$`^{\text{14}}`$. These high energy experiments address the unknown physics in a direct manner, and they are therefore most important.
$`b`$) On rare occasions, new results may also be expected from precision experiments at lower energies. There have been a number of interesting examples in the recent past:
— Atomic parity violations could be measured with better than one percent precision, in spite of the fact that these are minute effects$`^{\text{15}}`$, yielding independent confirmation of the effects due to $`W`$ and $`Z`$ exchanges within atomic nuclei.
— The $`K_L/K_S`$ system is a beautiful laboratory. Precision measurements can be made of the parameter $`\epsilon ^{}/\epsilon `$, which may reveal features from deeply inside, or possibly beyond, the Standard Model, as we will learn at this School$`^{\text{16}}`$.
— Other known fundamental principles of Theoretical Physics can be put to a test, such as $`CPT`$ invariance, relativity tests, the ratios $`Q/M`$ can be compared between particles and antiparticles$`^{\text{17}}`$, the Quantum Mechanics of gravitating systems can be investigated, etc.
— New ideas were launched suggesting that Newton’s law of the gravitational force might change at scales below a mm. This can be experimentally tested.$`^{\text{18}}`$
— Tiny mass terms that produce mixing between various neutrino species can be detected in dedicated experiments.$`^{\text{19}}`$
— And there are doubtlessly many more subtle effects that may be discovered and that will alter our views concerning the fundamental interactions.$`c`$) A third source of information is cosmology. It used to be well within the domain of Science-Fiction, but nowadays cosmological models are becoming more mature. They yield precise predictions that can be verified by astrophysical observations. Models of the inflationary universe probe deeply into regions at extremely high energy, and so the information they deliver is unique:
— Structures in the spectrum of the cosmic background radiation are predicted and more detailed observations are to be expected.
— The distribution of galaxies is speculated to be due to quantum fluctuations in a very early universe. They will be calculated and compared to what is observed.
— The search for dark matter continues. The outcome will deeply affect our thinking about the fundamental interactions.
— Statistical analysis of distant galaxies may finally also reveal the presence of a cosmological constant term in the Einstein-Hilbert action.
— Other tests of the models, for instance the baryon-antibaryon asymmetry and $`CP`$ violation.
In spite of this long list, there are reasons to worry about the increasingly difficult barriers from behind which we are trying to understand the small-distance structure of our world. Which purely theoretical approaches will help us find the answers? We have to concentrate on fundamental inconsistencies in our present picture. There are many of these:
— As it was already mentioned, the hierarchies seen in the distance scales are not properly explained by what is presently known.
— The apparent absence of a cosmological constant is at odds with what we understand about Quantum Gravity.
— Indeed, quantizing gravity is still a deep problem. Superstring theory is vigorously trying to bring the gravitational force under control, but it surely is a wild animal. From superstrings came $`D`$-branes, from $`D`$-branes came “$`M`$-theory”, but it has as yet not been possible to even come close to an accurate formulation of the laws. These ideas are of extreme importance, but new avenues must still be found. What is known for sure is that Quantum Mechanics works, that the gravitational force exists, and that General Relativity works. The approach advocated by me during the last decades is to consider in a direct way the problems that arise when one tries to combine these theories, in particular the problem of gravitational instability. These considerations have now led to what is called “the Holographic Principle”, and it in turn led to the more speculative idea of deterministic quantum gravity. This theory, and the effects of dissipation of information, will be discussed in a separate lecture. Our central issue is: What is Nature’s bookkeeping system at the Planck scale?
2. The inevitable existence of black holes In sufficiently large amounts of matter, gravitational collapse is inevitable.$`^{\text{20}}`$ There are various ways to derive this fact. First, one may consider a stationary, spherically symmetric configuration of matter, held together by gravity. Near the surface, we assume that there is a region $`r=r_1`$ where the temperature is low enough so that the density $`\varrho _1`$ there is sufficiently high, say that of water. Assume that inside the sphere $`r=r_1`$, there is a certain amount of mass $`M_1`$. The pressure $`p`$ at this surface is still negligible. It is now easy to argue that $`M_1`$ must be subject to an upper limit. In any case, the pressure $`p`$ rises if we look at smaller distances $`r`$ from the centre. If the density were constant, and the general relativistic effects negligible, then one could readily compute the pressure at the centre. However, the density is likely to increase if we go down. In fact, if we assume our material to be non-exotic, then a finite compressibility follows. Matter is defined to be non-exotic if the speed of sound $`v_s`$ is less than the speed of light, $`c`$. This means that the gradient of the gravitational field will rise, and hence the gradient in the pressure will become steeper, and this will cause an instability. If $`M_1`$ was chosen large enough, there will be a point $`r=r_2<r_1`$ where the pressure becomes infinite. Even without any other relativistic arguments, this gives us as a limit: $`\frac{2}{3}\pi G_N\varrho _1r_1^2<1`$. Adding general relativistic effects correctly will give a more stringent limit, as I will briefly explain later.
But first: does matter have to be non-exotic? Suppose the speed of sound exceeds the speed of light. Would there be an immediate contradiction? Special relativity would normally demand that no signal can go faster than light. However, the reason why we demand this is causality: no signal should be able to propagate backwards in time. This is then combined with demanding Lorentz-invariance. However, matter in equilibrium represents a preferred Lorentz frame, so we could drop the latter demand. Still, there must be restrictions. Consider two regions in which matter has different local velocities. Imagine two adjacent pipes in which matter streams in opposite directions, and in both, sound goes faster than light, also in opposite directions. Due to the time shifts when Lorentz transforming, an outside observer may see both signals move backwards in time. This, in principle, could then generate a closed loop of information transport in space-time, which is an undesirable situation. This we must forbid.
Even exotic matter, however, will not be able to stop black holes from being formed. This is seen if we insert the complete general relativistic equations instead of our above pseudo-relativistic argument. These are the so-called Tolman-Oppenheimer-Volkoff equations$`^{\text{21}}`$ These equations tell us how density and pressure increase when followed inwards, given some equation of state. It is an elementary exercise to solve these equations for constant density. Even then, one finds that, due to space-time curvature, the pressure diverges to infinity at a finite radius, if we start with too much mass and a too high density at a too small radius on the surface.
We can however also produce a black hole from ordinary matter at zero pressure. Consider a spherically symmetric arrangement of matter in the form of a shell, with some finite thickness. We allow the shell to contract due to its own gravitational field. Inside the shell, there is no gravitational field at all, something that one can understand using the same arguments that tell us that inside a conducting metal sphere there is no electric or magnetic field. If the original amount of material was big enough, the contraction will proceed, and, in the limit of zero pressure and purely radial, spherically symmetric motion, the equations can easily be solved exactly. We obtain flat space-time inside, and a pure Schwarzschild metric outside. As the ball contracts, a moment will arrive when the Schwarzschild horizon appears. From that moment on, an outside observer will no longer detect any radiation from the shell, but a black hole instead. 3. Hawking radiation and quantum states The standard generally relativistic black hole solution has as a special feature that shortly after its formation, no signals will be seen coming out. It should be truly black. As is well-known, this picture changed when Hawking$`^{\text{22}}`$ discovered an elementary consequence of quantum field theory when applied to fields living in the black hole metric. The rearrangement of creation and annihilation operators is such that the states near the horizon are not truly vacuum, but they contain a precisely computable density of particles, which are emitted as black body radiation at a temperature given by$`^{\text{22}}`$
$$kT_H=\frac{\mathrm{}c^3}{8\pi GM_{\mathrm{BH}}}.$$
$`(3.1)`$
This result allows us to compute the density of quantum states of a black hole. The easiest way to do this is by using thermodynamics. However, one could object that a black hole is not truly in thermodynamic equilibrium; if energy is added to a black hole, its mass and its size will increase, and consequently its temperature will drop.
We can avoid thermodynamics by deriving the spectral density of a black hole directly from its Hawking temperature. All one needs is some form of time reversal invariance$`^{\text{23}}`$. We have at our disposal both the emission rate (the Hawking radiation intensity), and the capture probability, or the effective cross section of the black hole for infalling matter.
In units at which $`G=\mathrm{}=c=1`$, the cross section $`\sigma `$ is approximately:
$$\sigma \mathrm{\hspace{0.17em}2}\pi R^2=\mathrm{\hspace{0.17em}8}\pi M^2,$$
$`(3.2)`$
and slightly more for objects moving in slowly. The emission probability $`W\mathrm{d}t`$ for a given particle type, in a given quantum state, in a large volume $`V=L^3`$ is:
$$W\mathrm{d}t=\frac{\sigma (𝐤)v}{V}e^{E/kT}\mathrm{d}t,$$
$`(3.3)`$
where $`𝐤`$ is the wave number characterizing the quantum state, $`v`$ is the particle velocity, and $`E`$ is its momentum.
Now we assume that the process is also governed by a Schrödinger equation. This means that there exist quantum mechanical transition amplitudes,
$$\begin{array}{cc}\hfill 𝒯_{\mathrm{in}}& =_{\mathrm{BH}}M+E|𝒯|M_{\mathrm{BH}}|E_{\mathrm{in}},\hfill \\ \hfill \mathrm{and}𝒯_{\mathrm{out}}& =_{\mathrm{BH}}M|_{\mathrm{out}}E|𝒯|M+E_{\mathrm{BH}},\hfill \end{array}$$
$`(3.4)`$
where the states $`|M_{\mathrm{BH}}`$ represent black hole states with mass $`M`$, and the other states are energy eigenstates of particles in the volume $`V`$. In terms of these amplitudes, using the so-called Fermi Golden Rule, the cross section and the emission probabilities can be written as
$$\begin{array}{ccc}\hfill \sigma & =|𝒯_{\mathrm{in}}|^2\varrho (M+E)/v,\hfill & (3.5)\hfill \\ \hfill W& =|𝒯_{\mathrm{out}}|^2\varrho (M)\frac{1}{V}.\hfill & (3.6)\hfill \end{array}$$
where $`\varrho (M)`$ stands for the level density of a black hole with mass $`M`$. The factor $`v^1`$ in Eq. (3.5) is a kinematical factor, and the factor $`V^1`$ in $`W`$ arises from the normalization of the wave function.
Now, time reversal invariance relates $`𝒯_{\mathrm{in}}`$ to $`𝒯_{\mathrm{out}}`$. To be precise, all one needs is $`PCT`$ invariance, since the parity transformation $`P`$ and charge conjugation $`C`$ have no effect on our calculation of $`\sigma `$. Dividing the expressions (3.5) and (3.6), and using (3.3), one finds:
$$\frac{\varrho (M+E)}{\varrho (M)}=e^{E/kT}=e^{8\pi ME}.$$
$`(3.7)`$
This is easy to integrate:
$$\varrho (M)=e^{4\pi M^2+C}=e^S.$$
$`(3.8)`$
We can rewrite this as
$$\varrho (M)=\mathrm{\hspace{0.17em}2}^{A/A_0},$$
$`(3.9)`$
where $`A`$ is the horizon area and $`A_0`$ is a fundamental unit of area,
$$A_0=\mathrm{\hspace{0.17em}4}\mathrm{ln}2L_{\mathrm{Planck}}^2.$$
$`(3.10)`$
This suggests a spin-like degree of freedom on all surface elements of size $`A_0`$, see Fig. 1.
Figure 1. Information on a black hole horizon.
The importance of this derivation is the fact that the expressions used as starting points are the actual Hawking emission rate and the actual black hole absorption cross section. This implies that, if in more detailed considerations divergences are found near the horizon, these divergences should not be used as arguments to adjust the relation between entropy and level density by large renormalization factors.
4. The quantum information problem. It is tempting to conclude from the arguments presented above that the ‘black hole states’ form a natural extension of the spectrum of elementary particles. The lightest particles are known and have been identified as photons, neutrinos, electrons, muons, mesons, baryons, and onwards to the heavy leptons, the Higgs and so forth. The series could continue with as yet unknown particles in the ‘desert’ between 1 TeV and $`10^{19}`$ GeV, and beyond that region the first superstring recurrences could exist. The ‘most pointlike objects’ beyond the Planck mass must undoubtedly be black holes, simply because any sufficiently compact object with sufficiently high mass must carry a gravitational field and a horizon associated with that. Apparently, we now know the spectrum of the objects in this range, apart from the unknown multiplicative constant $`e^C`$ in Eq. (3.8).
It should be possible to handle these objects just as all quantum objects when we consider quantum mechanical amplitudes at high energies: they are represented as propagators describing intermediate states. Theoretical Physics should give us the computational rules, comparable to Feynman rules, for computing these amplitudes. What have we got?
The behaviour of quantum fields near the horizon of a black hole follows from the expression for $`\mathrm{d}s`$, the infinitesimal invariant distance element according to General Relativity:
$$\mathrm{d}s^2=\left(1\frac{2M}{r}\right)\mathrm{d}t^2+\frac{\mathrm{d}r^2}{12M/r}+r^2\mathrm{d}\mathrm{\Omega }^2,$$
$`(4.1)`$
where $`\mathrm{d}\mathrm{\Omega }^2`$ stands for $`\mathrm{d}\theta ^2+\mathrm{sin}^2\theta \mathrm{d}\varphi ^2`$. Writing
$$r2M=e^\sigma ,\mathrm{d}r=(r2M)\mathrm{d}\sigma ,$$
$`(4.2)`$
we see that
$$\mathrm{d}s^2=\left(1\frac{2M}{r}\right)\left(\mathrm{d}t^2+r^2\mathrm{d}\sigma ^2+\frac{r^2}{12M/r}\mathrm{d}\mathrm{\Omega }^2\right).$$
$`(4.3)`$
At high energies, the conformal factor has little effect on the wave equations for light particle species. Consider a wave equation close to the horizon, such as
$$_\mu \sqrt{g}g^{\mu \nu }_\nu \varphi m^2\sqrt{g}\varphi =0.$$
$`(4.4)`$
Since it contains the inverse of the metric tensor $`g_{\mu \nu }`$, the contribution from the angular part in Eq. (4.3) becomes insignificant, and likewise the mass term. Thus, the wave equations become 2-dimensional, and they generate plane waves in the $`\sigma `$-$`t`$ direction as if $`\sigma `$-$`t`$ space were flat. This means that one gets an unlimited number of oscillations near the horizon, as $`\sigma \mathrm{}`$. Quite unlike the case for systems such as the hydrogen atom, we see that the boundary condition at the horizon is ill-defined. The quantum states can generate an unlimited number of modes here. From this exercise, one would conclude that the density of quantum states for a black hole is not at all finite.
The physical origin of this divergence is not difficult to identify. Particles may move into a black hole, but, as long as we stick to linear field equations, the particles emerging from the black hole cannot at all be related to the ingoing ones. There cannot be any reflection against the horizon, since there should be an infinite time delay. Here we see that the situation with ingoing and outgoing spinors and vector particles will be equally hopeless.
Two other conundrums are closely related to the problem just signalled. First, we have the quantum decoherence problem. This problem becomes apparent when the Hawking effect is calculated explicitly. The initial state of elementary particles before the formation of a black hole is described in terms of various Fourier modes of their fields. All of these modes then are associated to observable operators. In a Heisenberg picture of the quantum states, these operators become time dependent. Part of the Fourier modes of the initial fields now enter into the black hole, and only the operators associated to the modes that emerge out of the black hole correspond to observables at a later time. The expectation values of these late-time observables turn out to be described by a thermal density matrix. In terms of the basis of states generated by the late-time observables alone, this density matrix turns out to have eigenvalues less than one, which is characteristic for a not fully coherent quantum state. This situation is similar to what one gets in a condensed matter system if one allows observable particles to escape, and subsequently omits the quantum states that they represent.
In the case of a black hole, the missing particles are the absorbed ones. If we were forced to keep these particles in our quantum description, an even worse infinity of quantum states would result.
A description of the ‘information problem’ that is easier to understand is the following. Choose a coordinate frame in which the formation of a black hole looks more or less regular. Ingoing particles are then seen to enter at rather late times. If now the returning particles were assumed to be not totally independent of the ingoing ones, one would have to accept the observation that, somehow, the information contained in the ingoing particles has been transferred to the outgoing ones. The outgoing ones, however, all belong to the Fourier modes that arose as quantum oscillations at the point where the black hole was formed, way back in the past. How could the required information have been imprinted on these particles, if they have already been there for such a long time?
A mathematically impeccable observation was made by Hawking$`^{\text{24}}`$: the black hole space-time describes two universes, not one, and these two are connected by a ‘wormhole’. After a black hole is formed, the quantum wave functions of elementary particles spread over these two universes, and they become intertwined. Cutting off the information concerning the contents of the ‘hidden’ universe will leave the other universe in a quantum mechanically decoherent state.
From a physical point of view, however, this argument is unsatisfactory. It implies that black holes are fundamentally different from all other forms of matter in the sense that they appear to produce decoherence. In all respects this result is equivalent to saying that the scattering matrix elements involving black holes are not fixed by our theory, but carry an uncertainty, distributed in some well-defined way. So, what we really have here is an ‘uncertain theory’. Our theory is incomplete. We should not be satisfied with that. Perhaps new physics can remove this uncertainty.
5. The Scattering Matrix Ansatz. How do we ‘improve’ our theory? Naturally, one may think of including more interactions; obviously, the procedures applied thus far assumed ingoing and outgoing fields not to interact — Eq. (4.4) is after all linear in $`\varphi `$. At first sight it seems that including interactions will resolve the paradox. As $`\sigma `$ approaches $`\mathrm{}`$, the plane $`\varphi `$ waves enter conformally into ever smaller regions of space and time. Effectively, the gravitational couplings increase rapidly, and as $`e^\sigma `$ approaches the Planck length, this effective coupling becomes super strong. An alternative way to verify this is by switching towards a coordinate frame that is locally regular near the horizon. Such a frame is given, e.g., by the Kruskal coordinates $`\{x,y\}`$:
$$xy=\left(\frac{r}{2M}1\right)e^{r/2M};x/y=e^{t/2M}.$$
$`(5.1)`$
Writing
$$x^0=xy;x^1=x+y,$$
$`(5.2)`$
one finds the metric to be regular near $`xy0`$. Particles sent in in the far past will align close to the axis $`x=0`$, and particles going out in the distant future align close to $`y=0`$. A boost in the Schwarzschild time parameter $`t`$ corresponds to a Lorentz transformation in $`(x,y)`$ space, where the scale is set by the mass parameter $`M`$. If we consider time lapses long compared to $`M`$, the Lorentz boosts separating ingoing and outgoing particles become horrendous. Thus we see that ingoing and outgoing particles meet each other near $`x=y=0`$ at tremendously large c.m. energies. Even if we could neglect Standard Model interactions at these energies, the gravitational interactions, which grow with the energy squared, can no longer be ignored from some point onwards.
This observation however does not resolve the decoherence problem. Even with the interactions in place, one may still argue that information is drained by the black hole, and a theory for pure states interacting with pure states without decoherence does not follow. A more powerful approach is wanted.
It is strongly advocated now to start from the other end: we must assume that there exists a quantum mechanically fully coherent scattering matrix $`S`$. The assumption is somewhat dogmatic; we cannot prove it from first principles, other than demanding the existence of a theory. Even if standard techniques at best only provide us with some ‘distribution’ for the physical scattering matrix elements, we assume that the ‘true’ scattering matrix elements are exactly defined. Even if no theory would exist to derive them, they could in principle be derived from experiment.
Demanding consistency with existing theory however gives us important constraints. Indeed, the scattering matrix can now almost be derived from the information we already have. The calculations have been presented at length elsewhere, so here we give a summary.
The dominant interaction is assumed to be the gravitational one, simply because the c.m. energies tend to infinity. Other interactions, such as in particular the electro-magnetic one, can be corrected for later (the effects from electro-magnetism are important, but they do not affect the main structure that will be obtained). The procedure then is as follows.$`^{\text{25}}`$
First, assume a black hole with some well-specified initial history of ingoing particles, for instance we specify the way in which a star imploded to give this black hole, and afterwards more objects may have fallen in at later times. We assume that all this leads to a pure quantum state, to be referred to as the state $`|1`$. It evolves and decays in some prescribed way. It leads to some superposition of many possible states for the outgoing particles, including states describing the final explosion.
Now, we consider the same state $`|1`$, but we either add or remove one ingoing particle, and we call this state $`|1,\delta p`$, where $`\delta p`$ stands for the momentum (and possible other details) of the extra ingoing object. What can be done now is a calculation, as detailed as possible, of the effects this extra ingoing particle has on the outgoing objects. Surely there is interaction. The gravitational one is most interesting. It leads to a shift of the outgoing wave functions. This shift is simply the Shapiro delay due to the gravitational field of the ingoing object. The calculation is in principle entirely straightforward, but has to be done with some care since the ingoing object goes essentially with the speed of light. The shift of the Kruskal $`y`$ coordinate is found to be
$$\delta y=p_{\mathrm{in}}G(|\stackrel{~}{x}_{\mathrm{in}}\stackrel{~}{x}_{\mathrm{out}}|),$$
$`(5.3)`$
where $`G`$ is a simple calculable function of the coordinates $`\stackrel{~}{x}`$ on the horizon. In the limit where the black hole is large and the separations on the horizon small, we can approximately view $`\stackrel{~}{x}`$ as flat coordinates, and in that limit, the function $`G`$ is proportional to $`\mathrm{log}|\stackrel{~}{x}_{\mathrm{in}}\stackrel{~}{x}_{\mathrm{out}}|`$. We then find that this shift obeys a Laplace equation:
$$\stackrel{~}{}^2\delta y=Cp_{\mathrm{in}}\delta ^2(\stackrel{~}{x}_{\mathrm{in}}\stackrel{~}{x}_{\mathrm{out}}).$$
$`(5.4)`$
$`G`$ is therefore a Green function. Because of this shift, the outgoing state $`|\psi _{\mathrm{out}}`$ turns into
$$\mathrm{exp}\left(i\mathrm{d}^2\stackrel{~}{x}P_{\mathrm{out}}^+(\stackrel{~}{x})\delta y(\stackrel{~}{x})\right)|\psi _{\mathrm{out}}.$$
$`(5.5)`$
Writing $`\delta y(\stackrel{~}{x})=\mathrm{d}^2\stackrel{~}{x}^{}G(\stackrel{~}{x}\stackrel{~}{x}^{})\delta p_{\mathrm{in}}(\stackrel{~}{x}^{})`$, the new state is
$$|\psi _{\mathrm{out}}^{}=\left(\mathrm{exp}i\mathrm{d}^2\stackrel{~}{x}\mathrm{d}^2\stackrel{~}{x}^{}P_{\mathrm{out}}^+(\stackrel{~}{x})G(\stackrel{~}{x}\stackrel{~}{x}^{})\delta p_{\mathrm{in}}(\stackrel{~}{x}^{})\right)|\psi _{\mathrm{out}}.$$
$`(5.6)`$
And now we can repeat this many times. Let $`P^{}(\stackrel{~}{x}^{})`$ stand for the total momentum (in Kruskal coordinates) for all particles added, or subtracted using a minus sign, from the state $`|1`$. Then we have
$$|\psi _{\mathrm{out}}=\left(\mathrm{exp}i\mathrm{d}^2\stackrel{~}{x}\mathrm{d}^2\stackrel{~}{x}^{}P_{\mathrm{out}}^+(\stackrel{~}{x})G(\stackrel{~}{x}\stackrel{~}{x}^{})P_{\mathrm{in}}^{}(\stackrel{~}{x}^{})\right)|1.$$
$`(5.6)`$
This way, the state $`|1`$ can be used as a universal reference state. The true state is then specified by giving $`P_{\mathrm{in}}^{}(\stackrel{~}{x})`$.
After some simple manipulations, we find that both the initial and the final state could be described by giving the transverse coordinates $`\stackrel{~}{x}^{(i)}`$ and the radial momenta $`p^{(i)}`$ for all particles. We get
$$\begin{array}{cc}\hfill _{\mathrm{out}}& q^1,\stackrel{~}{y}^1,q^2,\stackrel{~}{y}^2,\mathrm{}||p^1,\stackrel{~}{x}^1,p^2,\stackrel{~}{x}^2,\mathrm{}_{\mathrm{in}}=\hfill \\ & =𝒟X^+(\stackrel{~}{x})𝒟X^{}(\stackrel{~}{x})\mathrm{exp}(\mathrm{d}^2\stackrel{~}{x}[i\stackrel{~}{}X^+\stackrel{~}{}X^{}+\hfill \\ & +\underset{i}{}i\delta ^2(\stackrel{~}{x}\stackrel{~}{x}^i)p^iX^+(\stackrel{~}{x})\underset{i}{}i\delta ^2(\stackrel{~}{x}\stackrel{~}{y}^i)q^iX^{}(\stackrel{~}{x})]),\hfill \end{array}$$
$`(5.7)`$
and if the in- and outgoing particles are described by transverse wave functions $`e^{i\stackrel{~}{p}^i\stackrel{~}{x}^i}`$ and $`e^{i\stackrel{~}{q}^i\stackrel{~}{y}^i}`$, then another set of integrations has to be performed, over the transverse coordinates $`\stackrel{~}{x}^i`$ and $`\stackrel{~}{y}^i`$. All of this yields an amplitude that is very much reminiscent of a string amplitude, with the exception of the $`i`$ in front of the ‘kinetic’ term $`(X)^2`$ in Eq. (5.7). In Eq. (5.7), Newton’s constant $`G_N`$ has been normalized according to
$$8\pi G_N=1.$$
$`(5.8)`$
6. Fock space. The result of the previous section appears to be beautiful. We managed to construct the $`S`$-matrix using only known facts about the gravitational interaction between fast moving objects. In addition, it appears not to be too difficult to impose unitarity for this scattering matrix. Unitarity just fixes the measure of the functional integration in (5.7). Only the phase then remains undetermined, but it was arbitrary anyway since the amplitudes in question violate many of the conventional conservation laws such as all combinations of baryon and lepton number.
There are, however, two problems, both having to do with the Hilbert space in terms of which this scattering matrix appears to be defined.
Problem # 1: the space of all momenta, $`\{P^\pm (\stackrel{~}{x})\}`$, is infinite dimensional, even for small black holes, whereas we expected a finite total number of states (the entropy was supposed to be finite). So, this Hilbert space is far too large.
Problem # 2: the space of all momentum distributions, $`\{P^\pm (\stackrel{~}{x})\}`$, is far too small to accommodate for all possible particle configurations. If two or more particles enter at the same transverse point $`\stackrel{~}{x}`$, then, in our expressions, only the total momentum counts. States for which the total momentum distributions are identical will be indistinguishable, and since we want our scattering matrix to be unitary, these states must be identical. This is not Fock space for elementary particles as we are used to.
It is important to note, on the other hand, that string amplitudes, which are like Eq. (5.7) but without the $`i`$ in the kinetic term, share the same feature: states with two or more particles entering the string world sheet at one point, are indistinguishable from states with just a single particle entering at that point. Distinctions only come after the $`\stackrel{~}{x}`$ integrations, at which the particle number becomes unambiguous.
7. The holographic principle.$`^{\text{26}}`$ We have reached a point where, for a proper description of the particle states in the vicinity of a black hole, a two-dimensional function is required: the momentum distribution over a two-dimensional coordinate on the horizon. In addition, this function must be further reduced, since it must effectively contain not more than one $`Z(2)`$ variable per surface element $`A_0`$ (see Eq. (3.10)). A comparison with a holographic photograph is quickly made. In a holographic set-up, a laser beam shines onto some three-dimensional object, and the reflected light interferes with an unperturbed laser beam. The interference pattern is registered on a photographic plate. In turn, after having developed the plate, we can shine a laser beam on it. An image of the three-dimensional object re-emerges. This appears to be a way to register three-dimensional objects on a two-dimensional photographic plate.
Now imagine that the photographic plate has a limited resolution, and that its colouring can only be black-and-white, no gray tones. In that case, the image we see of the original object will be blurred somewhat, since information went astray. This must actually be the situation in our description of particles entering a black hole: the momentum distribution cannot represent as many details as a fully three-dimensional description: our image of the universe is blurred. Of course, since it is the Planck scale where this limit is attained, in practice we perceive our universe very sharply.
Although this holographic nature of our description of the particles appears to apply only for particles entering a black hole, one may argue that it must have a much more universal validity. According to general relativity, there should exist a direct mapping that relates physical phenomena in one setting (with a gravitational field present) to another one (freely falling coordinates). Normally, the mapping goes both ways. It is indeed unlikely that freely falling particles can be described in more detail than the limits set by the holographic principle: one bit of information per surface element of size $`A_0`$. It can be computed that the energy needed to detect more details would be so large that gravitational collapse would be inevitable; the entire scene would be absorbed by a black hole – and indeed be impossible to observe at all!
This is what we found out about Nature’s book keeping system: the data can be written onto a surface, and the pen with which the data are written has a finite size.
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26. G. ’t Hooft, “Dimensional Reduction in Quantum Gravity”, Essay dedicated to Abdus Salam, Utrecht preprint THU-93/26 (gr-qc/9310026); id., “Black holes and the dimensionality of space-time”, in Proceedings of the Symposium “The Oskar Klein Centenary”, 19-21 Sept. 1994, Stockholm, Sweden. Ed. U. Lindström, World Scientific 1995, p. 122; L. Susskind, L. Thorlacius and J. Uglum, Phys. Rev. D48 (1993) 3743 (hep-th 9306069).
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# Meron-cluster algorithms and chiral symmetry breaking in a (2+1)-d staggered fermion model 11footnote 1This work is supported in part by funds provided by the U.S. Department of Energy (D.O.E.) under cooperative research agreement DE-FC02-94ER40818 and by the Schweizerischer Nationalfonds.
## 1 Introduction
There are a number of models of interest which suffer from a very severe sign problem. This includes QCD and other field theories with a non-zero chemical potential or a non-zero vacuum angle or odd numbers of fermion flavors, frustrated quantum spin systems, like the quantum antiferromagnet in an external magnetic field, and models for strongly-correlated electrons, like the Hubbard model for high-temperature superconductivity. These models have a Boltzmann weight which can be negative or even complex and so cannot be interpreted as a probability. This difficulty can be overcome in numerical simulations by including the sign or phase of the Boltzmann weight with observables. Unfortunately, this leads to large cancellations and gives exponentially small observables. This requires exponentially large statistics, which makes it in practice impossible to simulate these models numerically.
Recently, a new technique has been developed, called Meron-Cluster algorithms , which completely solves the sign problem for some of these models . It identifies the origin of the sign problem with properties of the clusters, which enables it to be eliminated. Cluster algorithms in general are extremely efficient at exploring configuration spaces and very often do not suffer from critical slowing down as a phase transition is approached, unlike many other algorithms. For example, in this and previous papers, we can work directly in the chiral limit with massless fermions. Combined with the ability to construct improved estimators, we can perform a high precision study of these models with only modest statistics.
In this paper, we explore a model of $`N=1`$ flavor of staggered fermions in (2+1)-dimensions with a four-fermion interaction. This model has a very severe sign problem and cannot be simulated with standard techniques. We build a Meron-Cluster algorithm, which we use to perform a high-precision study. We find that the $`𝖹𝖹(2)`$ chiral symmetry of this model is spontaneously broken at low temperatures and, using finite-size scaling analysis, we verify that the finite-temperature chiral phase transition is in the universality class of the 2-d Ising model. A recent study of the same model with $`𝖹𝖹(2)`$ chiral symmetry in (3+1)-dimensions has shown, using a Meron-Cluster algorithm, that a finite temperature chiral phase transition occurs which has the universal behavior of the 3-d Ising model . The work presented in this paper concerns a different universality class and also constructs more observables than were previously considered.
The identification of the finite temperature critical behavior is not entirely straightforward . A model of $`N`$ fermion flavors with a four-fermion interaction shows mean-field behavior in the $`N=\mathrm{}`$ limit. On the other hand, at finite $`N`$ one finds the non-trivial critical behavior that one expects based on dimensional reduction and standard universality arguments. For example, in it has been verified that the chiral phase transition in a $`(2+1)`$-d four-fermion interaction model with $`N=4`$ flavors and $`𝖹𝖹(2)`$ chiral symmetry is in the universality class of the 2-d Ising model. Due to the fermion sign problem, standard fermion simulation methods often do not work in models with too small a number of flavors. The work presented in this paper shows that the same universal behavior holds for $`N=1`$ flavor.
The standard technique to deal with fermions in Monte-Carlo simulations is to integrate them out, resulting in a non-local bosonic theory. The Meron-Cluster algorithm does not integrate out the fermions, instead we describe them with a local theory using a Fock space basis of occupation number. The fermion sign arises as a non-local property due to the permutation of fermion world lines. Using probabilistic rules, we connect neighboring lattice sites, producing closed loops, which are the clusters. A cluster is flipped by making all of its occupied sites empty and its empty ones occupied. Such a cluster flip can change the fermion sign by changing the permutation of fermion world lines. A cluster whose flip changes the sign we call a meron. We can tell if a cluster is a meron simply from its structure. A typical configuration contains many merons, yet the observables of interest are only non-zero for configurations with very few or no merons. The signals from standard Monte Carlo algorithms are so exponentially small because the Markov chain explores a vast configuration space, yet only an exponentially small sub-space makes any contribution to measurables. By restricting ourselves to only explore the relevant sub-space, we completely solve the sign problem.
This paper is organized as follows. In section 2, we present the fermionic model which we have studied, calculate its partition function using the Hamiltonian formulation and find that there is a sign problem. In section 3, we describe briefly the Meron-Cluster algorithm which we have used to perform numerical simulations of this model. We present the results of the simulations in section 4 and give our conclusions in section 5.
## 2 The Staggered Fermion Model
We consider staggered fermions in the Hamiltonian formulation on a 2-dimensional spatial lattice of extent $`L`$, which is even. The Hamiltonian operator is
$`H`$ $`=`$ $`{\displaystyle \underset{x,i}{}}h_{x,i}+m{\displaystyle \underset{x}{}}(1)^{x_1+x_2}\mathrm{\Psi }_x^+\mathrm{\Psi }_x`$
$`h_{x,i}`$ $`=`$ $`\eta _{x,i}(\mathrm{\Psi }_x^+\mathrm{\Psi }_{x+\widehat{i}}+\mathrm{\Psi }_{x+\widehat{i}}^+\mathrm{\Psi }_x)+G(\mathrm{\Psi }_x^+\mathrm{\Psi }_x{\displaystyle \frac{1}{2}})(\mathrm{\Psi }_{x+\widehat{i}}^+\mathrm{\Psi }_{x+\widehat{i}}{\displaystyle \frac{1}{2}}),`$ (2.1)
where $`\eta _{x,1}=1`$ and $`\eta _{x,2}=(1)^{x_1}`$ are the standard Kawamoto-Smit phases for staggered fermions and $`G`$ is a constant. The fermionic operators satisfy the usual anticommutation relations $`\{\mathrm{\Psi }_x,\mathrm{\Psi }_y\}=\{\mathrm{\Psi }_x^+,\mathrm{\Psi }_y^+\}=0,\{\mathrm{\Psi }_x^+,\mathrm{\Psi }_y\}=\delta _{xy}`$. The same model in (3+1)-dimensions was explored in . We refer the reader to this paper, where various features of the model and the Meron-Cluster algorithm are discussed in more detail than we give here.
In the Hamiltonian formulation of the theory, fermion doubling on the lattice occurs only in the spatial dimensions. Using staggered fermions, the Dirac components of a spinor are distributed spatially, reducing the number of fermion flavors by a factor of four. Thus this (2+1)-dimensional model contains $`N=1`$ fermion flavor. The model has a global $`U(1)`$ symmetry corresponding to conserved particle number, as the total particle number operator commutes with the Hamiltonian
$$N=\underset{x}{}\mathrm{\Psi }_x^+\mathrm{\Psi }_x,[H,N]=0.$$
(2.2)
Furthermore the Hamiltonian has, for $`m=0`$, a discrete $`𝖹𝖹(2)`$ symmetry corresponding to shifts by one lattice spacing. However the mass term breaks that symmetry explicitly. For a single flavor of massless fermions, the symmetry of the lattice model is $`U(1)𝖹𝖹(2)`$ and we refer to the discrete symmetry as chiral symmetry. In the continuum, a single massless fermion flavor has a $`U(1)`$ axial symmetry (there is no gauge interaction, so this symmetry is not anomalously broken). The discrete $`𝖹𝖹(2)`$ symmetry is the lattice remnant of this continuous symmetry. From now on, we set $`m=0`$ and explore the behavior of the chiral symmetry of this model. The symmetries of staggered fermions are discussed in detail in Ref. . If the $`𝖹𝖹(2)`$ chiral symmetry is spontaneously broken at some finite temperature, from universality we expect this to be a second-order phase transition. As the critical point is approached, the correlation length $`\xi `$ diverges and the system becomes insensitive to the time extent. Due to dimensional reduction, we expect a finite-temperature chiral phase transition in this model to belong to the 2-d Ising universality class.
The partition function of the model is
$`Z=\text{Tr}[\mathrm{exp}(\beta H)]=\underset{M\mathrm{}}{lim}\text{Tr}[\mathrm{exp}(ϵH)]^M`$ (2.3)
$`=\underset{M\mathrm{}}{lim}\text{Tr}[\mathrm{exp}(ϵH_1)\mathrm{exp}(ϵH_2)\mathrm{exp}(ϵH_3)\mathrm{exp}(ϵH_4)]^M,`$
where we use the Suzuki-Trotter decomposition to divide the Euclidean time extent $`\beta `$ into $`4M`$ time slices, the lattice spacing in the time direction being $`ϵ=\beta /M`$. The Hamiltonian operator is decomposed into four parts $`H=H_1+H_2+H_3+H_4`$. All of the terms that contribute to a particular $`H_i`$ commute with one another, as each term is an interaction between nearest-neighbors and each lattice site appears in only one such nearest-neighbor pair. However, the $`H_i`$ do not commute with one another. We note that it is not actually necessary to discretize the time direction, as it is possible to work directly in the Euclidean time continuum .
We can equivalently describe this model with bosonic operators, using a transformation by Jordan and Wigner . We order the lattice sites on each time slice arbitrarily into a chain, which can be done in any number of spatial dimensions. For example, a possible ordering of points in two spatial dimensions is by an index $`l=x_1+(x_21)L`$. The fermionic operators are now represented by a chain of Pauli matrices
$`\mathrm{\Psi }_x^+=\sigma _1^3\sigma _2^3\mathrm{}\sigma _{l1}^3\sigma _l^+,\mathrm{\Psi }_x=\sigma _1^3\sigma _2^3\mathrm{}\sigma _{l1}^3\sigma _l^{},\mathrm{\Psi }_x^+\mathrm{\Psi }_x={\displaystyle \frac{1}{2}}(\sigma _l^3+1)`$ (2.4)
$`\sigma ^\pm ={\displaystyle \frac{1}{2}}(\sigma ^1\pm i\sigma ^2),[\sigma _l^i,\sigma _m^j]=2i\delta _{lm}ϵ^{ijk}\sigma _l^k,`$
where the spatial position $`x`$ is denoted by the index $`l`$ and the Pauli matrices satisfy the usual commutation relations. To calculate the partition function of the theory, we use the Fock space basis of occupation number $`n_x=0,1`$ i.e. the eigenstates of $`\sigma ^3`$. The occupied and empty states are respectively $`|1`$ and $`|0`$, which satisfy $`\sigma ^3|1=|1`$ and $`\sigma ^3|0=|0`$.
The time evolution operator $`\mathrm{exp}(ϵH_i)`$ acts on a time slice of occupation number states, producing the next time slice. This is decomposed into the product of operators $`\mathrm{exp}(ϵh_{x,i})`$ acting on nearest-neighbor occupation states. The transfer matrix is
$$\mathrm{exp}(ϵh_{x,i})=\mathrm{exp}(\frac{ϵG}{4})\left(\begin{array}{cccc}\mathrm{exp}(\frac{ϵG}{2})& 0& 0& 0\\ 0& \mathrm{cosh}\frac{ϵ}{2}& \mathrm{\Sigma }\mathrm{sinh}\frac{ϵ}{2}& 0\\ 0& \mathrm{\Sigma }\mathrm{sinh}\frac{ϵ}{2}& \mathrm{cosh}\frac{ϵ}{2}& 0\\ 0& 0& 0& \mathrm{exp}(\frac{ϵG}{2})\end{array}\right),$$
the basis being $`|00,|01,|10`$ and $`|11`$, where e.g. $`|01`$ represents state $`|0`$ at $`x`$ and $`|1`$ at $`x+\widehat{i}`$. If these nearest-neighbors are labelled $`l`$ and $`m`$, the off-diagonal transfer matrix elements have a factor $`\mathrm{\Sigma }=\eta _{x,i}\sigma _{l+1}^3\sigma _{l+2}^3\mathrm{}\sigma _{m1}^3`$. Note that this operator is diagonal in the occupation number basis.
The partition function of the theory is given as a path integral
$$Z_f=\underset{n}{}\text{Sign}[n]\mathrm{exp}(S[n]),$$
(2.5)
where we sum over all possible configurations of occupation numbers $`n(x,t)=0,1`$ on a $`(2+1)`$-d space-time lattice of points $`(x,t)`$. The Boltzmann factor $`\mathrm{exp}(S[n])`$ for a configuration is the product of the Boltzmann factors for each space-time plaquette $`\mathrm{exp}(s[n(x,t),n(x+\widehat{i},t),n(x,t+1),n(x+\widehat{i},t+1)])`$, which are
$`\mathrm{exp}(s[0,0,0,0])=\mathrm{exp}(s[1,1,1,1])=\mathrm{exp}({\displaystyle \frac{ϵG}{2}}),`$
$`\mathrm{exp}(s[0,1,0,1])=\mathrm{exp}(s[1,0,1,0])=\mathrm{cosh}{\displaystyle \frac{ϵ}{2}},`$
$`\mathrm{exp}(s[0,1,1,0])=\mathrm{exp}(s[1,0,0,1])=\mathrm{sinh}{\displaystyle \frac{ϵ}{2}}.`$ (2.6)
All other plaquettes are illegal and have Boltzmann weight zero, as they represent non-conservation of fermion number. Any configuration which contains illegal plaquettes has itself Boltzmann weight zero and makes no contribution to the partition function. We are only interested in legal configurations, which have to satisfy several constraints. Note that here we have dropped the overall factor $`\mathrm{exp}(ϵG/4)`$ that appeared in eq.(2). The sign of a configuration, $`\text{Sign}[n]`$, is also a product of space-time plaquette contributions $`\text{sign}[n(x,t),n(x+\widehat{i},t),n(x,t+1),n(x+\widehat{i},t+1)]`$ with
$`\text{sign}[0,0,0,0]=\text{sign}[0,1,0,1]=\text{sign}[1,0,1,0]=\text{sign}[1,1,1,1]=1,`$
$`\text{sign}[0,1,1,0]=\text{sign}[1,0,0,1]=\mathrm{\Sigma }.`$ (2.7)
The occupied lattice sites define world-lines of fermions, which close due to the periodicity of the Euclidean time direction. The world-lines are free to permute during their time evolution as the fermions interchange position and each configuration has a well-defined permutation of fermions. The Pauli exclusion principle tells us that the sign of a configuration is the permutation sign of the fermions, hence $`\text{Sign}[n]=\pm 1`$. This non-local effect is contained in the factors $`\mathrm{\Sigma }`$ of each space-time plaquette.
The expectation value of a fermionic observable $`A[n]`$ is given by
$`A_f={\displaystyle \frac{1}{Z_f}}{\displaystyle \underset{n}{}}A[n]\text{Sign}[n]\mathrm{exp}(S[n])={\displaystyle \frac{A\text{Sign}}{\text{Sign}}},`$ (2.8)
$`\text{Sign}={\displaystyle \frac{1}{Z_b}}{\displaystyle \underset{n}{}}\text{Sign}[n]\mathrm{exp}(S[n]),`$
where $`\mathrm{}`$ means a measurement made in the bosonic ensemble, whose partition function is $`Z_b=_n\mathrm{exp}(S[n])`$. To measure one fermionic observable requires two bosonic measurements. The quantities of physical interest which we measure are the chiral condensate $`\overline{\mathrm{\Psi }}\mathrm{\Psi }`$, the chiral susceptibility $`\chi `$ and a Binder cumulant $`U`$ of the chiral condensate, respectively
$`\overline{\mathrm{\Psi }}\mathrm{\Psi }[n]={\displaystyle \frac{ϵ}{4}}{\displaystyle \underset{x,t}{}}(1)^{x_1+x_2}(n(x,t){\displaystyle \frac{1}{2}}),`$
$`\chi ={\displaystyle \frac{1}{\beta V}}(\overline{\mathrm{\Psi }}\mathrm{\Psi })^2_f,`$ $`U=1{\displaystyle \frac{(\overline{\mathrm{\Psi }}\mathrm{\Psi })^4_f}{3[(\overline{\mathrm{\Psi }}\mathrm{\Psi })^2_f]^2}}.`$ (2.9)
## 3 The Meron-Cluster Algorithm
We now describe briefly the Meron-Cluster algorithm which we used to sample the bosonic ensemble corresponding to the fermionic model without the sign factor. We set $`G=1`$, for which the bosonic model is the isotropic antiferromagnetic quantum Heisenberg model, whose Hamiltonian is $`H=_{x,i}(S_x^1S_{x+\widehat{i}}^1+S_x^2S_{x+\widehat{i}}^2+S_x^3S_{x+\widehat{i}}^3)`$, where $`S_x^i=\frac{1}{2}\sigma _l^i`$ is a spin $`1/2`$ operator at the lattice site $`x`$, labelled by $`l`$ in the Jordan-Wigner chain. There already exist extremely efficient cluster algorithms to simulate bosonic quantum spin systems , and the first cluster algorithm for lattice fermions was constructed in . These algorithms can be implemented directly in the time continuum , i.e. the Suzuki-Trotter time discretization is not even necessary. In this study, we discretize the time direction.
We use the same algorithm that was used in . Each configuration is decomposed into a set of clusters, which consist of connected lattice sites. A new configuration is generated by flipping the clusters. When a cluster is flipped, all lattice sites contained in that cluster change occupation number from $`n(x,t)`$ to $`1n(x,t)`$, i.e. the occupied sites become empty and the empty ones occupied. To build the clusters, a probabilistic choice is made in each space-time interaction plaquette $`[n(x,t),n(x+\widehat{i},t),n(x,t+1),n(x+\widehat{i},t+1)]`$ as to which neighboring lattice sites are connected to one another. A cluster is a sequence of connected sites. In this algorithm, the clusters are closed loops. The probabilistic choices (called cluster break-ups) which build the clusters are designed to obey detailed balance and we only allow break-ups which generate legal plaquettes under cluster flips. The cluster rules are illustrated in Table 1. For plaquette configurations $`[0,0,0,0]`$ and $`[1,1,1,1]`$, i.e. entirely empty or entirely occupied, we always connect sites with their time-like neighbors. For configurations $`[1,0,0,1]`$ and $`[0,1,1,0]`$ where a fermion hops to a neighboring site, we always connect sites with their space-like neighbors. For configurations $`[1,0,1,0]`$ and $`[0,1,0,1]`$, i.e. a static fermion next to an empty site, we connect the sites with their time-like neighbors with probability $`p=2/[1+\mathrm{exp}(ϵ/2)]`$ and with their space-like neighbors with probability $`1p`$. This algorithm was also used in . It is extremely efficient, has almost no detectable autocorrelations and its dynamical exponent for critical slowing down is compatible with zero.
Each cluster has two orientations, with lattice site occupancies $`n(x,t)`$ and $`1n(x,t)`$. When a cluster is flipped, the new configuration which is generated may have a different sign from the previous one, depending on whether or not the permutation of fermion world-lines is changed. A cluster whose flip changes $`\text{Sign}[n]`$ we call a meron, those which leave $`\text{Sign}[n]`$ unchanged we call non-merons. Flipping a meron changes the topology of the fermion world-lines. The term meron has been used before to denote half-instantons , such as in the 2-d $`O(3)`$ model at non-zero vacuum angle $`\theta `$ . The number of merons in a configuration is always even, as flipping all clusters leaves the sign unchanged. An example of a meron-cluster is given in Figure 1. When the meron-cluster is flipped the first configuration with $`\text{Sign}[n]=1`$ turns into the second configuration with $`\text{Sign}[n]=1`$. For cluster algorithms more general than the one described here, it is not always possible to identify certain clusters as merons .
The meron concept alone gives us an exponential gain in statistics. Starting from a configuration containing $`N_C`$ clusters, we consider the ensemble of $`2^{N_C}`$ configurations where we allow all possible cluster orientations. If a configuration contains no merons, all configurations in the ensemble have $`\text{Sign}[n]=1`$. However, if it contains merons, half the ensemble has $`\text{Sign}[n]=1`$ and the other half $`\text{Sign}[n]=1`$, which exactly cancel, giving a contribution 0. The improved estimator gives $`\text{Sign}=\delta _{N,0}`$, i.e. the probability that a configuration contains $`N=0`$ merons, which is an exponential improvement on standard algorithms, which measure a statistical average of $`\pm 1`$. As explained in , this solves half of the sign problem.
We also construct improved estimators for observables. The chiral susceptibility is
$$\chi =\frac{1}{\beta V}(\overline{\mathrm{\Psi }}\mathrm{\Psi })^2_f=\frac{1}{\beta V}\frac{(\overline{\mathrm{\Psi }}\mathrm{\Psi })^2\text{Sign}}{\text{Sign}}.$$
(3.1)
The total chiral condensate for a given configuration, $`\overline{\mathrm{\Psi }}\mathrm{\Psi }[n]=_C\overline{\mathrm{\Psi }}\mathrm{\Psi }_C`$, is a sum of cluster contributions. Averaging $`\chi `$ over the ensemble of $`2^{N_C}`$ configurations gives
$$\chi =\frac{_C|\overline{\mathrm{\Psi }}\mathrm{\Psi }_C|^2\delta _{N,0}+2|\overline{\mathrm{\Psi }}\mathrm{\Psi }_{C_1}||\overline{\mathrm{\Psi }}\mathrm{\Psi }_{C_2}|\delta _{N,2}}{\beta V\delta _{N,0}},$$
(3.2)
This only gets contributions from configurations with $`N=0`$ or $`N=2`$ merons ($`C_1`$ and $`C_2`$ are the two merons). The vast majority of configurations contain many merons, but they make no contribution to observables. The zero- and two-meron sectors of configuration space are exponentially small, but they contain all of the contributions to $`\chi `$. Restricting ourselves to only explore this sub-space, we exponentially enhance both the numerator and denominator of eq.(3.2), leaving the ratio invariant. This solves the remaining half of the sign problem.
For the Binder cumulant $`U`$, we need to measure $`(\overline{\mathrm{\Psi }}\mathrm{\Psi })^4_f`$ and hence
$$\text{Sign}(\overline{\mathrm{\Psi }}\mathrm{\Psi })^4=\text{Sign}\underset{C_i,C_j,C_k,C_l}{}\overline{\mathrm{\Psi }}\mathrm{\Psi }_{C_i}\overline{\mathrm{\Psi }}\mathrm{\Psi }_{C_j}\overline{\mathrm{\Psi }}\mathrm{\Psi }_{C_k}\overline{\mathrm{\Psi }}\mathrm{\Psi }_{C_l}.$$
(3.3)
A cluster’s condensate contribution $`\overline{\mathrm{\Psi }}\mathrm{\Psi }_C`$ changes sign when the cluster is flipped. When a meron-cluster is flipped, Sign is changed. The non-zero terms in $`\text{Sign}(\overline{\mathrm{\Psi }}\mathrm{\Psi })^4`$ do not change sign if any cluster in the configuration is flipped. These non-zero terms must contain odd powers of $`\overline{\mathrm{\Psi }}\mathrm{\Psi }_C`$ for all merons $`C`$ in the configuration and even powers of $`\overline{\mathrm{\Psi }}\mathrm{\Psi }_C^{}`$ for all non-merons $`C^{}`$. The average over the ensemble of $`2^{N_C}`$ configurations is
$`\text{Sign}(\overline{\mathrm{\Psi }}\mathrm{\Psi })^4_{2^{N_C}}=\delta _{N,0}\left[{\displaystyle \underset{C}{}}|\overline{\mathrm{\Psi }}\mathrm{\Psi }_C|^4+6{\displaystyle \underset{C,C^{}}{}}|\overline{\mathrm{\Psi }}\mathrm{\Psi }_C|^2|\overline{\mathrm{\Psi }}\mathrm{\Psi }_C^{}|^2\right]`$ (3.4)
$`+\delta _{N,2}\left[4|\overline{\mathrm{\Psi }}\mathrm{\Psi }_{C_1}|^3|\overline{\mathrm{\Psi }}\mathrm{\Psi }_{C_2}|+4|\overline{\mathrm{\Psi }}\mathrm{\Psi }_{C_2}|^3|\overline{\mathrm{\Psi }}\mathrm{\Psi }_{C_1}|+12{\displaystyle \underset{C}{}}|\overline{\mathrm{\Psi }}\mathrm{\Psi }_C|^2|\overline{\mathrm{\Psi }}\mathrm{\Psi }_{C_1}||\overline{\mathrm{\Psi }}\mathrm{\Psi }_{C_2}|\right]`$
$`+\delta _{N,4}\left[24|\overline{\mathrm{\Psi }}\mathrm{\Psi }_{C_1}||\overline{\mathrm{\Psi }}\mathrm{\Psi }_{C_2}||\overline{\mathrm{\Psi }}\mathrm{\Psi }_{C_3}||\overline{\mathrm{\Psi }}\mathrm{\Psi }_{C_4}|\right],`$
where $`N`$ is the number of merons in the configurations, $`C_1,C_2,C_3`$ and $`C_4`$ are the merons and all sums in eq.(3.4) are over non-meron clusters. This average only gets contributions from the zero-, two- and four-meron sectors and so we need only explore this sub-space. We average this quantity over the complete bosonic ensemble to measure $`\text{Sign}(\overline{\mathrm{\Psi }}\mathrm{\Psi })^4`$ and hence $`U`$.
Consider the case of measuring $`\chi `$. We expect that $`p(0)/p(2)(|C|/V\beta )^2`$, where $`p(0)`$ and $`p(2)`$ are the probabilities that a configuration has zero or two merons and $`|C|`$ is the average cluster size. In large volumes, the majority of configurations has two merons, contributing $`0`$ to $`\text{Sign}`$. For even greater accuracy, we reweight the meron-sectors with trial probabilities $`p_t(0)`$ and $`p_t(2)`$, so that they appear with roughly equal frequency. This gives
$$\chi =\frac{_C|\overline{\mathrm{\Psi }}\mathrm{\Psi }_C|^2\delta _{N,0}p_t(0)+2|\overline{\mathrm{\Psi }}\mathrm{\Psi }_{C_1}||\overline{\mathrm{\Psi }}\mathrm{\Psi }_{C_2}|\delta _{N,2}p_t(2)}{\beta V\delta _{N,0}p_t(0)}.$$
(3.5)
The reweighting probabilities can be adjusted to minimize the statistical error. This technique was previously used in . To measure the Binder cumulant $`U`$, we use reweighting probabilities $`p_t(0),p_t(2)`$ and $`p_t(4)`$.
## 4 Numerical Results
We have performed simulations of the staggered fermion model on lattices with antiperiodic spatial boundary conditions from $`L=4`$ up to $`L=30`$ at inverse temperatures in the range $`\beta [1.0,3.0]`$, which includes the critical temperature where the chiral symmetry is spontaneously broken. We have made separate runs with either a fixed number of time slices (typically $`M=10`$, i.e. 40 time slices) or with fixed lattice spacing in the time direction ($`ϵ=0.1`$). In each simulation, we have made at least 1000 thermalization sweeps followed by 10000 measurements, with these numbers increased by a factor of 10 for $`L10`$. In one sweep of the lattice, a new cluster connection is proposed on each interaction plaquette and each cluster is flipped with probability $`1/2`$. To find the optimal reweighting probabilities $`p_t(N)`$ which minimize the statistical error, we first make a sample run without reweighting, only exploring the relevant meron-sectors. The observed relative weights are then used in production runs, where the sectors appear with equal probability. The major part of the sign problem is removed by the improved estimators, but reweighting is necessary for accurate measurements in large volumes.
A sample of the data measured is given in Table 2. The Table contains $`\text{Sign}`$ and the susceptibility $`\chi `$ measured over all meron-sectors, and the reweighted $`\text{Sign}_r`$ and $`\chi _r`$ measured over the zero- and two-meron sectors only with the reweighting factor $`p_t(0)/p_t(2)`$. All of these data are produced with 1000 thermalization sweeps and 10000 measurements. As all of the contributions to $`\chi `$ come from the zero- and two-meron sectors, $`\chi `$ and $`\chi _r`$ should be identical. Note that $`\text{Sign}_r`$, the fraction of zero-meron configurations generated by sampling the zero- and two-meron sectors only, is typically a lot bigger than $`\text{Sign}`$, the fraction of zero-meron configurations generated over all meron sectors. In small space-time volumes, $`\chi `$ can be accurately measured even when sampling all meron sectors. However, in large space-time volumes, $`\text{Sign}`$ is too small to be measured and we can only determine the susceptibility by restricting ourselves to the zero- and two-meron sectors. The staggered fermion model suffers from a very severe sign problem which is solved by the Meron-Cluster algorithm.
Figure 2 shows the meron number probability distribution in an algorithm that samples all meron sectors without reweighting. For small volumes the zero-meron sector and hence $`\text{Sign}`$ are relatively large, while multi-meron configurations are rare. On the other hand, in larger volumes the vast majority of configurations has a large number of merons and hence $`\text{Sign}`$ is exponentially small. For example, an extrapolation from smaller volumes gives a rough estimate for the non-reweighted $`\text{Sign}10^9`$ on the $`L=28`$ lattice at $`\beta =2.4`$., while the reweighted $`\text{Sign}_r=0.329(9)`$. Even if the configurations are entirely uncorrelated, to achieve a similar accuracy without the meron-cluster algorithm one would have to increase the statistics by a factor $`10^{18}`$, which is obviously impossible. In fact, at present there is no other method that can be used to simulate this model.
Figure 3 shows the chiral susceptibility $`\chi `$ as a function of $`\beta `$ for various spatial sizes $`L`$. At high temperatures (small $`\beta `$) $`\chi `$ is almost independent of the volume, indicating that chiral symmetry is intact. On the other hand, at low temperatures (large $`\beta `$) $`\chi `$ increases with the volume, which implies that chiral symmetry is spontaneously broken. To study the critical behavior in detail, we have performed a finite-size scaling analysis for $`\chi `$ focusing on the range $`\beta [2.2,2.6]`$ around the critical point. Since a $`𝖹𝖹(2)`$ chiral symmetry is spontaneously broken at finite temperature in this $`(2+1)`$-d model, one expects to find the critical behavior of the 2-d Ising model. The corresponding finite-size scaling formula valid close to $`\beta _c`$ is
$`\chi (L,\beta )=a(x)+b(y)L^{\gamma /\nu },`$
$`a(x)=a_0+a_1x+a_2x^2+\mathrm{},x=\beta \beta _c,`$
$`b(y)=b_0+b_1y+b_2y^2+\mathrm{},y=(\beta \beta _c)L^{1/\nu }.`$ (4.1)
For the 2-d Ising model the critical exponents are given by $`\nu =1.0`$ and $`\gamma /\nu =1.75`$. Assuming these values for the exponents, we obtain $`\beta _c=2.43(1)`$ for fixed $`ϵ=0.1`$ from the finite-size scaling fit, with a chi squared per degree of freedom of 0.84. The fit of the data is plotted in Figure 4. The value of $`\beta _c`$ is slightly dependent on $`ϵ`$.
In the finite-size scaling equation (4), for large enough $`L`$ one can neglect the term $`a(x)`$. Then $`\chi /L^{\gamma /\nu }`$ is a function of $`y=(\beta \beta _c)L^{1/\nu }`$ alone, i.e. the susceptibility data in various volumes at various $`\beta `$ can be described by one universal function. We have varied the value $`\beta _c`$ to find if all the data can be collapsed onto one universal curve. In Figure 5, we plot the universal curve obtained by taking $`\beta _c=2.43`$. The excellent agreement over a large range of spatial volumes $`L`$ and inverse temperatures $`\beta `$ is an indication of the quality of the finite-size scaling fit.
We also measure $`U_L`$, the Binder cumulant in volumes of extent $`L`$. In Figure 6, we plot the expected behavior of $`U_L`$ as $`L`$ increases for different temperatures. For $`T>T_c`$, the chiral symmetry is intact and $`U_L`$ flows into the $`T=\mathrm{}`$ fixed point $`U=0`$. For $`T<T_c`$, the chiral symmetry is spontaneously broken and $`U_L`$ flows into the $`T=0`$ fixed point $`U=2/3`$. If the universality class has a non-trivial fixed point $`U=U_{}`$, then $`U_L`$ flows into this value at $`T=T_c`$. By measuring $`U_L`$ in various volumes at many different temperatures, we determine this flow numerically. We have measured the Binder cumulant values in volumes up to $`L=30`$ and we plot some of these values as a function of $`1/L`$ in Figure 7. These measurements are made with the number of time slices fixed at $`40`$. Each curve in the figure represents some fixed temperature. In Figure 7, for small $`\beta `$ (i.e. high temperatures), $`U_L`$ clearly flows into the infinite temperature fixed point $`U=0`$, while for $`\beta `$ large (low temperatures), $`U_L`$ flows into the zero temperature fixed point $`U=2/3`$. For $`\beta `$ close to $`\beta _c`$, we have to go to larger volumes to see this behavior. Near $`\beta =2.35`$, the cumulant values appear to flow into a non-trivial fixed point $`U_{}`$. Examining this region closely, we estimate the critical inverse temperature as $`\beta _c=2.36(2)`$ and the fixed point value $`U_{}=0.60(1)`$. The finite-size scaling fit of $`\chi `$ measured at this $`ϵ0.24`$ gives the same value of $`\beta _c`$. Note that this deviates slightly from the critical temperature measured at $`ϵ=0.1`$. The universal fixed point value for the 2-d Ising model is estimated as $`U_{}0.58`$ . This is further evidence that the chiral phase transition belongs to the 2-d Ising universality class.
## 5 Conclusions
The Meron-Cluster algorithm has recently been developed to allow numerical simulations in models which suffer from a very severe sign problem. In this paper, we have applied this technique to investigate a model of staggered fermions. Unlike standard methods, which integrate out the fermions, resulting in a non-local bosonic action, we use a Fock space of occupation number to describe the fermions. We have a local bosonic action, with an additional non-local sign factor which contains the Fermi statistics. Due to the Pauli exclusion principle, configurations which have an odd permutation of fermion world lines have a negative sign. This sign leads to very large cancellations in observables and usually makes it impossible to make accurate measurements in numerical simulations. The Meron-Cluster algorithm decomposes every configuration into closed loops of connected sites, each loop being a cluster. Loops which change the fermion sign when flipped are identified as meron-clusters. A meron-cluster identifies a pair of configurations with equal weight and opposite sign. This results in an exact cancellation of two contributions $`\pm 1`$ to the path integral, such that only configurations without merons contribute to the partition function. Observables only receive contributions from configurations which contain very few or no merons, whereas the vast majority of configurations contain many merons. By only exploring the sectors of configuration space with the relevant numbers of merons, one makes an exponential gain in statistics. Combined with efficient re-weighting of the remaining meron sectors, this completely solves the sign problem. Cluster algorithms are extremely efficient at exploring configuration spaces and generating uncorrelated configurations and generally do not suffer from critical slowing down. Even in models without a sign problem, the Meron-Cluster algorithm is more efficient than standard fermion simulation methods.
In this paper, we examined a model of $`N=1`$ flavor of staggered fermions in $`(2+1)`$-dimensions, which has a $`𝖹𝖹(2)`$ chiral symmetry. The model has a very severe sign problem and cannot be solved by standard fermion simulation algorithms. Using a Meron-Cluster algorithm, we were able to make high-precision measurements of the chiral susceptibility and Binder cumulant even in very large volumes and low temperatures. In order to perform an accurate and reliable finite-size scaling analysis, it was necessary to go to volumes so large, where the sign problem is so severe, that a standard algorithm would require statistics on the order of $`10^{18}`$ to attain a similar accuracy. We were able to verify that the model undergoes a finite-temperature chiral phase transition, which belongs to the universality class of the $`2`$d Ising model. This is the behavior expected from dimensional reduction and universality. The same universal behavior was observed in the $`N=4`$ flavor case . However, the standard fermion algorithm that was used in that study does not work for $`N<4`$ due to the fermion sign problem.
It is quite natural to use cluster algorithms in models of discrete variables. A future possible application of the Meron-Cluster algorithm is in exploring quantum link models which are used in the D-theory formulation of QCD . In D-theory, a model of discrete quantum variables undergoes dimensional reduction, resulting in an effective theory of continuous classical variables. In quantum link QCD the quarks arise as domain wall fermions. The application of meron-cluster algorithms to domain wall fermions is in progress. Also there are many applications to sign problems in condensed matter physics. Investigations of antiferromagnets in a magnetic field and of systems in the Hubbard model family are given in Ref..
At present, the Meron-Cluster algorithm is the only method that allows us to solve the fermion sign problem. A severe sign problem arises in lattice QCD calculations at non-zero baryon number due to a complex action. It is therefore natural to ask if our algorithm can be applied to this case. At non-zero chemical potential the 2-d $`O(3)`$ model, which is a toy model for QCD, also suffers from a sign problem due to a complex action. When applied to the D-theory formulation of this model, the Meron-Cluster algorithm solves the sign problem completely . It is an open question if such progress can be made in investigations of QCD.
## Acknowledgements
We would like to thank Shailesh Chandrasekharan and Uwe-Jens Wiese for helpful discussions.
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# On a generalization of the logistic map
## I Introduction
In the theory of nonlinear systems, the logistic map
$$x_{n+1}=f_r(x_n)=rx_n(1x_n),$$
(1)
with $`0r4`$, $`x_n[0,1]`$, and $`n=0,1,2,\mathrm{}`$, is well-known to provide one of the simplest example of what is referred to as a chaotic system. That is, a system which, for some range of its parameters, possesses at least one bounded orbit $`\{x_0,x_1,\mathrm{}\}`$ such that (i) $`\{x_0,x_1,\mathrm{}\}`$ is not asymptotically periodic and (ii) the Lyapunov exponent, defined in the context of Eq.(1) by the usual relation
$$\lambda (x_0)=\underset{N\mathrm{}}{lim}\underset{i=0}{\overset{N1}{}}\mathrm{ln}\left|_xf_r(x_i)\right|,$$
(2)
is greater than zero . For the logistic map at $`r=4`$, specifically, $`\lambda (x_0)`$ is positive and equals $`\mathrm{ln}2`$ for almost all $`x_0`$. This value can be easily calculated by noting that the logistic map is conjugate to the tent map or by having recourse to the fact that the dynamical system defined by Eq.(1) is ergodic for this particular value of $`r`$ . In this latter case, one is justified to calculate analytically the Lyapunov exponent as an ensemble average
$$\lambda =_0^1\rho (x)\mathrm{ln}\left|_xf_r(x)\right|dx,$$
(3)
the function
$$\rho (x)=\frac{1}{\pi \sqrt{x(1x)}},x[0,1]$$
(4)
being the invariant density of the map satisfying $`\rho (A)=\rho (f^1(A))`$, where $`f^1(A)`$ is the preimage of an arbitrary subset $`A`$ of the unit interval .
In this paper, we shall focus on another property of the logistic map at $`r=4`$, namely the existence of the following closed-form solution
$$x_n=\mathrm{sin}^2(2^n\pi \theta _0),$$
(5)
where $`\theta _0=\pi ^1\mathrm{arcsin}(\sqrt{x_0})`$. The importance of this formula relies evidently on the fact that $`x_n`$ can be evaluated directly, for any initial point $`x_0`$, without actually computing the intermediate values $`x_1,x_2,\mathrm{},x_{n1}`$. In the remaining of this work, we follow the steps used in the derivation of the above formula to construct similar expressions involving general terms of the form $`\mathrm{sin}^2(N^n\theta )`$. In doing so, we shall see that a whole new family of polynomial maps can be defined which, essentially, generalize the logistic map by preserving the density Eq.(4), and whose Lyapunov exponents can be calculated easily. Other properties of these newly defined dynamical maps are also studied. In particular, we shall point out some similarities between the maps and the Tschebysheff polynomials, and prove, finally, that the family of maps as a whole is closed under the composition of functions.
## II Definition of the logistic family
The explicit solution of the logistic map at the particular value $`r=4`$ can be derived by substituting in $`f_4(x_n)=4x_n(1x_n)`$ the change of coordinates $`x_n=\mathrm{sin}^2(\pi \theta _n)`$, valid for $`x_n[0,1]`$, and by using the trigonometric identity
$$4\mathrm{sin}^2\theta (1\mathrm{sin}^2\theta )=\mathrm{sin}^2(2\theta )$$
(6)
in order to obtain $`\mathrm{sin}^2(\pi \theta _{n+1})=\mathrm{sin}^2(2\pi \theta _n)`$. This is equivalent to the map $`\theta _{n+1}=2\theta _nmod1`$, which has the explicit solution
$$\theta _n=2^n\theta _0mod1.$$
(7)
Hence, the complete solution in terms of the coordinate $`x`$ must correspond to Eq.(5).
Following these steps, one can imagine to construct a map having the solution $`x_n=\mathrm{sin}(3^n\pi \theta _0)`$ by expressing $`\mathrm{sin}^2(3\theta )`$ in terms of $`\mathrm{sin}^2(\theta )`$ with the identity
$$\mathrm{sin}^2(3\theta )=16\mathrm{sin}^6\theta 24\mathrm{sin}^4\theta +9\mathrm{sin}^2\theta .$$
(8)
In this case, one gets $`g(x)=16x^324x^2+9x`$ as a possible dynamical map on $`[0,1]`$ having the required solution. More generally, one can use the recurrence formula
$$\mathrm{sin}(N\theta )=2\mathrm{cos}\theta \mathrm{sin}[(N1)\theta ]\mathrm{sin}[(N2)\theta ]$$
(9)
to define a whole set of maps which express $`\mathrm{sin}^2(N\theta )`$ in terms of $`\mathrm{sin}^2\theta `$. This set, the sine functions set, is defined specifically as
$$𝒮=\{SN(x):N=1,2,3,\mathrm{}\},$$
(10)
where $`SN(x)=[s_N(\sqrt{x})]^2`$, and
$`s_1(x)`$ $`=`$ $`x`$ (11)
$`s_2(x)`$ $`=`$ $`2x\sqrt{1x^2}`$ (12)
$`\mathrm{}`$ (13)
$`s_N(x)`$ $`=`$ $`2\sqrt{1x^2}s_{N1}(x)s_{N2}(x).`$ (14)
Note that the intermediate functions $`s_N(x)`$ are expressed as such in order to verify Eq.(9) with the variable change $`x=\mathrm{sin}\theta `$.
| $`SN(x)`$ |
| --- |
| $`S1(x)=x`$ |
| $`S2(x)=4x^2+4x`$ |
| $`S3(x)=16x^324x^2+9x`$ |
| $`S4(x)=64x^4+128x^380x^2+16x`$ |
| $`S5(x)=256x^5640x^4+560x^3200x^2+25x`$ |
Table 1. First five members of $`𝒮`$.
Using these definitions, we calculate the first five functions of $`𝒮`$ listed in table 1. Obviously, by construction of $`𝒮`$, $`S2(x)`$ is the logistic equation itself with parameter $`r=4`$. From a more general perspective, it can also be seen that the maps $`SN(x)`$ are degree $`N`$ polynomials whose leading coefficient, i.e., the coefficient of the highest degree term, is equal to $`4^{N1}`$ in absolute value. These two results, satisfied by any function $`SN(x)`$, is proven more formally in ref.. Note that the latter property allows us to extend the similarity with the logistic map by parameterizing the functions of $`𝒮`$ in the following manner
$$SN_r(x)=r\frac{SN(x)}{4^{N1}},$$
(15)
with $`0r4^{N1}`$. We call the set of functions $`\{SN_r(x):0r4^{N1}\},`$ the family of $`SN(x)`$, which can be characterized numerically by bifurcation diagrams and Lyapunov spectrums such as the ones shown in figure 1.
## III Particularities of $`𝒮`$
Many of the interesting properties of the logistic map at $`r=4`$ can be investigated more intuitively by making explicit the fact that Eq.(7) is equivalent to a shift map $`S`$ on the binary expression of $`\theta _0`$ . Indeed, if we express $`\theta _0`$ as a binary number
$$\theta _0=0.b_0b_1b_2\mathrm{}=\underset{i=0}{\overset{\mathrm{}}{}}\frac{b_i}{2^{i+1}},b_i\{0,1\},$$
(16)
then applying Eq.(7) to $`\theta _0`$ is equivalent to shifting all the bits of $`\theta _0`$ to the left and dropping the integer part. In other words,
$`\theta _n`$ $`=`$ $`0.b_nb_{n+1}b_{n+2}\mathrm{}`$ (17)
$`=`$ $`S^n(0.b_0b_1b_2\mathrm{}),`$ (18)
where $`S^n=SS^{n1}`$ for $`n>1`$, and $`S^1=S`$. Not surprisingly, the same is true for the maps $`SN(x)`$, since Eq.(7) was the guideline in defining the family $`𝒮`$. However, in the case of $`SN(x)`$, the shift map to consider takes effect on $`\theta _0`$ written in base $`N`$. This follows from the following result which generalizes effectively the solution of Eqs.(5) and (7).
Theorem 1. Let $`\{x_0,x_1,x_2,\mathrm{}\}`$ be the orbit of $`x_0`$ under $`SN(x)`$. If we write $`x_n=\mathrm{sin}^2(\pi \theta _n)`$, then we have that
$$\theta _{n+1}=N^n\theta _0mod1,$$
(19)
where, as usual, $`\theta _0=\pi ^1\mathrm{arcsin}(\sqrt{x_0})`$.
We omit the proof of this theorem as it follows directly from the next lemma.
Lemma 1. Consider $`SN(x)`$ as defined previously. We have that
$$SN(\mathrm{sin}^2\theta )=\mathrm{sin}^2(N\theta ).$$
(20)
Proof: The result is obvious for $`N=1`$ and $`N=2`$. Suppose Eq.(20) true for $`N1`$ and $`N2`$, that is to say
$`s_{N1}^2(\mathrm{sin}\theta )`$ $`=`$ $`S[N1](\mathrm{sin}^2\theta )`$
$`=`$ $`\mathrm{sin}^2[(N1)\theta ],`$
and $`s_{N2}(\mathrm{sin}\theta )=\mathrm{sin}[(N2)\theta ]`$. Then, for $`\theta ,`$
$`SN(\mathrm{sin}^2\theta )`$ $`=`$ $`s_N^2(\mathrm{sin}\theta )`$
$`=`$ $`[2\sqrt{1\mathrm{sin}^2\theta }s_{N1}(\mathrm{sin}\theta )s_{N2}(\mathrm{sin}\theta )]^2`$
$`=`$ $`[2\mathrm{cos}\theta \mathrm{sin}((N1)\theta )\mathrm{sin}((N2)\theta )]^2`$
$`=`$ $`\mathrm{sin}^2(N\theta ),`$
where we have used the identity (9). $`\mathrm{}`$
Note that we could have proceeded to a similar generalization of the logistic map using cosine functions instead of sine functions, while preserving its shift property. One possible way of achieving this is to define the cosine functions set
$$𝒞=\{CN(x):N=1,2,3,\mathrm{}\},$$
(21)
where $`CN(x)=[c_N(\sqrt{x})]^2`$, and
$`c_1(x)`$ $`=`$ $`x`$ (22)
$`c_2(x)`$ $`=`$ $`2x^21`$ (23)
$`\mathrm{}`$ (24)
$`c_N(x)`$ $`=`$ $`2xc_{N1}(x)c_{N2}(x).`$ (25)
Contrary to the $`s_N(x)`$’s, the functions $`c_N(x)`$ have the interesting property that they are polynomials of degree $`N`$. In fact, the set $`\{c_N(x):N\}`$ coincides with the set of Tschebysheff polynomials on the unit interval , the latter set satisfying the exact same recurrence formula as Eq.(25). We thus have that $`\{c_N(x)\}`$ must constitute a set of orthogonal polynomials, i.e.,
$$_0^1c_N(x)c_N^{}(x)𝑑x=\delta _{N,N^{}},$$
(26)
for all integers $`N`$ and $`N^{}`$, where $`\delta _{i,j}`$ is the delta-Kronecker function. This fact can be further proved using the property $`c_N(\mathrm{cos}\theta )=\mathrm{cos}(N\theta )`$, well-known to be satisfied by the Tschebysheff functions. Note that $`\{s_N(x)\}`$ is also a set of orthogonal functions; its members satisfy indeed the relation $`s_N(\mathrm{sin}\theta )=\mathrm{sin}(N\theta )`$. In the remaining of this work, we shall restrain our study to the set $`𝒮`$, since the maps $`SN(x)`$ are directly related to $`CN(x)`$ by the expression
$$SN(x)=\{\begin{array}{ccc}CN(x),\hfill & & \text{for }N\text{ odd}\hfill \\ 1CN(x),\hfill & & \text{for }N\text{ even.}\hfill \end{array}$$
(27)
Hence, as far as their dynamics are concerned, the functions $`SN(x)`$ and $`CN(x)`$ are totally equivalent.
## IV Conjugacies
The analysis of the chaoticity properties of a map $`f`$ is greatly simplified by studying conjugate maps of $`f`$ which are obtained by applying a global change of variables. Recall that two maps $`f:II`$ and $`g:JJ`$ are conjugate if there exists a homeomorphism, i.e., a bijective and continuous map $`H:IJ`$ such that $`Hf=gH`$. The function $`H`$ is called a conjugacy. In the context of $`𝒮`$, a possible conjugate function of $`SN(x)`$ can be constructed as follows. Let $`TN(x)`$ be a piecewise linear function (a generalized tent map) defined on subintervals $`[k/N,(k+1)/N]`$ of $`[0,1]`$ by setting
$$TN(x)=\{\begin{array}{cc}Nxk,\hfill & \text{for }k\text{ even}\hfill \\ Nx+k+1,\hfill & \text{for }k\text{ odd,}\hfill \end{array}$$
(28)
with $`k=0,1,\mathrm{},N1`$.
Theorem 2. $`SN(x)`$ is conjugate to $`TN(x)`$ with conjugacy $`H(x)=\mathrm{sin}^2(\pi x/2)`$.
Proof: First note that $`H(x)`$ is both continuous and bijective on the interval $`[0,1]`$. Now, on the first hand we have that $`SN(H(x))=SN(\mathrm{sin}^2(\pi x/2))=\mathrm{sin}^2(N\pi x/2)`$. On the other hand,
$`H(TN(x))`$ $`=`$ $`\{\begin{array}{cc}\mathrm{sin}^2(N\pi x/2k\pi /2),\hfill & \text{for }k\text{ even}\hfill \\ \mathrm{sin}^2(N\pi x/2+(k+1)\pi /2),\hfill & \text{for }k\text{ odd}\hfill \end{array}`$ (31)
$`=`$ $`\left(\pm \mathrm{sin}(N\pi x/2)\right)^2`$ (32)
$`=`$ $`\mathrm{sin}^2(N\pi x/2).`$ (33)
Thus, we have proved that $`SNH=HTN`$ for all integers $`N`$. $`\mathrm{}`$
Figure 2 depicts the graphs of $`SN(x)`$ and the corresponding $`TN(x)`$ for $`N=3,4`$. ($`T2(x)`$ is only the tent map since, as we mentioned, $`S2(x)`$ is the logistic map.) In each case, note that $`SN(x)`$ is a “smooth” version of $`TN(x)`$.
Remark. The functions $`CN(x)`$ are pairwise conjugate to the functions $`TN(x)`$ with conjugacy map $`H(x)=\mathrm{cos}^2(\pi x/2)`$. Evidently, since the conjugacy is an equivalence relation, this has the consequence that any $`CN(x)`$ is conjugate to $`SN(x)`$.
## V Ergodic properties of $`𝒮`$
In this section, we evaluate the invariant densities of the maps $`SN(x)`$ and their Lyapunov exponents based on the conjugates found in the previous section. As one would expect from the similarity of the graphs of $`T3(x)`$ and $`T4(x)`$ in figure 2, the functions $`TN(x)`$ should exhibit a unique and constant invariant density
$$\rho _{TN}(x)=1,x[0,1],$$
(34)
just as it is the case for the tent map. This can be proven using Adler and Bowen’s results on Markov transformations , (see also ). At this point, the invariant density of the functions $`SN(x)`$ can be found using the transformation formula
$$\rho _{SN}(x)=\rho _{TN}(y)\left|\frac{dy}{dx}\right|,$$
(35)
where $`y=\frac{2}{\pi }\mathrm{arcsin}(\sqrt{x}),`$ so as to obtain
$$\rho _{SN}(x)=\frac{1}{\pi \sqrt{x(1x)}}.$$
(36)
The unicity of $`\rho _{SN}(x)`$ is assured by the fact that $`\rho _{TN}(x)`$ is also unique, which means that both functions $`SN(x)`$ and $`TN(x)`$ must be ergodic .
Now, to evaluate the Lyapunov exponent of $`SN(x)`$ we may use the fact that $`\left|_xTN(x)\right|=N`$ for almost all $`x[0,1]`$ to infer that $`\lambda (x_0)=\mathrm{ln}N`$ almost everywhere in the case of $`TN(x)`$. Accordingly, since Lyapunov exponents are invariant under smooth and differentiable coordinate transformations , we have the following theorem. (A more extensive proof of this result, which takes care of the pathological points where $`_xTN(x)`$ is not defined, is contained in ref..)
Theorem 3. The Lyapunov exponent of $`SN(x)`$ is $`\mathrm{ln}N`$ almost everywhere (with respect to the invariant measure $`\rho _{SN}(x)`$).
The above theorem shows that the members of $`𝒮`$ are non-conjugate to each other simply because they possess different Lyapunov exponents. It also shows that $`𝒮\backslash \{S1\}`$, and consequently the set of Tschebysheff polynomials, are sets of chaotic maps. Indeed, $`\lambda =\mathrm{ln}N>0`$ for $`N>1`$, and by using the shift property of $`SN(x)`$ we can choose $`x_0=\mathrm{sin}^2(\pi \theta _0)`$, with $`\theta _0`$ irrational, to build an orbit that is not asymptotically periodic. Another way to convince ourselves that all the polynomials in $`𝒮`$ have chaotic orbits is to use the celebrated result “period-3 implies chaos” , and find an initial point $`x_0`$ of period 3 for each $`SN(x)`$. For instance, for a $`N>1`$ let $`x_0=\mathrm{sin}^2(\pi \theta _0)`$ where
$$\theta _0=\frac{1}{N^31}=0.001001\mathrm{}\text{(in base }N\text{).}$$
(37)
Again, using the shift map property, we must have
$`x_0`$ $`=`$ $`\mathrm{sin}^2(\pi 0.001001\mathrm{})`$ (38)
$`x_1`$ $`=`$ $`\mathrm{sin}^2(\pi 0.010010\mathrm{})`$ (39)
$`x_2`$ $`=`$ $`\mathrm{sin}^2(\pi 0.100100\mathrm{})`$ (40)
$`x_3`$ $`=`$ $`\mathrm{sin}^2(\pi 0.001001\mathrm{})=x_0.`$ (41)
We thus extended the chaoticity properties of the logistic map to an infinite family of polynomials.
## VI Algebraic properties
To complete the study of the properties of $`𝒮`$, we now deduce that it is an abelian monoid with respect to the composition of functions ($``$). A monoid, precisely, is a non-empty set $`M`$ together with a binary associative operation, say $``$, such that $`xyM`$ for $`x,yM`$. There must also be an element $`eM`$, called the identity element, for which $`xe=ex=x`$ for all $`xM`$. Moreover, a monoid is called abelian if the binary operation is commutative . In $`𝒮`$, the identity element is $`S1(x)=x`$. Also, the composition of function is clearly associative. Now, to prove that $`𝒮`$ is indeed abelian monoid, we verify that it is closed under composition and that this composition is commutative, a condition that is not verified in the case of composition of general functions. However, before we do so, we present next a new expression of $`SN(x)`$ on the unit interval.
Lemma 2. For all $`SN(x)`$ and $`x[0,1]`$,
$`SN(x)=\mathrm{sin}^2(N\mathrm{arcsin}\sqrt{x}).`$
Proof: Let $`x[0,1]`$. There exists a $`\theta [0,\pi /2]`$ such that $`x=\mathrm{sin}^2\theta `$, and thus $`\theta =\mathrm{arcsin}\sqrt{x}`$. Now, from Lemma 1 we have $`SN(x)=SN(\mathrm{sin}^2\theta )=\mathrm{sin}^2(N\theta )=\mathrm{sin}^2(N\mathrm{arcsin}\sqrt{x})`$. $`\mathrm{}`$
Theorem 4. (Monoid property) Let $`N_1`$ and $`N_2`$ be any positive integers. We have that
$$SN_1SN_2=SN_2SN_1=S[N_1N_2].$$
(42)
Proof: For $`N_1`$ and $`N_2`$ given, consider $`SN_1(x)`$ and $`SN_2(x)`$. Then, for any $`x[0,1]`$, we obtain from Lemma 2
$`SN_1(SN_2(x))`$ $`=`$ $`\mathrm{sin}^2[N_1\mathrm{arcsin}(\sqrt{SN_2(x)})]`$ (43)
$`=`$ $`\mathrm{sin}^2[N_1\mathrm{arcsin}(\mathrm{sin}(N_2\mathrm{arcsin}\sqrt{x}))]`$ (44)
$`=`$ $`\mathrm{sin}^2(N_1N_2\mathrm{arcsin}\sqrt{x})`$ (45)
$`=`$ $`S[N_1N_2](x).`$ (46)
Obviously, $`S[N_1N_2](x)=S[N_2N_1](x)`$, so the composition is commutative. $`\mathrm{}`$
As a direct consequence of the monoid property, $`k`$-periodic points of a certain polynomial $`SN(x)`$ can be looked at as fixed points of the function $`SM(x)`$ where $`M=N^k`$. Furthermore, a polynomial $`SN(x)`$ of very high degree can be computed easily by decomposing its expression using lower degree polynomial of the family $`𝒮`$. Explicitly, consider $`SN(x)𝒮`$. We say that $`SN(x)`$ is a prime element of $`𝒮`$ if $`N`$ is a prime number. Using this definition, we have as a result of Theorem 4 and the Fundamental Theorem of Arithmetic that any polynomial $`SN(x)`$ must possess a unique decomposition in prime elements of $`𝒮`$.
## VII Final remarks
To conclude, note that our study of the sine functions, written in the form $`SN(x)=\mathrm{sin}^2(N\mathrm{arcsin}\sqrt{x})`$, have been restricted to positive integers $`N`$. In a similar manner, it could be interesting to investigate functions of the type $`S\alpha (x)=\mathrm{sin}^2(\alpha \mathrm{arcsin}\sqrt{x})`$ with $`\alpha `$ real. One observation about this extra generalization is that, as for $`S1(x)`$, the function $`S\alpha (x)`$ does not exhibit chaotic properties for $`0\alpha 1`$. The function $`S\frac{1}{2}(x)=\mathrm{sin}^2(\frac{1}{2}\mathrm{arcsin}\sqrt{x})`$, for example, is conjugate to $`g(x)=x/2`$, and has all of its orbits attracted to $`x=0`$. Yet, this is not surprising since the Lyapunov exponent of this map must be $`\mathrm{ln}(1/2)<0`$. This brings us to conjecture that $`S\alpha (x)`$ must admit chaotic behavior if and only if $`\left|\alpha \right|>1`$, considering that the Lyapunov exponents of $`S\alpha (x)`$ should be $`\mathrm{ln}\alpha `$. A complete proof of this result, however, cannot be given here using the same symbolic dynamic approach used for $`SN(x)`$, for the simple reason that the expression of a point in “base $`N`$” makes sense only if $`N`$ is an integer greater than $`1`$.
## Aknowledgements
V.P. would like to thank A. Mingarelli for helpful discussions. This work was supported in part by the National Sciences and Engineering Research Council of Canada (NSERC) through the ES A Scolarship program.
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# On the semiclassical expansion for 1-dim 𝑥^𝑁 potentials
## Acknowledgements
This project was supported by the Ministry of Science and Technology of the Republic of Slovenia and by the Rector’s Fund of the University of Maribor. VR acknowledges the support of the work by the grant of the Ministry of Science and Technology of the Republic of Slovenia and the Abdus Salam ICTP (Trieste) Joint Programme and also the support of the Foundation of Fundamental Research of the Republic of Belarus.
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# The volume and lengths on a three sphere
## 1 Introduction
In this paper we consider the volume, $`Vol(g)`$, of a Riemannian metric $`g`$ on the three sphere. We let $`L(g)`$ represent the length of the shortest nontrivial closed geodesic in $`(S^3,g)`$. By an “antipodal map”, $`A`$, on an $`n`$-sphere we will mean an order 2, fixed point free diffeomorphism. We will let $`D(g,A)=inf_{xS^3}d(x,Ax)`$, and $`D(g)=sup_AD(g,A)`$. The main result of this paper is:
Theorem 1.1. For any Riemannian metric on $`S^3`$ we have:
$$Vol(g)^{\frac{1}{3}}C_1\mathrm{min}\{L(g),2D(g)\}$$
where $`C_1`$ is a universal constant (which can be taken to be $`\frac{1}{3180}`$)
This theorem is an example of a “universal” inequality. This is a term introduced by Berger to describe inequalities on a Riemannian manifold $`(M,g)`$ (usually between minimizing objects in some topological class) which hold for all metrics $`g`$ on $`M`$. With the glaring exception of Gromov’s work \[Gr\] (which we discuss below) most such inequalities are known in 2 dimensions only. Many of these inequalities involve the systole, $`sys(g)`$, which is the length of the shortest closed noncontractible curve in $`M`$. Estimates of the form $`A(g)^{\frac{1}{2}}c(M)sys(g)`$, Where $`A(g)`$ represents the area, have been proved for all surfaces $`M`$. The first of these was proved by Loewner (unpublished) where he finds the sharp value of $`c(T^2)`$ (sharp for the flat regular hexagonal torus). The sharp values of $`c(RP^2)`$ and $`c(K^2)`$ for the projective plane (the round one is best) and Klein bottle (the best one is singular and not flat!) were proven by Pu \[Pu\] and Bavard \[Ba\] respectively. For surfaces of genus $`\gamma `$, nonsharp constants $`c(\gamma )`$ (which unfortunately went to 0 as $`\gamma `$ grew) were first proved independently by Accola \[Ac\] and Blatter \[Bl\]. Gromov greatly improve these constants in \[Gr\] so that they grow correctly with $`\gamma `$. However, the best constants are still unknown. Finding them is likely to be an extremely hard problem.
None of the above says anything about $`S^2`$ since it is simply connected. However there are two inequalities of this type on $`S^2`$. The first, due to Berger \[Be\] says that there is a universal constant $`c_2`$ such that for any antipodal map $`A:S^2S^2`$:
$$Vol(S^2,g)^{\frac{1}{2}}c_2D(g,A).$$
$`1.1`$
The sharp constant is not known, but is conjectured to be that of the round sphere. The other result (see \[Cr2\]) is that
$$Vol(S^2,g)^{\frac{1}{2}}c_3L(g).$$
$`1.2`$
Also in this case the sharp constant is not known. It is conjectured to be that for the singular metric one gets by gluing two flat equilateral triangles along their boundaries.
We now turn to higher dimensions. The most important results in the area are in Gromov’s paper “Filling Riemannian Manifolds” \[Gr\]. Gromov shows that for essential manifolds (which includes $`T^n`$, $`RP^n`$, and all compact $`K(\pi ,1)`$ spaces) we have $`V(g)^{\frac{1}{n}}c(n)sys(g)`$. However, this says nothing directly about $`S^n`$ (or any compact simply connected manifold). For example the natural generalizations of 1.1 and 1.2 above for $`S^2`$ are open questions for $`S^n`$, $`n3`$. Note that for any Riemannian manifold $`(M,g)`$ of injectivity radius $`inj(M)`$ and any $`r\frac{1}{2}inj(M)`$ the metric spheres $`S(x,r)`$ with their induced Riemannian metrics $`\overline{g}`$ satisfy $`D(\overline{g},A)2r`$, so an estimate $`Vol(g)^{\frac{1}{n}}c(S^n)D(g,A)`$ which was sharp for the round spheres would give sharp estimates for the volume of small metric spheres. See \[Cr1\] for nonsharp estimates for the volume of such metric spheres.
In the other direction, Ivanov (see \[I1\] or \[I2\]) has given examples of a sequence of metrics on $`S^3`$ that Gromov-Housdorff converge to the standard metric but whose volumes go to zero.
Although Theorem 1 does not yield either 1.1. or 1.2 for $`S^3`$ it does show that for any given metric one or the other inequality must hold.
For $`A:S^3S^3`$ an antipodal map we let $`L(g,A)`$ be the infimum of the lengths of curves, $`\gamma `$ such that $`\gamma `$ links $`A(\gamma )`$. Here we will say $`\gamma `$ links $`\tau `$ if either $`\gamma `$ intersects $`\tau `$ or the linking number of $`\gamma `$ and $`\tau `$ is nonzero. We consider two cases. If there is no closed curve $`\gamma `$ that links $`A(\gamma )`$ of length less than $`2D(g,A)`$ then $`L(g,A)=2D(g,A)`$ and the minimum length $`\gamma `$ is the union of two minimizing geodesics between $`x`$ and $`A(x)`$ for some $`x`$. Otherwise there is such a closed curve $`\gamma `$. In this case we can apply curve shortening to $`\gamma `$ to get a continuous family $`\gamma _t`$ each of which is shorter than $`\gamma `$ and hence does not intersect $`A(\gamma _t)`$. Thus we see that $`\gamma _t`$ links with $`A(\gamma _t)`$ and hence has length bounded below by $`L(g,A)`$. By looking at limits of such $`\gamma _t`$ we find nontrivial closed geodesics of length less than the length of $`\gamma `$. Since this is true for all such $`\gamma `$ we find that there is a closed geodesic of length $`L(g,A)`$. Thus the invariants $`L(g)`$ and $`D(g,A)`$ are thus related to $`L(g,A)`$ via:
$$L(g,A)min\{L(g),2D(g,A)\}.$$
Theorem 1 thus will follow from:
Theorem 1.2. For any Riemannian metric on $`S^3`$ and any antipodal map $`A`$, we have:
$$Vol(g)^{\frac{1}{3}}C_1L(A,g)$$
where $`C_1`$ is a universal constant (which can be taken to be $`\frac{1}{3180}`$)
The fundamental result in the proof of Gromov’s isosystolic inequality is his filling radius theorem: $`Fillrad(g)c_nVol(g)^{\frac{1}{n}}`$, which holds for all Riemannian $`n`$-manifolds (here $`c_3<265`$). We define the filling radius, $`Fillrad(g)`$, in the next section. We also make fundamental use of it since Theorem 1 will thus follow from
Theorem 1.3. For any Riemannian metric on $`S^3`$ and any antipodal map $`A`$, we have:
$$Fillrad(g)C_4L(g,A)$$
Where $`C_4`$ is a universal constant (which can be taken to be $`\frac{1}{12}`$)
The proof of this theorem is a generalization of an argument that works on $`S^2`$ to yield:
Theorem 1.4. For any Riemannian metric on $`S^2`$ we and any antipodal map $`A`$ we have:
$$Fillrad(g)\frac{1}{4}D(g,A)$$
This theorem along with the Filling radius theorem recovers Berger’s estimate equation 1.1 (although with a worse constant). We present the argument in section 3.
The author would like to thank Herman Gluck for helpful conversations.
## 2 Notation/Preliminaries
The main purpose of this section is to set up notation and remind the reader about the Filling radius (for details see \[Gr\]):
In \[Gr\] Gromov introduced the notion of the filling radius (we use integer coefficients) $`Fillrad(M)`$ of a closed $`n`$-dimensional manifold $`M`$ with a metric $`d`$ (not necessarily Riemannian). We will only consider the case where $`M`$ is homeomorphic to $`S^2`$ or $`S^3`$ and the metric is Riemannian. There is a natural isometric embedding (in the metric space sense!) $`i:ML^{\mathrm{}}(M)`$ defined by $`i(x)()=d_M(x,)`$. The filling radius is the infimum of $`r`$ such that $`i(M)`$ bounds in the tubular neighborhood $`T_r(i(M))`$ in the sense that $`i_{}(H_n(M;))`$ vanishes in $`H_n(T_r(i(M));)`$. We will represent a filling as continuous map $`\sigma :\mathrm{\Sigma }T_r(i(M))`$ from an $`n+1`$-dimensional simplicial complex $`\mathrm{\Sigma }`$ such that $`\sigma |_\mathrm{\Sigma }:\mathrm{\Sigma }i(M)`$ represents a generator in $`H_n(M;)`$.
We note that for any fixed $`ϵ>0`$ by taking Barycentric subdivisions as needed we may assume that the $`\sigma `$-image of any simplex has diameter less than $`ϵ`$ in $`L^{\mathrm{}}(X)`$.
We note that there can be no continuous map $`f:\mathrm{\Sigma }i(X)`$ which agrees with $`\sigma `$ on $`\mathrm{\Sigma }`$ since $`\sigma |_\mathrm{\Sigma }`$ represents a generator of the top homology (so is not a boundary). Our proof of Theorem 1 will be by contradiction. We assume that $`Fillrad(g)`$ is small, take a filling $`\sigma :\mathrm{\Sigma }L^{\mathrm{}}(M)`$ as above, and show that $`\sigma :\mathrm{\Sigma }i(M)`$ extends to a continuous $`f:\mathrm{\Sigma }i(M)`$ giving the desired contradiction. Our proof of Theorem 1 uses a similar contradiction.
We use the notation $`S_i(\mathrm{\Sigma })`$ to denote the $`i`$-skeleton of a simplicial complex $`\mathrm{\Sigma }`$. For each $`i=0,1,2,\mathrm{},n+1`$ We will let $`\{\mathrm{\Delta }_j^i\}`$ denote the set of $`i`$ simplices in the $`n+1`$ dimensional simplicial complex $`\mathrm{\Sigma }`$. For each i-simplex $`\mathrm{\Delta }_j^i`$ we let $`G_j^i`$ be the connected graph (i.e. 1-complex) $`G_j^i=\{S_1(\mathrm{\Delta }_k^{n+1})|\mathrm{\Delta }_j^i`$ is a face of $`\mathrm{\Delta }_k^{n+1}\}`$. By taking Barycentric subdivisions if needed we can assume that any two simplices intersect in a single (possibly empty) common face. Then we see that $`H_1(G_j^i)`$ is generated by the boundaries of the $`2`$-simplices that are faces of the $`n+1`$ simplices in the above union, since Van Kampen’s theorem implies that the two complex $`\{S_2(\mathrm{\Delta }_k^{n+1})|\mathrm{\Delta }_j^i`$ is a face of $`\mathrm{\Delta }_k^{n+1}\}`$ is simply connected.
## 3 Proof of Theorem 1
Proof. (of Theorem 1) Let $`ϵ>0`$ and choose a filling $`\sigma :\mathrm{\Sigma }L^{\mathrm{}}(S^2,g)`$ of $`(S^2,g)`$ in the $`Fillrad(g)+ϵ`$ tubular neighborhood of $`S^2L^{\mathrm{}}(S^2,g)`$ (we will confuse $`S^2`$ with $`i(S^2)`$ since $`i`$ is an isometric embedding). By taking subdivisions we can assume that $`diam(\sigma (\mathrm{\Delta }_k^3))<ϵ`$ for each $`\mathrm{\Delta }_k^3\mathrm{\Sigma }`$. We will prove the theorem as suggested in the previous section by showing that if $`Fillrad(g)<\frac{D(g,A)}{4}`$ then we can find a continuous map $`f:\mathrm{\Sigma }S^2`$ extending the map $`\sigma |_\mathrm{\Sigma }`$. We choose $`ϵ`$ so small that $`2Fillrad(g)+3ϵ\frac{D(g,A)}{2}`$
Step1: The 0 and 1 Skeletons
We define $`f`$ on the 0-skeleton, $`S_0(\mathrm{\Sigma })`$, of $`\mathrm{\Sigma }`$ by mapping each 0 simplex $`v`$ to a point $`f(v)`$ on $`i(S^2)`$ closest to $`\sigma (v)`$, hence
$$d_{L^{\mathrm{}}(S^2)}(f(v),\sigma (v))<Fillrad(g)+ϵ.$$
In particular, $`f`$ takes 0-simplices in the boundary to the same point as $`\sigma `$ does. We note that if two vertices $`v_1`$ and $`v_2`$ are the endpoints of an edge then
$$d_{S^2}(f(v_1),f(v_2))=d_{L^{\mathrm{}}(S^2)}(f(v_1),f(v_2))$$
$$d_{L^{\mathrm{}}(S^2)}(f(v_1),\sigma (v_1))+d_{L^{\mathrm{}}(S^2)}(\sigma (v_1),\sigma (v_2))+d_{L^{\mathrm{}}(S^2)}(\sigma (v_2),f(v_2))<$$
$$<2Fillrad(g)+3ϵ\frac{D(g,A)}{2}.$$
We define $`f`$ to map each nonboundary $`\mathrm{\Delta }_i^1S_1(\mathrm{\Sigma })`$ to a minimizing geodesic between the $`f`$ image of the endpoints, hence the length $`L(\mathrm{\Delta }_i^1)`$ satisfies $`L(\mathrm{\Delta }_i^1)<\frac{D(g,A)}{2}`$. For $`\mathrm{\Delta }_i^1`$ on the boundary, we let $`f|_{\mathrm{\Delta }_i^1}=\sigma |_{\mathrm{\Delta }_i^1}`$ and hence $`Diam(f(\mathrm{\Delta }_i^1))<ϵ`$.
Step 2: The 2 skeleton
We now extend $`f`$ to $`S_2(\mathrm{\Sigma })`$. We note that $`f(\mathrm{\Delta }_j^2)Af(G_j^2)=\mathrm{}`$, for otherwise there would be a point $`xf(G_j^2)`$ such that $`A(x)f(G_j^2)`$ but this cannot happen since step 1 guarantees that the diameter of $`G_j^2`$ is less than $`D(g,A)`$.
We can thus extend $`f`$ to $`S_2(\mathrm{\Sigma })`$ in such a way that we have
$$f:\mathrm{\Delta }_j^2S^2Af(G_j^2).$$
Note that for boundary simplices, $`\mathrm{\Delta }_j^2`$, this will hold when we take $`f|_{\mathrm{\Delta }_j^2}=\sigma |_{\mathrm{\Delta }_j^2}`$ by the triangle inequality.
Step 3: The 3 skeleton
Now for every $`\mathrm{\Delta }_k^3`$ in $`\mathrm{\Sigma }`$ the previous step guarantees that
$$f:\mathrm{\Delta }_k^3S^2Af(S_1(\mathrm{\Delta }_k^3)).$$
Hence, since $`f|_{\mathrm{\Delta }_k^3}`$ misses a point, we can extend $`f`$ to $`\mathrm{\Delta }_k^3`$.
This completes the proof.
## 4 Proof of Theorem 1
Throughout this section $`S^3`$ will be endowed with a fixed Riemannian metric $`g`$ and an antipodal map $`A:S^3S^3`$.
We will proceed analogously to the proof of Theorem 1. However, in this case we will get our contradiction by finding a singular chain in $`S^3`$ whose boundary is $`\sigma `$ restricted to $`\mathrm{\Sigma }`$. We will do this one skeleton at a time, associating to each simplex $`\mathrm{\Delta }_i^j`$ a singular simplicial chain $`c_i^j`$ whose boundary $`c_i^j`$ corresponds to the already defined chain associated to $`\mathrm{\Delta }_i^j`$ (i.e. if $`\mathrm{\Delta }_i^j=\mathrm{\Sigma }_k(1)^{\alpha (i,k)}\mathrm{\Delta }_k^{j1}`$ then $`c_i^j=\mathrm{\Sigma }_k(1)^{\alpha (i,k)}c_k^{j1}`$). We do this while associating simplices $`\mathrm{\Delta }_i^j`$ of $`\mathrm{\Sigma }`$ to the chain consisting only of $`\sigma `$ applied to $`\mathrm{\Delta }_i^j`$.
Let $`ϵ>0`$ and choose a filling $`\sigma `$ of $`(S^3,g)`$ in the $`Fillrad(g)+ϵ`$ tubular neighborhood of $`S^3L^{\mathrm{}}(S^3,g)`$. By taking subdivisions we can assume that $`diam(F(\mathrm{\Delta }_l^4))<ϵ`$ for each $`\mathrm{\Delta }_l^4\mathrm{\Sigma }`$. We will prove the theorem by showing that if $`Fillrad(g)<\frac{L(g,A)}{12}`$ then we find a chain as above extending $`\sigma |_\mathrm{\Sigma }`$. We choose $`ϵ`$ so small that $`2Fillrad(g)+3ϵ<\frac{L(g,A)}{6}`$. During the rest of the argument we will associate boundary simplices to themselves without explicitly mentioning this special case. The arguments for these simplices will always follow from the other arguments along with the fact that the diameters are bounded by $`ϵ`$.
Step1: The 0 and 1 Skeletons
We define the 0-chains and 1-chains, $`c_i^0`$ and $`c_i^1`$, associated to the 0-skeleton, $`S_0(\mathrm{\Sigma })`$, and the 1-skeleton, $`S_1(\mathrm{\Sigma })`$ of $`\mathrm{\Sigma }`$ just as before; i.e. by mapping each 0 simplex to a closest point on $`S^3`$, and mapping each edge in the one skeleton to a minimizing geodesic between the endpoints. Hence the length of the image of a 1 simplex is less than $`2Fillrad(g)+3ϵ<\frac{L(g,A)}{6}`$. We can assume (by small moves) that 0-chains of distinct vertices of $`\mathrm{\Sigma }`$ are distinct and that the geodesic segments only intersect each other at endpoints.
Step 2: The 2 skeleton
Now consider a two simplex $`\mathrm{\Delta }_j^2`$ of $`\mathrm{\Sigma }`$. Let $`\overline{G}_j^2`$ be the embedded geodesic graph in $`S^3`$ which is the union of the geodesic segments that correspond to the edges of $`G_j^2`$. We know that the support of $`c_j^2`$ (i.e. a geodesic triangle) does not intersect $`A((\overline{G}_j^2))`$ since in fact $`\overline{G}_j^2`$ does not intersect $`A((\overline{G}_j^2))`$, for if so there would be an $`x\overline{G}_j^2`$ such that $`A(x)`$ is also in $`\overline{G}_j^2`$. But this can’t happen because the diameter of $`\overline{G}_j^2`$ is $`<\frac{3}{2}\frac{L(g,A)}{6}<L(g,A).`$
We claim that $`c_j^2`$ represents zero in $`H_1(S^3A(\overline{G}_k^2))`$. Alexander duality along with the fact that $`H_1(\overline{G}_j^2)`$ is generated by $`\{c_k^2|\mathrm{\Delta }_k^2`$ and $`\mathrm{\Delta }_j^2`$ lie in a common 4-simplex$`\}`$ says that we need only show that $`(c_j^2,A(c_k^2))=0`$ for each such $`c_k^2`$. So assume that $`(c_j^2,A(c_k^2))0`$. Since $`\mathrm{\Delta }_k^2`$ and $`\mathrm{\Delta }_j^2`$ lie in a common 4-simplex they share at least one vertex and hence we let $`\gamma `$ be the closed simplicial curve (of combinatorial length 6) which is just $`\mathrm{\Delta }_1^2`$ followed by $`\mathrm{\Delta }_k^2`$. We let $`\overline{\gamma }`$ be the corresponding closed piecewise geodesic curve in $`\overline{G}_j^2`$ of length $`<6\frac{L(g,A)}{6}=L(g,A)`$ (which of course does not intersect $`A(\overline{\gamma })`$). By the definition of $`L(g,A)`$ we see that $`(\overline{\gamma },A(\overline{\gamma }))`$ is zero as are $`(c_j^2,A(c_j^2))`$ and $`(c_k^2,A(c_k^2))`$. On the other hand
$$0=(\gamma ,A(\gamma ))=(c_j^2,A(c_j^2))+(c_k^2,A(c_k^2))+(c_j^2,A(c_k^2))+(c_k^2,A(c_j^2))$$
and hence $`(c_j^2,A(c_k^2))=(c_k^2,A(c_j^2))`$. But since $`A^2=id`$ and $`A`$ preserves orientation (since it is fixed point free) we have
$$(c_j^2,A(c_k^2))=(A(c_j^2),A^2(c_k^2))=(A(c_j^2),c_k^2)$$
and the claim follows.
Thus we can define a 2-chain $`c_j^2`$ whose boundary is $`c_j^2`$ and whose support is contained in $`S^3A(\overline{G}_j^2)`$.
Step 3: the 3 and 4 skeleta
Let $`\mathrm{\Delta }_i^3`$ be a 3-simplex in $`\mathrm{\Sigma }`$. For each 2-face $`\mathrm{\Delta }_k^2`$ the support of $`c_k^2`$ lies in $`S^3A(\overline{G}_k^2)S^3A(\overline{G}_i^3)`$ (since $`\overline{G}_i^3\overline{G}_k^2`$). Thus we have already defined the cycle $`c_i^3`$ in such a way that its support is in $`S^3A(\overline{G}_i^3)`$ and since by Alexander duality $`H_3(S^3A(\overline{G}_i^3))=0`$ we can find a chain $`c_i^3`$ with boundary $`c_i^3`$ whose support also lies in $`S^3A(\overline{G}_i^3)`$. Thus we can extend to the 3 skeleton.
Now for $`\mathrm{\Delta }_i^4`$ a 4-simplex of $`\mathrm{\Sigma }`$ we have now defined the cycle $`c_i^4`$ in such a way that its support is in $`S^3A(\overline{G}_i^4)`$. Since the support of $`c_i^4`$ misses a point of $`S^3`$ there is a $`c_i^4`$ whose boundary is $`c_i^4`$.
$`\mathrm{}`$
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# All Teleportation and Dense Coding Schemes
## 1 Introduction
Teleportation and dense coding are two processes, which stood at the beginning of modern Quantum Information Theory. They both demonstrated radically new features of quantum information as opposed to classical information, in that both would be impossible without the assistance of entangled states. Indeed, the attempt of using the properties of a classically correlated system shared by sender and receiver to improve the transmission rate of a classical channel can easily be seen to be hopeless. But this is precisely what happens in teleportation and dense coding, and dramatically so, because without entanglement assistance teleportation, i.e., the transmission of quantum information on a classical channel, would not only be less efficient, but virtually impossible.
In the original papers \[BW, BB\] the new possibilities were demonstrated by giving an explicit example, based on qubits. It was clear early on that extensions to systems with higher dimensional Hilbert spaces were possible, not only to powers of $`2`$, by running the process several times, but to any dimension $`2d<\mathrm{}`$ \[BB\].
The task set in this paper is to do this systematically, and to classify all schemes for teleportation and dense coding. There are several reasons for doing this. The first is, of course, to take these miracle machines apart and to analyze what makes them work: what is the mathematical structure one really needs to set up such a scheme? For the present author one motivation of this kind was to understand the surprising observation that each of the published teleportation schemes also works as a dense coding scheme, and conversely: sender Alice and receiver Bob merely have to swap the equipment they use. An attempt at a direct proof of this failed, and indeed, as discussed below, the statement fails in general, but is true in the special case of “tight” schemes.
The second reason for attempting a complete classification of teleportation schemes is more practical. In spite of amazing progress in recent years, experiments in quantum information processing are still quite difficult. Hence, for realizing a teleportation scheme it is useful to have a systematic overview of the options, before going on to find the one which is the easiest to implement. This also goes for approximate realizations. And in order to find feasible approximate teleportation schemes it is probably once again necessary to understand the manifold of exact realizations.
The aim of determining all schemes is not quite achieved in this paper, in two respects. Firstly, we will only look at the case when dense coding and teleportation are realized optimally with minimal resources, in the sense of Hilbert space dimensions and number of distinguishable classical signals. As in the well-known qubit case, this means that an entangled state between systems of the same dimension $`d`$ as the input systems is used, and the classical channel distinguishes $`d^2`$ signals. That is, the classical capacity of the quantum channel is exactly doubled by dense coding, and teleportation requires twice as much classical channel capacity as the quantum capacity of the channel set up by this scheme. We will call schemes with these dimension parameters tight. As mentioned above, for these dimensions the symmetry between teleportation and dense coding holds perfectly. Classifying all schemes beyond the tight case appears to be more difficult because there is too much freedom, which cannot be parametrized in a simple way (see, however, \[BD\]).
The second respect in which this paper falls short of a complete classification is that we can only reduce it to another “standard” problem, namely the construction of orthonormal bases of unitary operators with respect to the scalar product $`(A,B)d^1tr(A^{}B)`$. In the last section we provide a fairly general construction for such bases. However, even this construction has to rely on other well-known but not completely classified combinatorial designs, namely Latin squares, and complex Hadamard matrices. This suggests that a complete construction procedure for all unitary bases would be at least as difficult as a complete classification of Latin squares or Hadamard matrices, and hence hardly a promising task.
The paper is organized as follows: in Section 2 the Main Theorem is stated: an equivalence in the tight case between teleportation schemes, dense coding schemes, orthonormal unitary bases, bases of maximally entangled vectors, and so-called unitary depolarizers. Basic consequences of the Theorem are discussed. Section 3 contains the proof, divided into subsections, each devoted to some implication in the big equivalence. In writing the proof an attempt was made to include also simple steps explicitly, and to make as transparent as possible why the tightness condition is crucial. Finally, in Section 4 we present the “Shift and Multiply” construction of unitary bases, which are then classified in terms of Latin squares and Hadamard matrices.
## 2 Main result
In order to state our result we use the following notation and terminology: When $``$ is a Hilbert space, we denote by $`()`$ the space of bounded linear operators on $``$. A channel converting quantum systems with Hilbert space $`_{\mathrm{in}}`$ into systems with Hilbert space $`_{\mathrm{out}}`$ is a linear operator $`T:(_{\mathrm{out}})(_{\mathrm{in}})`$, which is completely positive \[Da, Pa\] and normalized as $`T(1\mathrm{I})=1\mathrm{I}`$. A (discrete) observable $`F`$ on $``$ over an output parameter space $`X`$ is a collection of positive operators $`F_x()`$ such that $`_xF_x=1\mathrm{I}`$. A density operator on $``$ is a positive operator with trace $`1`$. The basic probabilistic interpretation of these objects is fixed by the prescription that $`tr(\omega T(F_x))`$ is the probability to get the measuring result “$`x`$” on systems prepared according $`\omega `$, before passing through the channel $`T`$. Finally, we call a vector $`\mathrm{\Psi }`$ maximally entangled, if it is normalized, and its reduced density operator is maximally mixed, i.e., a multiple of $`1\mathrm{I}`$:
$$\mathrm{\Psi }|(A1\mathrm{I})\mathrm{\Psi }=(dim)^1tr(A).$$
(1)
Let us set up the equations describing dense coding and teleportation in this language. In both cases, the beginning of each transmission is to distribute the parts of an entangled state $`\omega `$ between sender Alice and receiver Bob. Only then Alice is given the message she is supposed to send, which is a quantum state in the case of teleportation and a classical value in case of dense coding. She codes this in a suitable way, and Bob reconstructs the original message by evaluating Alice’s signal jointly with his entangled subsystem. For dense coding, assume that $`xX`$ is the message given to Alice. She encodes it by transforming her entangled system by a channel $`T_x`$, and sending the resulting quantum system to Bob, who measures an observable $`F`$ jointly on Alice’s particle and his. The probability for getting $`y`$ as a result is then $`tr\left(\omega (T_x\mathrm{id})(F_y)\right)`$, where the “$`\mathrm{id}`$” expresses the fact that no transformation is done to Bob’s particle while Alice applies $`T_x`$ to hers. If everything works correctly, this expression has to be $`1`$ for $`x=y`$, and $`0`$ otherwise (see eq. (3)).
Let us take a similar look at teleportation. Here three quantum systems are involved: the entangled pair in state $`\omega `$, and the input system given to Alice, in state $`\rho `$. Thus the overall initial state is $`\rho \omega `$. Alice measures an observable $`F`$ on the first two factors, obtaining a result $`x`$ sent to Bob. Bob applies a transformation $`T_x`$ to his particle, and makes a final measurement of an observable $`A`$ of his choice. Thus the probability for Alice measuring $`x`$ and for Bob getting a result “yes” on $`A`$, is $`tr(\rho \omega )(F_xT_x(A))`$. Note that the tensor symbols in this equation refer to different splittings of the system ($`123`$ and $`123`$, respectively). Teleportation is successful, if the overall probability for getting $`A`$, computed by summing over all possibilities $`x`$, is the same as for an ideal channel, i.e., $`tr(\rho A)`$, as in eq. (2).
The only relationship between the Hilbert spaces involved, which this description requires, is that the input and output spaces of the teleportation line are the same, since the whole teleportation process is equivalent to the identity. In some sense the best results (minimal dimension for the Hilbert spaces carrying the entangled state, best ratio of achieved capacity to capacity used) are obtained in the special case, where all Hilbert spaces have the same dimension $`d`$, and exactly $`|X|=d^2`$ signals are distinguished. We call this the tight case, and the main Theorem refers only to this case.
###### Theorem 1
Let $``$ be a $`d`$-dimensional Hilbert space ($`d<\mathrm{}`$), and $`X`$ a set of $`d^2`$ elements. Consider the following types of objects:
1. Teleportation schemes consisting of
* a density operator $`\omega `$ on $``$
* a collection of channels $`T_x:()()`$, $`xX`$
* an observable $`F_x`$, $`xX`$ on $``$
such that, for all density operators $`\rho `$ on $``$, and $`A()`$:
$$\underset{xX}{}tr(\rho \omega )(F_xT_x(A))=tr\rho A.$$
(2)
2. Dense coding schemes, consisting of the same objects as a Teleportation Scheme, but satisfying, instead of (2), the equation
$$tr\left(\omega (T_x\mathrm{id})(F_y)\right)=\delta _{xy}.$$
(3)
3. Bases of maximally entangled vectors, i.e. families of maximally entangled vectors $`\mathrm{\Phi }_x`$, $`xX`$ such that
$$\mathrm{\Phi }_x|\mathrm{\Phi }_y=\delta _{xy}$$
(4)
4. Bases of unitary operators, i.e., collections of unitary operators $`U_x()`$, $`xX`$ such that
$$tr(U_x^{}U_y)=d\delta _{xy}.$$
(5)
5. Unitary depolarizers,i.e., collections of unitary operators $`U_x()`$, $`xX`$ such that for any $`A()`$:
$$\underset{x}{}U_x^{}AU_x=dtr(A)\mathrm{\hspace{0.17em}1}\mathrm{I}.$$
(6)
Then, given any object of any one of these types, one can construct an object of each of the types, using the following equations:
$`\omega `$ $`=`$ $`|\mathrm{\Omega }\mathrm{\Omega }|,\text{with }\mathrm{\Omega }\text{ maximally entangled.}`$ (7)
$`F_x`$ $`=`$ $`|\mathrm{\Phi }_x\mathrm{\Phi }_x|`$ (8)
$`T_x(A)`$ $`=`$ $`U_x^{}AU_x`$ (9)
$`\mathrm{\Phi }_x`$ $`=`$ $`(U_x1\mathrm{I})\mathrm{\Omega }.`$ (10)
The logical structure of this result is maybe slightly unusual, so we begin by giving some examples how it is used. We can use it, for example, as a construction procedure: once we are given a unitary basis, we can get from the equations (7) to (10) a teleportation scheme and a dense coding scheme. Moreover, since we could also start with these schemes, ending up with the unitary basis we are assured that every teleportation or dense coding scheme is obtained in this way, i.e., this construction is exhaustive. In particular, we learn that any tight teleportation scheme is necessarily of a very special form: the entangled state $`\omega `$ must be pure and maximally entangled, the channels $`T_x`$ must be unitarily implemented, and the observable $`F`$ must be a complete von Neumann measurement.
Another result contained in this Theorem is the amazing equivalence between (1) and (2): any teleportation scheme works as a dense coding scheme and conversely. Alice and Bob merely have to swap their equipment to convert one into the other. We must emphasize, however, that the tightness condition is absolutely crucial for this equivalence. For simplicity, we will discuss this only in the case that $`|X|=n`$ is not fixed to be $`d^2`$, leaving aside the more difficult question what kind of trade-off between resources becomes possible, when $`\omega `$ lives on $`_1_2`$, with dimensions other than $`dd`$.
The basic difference between teleportation and dense coding is that the parameters $`d`$ and $`n`$ have opposite roles: For teleportation $`d`$ describes the size of the signal to be sent, and $`n`$ describes a resource, so the problem becomes more difficult when we increase $`d`$ and decrease $`n`$. For dense coding, it is exactly the opposite. Therefore, it is easy to show that teleportation (resp. dense coding) schemes exist whenever $`nd^2`$ (resp. $`nd^2`$). In fact, for teleportation one can take $`X`$ to be a continuum, and replace the sum in the teleportation equation by an integral \[BD\], but the dense coding equation would make no sense then. The optimality of these dimension inequalities, i.e., that no teleportation (resp. dense coding) scheme exists with $`n<d^2`$ (resp. $`n>d^2`$), is also a corollary of Theorem 1. To prove it, suppose we had a teleportation scheme with $`n<d^2`$. Then we could add $`(d^2n)`$ irrelevant classical signals happening with probability zero ($`F_x=0`$), and apply the Theorem, which says that all $`F_x`$ must be non-zero after all. The same reasoning works for dense coding with the operation of throwing in a few unused Hilbert space dimensions.
Of course, our Theorem is efficient as a construction procedure for dense coding and teleportation schemes only to the extent that unitary bases can be generated. After giving the proof of the Theorem, we will therefore describe the most general construction for such bases known to us.
## 3 Proof of Theorem 1
### Proof of the Implications “3$``$4”
Implicit in the formulation of the Theorem is the claim that the equation (10) $`\mathrm{\Phi }_x=(U_x1\mathrm{I})\mathrm{\Omega }`$ not only determines $`\mathrm{\Phi }_x`$ in terms of $`U_x`$ but also, conversely, determines $`U_x`$ in terms of $`\mathrm{\Phi }_x`$. This connection is based on a general construction, by which the $`d^2`$ matrix elements of an operator $`A:`$ are identified with the $`d^2`$ components of a vector $`\mathrm{\Psi }`$. This identification depends on the choice of a maximally entangled vector $`\mathrm{\Omega }`$. By choosing appropriate orthonormal bases $`e_k`$, $`k=1,\mathrm{},d`$ in the first and second tensor factor, such a vector can be written in “Schmidt form” as
$$\mathrm{\Omega }=\frac{1}{\sqrt{d}}\underset{k}{}e_ke_k.$$
(11)
Then a one-to-one correspondence between operators $`A()`$ and $`\mathrm{\Psi }`$ is given by the equation $`e_k|Ae_{\mathrm{}}=\sqrt{d}e_ke_{\mathrm{}}|\mathrm{\Psi }`$. We will use this in the form
$$\mathrm{\Psi }=(A1\mathrm{I})\mathrm{\Omega }=(1\mathrm{I}A^T)\mathrm{\Omega },$$
(12)
where the transpose operation $`AA^T`$ is defined in the basis $`e_k`$. Then if $`A`$ and $`\mathrm{\Psi }`$ and, similarly, $`A^{}`$ and $`\mathrm{\Psi }^{}`$ are related in this way,
$$\mathrm{\Psi }|(B1\mathrm{I})\mathrm{\Psi }^{}=\frac{1}{d}tr(A^{}BA^{}),$$
(13)
for arbitrary $`B()`$. Thus $`\mathrm{\Psi }`$ is maximally entangled iff this expression (for $`A=A^{}`$) is equal to $`d^1tr(B)`$, i.e., iff $`A`$ is unitary. Moreover, setting $`B=1\mathrm{I}`$, the scalar product of vectors $`\mathrm{\Psi },\mathrm{\Psi }^{}`$ is translated to $`d^1tr(A^{}A^{})`$ in terms of $`A,A^{}`$. Taking all this together, we get the one-to one correspondence between unitary bases and bases of maximally entangled vectors, as claimed. Note, however, that this correspondence depends on the choice of the reference maximally entangled vector $`\mathrm{\Omega }`$.
### Proof of the Implications “4$``$5”
This proof is relatively straightforward, since we are talking about only one type of objects, collections of $`d^2`$ unitaries $`U_x()`$. It is, however, also a crucial step for the entire proof, since it is here that the consequences of the tightness condition are seen. We will prove this in a form, which is also needed later to establish that the state $`\omega `$ in teleportation and dense coding schemes is necessarily maximally entangled.
The basic observation concerning matching dimensions is the following.
###### Lemma 2
$`D`$ vectors $`\varphi _1,\mathrm{},\varphi _D`$ in a $`D`$-dimensional Hilbert space form an orthonormal basis if and only if
$$\underset{k=1}{\overset{D}{}}|\varphi _k\varphi _k|=1\mathrm{I}.$$
(14)
Of course, this is false when there are more vectors than the dimension of the Hilbert space. Such families of vectors are called “overcomplete”. They exist and are an interesting mathematical structure of their own. On the other hand, fewer vectors than the dimension can never satisfy eq. (14), because the rank (dimension of the range) of the operator on the left hand side is at most the number of vectors.
Proof: It is a well-known fact that eq. (14) holds for any orthonormal basis. Conversely, we find from eq. (14) that, for each $`k`$, $`|\varphi _k\varphi _k|1\mathrm{I}`$, which is the same as $`\varphi _k^21`$. On the other hand, taking the trace of (14), we get $`_k\varphi _k^2=tr(1\mathrm{I})=D`$. This is only possible, when $`\varphi _k^2=1`$ for all $`k`$. Hence the operators $`|\varphi _k\varphi _k|`$ are hermitian projections, and we can invoke the observation that hermitian projections $`p_1,p_2`$ with $`p_1+p_21\mathrm{I}`$ are necessarily orthogonal. (For a quick proof, sandwich the inequality between factors $`p_1`$, finding $`p_1+p_1p_2p_1p_1`$, i.e., $`p_1p_2p_1=(p_2p_1)^{}(p_2p_1)0`$, and hence $`p_2p_1=0`$).
We now apply this Lemma to a collection of $`D=d^2`$ operators in $`()`$, where this space is considered as a Hilbert space with a suitable scalar product.
###### Proposition 3
Consider $`d^2`$ operators $`K_1,\mathrm{},K_{d^2}`$ on a $`d`$-dimensional Hilbert space $``$, and let $`R>0`$ be an invertible operator on $``$.
Then the following conditions are equivalent
1. $`tr(K_x^{}R^1K_y)=\delta _{xy}`$, for $`x,y=1,\mathrm{},d^2`$
2. $`_xK_x^{}CK_x=tr(RC)1\mathrm{I}`$ for all $`C()`$.
Proof: Let us define a scalar product $`|_R`$ on $`()`$ by
$$A|B_R=tr(A^{}R^1B).$$
(15)
Since $`R`$ is positive and invertible, this is indeed a scalar product, satisfying $`A|A_R=0`$ only for $`A=0`$. Condition 1 then simply says that the $`K_x`$ are an orthonormal basis. By the previous Lemma this is equivalent to the completeness relation (14), so all we have to do is to show that this relation, adapted to the special scalar product at hand, is equivalent to Condition 2 of the present Lemma. The completeness relation is that, for any $`A,B()`$,
$$A|B_R=\underset{x}{}A|K_x_RK_x|B_R.$$
($``$)
It suffices to evaluate this on rank one operators $`A,B()`$, since these span the whole space. We take $`A=|\varphi _1\varphi _2|`$ and $`B=|\psi _1\psi _2|`$. Then the left hand side of equation ($``$) becomes
$$\varphi _1|R^1\psi _1\psi _2|\varphi _2,$$
($``$LHS)
whereas the right hand side is
$$\underset{x}{}\varphi _1|R^1K_x\varphi _2\psi _2|K_x^{}R^1\psi _1=\psi _2|M\varphi _2,$$
($``$RHS)
with
$$M=\underset{x}{}K_x^{}R^1|\psi _1\varphi _1|R^1K_x\underset{x}{}K_x^{}CK_x$$
where we have interchanged the two factors in each term, and introduced the abbreviation $`C`$. Since ($``$LHS)=($``$RHS) for every $`\psi _2,\varphi _2`$, we find $`M=\varphi _1|R^1\psi _1\mathrm{\hspace{0.33em}1}\mathrm{I}`$. The factor is readily identified as
$`\varphi _1|R^1\psi _1=tr(RC)`$. Since operators of the form $`C`$ span $`()`$, the completeness relation thus becomes equivalent to $`_xK_x^{}CK_x=tr(RC)1\mathrm{I}`$ for all $`C`$, which completes the proof.
The special case of this Proposition, where each $`K_x`$ is unitary and $`R=\frac{1}{d}1\mathrm{I}`$, is exactly the relationship between items 4 and 5 of Theorem 1. However, there is another consequence needed later on:
###### Corollary 4
Let $`U_1,\mathrm{},U_{d^2}()`$ be unitaries in a $`d`$-dimensional Hilbert space $``$, and $`\rho `$ a density operator such that $`tr(U_x^{}\rho U_y)=\delta _{xy}`$. Then $`\rho =d^11\mathrm{I}`$.
Proof: Since the $`U_x`$ are an orthonormal set whose cardinality is the dimension, there can be no null vectors of this scalar product, i.e., $`tr(A^{}\rho A)=0`$ implies $`A=0`$. Hence $`\rho `$ is invertible, and we can apply the previous Proposition with $`R=\rho ^1`$, finding that $`_xU_x^{}AU_x=tr(\rho ^1A)1\mathrm{I}`$. The trace of this equation is $`d^2tr(A)=dtr(\rho ^1A)`$. This holds for all $`A`$, i.e., $`\rho ^1=d1\mathrm{I}`$.
### Proof of “(3 or 4)$``$(1 and 2)”
Suppose now we are given either a basis of unitary operators or of maximally entangled vectors. Then we can choose a maximally entangled vector $`\mathrm{\Omega }`$ and use equation (10) as in in the proof of “3$``$4” to define the other kind of basis. Equations (8) and (9) then become explicit definitions of the observable $`F_x`$ and the transformations $`T_x`$, respectively, so all the objects needed for a teleportation or dense coding scheme are defined, and we only need to verify that equations (2) and (3) are indeed satisfied.
In the teleportation equation an expectation value is generated between a state on the first and an observable on the third factor of a triple tensor product. This is a consequence of a similar “teleportation equation” on the level of vectors, which we now state. For later use we prove a certain converse at the same time.
###### Lemma 5
Let $`\mathrm{\Omega }^d^d`$ be the maximally entangled vector $`\mathrm{\Omega }=d^{1/2}_ke_ke_k`$, where $`e_k`$, $`k=1,\mathrm{},d`$ is the standard basis of $`^d`$. Let $`M(^d)`$, and $`\mu `$. Then the equation
$$\varphi \mathrm{\Omega }|(1\mathrm{I}M1\mathrm{I})\mathrm{\Omega }\psi =\mu \varphi |\psi $$
holds for all $`\varphi ,\psi ^d`$, if and only if $`M=d\mu 1\mathrm{I}`$.
Proof: Inserting the sum defining $`\mathrm{\Omega }`$ we get
$`\varphi \mathrm{\Omega }|(1\mathrm{I}M1\mathrm{I})\mathrm{\Omega }\psi =`$
$`=`$ $`{\displaystyle \frac{1}{d}}{\displaystyle \underset{k\mathrm{}}{}}\varphi e_ke_k|(1\mathrm{I}M1\mathrm{I})e_{\mathrm{}}e_{\mathrm{}}\psi `$
$`=`$ $`{\displaystyle \frac{1}{d}}{\displaystyle \underset{k\mathrm{}}{}}\varphi |e_{\mathrm{}}e_k|Me_{\mathrm{}}e_k|\psi ={\displaystyle \frac{1}{d}}\varphi |M^T\psi ,`$
which is equal to $`\mu \varphi |\psi `$ for all $`\varphi ,\psi `$ iff $`M^T=d\mu 1\mathrm{I}`$.
Consider now the term with index $`xX`$ in the teleportation equation (2), with $`F_x`$ and $`T_x`$ defined via equations (8), (9), and (10). Without loss of generality we set $`\rho =|\varphi _1\varphi _2|`$, $`A=|\psi _1\psi _2|`$. Then
$$\mathrm{term}_x=\varphi _2\mathrm{\Omega }|\mathrm{\Phi }_xU_x^{}\psi _1\mathrm{\Phi }_xU_x^{}\psi _2|\varphi _1\mathrm{\Omega }.$$
The first scalar product can be rewritten by substituting $`\mathrm{\Phi }_x`$ from equation (10), using equation (12):
$`\varphi _2\mathrm{\Omega }|\mathrm{\Phi }_xU_x^{}\psi _1`$ $`=`$ $`\varphi _2((1\mathrm{I}U_x)\mathrm{\Omega })|((U_x1\mathrm{I})\mathrm{\Omega })\psi _1`$
$`=`$ $`\varphi _2((U_x^T1\mathrm{I})\mathrm{\Omega })|((1\mathrm{I}U_x^T)\mathrm{\Omega })\psi _1`$
$`=`$ $`(1\mathrm{I}U_x^T1\mathrm{I})\varphi _2\mathrm{\Omega }|(1\mathrm{I}U_x^T1\mathrm{I})\mathrm{\Omega }\psi _1`$
$`=`$ $`\varphi _2\mathrm{\Omega }|\mathrm{\Omega }\psi _1={\displaystyle \frac{1}{d}}\varphi _2|\psi _1,`$
where at the last equation we used Lemma 5 with $`\mu =1/d`$. Together with a similar computation for the second scalar product, we get $`\mathrm{term}_x=d^2\varphi _2|\psi _1\psi _2|\varphi _1=d^2tr(\rho A)`$, and equation (2) follows by summing over $`d^2`$ equal terms.
Similarly, for the dense coding equation (3) we get
$$tr\left(\omega (T_x\mathrm{id})(F_y)\right)=\mathrm{\Omega }|(U_x^{}1\mathrm{I})\mathrm{\Phi }_y\mathrm{\Phi }_y|(U_x1\mathrm{I})\mathrm{\Omega },$$
i.e., the absolute square of the scalar product
$$\mathrm{\Omega }|(U_x^{}1\mathrm{I})\mathrm{\Phi }_y=\mathrm{\Omega }|(U_x^{}U_y1\mathrm{I})\mathrm{\Omega }=\frac{1}{d}tr(U_x^{}U_y)=\delta _{xy},$$
where we have used, in turn equation (10), the maximal entangledness of $`\mathrm{\Omega }`$ (see equation (1)), and the orthogonality of the $`U_x`$. This completes the proof of the dense coding property.
### Proof of the Implications “2 $``$ Rest”
Let us now assume that a dense coding scheme is given. We have to conclude that it is of the special form given in equations (710).
Note first that if $`\omega =_\alpha \lambda _\alpha \omega _\alpha `$ ($`\lambda _\alpha >0`$) is a mixture of states satisfying the teleportation equation, then every $`\omega _\alpha `$ also satisfies it. Hence the assumption is also satisfied for each pure component $`\omega _\alpha `$, and we can first analyze the problem assuming $`\omega `$ to be pure. In order to show that $`\omega `$ indeed is pure, we only have to verify that the given $`F,T`$ are consistent only with one pure state. So for the moment we will assume that $`\omega =|\mathrm{\Omega }\mathrm{\Omega }|`$ is pure.
The next step is a simple general observation on the coding of classical information on quantum channels, which we isolate in a Lemma.
###### Lemma 6
Let $`𝒦`$ be a $`D`$-dimensional Hilbert space, and $`\sigma _x,F_x(𝒦)`$, for $`xX`$, a set with $`D`$ elements. Suppose that each $`\sigma _x`$ is a density operator, $`F`$ is an observable, and $`tr(\sigma _xF_y)=\delta _{xy}`$, for $`x,y=1,\mathrm{},D`$.
Then there is an orthonormal basis $`\mathrm{\Phi }_x`$ such that
$$\sigma _x=F_x=|\mathrm{\Phi }_x\mathrm{\Phi }_x|.$$
Proof: Let $`\mathrm{\Phi }_x`$ be one of the normalized eigenvectors of $`\sigma _x`$ with non-zero eigenvalue. Then since $`F_x1\mathrm{I}`$, and $`\mathrm{\Phi }_x|F_x\mathrm{\Phi }_x=1`$, $`\mathrm{\Phi }_x`$ must also be an eigenvector of $`F_x`$ with eigenvalue $`1`$. Similarly, for any $`yx`$ the $`F_x0`$, and the normalization $`_xF_x=1\mathrm{I}`$ forces $`F_y\mathrm{\Phi }_x=0`$. Hence the $`\mathrm{\Phi }_x`$ are orthonormal, and since their number is the dimension of the space, they must be a basis. Consequently we have jointly diagonalized the $`F_x`$ and the $`\sigma _x`$, with eigenvalues either $`0`$ or $`1`$.
We apply this Lemma with $`D=d^2`$ and $`\sigma _x`$ the state after application of $`T_x`$ to the first factor, i.e., $`tr(\sigma _xA)=tr(\omega (T_x\mathrm{id})(A))`$. This proves equation (8), although it remains to be seen that each $`\mathrm{\Phi }_x`$ is maximally entangled.
Since the $`\sigma _x`$ form a maximal set of pure states, there cannot be a non-zero projection $`P`$ such that, for all $`xX`$,
$`0`$ $`=`$ $`tr(\sigma _x(1\mathrm{I}P))=tr(\omega (T_x\mathrm{id})(1\mathrm{I}P))`$
$`=`$ $`tr(\omega (1\mathrm{I}P))=\mathrm{\Omega }|(1\mathrm{I}P)\mathrm{\Omega }.`$
Hence $`\mathrm{\Omega }`$ must have full Schmidt rank. We will need the consequence that the equation $`(A1\mathrm{I})\mathrm{\Omega }=(A^{}1\mathrm{I})\mathrm{\Omega }`$ implies $`A=A^{}`$.
Let $`T_x(A)=_\alpha K_{x,\alpha }^{}AK_{x,\alpha }`$ be the Kraus decomposition of $`T_x`$. Then the teleportation equation is
$$\underset{\alpha }{}|\mathrm{\Omega }|(K_{x,\alpha }^{}1\mathrm{I})\mathrm{\Phi }_y|^2=\delta _{xy}.$$
Therefore, $`(K_{x,\alpha }1\mathrm{I})\mathrm{\Omega }|\mathrm{\Phi }_y=0`$ for all $`yx`$, and for every $`x`$ there must be constants $`c_\alpha `$ such that
$$(K_{x,\alpha }1\mathrm{I})\mathrm{\Omega }=c_\alpha \mathrm{\Phi }_x.$$
(16)
Since $`\mathrm{\Omega }`$ has full Schmidt rank, this implies that all $`K_{x,\alpha }`$ are proportional to each other, i.e., that $`T_x`$ can be written with a single Kraus summand. Of course, the corresponding $`K_xU_x`$ must be unitary, and since both sides are normalized, equation (16) $`\mathrm{\Phi }_x=(U_x1\mathrm{I})\mathrm{\Omega }`$, possibly after fixing suitable phase factors (which influence neither $`T_x`$ nor $`F_x`$).
The orthonormality of the $`\mathrm{\Phi }_x`$ translates into $`tr(\rho U_x^{}U_y)=\delta _{xy}`$, where $`\rho `$ is the reduced density operator of $`\omega `$. But then Corollary 4 shows that $`\rho `$ must be a multiple of the identity, i.e., $`\mathrm{\Omega }`$ and each $`\mathrm{\Phi }_x`$ is maximally entangled.
Finally, we have to complete the argument for the purity of $`\omega `$ by showing that only one pure state is consistent with the other data $`T,F`$, encoded in $`U_x`$. But this is obvious from the explicit expression $`\mathrm{\Omega }=(U_x^{}1\mathrm{I})\mathrm{\Phi }_x`$.
### Proof of the Implications “1 $``$ Rest”
Let us now assume that a teleportation scheme is given. We have to conclude that it is of the special form given in equations (710).
The crucial input for this proof is the principle that in quantum mechanics there is no measurement without perturbation. It enters in the following form, a corollary of the so-called Radon-Nikodym Theorem for completely positive maps. We state it here as a Lemma.
###### Lemma 7
Let $``$ be a finite dimensional Hilbert space, and let $`T_\alpha :()()`$ be completely positive maps such that $`_\alpha T_\alpha =\mathrm{id}`$. Then there are positive numbers $`t_\alpha `$ such that $`T_\alpha =t_\alpha \mathrm{id}`$.
Proof: For readers less familiar with dilation theory of cp-maps we include a quick proof based on the Kraus decomposition $`T(A)=_\beta K_\beta ^{}AK_\beta `$, which exists for every completely positive map. Note that by decomposing each $`T_\alpha `$ in Kraus form, we get a finer decomposition of $`\mathrm{id}`$, so we may as well prove the Lemma for the case that each $`T_\alpha `$ is of the form $`T_\alpha (A)=K_\alpha ^{}AK_\alpha `$. With $`A=|\psi \psi |`$,
$$|K_\alpha ^{}\psi K_\alpha ^{}\psi |\underset{\alpha }{}|K_\alpha ^{}\psi \psi |K_\alpha =|\psi \psi |.$$
Hence $`K_\alpha ^{}\psi =\lambda (\psi )\psi `$, with a factor $`\lambda (\psi )`$. But then every vector $`\psi `$ is an eigenvector of the linear operator $`K_\alpha ^{}`$, which is only possible, if $`K_\alpha ^{}`$ is a multiple of the identity.
A collection of completely positive maps adding up to a normalized one should be understood as an “instrument” in the terminology of Davies \[Da\], i.e., a device which produces classical measurement results “$`k`$”, such that the probability for obtaining this result and a response to a subsequent measurement $`F`$ on an input state $`\rho `$ is $`tr(\rho T_k(F)))`$. The channel $`_kT_k`$ then describes the overall state change, when the measuring results are ignored. In this language the hypothesis of the Lemma says that there is no overall state change through the device, i.e., “no perturbation” of the system. The conclusion is that in that case the output probabilities are $`t_k`$, and independent of the input state, i.e., no information about the system is obtained.
As a first application, we conclude exactly as in the previous subsection that each convex component of the state $`\omega `$ again satisfies the teleportation equation. Hence we can once more assume that $`\omega =|\mathrm{\Omega }\mathrm{\Omega }|`$ is a pure state. The argument that $`\mathrm{\Omega }`$ is then uniquely determined by the other data, and hence that $`\omega `$ is pure is the same as in the dense coding case.
Clearly, this kind of argument is also useful for decompositions of $`T_x`$ or $`F_x`$ into sums of (completely) positive terms. To do this systematically, fix a maximally entangled unit vector $`\mathrm{\Xi }`$, so that vectors in $``$ become expressed as $`\mathrm{\Phi }=(A1\mathrm{I})\mathrm{\Xi }`$ for a uniquely determined operator $`A`$ (see equation (12)). In particular, we can write $`\mathrm{\Omega }=(W1\mathrm{I})\mathrm{\Xi }`$, and the Kraus decomposition and spectral decomposition of each $`F_x`$ in the form
$`T_x(A)`$ $`=`$ $`{\displaystyle \underset{\alpha }{}}K_{x,\alpha }^{}AK_{x,\alpha }`$ (17)
$`F_x`$ $`=`$ $`{\displaystyle \underset{\beta }{}}(A_{x,\beta }1\mathrm{I})|\mathrm{\Xi }\mathrm{\Xi }|(A_{x,\beta }1\mathrm{I})^{}.`$ (18)
Inserting this into the teleportation equation (2) we find a sum over $`x,\alpha ,\beta `$, in which each term represents a completely positive operator, and which sum up to the identity. Hence by Lemma 7, each term has to be multiple of the identity, $`\stackrel{~}{\mu }_{x,\alpha ,\beta }\mathrm{id}`$, say. This can be written in terms of scalar products, if we take $`\rho =|\varphi _1\varphi _2|`$ and $`A=|\psi _1\psi _2|`$:
$`\stackrel{~}{\mu }_{x,\alpha ,\beta }tr(\rho A)`$ $`=`$ $`\stackrel{~}{\mu }_{x,\alpha ,\beta }\varphi _2,\psi _1\psi _2,\varphi _1`$
$`=`$ $`\varphi _2\mathrm{\Omega },(A_{x,\beta }1\mathrm{I}K_{x,\alpha }^{})\mathrm{\Xi }\psi _1`$
$`\mathrm{\Xi }\psi _2,(A_{x,\beta }^{}1\mathrm{I}K_{x,\alpha })\varphi _1\mathrm{\Omega }.`$
Note that the two scalar products on the right hand side are complex conjugates of each other apart from a swapping of the arguments $`(\varphi _2,\psi _1)`$ and $`(\psi _2,\varphi _1)`$, which exactly matches the variable pairing on the left hand side. Since the equation is to hold for arbitrary vectors $`\varphi _1,\varphi _2,\psi _1,\psi _2`$, we can hold one pair fixed and find that
$$\varphi _2\mathrm{\Omega },(A_{x,\beta }1\mathrm{I}K_{x,\alpha }^{})\mathrm{\Xi }\psi _1=\mu _{x,\alpha ,\beta }\varphi _2,\psi _1,$$
(19)
where $`\mu _{x,\alpha ,\beta }`$ is a factor determined in terms of $`\stackrel{~}{\mu }`$, and the scalar products involving $`(\psi _2,\varphi _1)`$. With $`\mathrm{\Omega }=(W1\mathrm{I})\mathrm{\Xi }`$, and equation (12) we get
$`\varphi _2\mathrm{\Omega },(A_{x,\beta }1\mathrm{I}K_{x,\alpha }^{})\mathrm{\Xi }\psi _1`$
$`=`$ $`\varphi _2(1\mathrm{I}K_{x,\alpha })\mathrm{\Xi },(1\mathrm{I}W^{}1\mathrm{I})\left((A_{x,\beta }1\mathrm{I})\mathrm{\Xi }\psi _1\right)`$
$`=`$ $`\varphi _2(K_{x,\alpha }^T1\mathrm{I})\mathrm{\Xi },(1\mathrm{I}W^{}1\mathrm{I})\left((1\mathrm{I}A_{x,\beta }^T)\mathrm{\Xi }\psi _1\right)`$
$`=`$ $`\varphi _2\mathrm{\Xi },(1\mathrm{I}\overline{K_{x,\alpha }}W^{}A_{x,\beta }^T1\mathrm{I})\mathrm{\Xi }\psi _1`$
$``$ $`\mu _{x,\alpha ,\beta }\varphi _2,\psi _1,`$
where we have used the notation $`\overline{K}=(K^{})^T`$ for the matrix element-wise complex conjugation in the Schmidt basis belonging to the maximally entangled state $`\mathrm{\Xi }`$. Since the above equation holds for all $`\varphi _2`$ and $`\psi _1`$, Lemma 5 implies that
$$\overline{K_{x,\alpha }}W^{}A_{x,\beta }^T=d\mu _{x,\alpha ,\beta }1\mathrm{I},$$
(20)
for all $`x,\alpha ,\beta `$.
Let us say that a label $`xX`$ contributes to teleportation, if the corresponding term in the teleportation equation does not vanish for all $`\rho `$ and $`A`$. This is equivalent to saying that for some $`\alpha ,\beta `$ the factor $`\mu _{x,\alpha ,\beta }`$ is non-zero. For such triples $`(x,\alpha ,\beta )`$ all three operators on the left hand side of equation (20) have to be invertible.
Now since there has to be at least one contributing label, $`W`$ has to be non-singular, which means that $`\mathrm{\Omega }`$ has full Schmidt rank. Equivalently, the reduced density operator $`\omega _1`$ for the first factor has no zero eigenvalues. From this we conclude that the non-contributing labels are precisely those for which $`F_x=0`$. Indeed, we may set $`A=\rho =1\mathrm{I}`$, and use the normalization of $`T_x`$ to find
$$0=tr\left((1\mathrm{I}\omega )(F_x1\mathrm{I})\right)=tr\left((1\mathrm{I}\omega _1)F_x\right)$$
Since $`F_x0`$, and $`1\mathrm{I}\omega _1`$ has only strictly positive eigenvalues, this implies $`F_x=0`$.
Now let $`x`$ be a contributing index, and choose some triple $`(x,\alpha ,\beta )`$ with $`\mu _{x,\alpha ,\beta }0`$. If we now look at equation (20) for triples $`(x,\alpha ^{},\beta )`$ with arbitrary $`\alpha ^{}`$, we get $`\overline{K_{x,\alpha ^{}}}=(\mu _{x,\alpha ^{},\beta }/\mu _{x,\alpha ,\beta })\overline{K_{x,\alpha }}`$, i.e., all Kraus operators of $`T_x`$ are proportional, and hence $`T_x`$ can be written with a single Kraus summand, $`T_x(A)=U_x^{}AU_x`$, with a unitary $`U_x`$.
Similarly, we find that all $`A_{x,\beta ^{}}`$ are proportional, which means that $`F_x=|\mathrm{\Phi }_x\mathrm{\Phi }_x|`$ with $`\mathrm{\Phi }_x=(A_x1\mathrm{I})\mathrm{\Xi }`$.
We can now apply Lemma 2 to these vectors $`\mathrm{\Phi }_x`$, setting $`\mathrm{\Phi }_x=0`$ for non-contributing labels. The conclusion is that the $`\mathrm{\Phi }_x`$ are an orthonormal basis. In particular, all indices do contribute after all.
Equation (20) and the unitarity of $`U_x`$ allow us to express $`A_x`$ in terms of $`U_x`$:
$$A_x=d\mu _xU_x\overline{W}^1$$
(21)
Orthonormality of the $`\mathrm{\Phi }_x`$ becomes
$$\delta _{xy}=\frac{1}{d}tr(A_x^{}A_y)=d\overline{\mu _x}\mu _ytr(U_y\overline{W}^1(\overline{W}^1)^{}U_x^{}).$$
(22)
For $`x=y`$ we find that $`|\mu _x|^2`$ is independent of $`x`$, hence the operators $`(\overline{\mu _x}/|\mu _x|)U_x^{}`$ are unitary, and satisfy the hypothesis of Corollary 4 with $`\rho `$ a positive multiple of $`\overline{W}^1(\overline{W}^1)^{}`$. Hence this operator is a multiple of the identity, $`W`$ is unitary up to a factor, and $`\mathrm{\Omega }=(W1\mathrm{I})\mathrm{\Xi }`$ is maximally entangled. Moreover, we see from equation (22) and the $`U_x`$ form a unitary basis.
Since $`\mathrm{\Xi }`$ was an arbitrary maximally entangled vector, we may just as well take $`\mathrm{\Xi }=\mathrm{\Omega }`$, so equation (21) holds with $`W=1\mathrm{I}`$. Hence, $`\mathrm{\Phi }_x=c(U_x1\mathrm{I})\mathrm{\Omega }`$, where $`c`$ is a factor which has to be of modulus $`1`$, because $`\mathrm{\Omega }`$ and $`\mathrm{\Phi }_x`$ are normalized, and $`U_x`$ is unitary, and which can be chosen to be $`1`$ by adjusting the phase of $`\mathrm{\Phi }_x`$. This completes the proof.
## 4 Constructing bases of unitaries
It is not a priori clear that bases of unitary operators should exist in any dimension. Indeed, the system equation (5) of equations is formally overdetermined, according to the following rough dimension count. The variables in this system are the unitaries $`U_x`$, each of which we can take in the $`(d^21)`$-dimensional manifold $`SU_d`$, i.e., with $`det\left(U_x\right)=1`$, by fixing a phase factor. Since the transformations $`U_xV_1U_xV_2`$, for arbitrary $`V_1,V_2SU_d`$ leave the set of solutions invariant, we may fix $`U_1=1\mathrm{I}`$, and take $`U_2`$ diagonal without loss of generality. This reduces the number of variables to $`(d1)+(d^22)(d^21)`$. On the other hand, orthogonality introduces one complex constraint for every pair $`xy`$. None of these is trivially satisfied due to the special choices we made, so we have to take $`d^2(d^21)`$ constraints into account. This leaves, formally,
$$\text{\#variables-\#equations}=(d1)(2d+1)<0.$$
Of course, we know that this count is somehow too crude, because, after all, many inequivalent unitary bases are constructed below. But it is not so easy to spot the dependences among the constraints. Note also that the dimension count is essentially the same for bases orthogonal with respect to a weight $`\rho d^11\mathrm{I}`$, but in that case Corollary 4 shows that there is no solution at all.
In order to describe the best known construction for unitary bases \[VW\], let us introduce some terminology. We say that a (single) unitary matrix is of shift and multiply type, if it is the product of a permutation operator and a diagonal unitary. In other words, every row or column contains $`(d1)`$ zero entries, and one entry of modulus $`1`$. The bases we will construct not only have the property that each element is of this type, but also that the $`d^2`$ values for $`x`$ can be split into $`d`$ options for “shift” and $`d`$ options for “multiply”.
###### Definiton 8
A shift and multiply basis of unitary matrices in $`^d`$ is a collection of $`d^2`$ unitary operators $`U_{ij}`$, $`i,jI_d\{1,\mathrm{},d\}`$, satisfying the orthogonality relation $`tr(U_{ij}^{}U_k\mathrm{})=d\delta _{ik}\delta _j\mathrm{}`$, and acting on the basis vectors $`|k`$ as
$$U_{ij}|k=H_{ik}^j|\lambda (j,k),$$
(23)
where the $`H_{ik}^j`$ are complex numbers, and $`\lambda :I_d\times I_dI_d`$.
###### Proposition 9
The parameters and $`\lambda :I_d\times I_dI_d`$ define a shift and multiply basis of unitary matrices if and only if the following two conditions are satisfied
1. Each $`H^j`$ is a Hadamard matrix, i.e. $`|H_{ik}^j|=1`$ for all $`i,k`$, and $`H^j(H^j)^{}=d\mathrm{\hspace{0.33em}1}\mathrm{I}`$.
2. $`\lambda `$ is a Latin square, i.e., the maps $`k\lambda (k,\mathrm{})`$ and $`k\lambda (\mathrm{},k)`$ are injective for every $`\mathrm{}`$.
Proof: For $`U_{ij}`$ to be unitary, it is necessary and sufficient that the $`H_{ik}^j`$ are phases, and that $`k\lambda (j,k)`$ is injective (hence bijective) for every $`j`$. For the orthogonality we have to evaluate
$$tr(U_{ij}^{}U_{i^{}j^{}})=\underset{k}{}\overline{H_{ik}^j}H_{i^{}k}^j^{}\lambda (j,k)|\lambda (j^{},k).$$
We consider first the case $`j=j^{}`$. Then the scalar products in the sum are all equal to $`1`$, and equating this expression to $`\delta _{ii^{}}`$ we find that $`H^j`$ is Hadamard.
Now let $`jj^{}`$, and consider the “coincidence set” $`C=\{k\lambda (j,k)=\lambda (j^{},k)\}`$. Then orthogonality requires, for every $`i,i^{}`$, that
$$0=\underset{kC}{}\overline{H_{ik}^j}H_{i^{}k}^j^{}=\underset{k=1}{\overset{d}{}}H_{i^{}k}^j^{}\chi _C(k)(H^j)_{ki}=(H^j^{}\chi _CH^j)_{i^{}i},$$
(24)
where $`\chi _C(k)=1`$ for $`kC`$, and zero otherwise, and in the last line $`\chi _C`$ denotes the projection $`\chi _C|k=\chi _C(k)|k`$. But since $`H^j^{}`$ and $`H^j`$ are Hadamard, and in particular invertible, this implies $`\chi _C=0`$. Hence $`C`$ is empty, and the second injectivity of $`\lambda `$ is proved.
In order to construct unitary bases of this form, we must now construct Hadamard matrices and Latin squares of the appropriate dimension. For both of these tasks there is a rich literature, and below we will give a brief summary on what is known for each.
It is useful to note that each of the structures ‘unitary bases’, ‘Hadamard matrices’, and ‘Latin squares’ has a natural notion of equivalence, and to some extent these equivalences are related. We call two unitary bases $`U,U^{}`$ equivalent, if $`U_x^{}=V_1U_x^{}V_2`$, for some unitaries $`V_1,V_2`$, and a re-labelling $`xx^{}`$. Hadamard matrices are called equivalent, if one is obtained from the other by permuting rows or columns, or multiplying rows or columns with phases. Finally, a Latin square $`\lambda :I_d\times I_dI_d`$ is equivalent to any other obtained by applying a permutation on each of the three copies of $`I_d`$ involved. In each case there are also discrete transformations, such as transposition or complex conjugation (where applicable). It should be noted that replacing each $`H^j`$ by an equivalent one, typically only leads to an equivalent unitary basis, if the equivalence operation is the same for each $`j`$. With $`j`$-dependent equivalence transformations it is possible to construct inequivalent unitary bases in $`d=3`$, although in this dimension there is only one Hadamard matrix and only one Latin square – up to equivalence. Of course, in $`d=2`$ all three structures, including the unitary bases are unique up to equivalence \[VW\]. The unique unitary basis is then given by the three Pauli matrices and the identity and, of course generates via the Theorem 1, the usual two qubit examples of teleportation and dense coding.
For each of the three structures we furthermore have an obvious notion of tensor product, allowing the construction of a unitary basis (resp. a Hadamard matrix, or Latin square) in dimension $`d=d_1d_2`$, if counterparts in dimension $`d_1`$ and dimension $`d_2`$ are given.
In order to show that unitary bases exist in any dimension it is easiest to use group theory based constructions: the Latin square can be taken as the multiplication table of any group of order $`d`$, for example the cyclic group. The Hadamard matrix can be taken as the matrix implementing the Fourier transform on an abelian group of order $`d`$, the standard example being given once again by the cyclic group of order $`d`$. Thus $`H_k\mathrm{}=\mathrm{exp}(\frac{2\pi i}{d}k\mathrm{})`$, where $`k`$ and $`\mathrm{}`$ are taken modulo $`d`$. If we combine these data into a unitary basis we get an instance of what we propose to call a unitary basis of group type (“nice error basis” in \[Kn\]). These are orthonormal unitary bases with the additional property that the operator product of any two elements is a third, up to a phase. That is to say, the index set $`X`$ is a group, and
$$U_xU_y=\mu (x,y)U_{xy},$$
(25)
with $`|\mu (x,y)|=1`$. In the special case of an abelian group $`X`$ this is a discrete version of Weyl systems of unitary operators, named after their continuous variable counterpart, well-known from quantum optics and non-relativistic “phase space” quantum mechanics.
Latin squares are not completely classified, nor does there seem to be a realistic hope to do so. A standard work on the subject is \[DK\], a useful net resource is \[Ri\]. Counts of squares are usually done for “normalized squares”, in which the first row and column are in natural order, thus eliminating some trivial freedom. In $`d=5`$ Euler counted $`56`$ of these, but only $`2`$ are inequivalent, because the symbols themselves can also be permuted. Counts of normalized squares have now gone all the way up to $`d=10`$, but are no longer done by hand (there are roughly $`7.5\times 10^{24}`$ \[MR\]). It is also clear from these numbers that group based constructions exhaust only a tiny fraction of the possible unitary bases.
Hadamard matrices are also a standard subject in coding theory. However, usually only the real case (orthogonal matrices with entries $`\pm 1`$) is considered. It is easy to see that real Hadamard matrices exist only in dimension two and multiples of four. Again, the possibilities for such designs by far exceed the group based possibilities (the characters of an abelian group are real only if $`d=2^n`$). A standard reference is \[Ag\].
For complex Hadamard matrices the Fourier matrices show that there is no constraint on dimension. The uniqueness in $`d=3`$ is easy to get. The general form in $`d=4`$ is, up to equivalence
$$\left(\begin{array}{cccc}1& 1& 1& 1\\ 1& 1& 1& 1\\ 1& 1& u& u\\ 1& 1& u& u\end{array}\right),$$
(26)
where $`u`$ is an arbitrary phase. For $`u=1`$ this is equivalent to the Fourier matrix of “Klein’s Four Group”, the product of two copies of the two-element group, and for $`u=i`$ it is equivalent to the Fourier matrix of the cyclic group. The possibility of embedding the cyclic group Fourier matrix into a higher dimensional manifold can be generalized to arbitrary composite numbers $`d=pq`$: whenever $`V_k\mathrm{}`$ is a matrix of phases satisfying the periodicity conditions $`V_{k,\mathrm{}}=V_{k+p,\mathrm{}}=V_{k,\mathrm{}+q}`$, we get a Hadamard matrix as
$$H_k\mathrm{}=V_k\mathrm{}\mathrm{exp}\left(\frac{2\pi i}{d}k\mathrm{}\right).$$
(27)
One might conjecture from this that for prime orders $`d`$ the Hadamard matrix is unique. This problem was discussed by Haagerup \[Ha\] on the basis of a completely different motivation (theory of von Neumann algebras). There it is shown that $`d=5`$ there is uniqueness, but for $`d=7`$ there are at least $`5`$ solutions. For some primes, uncountably many inequivalent Hadamard matrices are known.
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# Accretion by Galaxies
## 1. Introduction
The Milky Way is an extremely typical galaxy in the sense that most of the luminosity in the Universe comes from similar systems. Therefore, indications that the Milky Way is accreting at a significant rate would imply that accretion is an important phenomenon within the Universe at large.
Several observations suggest that material is falling into the Milky Way. First and foremost, our nearest large neighbour, M31, is approaching us at $`140\mathrm{km}\mathrm{s}^1`$. Second, the so-called ‘high-velocity clouds’ have, on the average, negative line-of-sight velocities. Third, the Magellanic Clouds are on an orbit that must be decaying through dynamical friction, and I shall argue that the orbit of the Sagittarius Dwarf Galaxy is in a still more advanced state of decay. Fourth, beyond $`R_0`$ the Galactic disk has a complex warped structure, and the most plausible explanation of this phenomenon invokes infall as the driver. Fifth, at a given radius the ISM appears to vary significantly in metallicity from place to place, and this phenomenon is most naturally explained by persistent accretion of low-metallicity material.
In this talk I shall review each of these lines of evidence in turn.
## 2. Collapse of the Local Group
The luminosity of the Local Group is dominated by the Milky Way and M31. Since M31 is approaching us, the Local Group is either virialized or collapsing. The simplest dynamical model in which the Group’s mass is concentrated in point particles coincident with the Milky Way and M31 suffices to show that the age of the Universe, combined with the present distance and Galactocentric velocity of M31, implies that the Local Group is collapsing for the first time (Kahn & Woltjer 1959).
When a cosmological structure collapses for the first time, its constituent parts move on highly elongated orbits towards the barycentre. Some of these orbits must enter the Galactic halo and become trapped. How much matter is the Galaxy likely to be accreting in this way?
The answer depends on two related quantities: how far out the Galactic halo reaches, and what fraction of the Local group lies outside M31 and the Galaxy. The mass of the Local Group has been estimated many times and found to lie in the range<sup>1</sup><sup>1</sup>1The Local Group’s $`V`$-band luminosity is $`4.2\times 10^{10}\mathrm{L}_{}`$ (e.g., Tables 2.1 and 4.3 of Binney & Merrifield, 1998), while the mean mass-to-light ratio of the Universe is $`400h(\mathrm{\Omega }/0.25)`$ (§10.3.1 of Binney & Tremaine, 1987), so for $`h=0.65`$ one expects the mass of the Local Group to be $`1.1\times 10^{13}`$. $`4`$ to $`8\times 10^{12}\mathrm{M}_{}`$ (e.g., Schmoldt & Saha 1998). If all this matter were in M31 and the Galaxy, with a third of the matter being in the Galaxy, the halo would have to extend at constant $`v_c=220\mathrm{km}\mathrm{s}^1`$ to $`120`$ to $`210\mathrm{kpc}`$. Consequently, material that comes within $`100\mathrm{kpc}`$ of the Galactic centre is likely to be accreted.
Blitz et al. (1999) describe a simulation of the formation of the Local Group in which $`10^6`$ test particles move in the field of two point masses representing M31 and the Galaxy, plus the tidal field generated by external galaxies from Raychaudhury & Lynden-Bell (1989). Particles that come within $`100`$ comoving kpc of either M31 or the Galaxy are captured. In the simulation the Hubble flow currently reverses $`1.5\mathrm{Mpc}`$ from the centre of the Local Group and the Galaxy’s current accretion rate is $`7.5\mathrm{M}_{}\mathrm{yr}^1`$, of which $`0.8\mathrm{M}_{}\mathrm{yr}^1`$ takes the form of neutral hydrogen.
Can we see infalling material? One school of thought has long held that the so-called high-velocity clouds (HVCs) are made of infalling gas (Oort 1966, 1970). HVCs were first identified when Muller, Oort & Raimond (1963) detected 21-cm emission in many directions at velocities that are incompatible with circular rotation. The location of the objects responsible for this emission has been controversial, however. One possibility is that they are small, nearby clouds that have been accelerated to large peculiar velocities by supernovae, stellar winds, and the like. Alternatively, they might be at distances in excess of a Mpc and be tracing the (disturbed) Hubble flow around the Local Group.
It is likely that no single explanation applies to all HVCs. Some of these objects are almost certainly associated with Local-Group galaxies such as M31, the Magellanic Clouds and the Phoenix dwarf spheroidal (St-Germain et al, 1999). Others, such as Complex A (van Woerden et al., 1998) and Complex M (Danly et al., 1993) are small clouds in the Galactic halo. But a powerful case can be mounted that many HVCs are systems over $`10\mathrm{kpc}`$ in diameter that lie at distances $`1\mathrm{Mpc}`$.
Braun & Burton (1999) identified a sample of 66 compact, isolated HVCs, 23 of which had not previously been catalogued. They showed that these objects are distributed fairly uniformly over the sky, define a mean velocity that, within the errors, agrees with the Local Group’s mean velocity (Karachentsev & Makarov, 1966), and with respect to this mean have a velocity dispersion of $`69\mathrm{km}\mathrm{s}^1`$ and an infall velocity $`100\mathrm{km}\mathrm{s}^1`$. They infer that these objects are typically $`15\mathrm{kpc}`$ in diameter and contain $`\mathrm{}>10^7\mathrm{M}_{}`$ of HI.
Blitz et al. (1999) point out first that HVCs are seen within the Galactic plane where local objects would collide at high velocity with gas in circular motion, yet there is no evidence for shock-excited gas. Second, attempts to detect HVCs in absorption against distant stars have been largely fruitless – only two HVCs have been detected in absorption, and there are other reasons for believing that these are two of three HVCs that are unusually nearby. Third, nearby clouds would be illuminated by UV photons from the Galaxy and be brighter H$`\alpha `$ sources than they actually are. Fourth, the velocities of the HVCs make more sense when referred to the barycentre of the Local Group than when referred to either the velocity of the Galactic centre or the Local Standard of Rest. Fifth, the metallicities of HVCs are low, which is inconsistent with their having been ejected from the Galactic ISM.
Blitz et al. use their simulation of the formation of the Local Group to show that the non-uniform distribution of HVCs on the sky arises naturally if HVCs lie at distances $`\mathrm{}>1\mathrm{Mpc}`$. Moreover, at such distances the HVCs would have a distribution over column density which resembles that of Ly$`\alpha `$ clouds, which are now known to be floating in intergalactic space (Theuns & Efstathiou, 1998; Davé et al. 1999). Hence, when the HVCs are interpreted as extragalactic objects, they fit naturally into the picture of intergalactic space that has emerged from a combination of theoretical and observational work in cosmology. Blitz et al. argue that objects lying near the top of the HVC mass spectrum have already been detected in external groups of galaxies, and that less massive objects are detectable with feasibly deep observations. They estimate that a typical cloud contains $`3\times 10^7\mathrm{M}_{}`$ of HI, so the population as a whole contains $`10^{10}\mathrm{M}_{}`$ of HI.
Once one accepts that many HVCs are extragalactic and associated with cosmic infall, it is inevitable that their detected hydrogen is associated with a very much larger mass of dark matter. The estimate of the Galaxy’s accretion rate cited above is based on the assumption that there is 10 times as much dark matter as neutral hydrogen. There could easily be more dark matter by a factor of a few, since $`\mathrm{\Omega }_{\mathrm{baryon}}`$ is thought to be $`\mathrm{}<0.05`$ and only a fraction of a HVC’s hydrogen content will be neutral.
## 3. Warps
Both theoretical and observational developments over the last few years have strengthened the argument that warps such as that of the Galactic disk are a reflection of cosmic infall, and the associated reorientation of galactic angular-momentum vectors. Since the discovery of the Galactic warp by Burke (1957) and Kerr (1957), there has been a debate as to whether warps are self-consistent intrinsic structures, or externally driven. Kahn & Woltjer (1959) suggested that the Galactic warp was driven by an intergalactic wind, while Lynden-Bell (1965) suggested that a warp can be an intrinsic structure. The latter proposal was shot down by Hunter & Toomre (1969). After the discovery of dark matter, Toomre (1983) and Dekel & Shlosman (1983) argued that warps might after all be intrinsic structures, and Sparke & Casertano (1988) developed this proposal to the point that it became very attractive. Recently Nelson & Tremaine (1995) and Binney, Jiang & Dutta (1998) have demonstrated that the Toomre–Dekel proposal, as elaborated by Sparke & Casertano, is not viable. The flaw in the approach of Sparke & Casertano is the treatment of the dark halo as a rigid, unflexing thing. The semi-analytic work of Nelson & Tremaine and the numerical simulations of Binney et al. demonstrate that a dark halo responds rapidly to any warp in an embedded disk in such a manner that warps predicted to endure for ever by Sparke & Casertano actually wind up within a few dynamic times.
Binney et al. did not demonstrate that intrinsic warps are impossible, but only that any long-lived warp would have to be a manifestation of a cooperative distortion of both disk and halo. However, a priori it is not clear that such configurations exist, and extensive numerical experimentation by several groups has failed to find one. Hence, the theoretical case for intrinsic warps must be considered at best doubtful. Moreover, new observations (Shang et al. 1998) of the classic example of an isolated galaxy with an integral-sign warp, NGC 5907, have shattered the observational case for the existence of warps in isolated galaxies: NGC 5907 possesses both a dwarf satellite $`37\mathrm{kpc}`$ from its centre, and an elliptical ring of luminosity that probably contains the debris of another, now shredded, companion.
What is the theoretical status of externally driven warps? The original proposal of Kahn & Woltjer (1959) is not viable because it relies upon a massive, subsonic wind past the Milky Way, for which there is no evidence. Similarly, proposals involving magnetic fields (Battaner, Florido & Sánchez-Saavedra, 1990) are for various reasons not in serious contention. An idea that remains plausible is that warps are a response to the accretion by a galaxy of material laden with angular momentum about an axis that is inclined to the galaxy’s original spin axis (Binney & May, 1986; Ostriker & Binney, 1989). Jiang & Binney (1999) followed the dynamics of a disk embedded in a live halo by decomposing the disk into a series of massive rings that interact gravitationally with one another and with the $`\mathrm{100\hspace{0.17em}000}`$ particles that represent the halo. The disk is exponential out to $`R=3.5R_\mathrm{d}`$ and is then smoothly tapered to zero surface density at $`R=4R_\mathrm{d}`$. The halo initially has ten times the mass of the disk and ensures that the overall rotation curve is very nearly flat out to $`5R_\mathrm{d}`$. The accretion of material by the halo is simulated by injecting particles into a torus of major radius $`8.9R_\mathrm{d}`$ whose spin axis is inclined by $`15^{}`$ to the original spin axis of the disk. In response to the accretion, the disk develops an integral-sign warp that carries the outermost ring $`0.3R_\mathrm{d}`$ above and below the plane of the innermost ring. Hence the warp is very comparable in magnitude to the warp of the Milky Way. This simulation shows very clearly that the warp is essentially a halo phenomenon: accretion causes the outer halo to tip with respect to the inner halo. The disk largely acts as a tracer of the internal dynamics of the halo.
How much infalling matter is needed to generate a warp? In this simulation, the reorientation of the galactic angular momentum is driven by the crude expedient of adding mass to a tilted annulus of fixed radius. This procedure is well defined and easy to describe, but it exaggerates the mass of infalling material that is required to slew a galaxy’s angular momentum by a given amount because neither the angular-momentum per unit mass of the infalling material, nor the angle between its spin axis and the galaxy’s original spin axis increases with time, as they would in a more realistic situation. Indeed, if the simulation were continued for longer, the galaxy’s angular momentum would everywhere become parallel to the symmetry axis of the torus, and the warp would fade away. In the real world, by contrast, the angular momentum vector of infalling material will be constantly shifting its direction, and, moreover, less and less material will be required to import a given quantity of angular momentum. Hence the simulation’s plausibility hinges on whether the rate of angular-momentum slewing in it is plausible, rather than on the likelihood that a galaxy will accrete as much material at $`10R_\mathrm{d}`$ as is crudely assumed.
Several authors have studied how tidal interactions endow protogalactic regions with angular momentum at early times (e.g. Heavens & Peacock 1988) and these studies are generally in good agreement with numerical simulations (Barnes & Efstathiou 1987). Ryden (1988) and Quinn & Binney (1992) investigated the rate at which the direction of infalling angular momentum should slew by studying the angular momenta of individual spherical shells of protogalactic material. Quinn & Binney found that the angular momenta acquired by shells that differ in radius by a factor 2 have a clear tendency to be antiparallel. Moreover, the angular momentum per unit mass of a shell rises strongly with radius, with the consequence that the net angular momentum of a galaxy tends to be aligned with the angular momentum of the most recently accreted shell. These two results together imply that the net spin axis of a halo tends to slew through more than $`90^{}`$ in the time required for the radius of the currently accreting shell to double. This time depends on the cosmology and the initial density profile (e.g. Fillmore & Goldreich 1984). For critical cosmic density, $`\mathrm{\Omega }=1`$, it is typically comparable to the current Hubble time and the halo’s spin axis is likely to slew by $`\mathrm{}>7^{}`$ in the $`0.9\mathrm{Gyr}`$ that the Jiang & Binney simulation lasted.
## 4. Infalling satellites
It has long been recognized that the Magellanic Clouds are spiralling into the Milky Way. Dynamical models of this process (Murai & Fujimoto 1980; Gardiner et al. 1994; Moore & Davis, 1994; Lin et al. 1995) have been successful in predicting the proper motions of the Clouds, which have now been measured to reasonable accuracy (Jones, Klemola & Lin, 1994; Kroupa & Bastian, 1997). These models are based on the assumption that the Magellanic Stream comprises material that has been tidally stripped out of the Clouds. Hence the success of these models implies that as the Clouds sink deeper into the Galactic halo, a polar ring will form. Since there will be more angular momentum in this ring than there is in the Galactic disk, the orbit of the Clouds constitutes direct evidence for the accretion of mis-aligned angular momentum.
Ibata, Gilmore & Irwin (1994) discovered that the Galaxy has a much nearer satellite than the Clouds, namely the Sagittarius Dwarf galaxy, that is almost hidden from us by the Galactic centre even though it is at a Galactocentric distance of only $`16\mathrm{kpc}`$. Like the Clouds, the Sgr Dwarf is in a nearly polar orbit (Ibata et al. 1997), but the poles of the two orbits make an angle of $`90^{}`$ with one another. It seems that the Dwarf’s orbit has a remarkably short period $`1\mathrm{Gyr}`$ (Ibata & Lewis 1998).
On such a short-period orbit the Dwarf is severely tidally limited by the Galaxy, and there is direct observational evidence that the Dwarf is being tidally shredded: an arc of associated material has now been detected that extends over $`\mathrm{}>60^{}`$ on the sky (Mateo, Olszewski & Morrison, 1999). Several authors have concluded that a Dwarf that contained only the observed luminous material could not have survived tidal shredding for a Hubble time on its present orbit. By enveloping the observed dwarf in a rather homogeneous cloud of dark matter, Lewis & Ibata (1998) were able to construct a model of the Dwarf which retained 46% of its original mass after a Hubble time on its present orbit.
Jiang & Binney (this volume) ask how the Dwarf could have got into its present configuration: galaxies are not likely to form in regions where an external tidal field is strong, because the field would shear away protogalactic material and prevent it accumulating locally. Moreover, the configuration proposed by Lewis & Ibata is finely tuned in the sense that the dark halo has to be very homogeneous and sharp-edged. For these reasons it seems likely that the Dwarf was formed at a considerable distance from the Milky Way, and has subsequently moved onto its present tight orbit by one of two mechanisms. Zhao (1998) suggested that the Dwarf was recently scattered onto its orbit by an encounter with the Clouds. The problem with this proposal is the softness of the potentials of the Dwarf and the Clouds relative to the large velocity of any encounter between them: the required scattering is through a substantial angle. Dynamical friction is the other mechanism that could have moved the Dwarf onto a short-period orbit. The problem here is that the present mass of the Dwarf, even including the Lewis & Ibata dark halo, is too small ($`10^9\mathrm{M}_{}`$) for dynamical friction to be effective.
Jiang & Binney combine $`N`$-body simulations, in which both Dwarf and Galaxy are live systems, with a semi-analytic model, that incorporates dynamical friction and tidal limitation, to explore the rather large parameter space of possible Dwarf histories. They identify a one-parameter family of initial configurations in which the Dwarf starts out at ever greater Galactocentric distances. In the most distant initial configuration explored, the Dwarf starts from $`R=250\mathrm{kpc}`$ as an object of mass $`10^{11}\mathrm{M}_{}`$, largely in the form of a dark halo with central velocity dispersion $`50\mathrm{km}\mathrm{s}^1`$ and roughly constant circular speed. At the other extreme, the Dwarf starts from $`R60\mathrm{kpc}`$ with mass $`1.4\times 10^{10}\mathrm{M}_{}`$. The off-axis angular momentum that the Dwarf brings to the Galaxy, which varies by a factor $`30`$ between extreme configurations, would have a pronounced effect of the Galactic warp at the upper end of the range. Hence, it should be possible to constrain the history of the Dwarf by modelling the combined effect of the Dwarf and the Clouds on the outer disk.
## 5. Conclusions
From the kinematics of M31 we know that the Local Group has yet to virialize. Hence, there is an a-priori expectation that out towards the edges of the Local Group there is a reservoir of material from which M31 and the Galaxy are accreting at a significant rate.
There is a real possibility that this reservoir can be traced through high-velocity clouds (HVCs). This proposition is controversial because the HVCs form a heterogeneous group. Some are clearly associated with M31, the Magellanic Clouds, and other, lower-mass systems such as Phoenix (St-Germain et al, 1999). Others are material in the Galactic halo, but one can identify substantial subset of HVCs that are most naturally interpreted as rather massive objects of order $`1\mathrm{Mpc}`$ from the Galaxy that have yet to fall to the centre of the Local Group. These clouds may together contain $`10^{10}\mathrm{M}_{}`$ of HI and enable the Galaxy to accrete HI at a rate of $`0.8\mathrm{M}_{}\mathrm{yr}^1`$ and dark matter at a rate ten times higher.
Currently we know of no viable mechanism by which warps could survive long in an isolated galaxy, and there is no longer observational support for the proposition that warps exist in isolated galaxies. It seems likely that warps are a manifestation of a galaxy accreting material whose net angular momentum is not parallel to that of the galaxy. To generate a typical warp, off-axis angular momentum must be accreted at such a rate that the spin axis of the inner galaxy slewed by $`10^{}`$ per Gyr.
Observations of the Magellanic Stream and the Sgr Dwarf Galaxy tell us that the Milky Way is certainly accreting off-axis angular momentum. Whether this accretion is fast enough to slew the inner Galaxy’s spin axis by $`10^{}`$ per Gyr depends on how much dark-matter is being accreted alongside the observed luminous matter. If there is 10 times as much dark as luminous matter, it should be possible to explain the Galactic warp, although the very peculiar morphology of the Galactic warp has yet to emerge from a model of Galactic accretion.
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# Exponential and Moment Inequalities for U-statistics
## 1. Introduction
Exponential inequalities, such as Bernstein’s and Prohorov’s, and moment inequalities, such as Rosenthal’s and Hoffmann-Jørgensen’s, are among the most basic tools for the analysis of sums of independent random variables. Our object here consists in developing analogues of such inequalities for generalized $`U`$-statistics, in particular, for $`U`$-statistics and for multilinear forms in independent random variables.
Hoffmann-Jørgensen type moment inequalities for canonical (that is, completely degenerate) $`U`$-statistics of any order $`m`$ were first considered by Giné and Zinn (1992), and their version for $`U`$-statistics with nonnegative kernels turned out to be useful for obtaining best possible necessary integrability conditions in limit theorems for $`U`$-statistics. (By Khinchin’s inequality it is irrelevant whether one considers canonical or nonnegative kernels in moment inequalities, at least if multiplicative constants are not at issue). Klass and Nowicki (1997) also obtained moment inequalities for nonnegative generalized $`U`$-statistics, but only for order $`m=2`$, and their decomposition of the moments is more complete than that in Giné and Zinn (1992). Ibragimov and Sharakhmetov (1998, 1999) recently obtained analogues of Rosenthal’s inequality for nonnegative and for canonical $`U`$-statistics. The moment inequalities we present in the first part of this article, valid for canonical and for nonnegative generalized $`U`$-statistics of any order $`m`$, when specialized to $`m=2`$, represent the same level of moment decomposition as the Klass-Nowicki inequalities, coincide with theirs for powers $`p>1`$ (except for constants) and are expressed in terms of different, simpler quantities for powers $`p<1`$. Proposition 2.1 below, which constitutes the first step towards more elaborate bounds such as those in Theorem 2.3 below, has also been obtained, up to constants, by Ibragimov and Sharakhmetov. Our proofs consist of simple iterations of the classical moment inequalities for sums of independent random variables.
The moment inequalities in the first part of this article do imply exponential bounds for canonical $`U`$-statistics of any order and with bounded kernels which are sharper than those in Arcones and Giné (1993); however, they are not of the best kind as they do not exhibit Gaussian behavior for part of the tail, which they should in view of the tail behavior of Gaussian chaos.
In the second part of this article we improve the moment inequalities from the first part in the case of generalized canonical $`U`$-statistics of order 2, and for moments of order $`p2`$ (Theorem 3.2). The bounds not only involve moments but also the $`L_2`$ operator norm of the matrix of kernels. Then we show how these improved moment inequalities imply what we believe is the correct analogue (up to constants) of Bernstein’s exponential inequality for generalized canonical $`U`$-statistics of order 2 (Theorem 3.3). This exponential inequality, which does exhibit Gaussian behavior for small values of $`t`$, is strong enough to imply the law of the iterated logarithm for canonical $`U`$-statistics under conditions which are also necessary. The main new ingredient in this part of the paper is Talagrand’s (1996) exponential bound for empirical processes, which gives a Rosenthal-Pinelis type inequality for moments of empirical processes (Proposition 3.1) basic for the derivation of the moment inequality for $`U`$-statistics of order 2.
Because of the decoupling results of de la Peña and Montgomery-Smith (1995), we can work with decoupled $`U`$-statistics, and this allows us to proceed by conditioning and iteration.
## 2. Moment inequalities
We consider estimation of moments of generalized decoupled $`U`$-statistics, defined as
(2.1)
$$\underset{1i_1,\mathrm{},i_mn}{}h_{i_1,\mathrm{},i_m}(X_{i_1}^{(1)},\mathrm{},X_{i_m}^{(m)}),$$
where the random variables $`X_i^{(j)}:1in,1jm`$, $`mn`$, are independent (not necessarily with the same distribution) and take values in a measurable space $`(S,𝒮)`$, and $`h_{i_1,\mathrm{},i_m}`$ are real valued measurable functions on $`S^m`$. For short, this sum is denoted by $`_𝐢h_𝐢`$.
Given $`J\{1,\mathrm{},m\}`$ ($`J=\mathrm{}`$ is not excluded), and $`𝐢=(i_1,\mathrm{},i_m)\{1,\mathrm{},n\}^m`$ we set $`𝐢_J`$ to be the point of $`\{1,\mathrm{},n\}^{|J|}`$ obtained from $`𝐢`$ by deleting the coordinates in the places not in $`J`$ (e.g., if $`𝐢=(3,4,2,1)`$ then $`𝐢_{\{1,3\}}=(3,2)`$). Also, $`_{𝐢_J}`$ indicates sum over $`1i_jn`$, $`jJ`$ (for instance, if $`m=4`$ and $`J=\{1,3\}`$, then
$$\underset{𝐢_J}{}h_𝐢=\underset{𝐢_{\{1,3\}}}{}h_{i_1,i_2,i_3,i_4}=\underset{1i_1,i_3n}{}h_{i_1,i_2,i_3,i_4}(X_{i_1}^{(1)},\mathrm{},X_{i_4}^{(4)}).)$$
By convention, $`_𝐢_{\mathrm{}}a=a`$.
Likewise, while $`E`$ will denote expected value with respect to all the variables, $`E_J`$ will denote expected value only with respect to the variables $`X_i^{(j)}`$ with $`jJ`$ and $`i\{1,\mathrm{},n\}`$. By convention, $`E_{\mathrm{}}a=a`$.
Rosenthal’s inequality is easiest to extend to $`U`$-statistics because it involves only moments of sums (as opposed to moments of maxima and quantiles for Hoffmann-Jørgensen’s inequality). So, we will first obtain analogues of Rosenthal’s inequality, and then we will transform these inequalities into analogues of Hoffmann-Jørgensen’s by first showing that some moments of sums can be replaced by moments of maxima, and then, that the lowest moment can in fact be replaced by a quantile. We will illustrate this three-steps procedure first in the case of nonnegative kernels and moments of order $`p1`$. Then we will see that this also solves, via Khinchin’s inequality, the case of canonical kernels and moments of order $`p2`$. Finally, we will consider the case of moments of order $`p<1`$ for positive kernels and $`p<2`$ for canonical, cases in which the inequalities are less neat, but still useful. We will pay some attention to the behavior of the constants as $`p\mathrm{}`$ in these inequalities since such behavior translates into (exponential) integrability properties.
2.1. Nonnegative kernels, moments of order $`p1`$. For nonnegative independent random variables $`\xi _i`$, we have the following two improvements of Rosenthal’s inequalities, valid for $`p1`$:
1) Latała’s, 1997:
$`(R_1)`$
$$E\left(\xi _i\right)^p(2e)^p\mathrm{max}[\frac{e}{p}p^pE\xi _i^p,e^p\left(E\xi _i\right)^p],p>1,$$
(see Pinelis (1994) for the corresponding inequality when the random variables are centered);
2) Johnson, Schechtman and Zinn’s, 1985:
$`(R_2)`$
$$E\left(\xi _i\right)^pK^p\left(\frac{p}{\mathrm{log}p}\right)^p\mathrm{max}[E\xi _i^p,\left(E\xi _i\right)^p],p>1,$$
where $`K`$ is a universal constant. See Utev (1985) and Figiel, Hitczenko, Johnson, Schechtman and Zinn (1997) for more precise inequalities of the same type.
And for general $`p>0`$, we have the following improved Hoffmann-Jørgensen inequality, that follows from Kwapień and Woyczyński (1992) and which can be obtained as in the proof of Theorem 1.2.3 in de la Peña and Giné (1999):
3)
$`(H)`$
$$E\xi _i^p2^{p2}2^{(p1)0}(p+1)^{p+1}\left[t_0^p+E\mathrm{max}\xi _i^p\right],p>0,$$
where
$$t_0:=inf[t>0:\mathrm{Pr}\{\xi _i>t\}\frac{1}{2}],$$
and where we write norm for absolute value in order to include not only independent nonnnegative real random variables, but also independent nonnegative random functions $`\xi _i`$ taking values in certain ‘rearrangement invariant spaces’ such as $`L_s(\mathrm{\Omega },\mathrm{\Sigma },\mu )`$, $`0<s<\mathrm{}`$, with $`\xi :=\left(|\xi |^s𝑑\mu \right)^{1/(s1)}`$, or $`\mathrm{}_{\mathrm{}}(L_s)`$. Note that, by Markov,
$$t_02^{1/r}\left(E\xi _i^r\right)^{1/r},$$
so that, $`(H)`$ becomes:
4) for $`0<r<p<\mathrm{}`$,
$$E\xi _i^p2^{p2}2^{(p1)0}(p+1)^{p+1}[2^{p/r}(E\xi _i^r)^{p/r}$$
$`(H_r)`$
$$\text{aaaaaaaaaaaaaaaaaaaaaaaaaaaaaaaaaaaaaa}+E\mathrm{max}\xi _i^p]$$
Inequalities $`(H)`$ and $`(H_r)`$ hold for spaces of functions which are quasinormed measurable linear spaces whose quasinorm $``$ has the property that $`xy`$ whenever $`0xy`$.
In the following proposition we extend inequalities $`(R_1)`$ and $`(R_2)`$ by means of an easy induction.
Proposition 2.1. Let $`m𝐍`$, $`p>1`$, and, for all $`𝐢\{1,\mathrm{},n\}^m`$, let $`h_𝐢`$ be a nonnegative function of $`m`$ variables whose $`p`$-th power is integrable for the law of $`𝐗_𝐢=(X_{i_1}^{(1)},\mathrm{},X_{i_m}^{(m)})`$. Then,
(2.2) $`\underset{J\{1,\mathrm{},m\}}{\mathrm{max}}\left[{\displaystyle \underset{𝐢_J}{}}E_J\left({\displaystyle \underset{𝐢_{J^c}}{}}E_{J^c}h_𝐢\right)^p\right]E\left({\displaystyle \underset{𝐢}{}}h_𝐢\right)^p`$
$`(2e^2)^{mp}{\displaystyle \underset{J\{1,\mathrm{},m\}}{}}\left[p^{|J|p}{\displaystyle \underset{𝐢_J}{}}E_J\left({\displaystyle \underset{𝐢_{J^c}}{}}E_{J^c}h_𝐢\right)^p\right],`$
and also, there exists a universal constant $`K<\mathrm{}`$ such that
$`(2.2^{})`$
$$E\left(\underset{𝐢}{}h_𝐢\right)^pK^{mp}\left(\frac{p}{\mathrm{log}p}\right)^{mp}\underset{J\{1,\mathrm{},m\}}{\mathrm{max}}\left[\underset{𝐢_J}{}E_J\left(\underset{𝐢_{J^c}}{}E_{J^c}h_𝐢\right)^p\right].$$
Proof. The proof of (2.2’) with sum over the subsets $`J`$ instead of maximum differs from that of (2.2) only in the starting point ($`(R_2)`$ instead of $`(R_1)`$); then, replacing sum by maximum simply increases the constant by a factor of $`2^m`$. The left side inequality in (2.2) follows by Hölder since $`p1`$. Consider the right hand side inequality. For $`m=1`$ this is just inequality ($`R_1`$) and we can proceed by induction. Suppose the result holds for $`m1`$. By applying the induction hypothesis to
$$E\left(\underset{𝐢}{}h_𝐢\right)^p=E_mE_{\{1,\mathrm{},m1\}}\left[\underset{𝐢_{\{1,\mathrm{},m1\}}}{}\left(\underset{i_m}{}h_𝐢\right)\right]^p,$$
we only have to consider the generic term in the decomposition (2.2) for the new kernels $`\left(_{i_m}h_𝐢\right)`$ with the $`X_i^{(m)}`$ variables fixed. In other words, letting $`J_{m1}`$ be any subset of $`\{1,\mathrm{},m1\}`$ and $`J_{m1}^c`$ its complement with respect to $`\{1,\mathrm{},m1\}`$, we must estimate
$`E_m{\displaystyle \underset{𝐢_{J_{m1}}}{}}E_{J_{m1}}\left({\displaystyle \underset{𝐢_{J_{m1}^c}}{}}E_{J_{m1}^c}\left({\displaystyle \underset{i_m}{}}h_𝐢\right)\right)^p`$
$`={\displaystyle \underset{𝐢_{J_{m1}}}{}}E_{J_{m1}}E_m\left({\displaystyle \underset{i_m}{}}\left(E_{J_{m1}^c}{\displaystyle \underset{𝐢_{J_{m1}^c}}{}}h_𝐢\right)\right)^p.`$
Rosenthal’s inequality ($`R_1`$) applied to the kernels $`E_{J_{m1}^c}_{𝐢_{J_{m1}^c}}h_𝐢`$ with the variables in $`J_{m1}`$ fixed, gives
$`E_m\left({\displaystyle \underset{i_m}{}}\left(E_{J_{m1}^c}{\displaystyle \underset{i_{J_{m1}^c}}{}}h_𝐢\right)\right)^p`$ $``$ $`(2e^2)^p[\left({\displaystyle \underset{i_m,i_{J_{m1}^c}}{}}E_mE_{J_{m1}^c}h_𝐢\right)^p`$
$`\text{aa}+p^p{\displaystyle \underset{i_m}{}}E_m\left(E_{J_{m1}^c}{\displaystyle \underset{𝐢_{J_{m1}^c}}{}}h_𝐢\right)^p].`$
Upon integrating each term with respect to $`E_{J_{m1}}`$ and summing over $`𝐢_{J_{m1}}`$, we then obtain
$`E_m{\displaystyle \underset{𝐢_{J_{m1}}}{}}E_{J_{m1}}\left({\displaystyle \underset{𝐢_{J_{m1}^c}}{}}E_{J_{m1}^c}\left({\displaystyle \underset{i_m}{}}h_𝐢\right)\right)^p`$
$`(2e^2)^p[{\displaystyle \underset{𝐢_{J_{m1}}}{}}E_{J_{m1}}\left({\displaystyle \underset{𝐢_{J_{m1}^c\{m\}}}{}}E_{J_{m1}^c\{m\}}h_𝐢\right)^p`$
$`\text{aaa}+p^p{\displaystyle \underset{𝐢_{J_{m1}\{m\}}}{}}E_{J_{m1}\{m\}}\left(E_{J_{m1}^c}{\displaystyle \underset{𝐢_{J_{m1}^c}}{}}h_𝐢\right)^p].`$
Multiplying by $`(2e^2)^{(m1)p}p^{|J_{m1}|}`$, this is the sum of two terms of the form $`(2e)^{mp}p^{|J|p}_{𝐢_J}E_J\left(_{𝐢_{J^c}}E_{J^c}h_𝐢\right)^p`$ (for $`J=J_{m1}`$ and for $`J=J_{m1}\{m\}`$), proving the proposition.
This proposition solves the problem of estimating, up to constants, the moments of a decoupled $`U`$-statistic by ‘computable’ expressions. For instance, if the functions $`h_𝐢`$ are all equal and if the variables $`X_i^{(j)}`$ are i.i.d., then the typical term at the right of (2.1) just becomes $`n^{|J|+p|J^c|}E_J(E_{J^c}h)^p`$, a ‘mixed moment’ of $`h`$. For $`m=2`$ the right hand side of inequality (2.2) is just:
$`E\left({\displaystyle \underset{i,j}{}}h_{i,j}(X_i^{(1)},X_j^{(2)})\right)^p`$ $``$ $`(2e^2)^{2p}[\left({\displaystyle \underset{i,j}{}}Eh_{i,j}(X_i^{(1)},X_j^{(2)})\right)^p`$
$`+p^p{\displaystyle \underset{i}{}}E_1\left({\displaystyle \underset{j}{}}E_2h_{i,j}(X_i^{(1)},X_j^{(2)})\right)^p`$
$`+p^p{\displaystyle \underset{j}{}}E_2\left({\displaystyle \underset{i}{}}E_1h_{i,j}(X_i^{(1)},X_j^{(2)})\right)^p`$
$`(2.2^{\prime \prime })\text{aaaaaaaaaaaaaaaaaaaaaaaaa}+p^{2p}{\displaystyle \underset{i,j}{}}Eh_{i,j}^p(X_i^{(1)},X_j^{(2)})].`$
We have been careful with the dependence on $`p`$ of the constants because it is of some interest to obtain constants of the best order as $`p\mathrm{}`$. In fact, (2.2’) exhibits constants of the best order as can be seen by taking the product of two independent copies of the example in Johnson, Schechtman and Zinn (1985), Proposition 2.9.
Next we replace the external sums of expected values at the right side of the above inequalities by expectations of maxima without significantly altering the order of the multiplicative constants. If $`\xi _i`$ are independent nonnegative random variables, then,
(2.3)
$$\frac{1}{2}\left[\delta _0^pE\xi _i^pI_{\xi _i>\delta _0}\right]E\mathrm{max}\xi _i^p\delta _0^p+E\xi _i^pI_{\xi _i>\delta _0},0<p<\mathrm{},$$
where
(2.4)
$$\delta _0=inf[t>0:\mathrm{Pr}\{\xi _i>t\}1]$$
(Giné and Zinn (1983); see also de la Peña and Giné (1999), page 22). The left hand side of (2.3) gives that, for $`0<r<p`$ and $`\xi _i`$ independent,
(2.5)
$$E|\xi _i|^p2E\mathrm{max}|\xi _i|^p+2\left(E|\xi _i|^r\right)\left(E\mathrm{max}|\xi _i|^p\right)^{(pr)/p}$$
(e.g., de la Peña and Giné (1999), page 48). This inequality, applied with $`r=1<p`$, yields
(2.6)
$$p^{\alpha p}E|\xi _i|^p2(1+p^\alpha )\mathrm{max}[p^{\alpha p}E\mathrm{max}|\xi _i|^p,\left(E|\xi _i|\right)^p]$$
for all $`\alpha 0`$. There are similar inequalities for other values of $`r`$; $`r=1`$ is adequate for $`\xi _i0`$, but $`r=2`$ is better for centered variables. If we use inequality (2.6) in (2.2”), iteratively for the last term, we obtain that, for a universal constant $`K`$ (easy but cumbersome to compute), $`h_{i,j}0`$, $`p>1`$,
(2.7) $`E\left({\displaystyle \underset{i,j}{}}h_{i,j}\right)^p`$ $``$ $`K^p(2e^2)^pp^4[\left({\displaystyle \underset{i,j}{}}Eh_{i,j}\right)^p+p^pE_1\underset{i}{\mathrm{max}}\left({\displaystyle \underset{j}{}}E_2h_{i,j}\right)^p`$
$`\text{aaaaaaaaa}+p^pE_2\underset{j}{\mathrm{max}}\left({\displaystyle \underset{i}{}}E_1h_{i,j}\right)^p+p^{2p}E\underset{i,j}{\mathrm{max}}h_{i,j}^p].`$
Inequality (2.7) was obtained, up to constants, by Klass and Nowicki (1997) (it is their inequality (4.14)). Our proof is different, and it is contained in the proof of the next corollary, which extends inequality (2.7) to any $`m`$.
Corollary 2.2. Under the same hypotheses as in Proposition 2.1, there exist universal constants $`K_m`$ such that
(2.8) $`\underset{J\{1,\mathrm{},m\}}{\mathrm{max}}\left[E_J\underset{𝐢_J}{\mathrm{max}}\left({\displaystyle \underset{𝐢_{J^c}}{}}E_{J^c}h_𝐢\right)^p\right]E\left({\displaystyle \underset{𝐢}{}}h_𝐢\right)^p`$
$`K_m^p{\displaystyle \underset{J\{1,\mathrm{},m\}}{}}\left[p^{|J|p}E_J\underset{𝐢_J}{\mathrm{max}}\left({\displaystyle \underset{𝐢_{J^c}}{}}E_{J^c}h_𝐢\right)^p\right],`$
and
$`(2.8^{})`$
$$E\left(\underset{𝐢}{}h_𝐢\right)^pK_m^p\left(\frac{p}{\mathrm{log}p}\right)^{mp}\underset{J\{1,\mathrm{},m\}}{\mathrm{max}}\left[E_J\underset{𝐢_J}{\mathrm{max}}\left(\underset{𝐢_{J^c}}{}E_{J^c}h_𝐢\right)^p\right].$$
Proof. The left side of (2.8) follows by Hölder. Inequality (2.8’) has a proof similar to that of the right hand side of (2.8), and therefore we only prove the latter. We will prove it by induction over $`m`$ simultaneously with the inequality
(2.9)
$$p^{mp}\underset{𝐢}{}Eh_𝐢^p\stackrel{~}{K}_m^p\underset{J\{1,\mathrm{},m\}}{}\left[p^{|J|p}E_J\underset{𝐢_J}{\mathrm{max}}\left(\underset{𝐢_{J^c}}{}E_{J^c}h_𝐢\right)^p\right].$$
Let us first note that the inequalities (2.9) for $`1,\mathrm{},m1`$ together with (2.2) imply (2.8). It is therefore enough to show that if (2.8) and (2.9) hold for $`1,\mathrm{},m1`$ then (2.9) is satisfied for $`m`$. We will follow the notation of the proof of Proposition 2.1. Inequality (2.9) for $`m=1`$ is just (2.6), and (2.8) for $`m=1`$ is $`(H_1)`$ (which also follows from $`(R_1)`$ and (2.6)). By the induction assumptions we have
(2.10)
$$p^{mp}\underset{𝐢}{}Eh_𝐢^p=p^p\underset{i_m}{}E_mp^{(m1)p}\underset{𝐢_{\{1,\mathrm{},m1\}}}{}E_{\{1,\mathrm{},m1\}}h_𝐢^p$$
$`\stackrel{~}{K}_{m1}^p`$
$$\times \underset{J_{m1}\{1,\mathrm{},m1\}}{}\left[p^{(|J_{m1}|+1)p}E_{J_{m1}}E_m\underset{i_m}{}\underset{𝐢_{J_{m1}}}{\mathrm{max}}\left(\underset{𝐢_{J_{m1}^c}}{}E_{J_{m1}^c}h_𝐢\right)^p\right].$$
Now, by (2.6), for any $`J_{m1}\{1,\mathrm{},m1\}`$ we have
(2.11) $`p^{(|J_{m1}|+1)p}E_{J_{m1}}E_m{\displaystyle \underset{i_m}{}}\underset{𝐢_{J_{m1}}}{\mathrm{max}}\left({\displaystyle \underset{𝐢_{J_{m1}^c}}{}}E_{J_{m1}^c}h_𝐢\right)^p`$
$`2(1+p)[p^{(|J_{m1}|+1)p}E_{J_{m1}\{m\}}\underset{𝐢_{J_{m1}\{m\}}}{\mathrm{max}}\left({\displaystyle \underset{𝐢_{J_{m1}^c}}{}}E_{J_{m1}^c}h_𝐢\right)^p`$
$`\text{aaaaaaaaaa}+p^{|J_{m1}|p}E_{J_{m1}}\left({\displaystyle \underset{i_m}{}}E_m\underset{𝐢_{J_{m1}}}{\mathrm{max}}{\displaystyle \underset{𝐢_{J_{m1}^c}}{}}E_{J_{m1}^c}h_𝐢\right)^p].`$
To estimate the last term we note that
(2.12) $`p^{|J_{m1}|p}E_{J_{m1}}\left({\displaystyle \underset{i_m}{}}E_m\underset{𝐢_{J_{m1}}}{\mathrm{max}}{\displaystyle \underset{𝐢_{J_{m1}^c}}{}}E_{J_{m1}^c}h_𝐢\right)^p`$
$`p^{|J_{m1}|p}E_{J_{m1}}\left({\displaystyle \underset{𝐢}{}}E_{J_{m1}^c\{m\}}h_𝐢\right)^p\text{aa}`$
$$\stackrel{~}{K}_{|J_{m1}|}^p\underset{JJ_{m1}}{}p^{|J|p}E_J\underset{𝐢_J}{\mathrm{max}}\left(\underset{𝐢_{(J_{m1}J)J_{m1}^c\{m\}}}{}E_{(J_{m1}J)J_{m1}^c\{m\}}h_𝐢\right)^p,$$
where in the last line we use the induction assumption (2.8) for $`|J_{m1}|<m`$. Finally (2.10), (2.11) and (2.12) imply (2.9) and complete the proof.
Remark. The proof of Proposition 2.6 below will use a version of Corollary 2.2 for nonnegative random functions taking values in $`L_r`$. The inequality is as follows: for $`p>1`$ there exists $`K_{m,p,r}<\mathrm{}`$ such that
$`(2.8^{\prime \prime })`$
$$E\underset{𝐢}{}h_𝐢^pK_{m,p,r}\underset{J\{1,\mathrm{},m\}}{\mathrm{max}}\left[E_J\underset{𝐢_J}{\mathrm{max}}(E_{J^c}\underset{𝐢_{J^c}}{}h_𝐢)^p\right].$$
The proof is similar to the previous ones and is omited: one takes $`(H_p)`$ as the starting point of the induction.
Finally we come to the third step, which will extend Hoffmann-Jørgen-sen’s inequality $`(H)`$ for $`p1`$. If we want to use the inequalities from Corollary 2.2 to obtain boundedness of moments from stochastic boundedness of a sequence of $`U`$-statistics, we need to replace the term corresponding to $`J=\mathrm{}`$ by the $`p`$-th power of a quantile of $`_𝐢h_𝐢`$. For this we use Paley-Zygmund’s inequality (e.g., Kahane (1968) or de la Peña and Giné (1999)): if $`A`$ is a nonnegative random variable and $`0<r<p<\mathrm{}`$, then, for all $`0<\lambda <1`$,
(2.13)
$$\mathrm{Pr}\left\{A>\lambda A_r\right\}\left[(1\lambda ^r)\frac{A_r}{A_p}\right]^{p/(pr)},$$
where $`A_r=\left(E|A|^r\right)^{1/r}`$ for $`0<r<\mathrm{}`$. Consider for instance inequality (2.8). It has the form
$$EA^pB+K_m^p(EA)^p,p>1,$$
with $`A=_𝐢h_𝐢`$. Then, either $`BK_m^p(EA)^p`$, in which case we have $`EA^p2B`$, or $`B<K_m^p(EA)^p`$, in which case we have $`EA^p2K_m^p(EA)^p`$ and we can apply Paley-Zygmund’s (2.13) with $`\lambda =1/2`$ and $`r=1`$. It gives
$$\mathrm{Pr}\left\{A>\frac{1}{2}EA\right\}\frac{1}{2^{(p+1)/(p1)}K_m^{p/(p1)}}.$$
Hence, if we define
(2.14)
$$t_0=inf[t0:\mathrm{Pr}\{A>t\}\frac{1}{2^{(p+1)/(p1)}K_m^{p/(p1)}}],$$
we obtain $`EA2t_0`$. So, in either case,
$$EA^p2B+2^{1+p}K_m^pt_0^p.$$
Also, by Markov’s inequality,
$$\frac{1}{2^{(p+1)/(p1)}K_m^{p/(p1)}}t_0^pEA^p.$$
We then have:
Theorem 2.3. Under the hypotheses of Proposition 2.1, there exist a universal constants $`K_m<\mathrm{}`$ such that, if $`t_0`$ is as defined by (2.14) for $`A=_𝐢h_𝐢`$, then
(2.20) $`{\displaystyle \frac{1}{(4K_m)^{p/(p1)}}}t_0^p\mathrm{max}[\underset{\begin{array}{c}J\{1,\mathrm{},m\}\\ J\mathrm{}\end{array}}{\mathrm{max}}\left[E_J\underset{𝐢_J}{\mathrm{max}}\left({\displaystyle \underset{𝐢_{J^c}}{}}E_{J^c}h_𝐢\right)^p\right]`$
$`E\left({\displaystyle \underset{𝐢}{}}h_𝐢\right)^p`$
$`(4K_m)^p\{2^{1+p}t_0^p+{\displaystyle \underset{\begin{array}{c}J\{1,\mathrm{},m\}\\ J\mathrm{}\end{array}}{}}\left[p^{|J|p}E_J\underset{𝐢_J}{\mathrm{max}}\left({\displaystyle \underset{𝐢_{J^c}}{}}E_{J^c}h_𝐢\right)^p\right]\}.`$
A similar inequality with different constants can be obtained from (2.8’). This is the most elaborate form we will give to our bounds for $`h0`$ and $`p>1`$.
The right hand side of (2.15) for $`m=2`$ becomes, disregarding constants,
$$E\left(\underset{i,j}{}h_{i,j}\right)^pC\mathrm{max}[E_1\underset{i}{\mathrm{max}}\left(\underset{j}{}E_2h_{i,j}\right)^p,E_2\underset{j}{\mathrm{max}}\left(\underset{i}{}E_1h_{i,j}\right)^p,$$
$`(2.15^{})`$
$$\text{aaaaaaaaaaaaaaaaaaaaaaaaaaaaaaaaaaaaaaaaaaa}E\underset{i,j}{\mathrm{max}}h_{i,j}^p,t_0^p].$$
So, we get the $`p`$-th moment of the double sum controlled by moments of partial maxima of conditional expectations plus a quantile. The Giné-Zinn (1992) inequality (for $`m=2`$),
$$E\left(\underset{i,j}{}h_{i,j}\right)^pC\mathrm{max}[E\underset{i}{\mathrm{max}}\left(\underset{j}{}h_{i,j}\right)^p,t_0^p],p1,$$
is slightly weaker in appearance than (2.15’) (actually, we only published the result for canonical $`U`$-statistics, but we applied it as well to nonnegative variables, for which the proof is the same: see, e.g., Giné and Zhang (1996)). For applications of this inequality in the asymptotic theory of $`U`$-statistics see Giné and Zhang (1996), Giné, Kwapień, Latała and Zinn (1999) and de la Peña and Giné (1999).
Remark. The constants in the definition of $`t_0`$ in (2.15) depend on $`p`$, hence, so does $`t_0`$. This is not the case when $`m=1`$ (as a consequence of the improved Hoffmann-Jørgensen’s inequality of Kwapień and Woyczyński -see, de la Peña and Giné (1999) p. 11-). But in most applications it does not matter whether the definition of the quantile depends on $`p`$.
2.2. Canonical kernels, moments or order $`p2`$. If $`\xi _i`$ are centered and independent and $`p2`$, then, by convexity and the Khinchin-Bonami inequality (e.g., de la Peña and Giné, 1999, p. 113), we have
(2.21) $`2^pE\left({\displaystyle \xi _i^2}\right)^{p/2}`$ $``$ $`2^pE\left|{\displaystyle \epsilon _i\xi _i}\right|^pE\left|{\displaystyle \xi _i}\right|^p`$
$``$ $`2^pE\left|{\displaystyle \epsilon _i\xi _i}\right|^p2^p(p1)^{p/2}E\left({\displaystyle \xi _i^2}\right)^{p/2},`$
where $`\epsilon _i`$ are independent identically distributed Rademacher random variables, independent from $`\{\xi _i\}`$. Suppose $`h_𝐢`$ is canonical for the variables $`\{X_i^{(j)}\}`$ given in the previous subsection, that is, suppose
(2.22)
$$E_jh(X_{i_1}^{(1)},\mathrm{},X_{i_m}^{(m)})=0\mathrm{a}.\mathrm{s}.\mathrm{for}\mathrm{all}j=1,\mathrm{},m,1i_1,\mathrm{},i_mn.$$
Let $`\epsilon _i^{(j)}`$ be an independent Rademacher array independent of $`\{X_i^{(j)}\}`$, and set
$$\epsilon _𝐢:=\epsilon _{i_1}^{(1)}\mathrm{}\epsilon _{i_m}^{(m)}.$$
Then, recursive application of inequality (2.16) gives
(2.23) $`2^{mp}E\left({\displaystyle \underset{𝐢}{}}h_𝐢^2\right)^{p/2}`$ $``$ $`2^{mp}E\left|{\displaystyle \underset{𝐢}{}}\epsilon _𝐢h_𝐢\right|^pE\left|{\displaystyle \underset{𝐢}{}}h_𝐢\right|^p`$
$``$ $`2^{mp}E\left|{\displaystyle \underset{𝐢}{}}\epsilon _𝐢h_𝐢\right|^p2^{mp}(p1)^{mp/2}E\left({\displaystyle \underset{𝐢}{}}h_𝐢^2\right)^{p/2}.`$
This inequality reduces estimation of moments of canonical $`U`$-statistics to estimation of moments of nonnegative ones (and conversely), at least if constants are not an issue. Combined with Proposition 2.1, it gives the analogue of Rosenthal’s inequality for centered variables and $`p>2`$, and if we apply it in conjunction with Corollary 2.2, we obtain the following inequality:
Proposition 2.4. If, for $`p>2`$ and all $`𝐢\{1,\mathrm{},n\}^m`$, $`h_𝐢(X_{i_1}^{(1)},\mathrm{},X_{i_m}^{(m)})`$ is $`p`$-integrable and $`E_jh_𝐢(X_{i_1}^{(1)},\mathrm{},X_{i_m}^{(m)})=0`$ a.s. for all $`j=1,\mathrm{},m`$, then
(2.24) $`2^{mp}\underset{J\{1,\mathrm{},m\}}{\mathrm{max}}\left[E_J\underset{𝐢_J}{\mathrm{max}}\left({\displaystyle \underset{𝐢_{J^c}}{}}E_{J^c}h_𝐢^2\right)^{p/2}\right]E|{\displaystyle \underset{𝐢}{}}h_𝐢|^p`$
$`K_m^p{\displaystyle \underset{J\{1,\mathrm{},m\}}{}}\left[p^{(m+|J|)p/2}E_J\underset{𝐢_J}{\mathrm{max}}\left({\displaystyle \underset{𝐢_{J^c}}{}}E_{J^c}h_𝐢^2\right)^{p/2}\right]`$
for universal constant $`K_m<\mathrm{}`$.
And, applying Paley-Zygmund with $`r=2`$, we finally have:
Theorem 2.5. Let $`h_𝐢`$ be as in Proposition 2.4, and let $`p>2`$. Then, there exist universal constants $`K_m<\mathrm{}`$ such that, if $`t_0`$ is defined as
$$t_0=inf[t0:\mathrm{Pr}\left\{\right|\underset{𝐢}{}h_𝐢|>t\}\left(\frac{3}{4}\right)^{p/(p2)}\frac{1}{\left(2K_m^pp^{mp/2}\right)^{1/(p2)}}],$$
then
$$\frac{1}{(4K_mp^{m/2})^{p/(p2)}}t_0^p\mathrm{max}[2^{mp}\underset{\begin{array}{c}J\{1,\mathrm{},m\}\\ J\mathrm{}\end{array}}{\mathrm{max}}\left[E_J\underset{𝐢_J}{\mathrm{max}}\left(\underset{𝐢_{J^c}}{}E_{J^c}h_𝐢^2\right)^{p/2}\right]$$
(2.25)
$$E\left|\underset{𝐢}{}h_𝐢\right|^p\text{aaaaaaaaaaaaaaaaaaaaaaaaaaaaaaaa}$$
$$2K_m^p\{(2p^{m/2})^pt_0^p+\underset{\begin{array}{c}J\{1,\mathrm{},m\}\\ J\mathrm{}\end{array}}{}\left[p^{(m+|J|)p/2}E_J\underset{𝐢_J}{\mathrm{max}}\left(\underset{𝐢_{J^c}}{}E_{J^c}h_𝐢^2\right)^{p/2}\right]\}.$$
If, instead of inequality (2.2), we wish to obtain an analogue of inequality (2.2’), that is, if we want to replace the constants at the right hand side of (2.19) by $`(Kp/\mathrm{log}p)^{mp}`$, then we cannot use Khinchin’s inequality and must proceed directly with an induction as in Proposition 2.1 with the following change: we must consider the variables $`_{𝐢_{J_{m1}^c}}h_𝐢`$ as taking values in $`L_2(J_{m1}^c)`$ and apply inequality (1.5) in Kwapień and Szulga (1991), which gives Rosenthal’s inequality with best constants for centered independent random variables in Banach spaces. We skip the details.
2.3. Nonnegative kernels, moments of order $`p1`$. It seems impossible to obtain inequalities as simple as in the previous section for this case. However, one can still obtain inequalities that may become useful when combined with Paley-Zygmund. Here is an analogue of Corollary 2.2 for $`h0`$ and $`p1`$. The method of proof is inefficient regarding constants as Hoffmann-Jørgensen is applied twice at each step. Hence, constants will not be specified.
Proposition 2.6. Let $`0<r<p1`$, $`m<\mathrm{}`$ and assume that the kernels $`h_𝐢0`$ have integrable $`p`$-th powers. Then
(2.26) $`\underset{J\{1,\mathrm{},m\}}{\mathrm{max}}[E_J\underset{𝐢_J}{\mathrm{max}}(E_{J^c}\left({\displaystyle \underset{𝐢_{J^c}}{}}h_𝐢\right)^r)^{p/r}]E\left({\displaystyle \underset{𝐢}{}}h_𝐢\right)^p`$
$`K_{r,p,m}\underset{J\{1,\mathrm{},m\}}{\mathrm{max}}[E_J\underset{𝐢_J}{\mathrm{max}}(E_{J^c}\left({\displaystyle \underset{𝐢_{J^c}}{}}h_𝐢\right)^r)^{p/r}],`$
where $`K_{r,p,m}`$ depends only on the parameters $`r,p,m`$.
Note that all the terms in this bound represent a reduction in the number of sums except for the term corresponding to $`J=\mathrm{}`$, which consists of a power of the $`r`$-th moment of a $`U`$-statistic of order $`m`$. We will deal later with this term by means of the Paley-Zygmund argument.
Proof. The inequality at the left side of (2.21) follows from Hölder. Inequality $`(H_r)`$ is just the right hand side of inequality (2.21) for $`m=1`$ and we can proceed by induction. We still use the notation from Proposition 2.1. By the induction hypothesis we have
$`\text{aaaaaaa}E\left({\displaystyle \underset{𝐢}{}}h_𝐢\right)^p=E_mE_{\{1,\mathrm{},m1\}}\left({\displaystyle \underset{𝐢_{\{1,\mathrm{},m1\}}}{}}{\displaystyle \underset{i_m}{}}h_𝐢\right)^p`$
$`K_{r,p,m1}{\displaystyle \underset{J_{m1}\{1,\mathrm{},m1\}}{}}E_{J_{m1}}E_m\underset{𝐢_{J_{m1}}}{\mathrm{max}}\left[E_{J_{m1}^c}\left({\displaystyle \underset{𝐢_{J_{m1}^c}}{}}{\displaystyle \underset{i_m}{}}h_𝐢\right)^r\right]^{p/r}.`$
Let us fix $`J_{m1}\{1,\mathrm{},m1\}`$ and note that, for fixed $`(X_j^{(i)})_{iJ_{m1}}`$, we have
$$\underset{𝐢_{J_{m1}}}{\mathrm{max}}E_{J_{m1}^c}\left(\underset{𝐢_{J_{m1}^c}}{}\underset{i_m}{}h_𝐢\right)^r:=\underset{i_m}{}\stackrel{~}{h}_{i_m}$$
for suitably chosen independent r.v.’s $`\stackrel{~}{h}_{i_m}`$ in $`l^{\mathrm{}}(L^r)`$. Therefore by ($`H_r`$), which still holds in this space (as the norm, restricted to nonnegative vectors, is monotone increasing), we have
(2.28) $`E_{J_{m1}}E_m\underset{𝐢_{J_{m1}}}{\mathrm{max}}\left[E_{J_{m1}^c}\left({\displaystyle \underset{𝐢_{J_{m1}^c}}{}}{\displaystyle \underset{i_m}{}}h_𝐢\right)^r\right]^{p/r}=E_{J_{m1}}E_m{\displaystyle \underset{i_m}{}}\stackrel{~}{h}_{i_m}^{p/r}`$
$`C_{p,r}E_{J_{m1}}\left[E_m\underset{i_m}{\mathrm{max}}\stackrel{~}{h}_{i_m}^{p/r}+\left(E_m{\displaystyle \underset{i_m}{}}\stackrel{~}{h}_{i_m}\right)^{p/r}\right]`$
$`=C_{p,r}[E_{J_{m1}\{m\}}\underset{𝐢_{J_{m1}\{m\}}}{\mathrm{max}}(E_{J_{m1}^c}({\displaystyle \underset{𝐢_{J_{m1}^c}}{}}h_𝐢)^r)^{p/r}`$
$`\text{aaaaaaaa}+E_{J_{m1}}\left(E_m\underset{𝐢_{J_{m1}}}{\mathrm{max}}E_{J_{m1}^c}\left({\displaystyle \underset{𝐢_{J_{m1}^c}}{}}{\displaystyle \underset{i_m}{}}h_𝐢\right)^r\right)^{p/r}].`$
Now, to estimate the last term, we note that
(2.29) $`E_{J_{m1}}\left[E_m\underset{𝐢_{J_{m1}}}{\mathrm{max}}E_{J_{m1}^c}\left({\displaystyle \underset{𝐢_{J_{m1}^c}}{}}{\displaystyle \underset{i_m}{}}h_𝐢\right)^r\right]^{p/r}`$
$`\text{aaaaaaa}E_{J_{m1}}\left[E_{J_{m1}^c\{m\}}\left({\displaystyle \underset{𝐢}{}}h_𝐢\right)^r\right]^{p/r}`$
$`K_{p/r,1,|J_{m1}|}`$
$$\times \underset{JJ_{m1}}{}E_J\underset{𝐢_J}{\mathrm{max}}\left[E_{J_{m1}JJ_{m1}^c\{m\}}\left(\underset{𝐢_{J_{m1}JJ_{m1}^c\{m\}}}{}h_𝐢\right)^r\right]^{p/r},$$
which follows by the version of Corollary 2.2 for $`L^r`$ ((2.8”) for $`p/r>1`$). Now (2), (2.28) and (2.29) complete the induction step.
To deal with the term corresponding to $`J=\mathrm{}`$ in Proposition 2.6 we apply Paley-Zygmund as above, but now with $`r<p`$ replacing $`1<p`$. The conclusion is:
Theorem 2.7. There is a constant $`K_{r,p,m}`$ such that for $`0<r<p1`$, $`m<\mathrm{}`$, and $`h_𝐢0`$ with integrable $`p`$-th powers, we have
(2.36) $`{\displaystyle \frac{1}{(2^{p+1}K_{r,p,m})^{1/(pr)}}}t_0^p{\displaystyle \underset{\begin{array}{c}J\{1,\mathrm{},m\}\\ J\mathrm{}\end{array}}{}}[E_J\underset{𝐢_J}{\mathrm{max}}(E_{J^c}\left({\displaystyle \underset{𝐢_{J^c}}{}}h_𝐢\right)^r)^{p/r}]`$
$`E\left({\displaystyle \underset{𝐢}{}}h_𝐢\right)^p`$
$`2K_{r,p,m}\{2^{p/r}t_0^p+{\displaystyle \underset{\begin{array}{c}J\{1,\mathrm{},m\}\\ J\mathrm{}\end{array}}{}}[E_J\underset{𝐢_J}{\mathrm{max}}\left(E_{J^c}\left({\displaystyle \underset{𝐢_{J^c}}{}}h_𝐢\right)^r)^{p/r}]\right\},`$
where
$$t_0=inf[t:\mathrm{Pr}\{\underset{𝐢}{}h_𝐢>t\}\frac{1}{2}(2^{p+1}K_{r,p,m})^{1/(pr)}].$$
Hence, the $`p`$-th moment of a $`U`$-statistic of order $`m`$ can be estimated by partial moments of maxima (or sums) of conditional moments of $`U`$-statistics of lower order plus the $`p`$-power of a quantile of the original $`U`$-statistic.
2.4. Canonical kernels and moments of order $`1p2`$, or kernels $`h`$ separately symmetric in each of the coordinates and $`0<p<1`$ . The canonical case reduces to the positive case by means of inequality (2.18), as before. The convexity part of inequality (2.18) fails for $`p<1`$, but in this case, if $`h`$ is symmetric separately in each of the coordinates, we can still randomize by products of independent Rademacher variables and recursive application of Khinchin’s inequality still reduces this case to nonnegative $`h`$. We leave the resulting statements to the reader in order to avoid repetition.
2.5. Regular (undecoupled) general $`U`$-statistics. If $`h_𝐢(𝐱)=h_{𝐢s}(𝐱s)`$ for any permutation $`s`$ of $`\{1,\mathrm{},m\}`$ and $`h_𝐢=0`$ if $`𝐢`$ has repeated indices, and if the sequences $`\{X_i^{(j)}:i=1,\mathrm{},n\}`$ are independent copies of each other, then the decoupling inequalities of de la Peña and Montgomery-Smith (1995), together with the decoupling inequality for maxima in Hitczenko (1988) in combination with the previous inequalities give moment inequalities for the generalized $`U`$-statistics
$$\underset{𝐢}{}h_{i_1,\mathrm{},i_m}(X_{i_1},\mathrm{},X_{i_m})$$
where $`\{X_i\}`$ is a sequence of independent random variables, at the cost of vastly increasing the numerical constants (see e.g. Giné and Zinn (1992) for a similar application of the decoupling inequalities). We omit the resulting statements.
2.6. Comparison with previous results. We have already noted, below the statement of Theorem 2.3, that the inequalities there are better than the Hoffmann-Jørgensen type inequalities for $`U`$-statistics in Giné and Zinn (1992) in that they represent a decomposition into simpler quantities. Also, as mentioned in the Introduction, Ibrahimov and Sharakhmetov (1998, 1999) obtained, except for constants, Proposition 2.1 and its analogue for canonical kernels for $`m=2`$ and announced the result for general $`m`$; the final results in the present article for $`p>1`$ in the nonnegative case (Theorem 2.3) and for $`p>2`$ in the canonical case (Theorem 2.5), replacing some sums by maxima and lower moments by quantiles, seem to be more useful. As mentioned above, Corollary 2.2 restricted to $`m=2`$ recovers inequalities (4.14) in Klass and Nowicki (1997). The inequalities in the last mentioned article for nonnegative kernels, $`p<1`$ and $`m=2`$ (the nonconvex case, inequalities (4.13) there) are different from our inequalities in Theorem 2.7 for $`m=2`$, although they represent a similar level of decomposition of the $`p`$-th moment of the $`U`$-statistic. Basically, the difference is that they use inverses of truncated conditional moments whereas we use inverses of tail probabilites together with partial moments. This can be better seen by comparing Hoffmann-Jørgensen, which is Theorem 2.7 for $`m=1`$, with their inequality for $`m=1`$. The result of Klass and Nowicki (1997) can be described as the iteration of an inequality that follows from Hoffmann-Jørgensen, Paley-Zygmund ((2.13)) and (2.3), as follows. Given $`\xi _i`$, $`i=1,\mathrm{},n`$, nonnegative, define $`v_0`$ as
(2.37)
$$v_0=sup\{v0:E\left(\frac{\xi _i}{v}1\right)1\}$$
or, what is the same, $`v_0`$ is the largest number satisfying
(2.38)
$$v_0=E\left(\xi _iv_0\right).$$
Then, the inequality in question is:
Corollary 2.8. (Klass and Nowicki, 1997, Cor. 2.7) Let $`\xi _i,`$ $`1=1,\mathrm{},n,`$ be independent nonnegative random variables. Then, for all $`p>0`$,
(2.39)
$$E\left(\xi _i\right)^pE\mathrm{max}\xi _i^p+v_0^p.$$
Proof. Since
$$E\left(\xi _i\delta _0\right)=E\xi _iI_{\xi _i<\delta _0}+\delta _0\mathrm{Pr}\{\xi _i\delta _0\}\delta _0,$$
it follows that $`\delta _0v_0`$. Therefore, if $`p1`$, inequality (2.3) and the definition of $`v_0`$ give
$$E\left(\xi _i\right)^p\left(E(\xi _iv_0)\right)^p+\underset{j}{}E\xi _i^pI_{\xi _i>v_0}v_0^p+2E\underset{i}{\mathrm{max}}\xi _i^p.$$
And if $`p>1`$, Hoffmann-Jørgensen ($`(H)`$) and the previous inequality (with $`p=1`$) give
$$E\left(\xi _i\right)^p\stackrel{<}{}\left(E\xi _i\right)^p+E\mathrm{max}\xi _i^p\stackrel{<}{}v_0^p+E\mathrm{max}\xi _i^p.$$
For the reverse inequality, if $`p>1`$,
$$v_0^p=\left(E(\xi _iv_0)\right)^pE\left(\xi _i\right)^p.$$
And if $`p<1`$, following the proof of Lemma 2.2 in Klass and Nowicki (1997), we first observe that Paley-Zygmund and the first part of this proof give that for some universal constant $`C`$,
$`\mathrm{Pr}\left\{{\displaystyle \xi _i}v_0>{\displaystyle \frac{v_0}{2}}\right\}`$ $``$ $`{\displaystyle \frac{1}{4}}{\displaystyle \frac{\left(E(\xi _iv_0)\right)^2}{E\left((\xi _iv_0)\right)^2}}={\displaystyle \frac{1}{4}}{\displaystyle \frac{v_0^2}{E\left((\xi _iv_0)\right)^2}}`$
$``$ $`{\displaystyle \frac{C}{4}}{\displaystyle \frac{v_0^2}{E\mathrm{max}(\xi _iv_0)^2+v_0^2}}{\displaystyle \frac{C}{8}};`$
therefore,
$`E\left({\displaystyle \xi _i}\right)^p`$ $``$ $`E\left({\displaystyle (\xi _iv_0)}\right)^p`$
$``$ $`E\left[\left({\displaystyle (\xi _iv_0)}\right)^pI_{{\scriptscriptstyle (\xi _iv_0)}>v_0/2}\right]{\displaystyle \frac{C}{8}}{\displaystyle \frac{v_0^p}{2^p}}.`$
In fact, if we bound $`t_0`$ by $`t_0^p2E\left((\xi _it_0)\right)^p`$ and apply the above proof to the variables $`\xi _it_0`$, Hoffmann-Jørgensen gives the following seemingly weaker inequality: letting $`\stackrel{~}{v}_0`$ be the parameter $`v_0`$ for the smaller variables $`\xi _it_0`$ (note $`\stackrel{~}{v}_0v_0`$), then
$`(2.22^{})`$
$$E\left(\xi _i\right)^pE\mathrm{max}\xi _i^p+\stackrel{~}{v}_0^p.$$
## 3. Improved moment inequalities and exponential inequalities for $`m=2`$
The right hand side of inequality (2.19) for $`m=1`$ is just
(3.1)
$$E\left|\xi _i\right|^pK^p\mathrm{max}[p^pE\mathrm{max}\xi _i^p,p^{p/2}\left(E\xi _i^2\right)^{p/2}],p2,$$
where $`\xi _i`$ are independent mean zero random variables. These inequalities were first obtained by Pinelis (1994). Part of their interest lie on the fact that they are basically equivalent to Bernstein’s inequality up to constants. Here is how (3.1) (for all $`p2`$) implies Bernstein’s inequality up to constants. Assume $`\xi _i_{\mathrm{}}A<\mathrm{}`$ for all $`i`$, and set $`C^2=E\xi _i^2`$. Then, (3.1) has the form
$$E\left|\xi _i\right|^pK^p\mathrm{max}[p^pA^p,p^{p/2}C^p],p2.$$
Let
$$p=\frac{x}{KeA}\left(\frac{x}{KeC}\right)^2$$
for any $`x`$ for which $`p2`$. Then, by Markov’s inequality, (3.1) gives, for these values of $`t`$,
$$\mathrm{Pr}\left\{\left|\xi _i\right|>x\right\}\{\begin{array}{cc}\frac{K^pp^pA^p}{x^p}e^p\hfill & \text{if }p^pA^pp^{p/2}C^p\hfill \\ & \\ \frac{K^pp^{p/2}C^p}{x^p}e^p\hfill & \text{otherwise.}\hfill \end{array}$$
Hence,
(3.2)
$$\mathrm{Pr}\left\{\left|\xi _i\right|>x\right\}e^2e^p=e^2\mathrm{exp}\left\{\frac{x}{KeA}\left(\frac{x}{KeC}\right)^2\right\}$$
for all $`x>0`$. Similarly, from the iteration (2.19) of the inequalities (3.1) we can obtain exponential inequalities for generalized decoupled $`U`$-statistics of any order. However, the inequalities we obtain, while better than the existing ones, are not of the best kind, as we will see below. We illustrate this comment by considering the case $`m=2`$. In this case, inequality (2.19) is as follows:
(3.3) $`E\left|{\displaystyle \underset{i,j}{}}h_{i,j}\right|^p`$ $``$ $`K^p\mathrm{max}[p^p\left({\displaystyle \underset{i,j}{}}Eh_{i,j}^2\right)^{p/2},p^{3p/2}E_1\underset{i}{\mathrm{max}}\left({\displaystyle \underset{j}{}}E_2h_{i,j}^2\right)^{p/2},`$
$`\text{aaaaaaaa}p^{3p/2}E_2\underset{j}{\mathrm{max}}\left({\displaystyle \underset{i}{}}E_1h_{i,j}^2\right)^{p/2},p^{2p}E\underset{i,j}{\mathrm{max}}|h_{i,j}|^p].`$
For bounded canonical kernels $`h_{i,j}`$ we define
(3.4)
$$A=\underset{i,j}{\mathrm{max}}h_{i,j}_{\mathrm{}},C^2=\underset{i,j}{}Eh_{i,j}^2,$$
$$B^2=\mathrm{max}[\underset{i}{}E_1h_{i,j}^2(X_i^{(1)},y)_{\mathrm{}},\underset{j}{}E_2h_{i,j}^2(x,X_j^{(2)})_{\mathrm{}}].$$
Then, we can proceed as in the deduction of (3.2) from (3.1), and easily obtain from (3.3) that there is a universal constant $`K`$ such that
(3.5)
$$\mathrm{Pr}\left\{\left|\underset{i,j}{}h_{i,j}\right|>x\right\}K\mathrm{exp}\left\{\frac{1}{K}\mathrm{min}[\frac{x}{C},\left(\frac{x}{B}\right)^{2/3},\left(\frac{x}{A}\right)^{1/2}]\right\}.$$
This inequality also holds for regular canonical $`U`$-statistics by the decoupling inequalities of de la Peña and Montgomery-Smith (1995).
Inequality (3.5) is better than the Bernstein type inequality in Arcones and Giné (1993) as it is better for $`xn^2A`$ and the probability is zero for $`xn^2A`$. Inequality (3.5) is suboptimal for small values of $`x`$, for which the exponent should be a constant times $`x^2`$, just as for chaos variables of order 2 (see Ledoux and Talagrand (1991) and Latała (1999)). This suggest that inequality (2.9) is not of the best kind, and can be improved.
Next we improve the Rosenthal type inequality (2.9) for $`m=2`$ (that is, (3.3)) and deduce from it an exponential inequality for canonical $`U`$-statistics of order two which does detect the Gaussian portion of the tail probability.
First we show how Talagrand’s (1996) extension of Prohorov’s inequality to empirical processes, actually in Massart’s (1999) version, produces an improved Rosenthal’s inequality for empirical processes. Then, we will use this inequality to estimate the terms resulting from conditionally applying inequality (3.1) to the $`U`$-statistic.
To describe Massart’s version of Talagrand’s inequality we must establish the setting and define some parameters. Let $`Z_i`$ be independent random variables with values in some measurable space $`(T,𝒯)`$, let $``$ be a countable class of measurable real functions on $`T`$, and define
$$S:=\underset{f}{sup}f(Z_i),\sigma ^2=\underset{f}{sup}E(f(Z_i))^2,a:=\underset{i}{\mathrm{max}}\underset{f}{sup}f(Z_i)_{\mathrm{}}.$$
Then,
(3.6)
$$\mathrm{Pr}\left\{|S|2E|S|+\sigma \sqrt{8x}+34.5ax\right\}e^x$$
for all $`x>0`$. It follows easily from inequality (3.6) that
(3.7)
$$E|S|^pK^p\left[(E|S|)^p+p^{p/2}\sigma ^p+p^pa^p\right]$$
for some universal constant $`K<\mathrm{}`$ and all $`p1`$, in fact, inequality (3.7) for all $`p`$ large enough and inequality (3.6) for all $`x>0`$ are equivalent up to constants. (We do not plan to keep track of constants in the derivation below and, therefore, we refrain from specifying a value for $`K`$ in (3.7).)
Proposition 3.1. Let $`\{Z_i\}`$ be as above, let $``$ be a countable class of functions such that $`Ef^2(Z_i)<\mathrm{}`$ and $`Ef(Z_i)=0`$ for all $`i`$. Then, in the notation from the previous paragraph,
(3.8)
$$E|S|^pK^p\left[(E|S|)^p+p^{p/2}\sigma ^p+p^pE\underset{i}{\mathrm{max}}\underset{f}{sup}\left|f(Z_i)\right|^p\right]$$
for all $`p1`$, where $`K`$ is a universal constant.
Proof. Set $`F:=sup_f|f|`$ and $`M^p:=83^pE\mathrm{max}_i|F(Z_i)|^p`$. Since the variables $`f(Z_i)`$ are centered, we can randomize by independent Rademacher variables $`\epsilon _i`$ independent of the $`Z`$ variables (at the price of increasing the value of the constant $`K`$). Set $`\stackrel{~}{S}:=sup_f\left|\epsilon _if(Z_i)\right|`$. Then,
$$|\stackrel{~}{S}|\underset{f}{sup}\left|\epsilon _if(Z_i)I_{F(Z_i)M}\right|+\underset{f}{sup}\left|\epsilon _if(Z_i)I_{F(Z_i)>M}\right|:=S_1+S_2,$$
and notice that, since $`ES_1^p2^{p+1}E|S|^p`$ (e.g., Lemmas 1.2.6 and 1.4.3 in de la Peña and Giné, 1999), inequality (3.7) gives
$$ES_1^pK^p\left[(E|S|)^p+p^{p/2}\sigma ^p+p^pM^p\right].$$
To estimate $`ES_2^p`$ we apply the original Hoffmann-Jørgensen inequality (from e.g., Ledoux and Talagrand (1991), (6.9) in page 156) to get
$$ES_2^p23^p\left(t_0^p+E\underset{i}{\mathrm{max}}F(Z_i)^p\right),$$
where $`t_0`$ is any number such that $`\mathrm{Pr}\{S_2>t_0\}(83^p)^1.`$ But the choice of $`M`$ implies that we can take $`t_0=0`$ because
$$\mathrm{Pr}\left\{S_2>0\right\}=\mathrm{Pr}\left\{\underset{i}{\mathrm{max}}F(Z_i)>M\right\}\frac{1}{83^p},$$
proving the proposition.
In what follows we will assume, just as above, that the kernels $`h_{i,j}`$, $`i,jn`$, are completely degenerate and define
(3.9)
$$D=(h_{i,j})_{L^2L^2}:=sup\{E\underset{i,j}{}h_{i,j}(X_i^{(1)},X_j^{(2)})f_i(X_i^{(1)})g_j(X_j^{(2)})$$
$$:E\underset{i}{}f_i^2(X_i^{(1)})1,E\underset{j}{}g_j^2(X_j^{(2)})1\}.$$
Theorem 3.2. There exists a universal constant $`K<\mathrm{}`$ such that, if $`h_{i,j}`$ are bounded canonical kernels of two variables for the independent random variables $`X_i^{(1)},X_j^{(2)}`$, $`i.j=1,\mathrm{},n`$, $`n𝐍`$, then
$`E|{\displaystyle \underset{1i,jn}{}}`$ $`h_{i,j}(X_i^{(1)},X_j^{(2)})|^pK^p[p^{p/2}\left({\displaystyle \underset{i,j}{}}Eh_{i,j}^2\right)^{p/2}+p^p(h_{i,j})_{L^2L^2}`$
$`+p^{3p/2}[E_1\underset{i}{\mathrm{max}}\left({\displaystyle \underset{j}{}}E_2h_{i,j}^2\right)^{p/2}+E_2\underset{j}{\mathrm{max}}\left({\displaystyle \underset{i}{}}E_1h_{i,j}^2\right)^{p/2}]`$
$`+p^{2p}E\underset{i,j}{\mathrm{max}}|h_{i,j}|^p]`$
for all $`p2`$.
Inequality (3.10) is strictly better than the right hand side inequality in (2.9) for $`m=2`$, that is, than (3.3).
Proof. Inequality (3.1) applied conditionally on the variables $`X_i^{(1)}`$ gives
(3.11)
$$E\left|\underset{i,j}{}h_{i,j}\right|^pK^pE_1\left(p^{p/2}\left[\underset{j}{}E_2\left(\underset{i}{}h_{i,j}\right)^2\right]^{p/2}+p^pE_2\underset{j}{}\left|\underset{i}{}h_{i,j}\right|^p\right).$$
To bound the first summand at the right hand side of (3.11) we first notice that
$`\left[{\displaystyle \underset{j}{}}E_2\left({\displaystyle \underset{i}{}}h_{i,j}\right)^2\right]^{1/2}`$
$`=sup[{\displaystyle \underset{i}{}}E_2{\displaystyle \underset{j}{}}h_{i,j}(X_i^{(1)},X_j^{(2)})f_j(X_j^{(2)}):E{\displaystyle \underset{j}{}}f_j^2(X_j^{(2)})1],`$
where in fact, the sup is taken only over a countable subset of mean zero vector functions $`(f_1,\mathrm{},f_n)`$ dense in the unit ball of $`L_2((X_1^{(2)}))\times \mathrm{}\times L_2((X_n^{(2)}))`$ for the seminorm $`|(f_j)_{jn}|=\left(Ef_j^2(X_j^{(2)})\right)^{1/2}`$. \[To see this, first apply duality in $`\mathrm{}_2^n`$ and then in $`L_2((X_j^{(2)}))`$ for each $`j`$.\] So we can apply (3.8) to $`Z_i=(h_{i,j})_{j=1}^n`$ with $`f(Z_i)=E_2_jh_{i,j}(X_i^{(1)},X_j^{(2)})f_j(X_j^{(2)})`$. In this case, the right hand side terms in (3.8) can be estimated as follows. The first term:
$$(E|S|)^2E|S|^2=E\left[\underset{j}{}E_2\left(\underset{i}{}h_{i,j}\right)^2\right]=E\underset{i,j}{}h_{i,j}^2=C^2.$$
For the second we see that, since, by the previous duality argument,
$$\underset{i}{}E_1\left(E_2\underset{j}{}h_{i,j}(X_i^{(1)},X_j^{(2)})f_j(X_j^{(2)})\right)^2(h_{i,j})_{L^2L^2}^2=D^2,$$
it follows that $`\sigma D`$. The third term:
$`E\underset{i}{\mathrm{max}}\underset{f}{sup}|f(Z_i)|^p`$ $`=`$ $`E_1\underset{i}{\mathrm{max}}\underset{E{\scriptscriptstyle f_j^2}1}{sup}\left[E_2{\displaystyle \underset{j}{}}h_{i,j}(X_i^{(1)},X_j^{(2)})f_j(X_j^{(2)})\right]^p`$
$``$ $`E_1\underset{i}{\mathrm{max}}\underset{E{\scriptscriptstyle f_j^2}1}{sup}\left[\left(E_2{\displaystyle \underset{j}{}}h_{i,j}^2\right)^{1/2}\left(E{\displaystyle \underset{j}{}}f_j^2\right)^{1/2}\right]^p`$
$`=`$ $`E_1\underset{i}{\mathrm{max}}\left(E_2{\displaystyle \underset{j}{}}h_{i,j}^2\right)^{p/2}.`$
Thus, inequality (3.8) gives
$`p^{p/2}`$ $`E\left({\displaystyle \underset{j}{}}E_2\left({\displaystyle \underset{i}{}}h_{i,j}\right)^2\right)^{p/2}`$
$`K^p[p^{p/2}C^p+p^pD^p+p^{3p/2}E_1\underset{i}{\mathrm{max}}\left(E_2{\displaystyle \underset{j}{}}h_{i,j}^2\right)^{p/2}].`$
To estimate the second summand at the right hand side of (3.11), we apply (3.1) once more and obtain
(3.13) $`p^pE_2{\displaystyle \underset{j}{}}E_1\left|{\displaystyle \underset{i}{}}h_{i,j}\right|^p`$
$`K^p\left[p^{3p/2}E_2{\displaystyle \underset{j}{}}\left({\displaystyle \underset{i}{}}E_1h_{i,j}^2\right)^{p/2}+p^{2p}E{\displaystyle \underset{i,j}{}}\left|h_{i,j}\right|^p\right].`$
Thus, to complete the proof of the theorem it suffices to replace the sum in $`j`$ and the sum in $`i,j`$ respectively by maxima in $`j`$ and in $`i,j`$ on the terms at the right hand side of this inequality. But this is an easy exercise of application of inequality (2.6). For completeness sake, here it is. Applying (2.6) with $`\alpha =3`$ and $`p/2`$ instead of $`p`$, the first term at the right of (3.13) bounds as:
$`p^{3p/2}E_2{\displaystyle \underset{j}{}}\left({\displaystyle \underset{i}{}}E_1h_{i,j}^2\right)^{p/2}`$
$`2^{1+3p/2}(1+(p/2)^3)[\left({\displaystyle \frac{p}{2}}\right)^{3p/2}E_2\underset{j}{\mathrm{max}}\left({\displaystyle \underset{i}{}}E_1h_{i,j}^2\right)^{p/2}+C^p],`$
which produces the conversion of the sum into a maximum without increasing the order of the multiplicative constant in front of $`C^p`$. The second term in (3.13) requires two steps. First, we apply (2.6) for $`p/2`$ and $`\alpha =4`$, conditionally on $`\{X_i^{(1)}\}`$:
(3.14) $`p^{2p}E{\displaystyle \underset{i,j}{}}\left|h_{i,j}\right|^p`$
$`2^{2p+1}(1+(p/2)^4)E_1{\displaystyle \underset{i}{}}\left[\left({\displaystyle \frac{p}{2}}\right)^{2p}E_2\underset{j}{\mathrm{max}}|h_{i,j}|^p+\left({\displaystyle \underset{j}{}}E_2h_{i,j}^2\right)^{p/2}\right].`$
We apply (2.6) with respect to $`E_1`$, for $`p/2`$ and $`\alpha =0`$, to the second term at the right hand side of (3.14) and we obtain the bound
$$2^{2p+3}(1+(p/2)^4)[E_1\underset{i}{\mathrm{max}}\left(\underset{j}{}E_2h_{i,j}^2\right)^{p/2}+C^p],$$
which is in terms of some of the quantities appearing at the right hand side of (3.10) and with coefficients of lower order. As for the first term at the right of (3.14), we apply (2.6) with respect to $`E_1`$, again for $`p/2`$ and $`\alpha =4`$, and get it bounded by
$$2^{4p+2}(1+(p/2)^4)^2\left[\left(\frac{p}{2}\right)^{2p}E\underset{i,j}{\mathrm{max}}|h_{i,j}|^p+E_2\left(\underset{i}{}E_1\underset{j}{\mathrm{max}}h_{i,j}^2\right)^{p/2}\right].$$
Here the first term coincides with the last one in (3.10), and the second is dominated by
$$K^pE_2\left[\underset{j}{}\left(\underset{i}{}E_1h_{i,j}^2\right)\right]^{p/2}.$$
Applying inequality $`(R_1)`$ with respect to $`E_2`$ this is in turn dominated by
$$K^p\left(\frac{p}{2}\right)^{p/2}E_2\underset{j}{}\left(\underset{i}{}E_1h_{i,j}^2\right)^{p/2}+K^pC^p,$$
and the first summand has alredy been handled above (first term at the right of (3.13)). Collecting terms we obtain inequality (3.10).
Theorem 3.2 gives the following moment inequality and exponential bound for bounded kernels.
Theorem 3.3. There exist universal constants $`K<\mathrm{}`$ and $`L<\mathrm{}`$ such that, if $`h_{i,j}`$ are bounded canonical kernels of two variables for the independent random variables $`X_i^{(1)},X_j^{(2)}`$, $`i.j=1,\mathrm{},n`$, and if $`A`$, $`B`$, $`C`$, $`D`$ are as defined in (3.4) and (3.9), then
(3.15)
$$E\left|\underset{1i,jn}{}h_{i,j}(X_i^{(1)},X_j^{(2)})\right|^pK^p\left[p^{p/2}C^p+p^pD^p+p^{3p/2}B^p+p^{2p}A^p\right]$$
for all $`p2`$ and, equivalently,
(3.16) $`\mathrm{Pr}\left\{\right|{\displaystyle \underset{i,jn}{}}h_{i,j}(X_i^{(1)},X_j^{(2)})|`$ $``$ $`x\}`$
$``$ $`L\mathrm{exp}\left[{\displaystyle \frac{1}{L}}\mathrm{min}({\displaystyle \frac{x^2}{C^2}},{\displaystyle \frac{x}{D}},{\displaystyle \frac{x^{2/3}}{B^{2/3}}},{\displaystyle \frac{x^{1/2}}{A^{1/2}}})\right]`$
for all $`x>0`$.
The moment inequality is immediate from Theorem 3.2 and the equivalence with the exponential inequality follows just like (3.2) follows from (3.1) in one direction, and, in the other, by integration of tail probabilities.
Next we comment on the exponential inequality. For comparison purposes, let $`h_{i,j}(X_i^{(1)},X_j^{(2)})=g_ig_j^{}x_{i,j}`$ with $`g_i,g_j^{}`$ independent standard normal. In this case,
$$C^2=\underset{i,j}{}x_{i,j}^2\mathrm{and}D=sup\{\underset{i,j}{}u_iv_jx_{i,j}:u_i^21,v_j^21\}$$
and the Gaussian chaos inequality in Latała (1999) yields the existence of universal constants $`0<k<K<\mathrm{}`$ such that
$$\mathrm{Pr}\left\{\left|\underset{i,j}{}h_{i,j}\right|K(Cx^{1/2}+Dx)\right\}e^x$$
and
$$\mathrm{Pr}\left\{\left|\underset{i,j}{}h_{i,j}\right|k(Cx^{1/2}+Dx)\right\}ke^x.$$
By the central limit theorem for canonical $`U`$-statistics, this implies that the coefficients of $`x^2`$ and $`x`$ in (3.16) are correct (except for $`K`$). It is natural to have terms in smaller powers of $`x`$ in (3.16) e.g., by comparison with Bernstein’s inequality for sums of independent random variables. In fact, the term in $`x^{1/2}`$ cannot be avoided, at least up to logarithmic factors. To see this, consider the product $`V`$ of two independent centered Poisson variables with parameter 1, which is the limit in law of $`V_n=_{i,jn}X_i^{(n)}Y_j^{(n)}`$ where $`X_i^{(n)}`$ and $`Y_j^{(n)}`$ are centered Bernoulli random variables with parameter $`p=1/n`$; then, for large $`x`$, the tail probabilities of $`V`$ are of the order of $`\mathrm{exp}(x^{1/2}\mathrm{log}x)`$, and therefore, so are those of $`V_n`$ for large $`n`$. Also, note that the term in $`x^{2/3}`$ in the exponent corresponds, up to logarithmic factors, to the tail probabilities of the product of two independent random variables, one normal and the other centered Poisson.
If $`X,Y,X_i^{(1)},X_j^{(2)}`$ are i.i.d., $`h_{i,j}=h`$ for all $`i,j`$ and $`h`$ is completely degenerate, then the parameters defined by (3.4) and (3.8) become:
$$A=h_{\mathrm{}},B^2=n\left(E_Yh^2(x,Y)_{\mathrm{}}+E_Xh^2(X,y)_{\mathrm{}}\right),C^2=n^2Eh^2$$
and
$`D`$ $`=`$ $`nsup\{Eh(X,Y)f(X)g(Y):Ef^2(X)1,Eg^2(Y)1\}`$
$`:=`$ $`nh_{L_2L_2},`$
where $`h_{L_2L_2}`$ is the norm of the operator of $`L_2((X))`$ with kernel $`h`$. Then, inequalities (3.15) and (3.16) become:
Corollary 3.4. Under the above assumptions, there exist universal constants $`K<\mathrm{}`$, $`L<\mathrm{}`$ such that, for all $`n𝐍`$ and $`p2`$,
$`E|{\displaystyle \underset{i,jn}{}}h(X_i^{(1)},X_j^{(2)})|^pK^p[p^{p/2}n^p(Eh^2)^{p/2}+p^pn^ph_{L_2L_2}^p`$
(3.17) $`+p^{3p/2}n^{p/2}(E_Yh^2_{\mathrm{}}+E_Xh^2_{\mathrm{}})^{p/2}+p^{2p}h_{\mathrm{}}^p]`$
and
$`\mathrm{Pr}\left\{\right|{\displaystyle \underset{i,jn}{}}h(X_i^{(1)},X_j^{(2)})|x\}K\mathrm{exp}[{\displaystyle \frac{1}{K}}\mathrm{min}({\displaystyle \frac{x^2}{n^2Eh^2}},`$
(3.18) $`\text{aaaaaaaa}{\displaystyle \frac{x}{nh_{L_2L_2}}},{\displaystyle \frac{x^{2/3}}{\left[n(E_Yh^2_{\mathrm{}}+E_Xh^2_{\mathrm{}})\right]^{1/3}}},{\displaystyle \frac{x^{1/2}}{h_{\mathrm{}}^{1/2}}})].`$
Inequality (3.18) provides an analogue of Bernstein’s inequality for degenerate $`U`$-statistics of order 2: note that inequalities (3.15), (3.16), (3.17) and (3.18) can all be ‘undecoupled’ using the result of de la Peña and Montgomery-Smith’s (1995). It should also be noted that this exponential inequality for canonical $`U`$-statistics is strong enough to imply the sufficiency part of the law of the iterated logarithm for these objects: this can be seen by applying it to the kernels $`h_n`$ in Steps 7 and 8 of the proof of Theorem 3.1 in Giné, Kwapień, Latała and Zinn (1999) (and using some of the computations there for the parameters $`C`$ to $`D`$). Neither inequality (3.5) nor any of the previously published inequalities for $`U`$-statistics can do this.
Acknowledgement. We thank Stanislaw Kwapień for several useful conversations.
References
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de la Peña, V. and Giné, E. (1999). Decoupling: ¿From Dependence to Independence. Springer-Verlag, New York.
de la Peña, V. and Montgomery–Smith, S. (1995). Decoupling inequalities for the tail probabilities of multivariate $`U`$-statistics. Ann. Probab. 23 806-816.
Figiel, T.; Hitczenko, P.; Johnson, W.B.; Schechtman, G.; and Zinn, J. (1997). Extremal properties of Rademacher functions with applications to the Khintchine and Rosenthal inequalities. Trans. Amer. Math. Soc. 349 997-1027.
Giné, E.; Kwapień, S.; Latała, R.; and Zinn, J. (1999). The LIL for canonical $`U`$-statistics of order two. To appear.
Giné, E. and Zhang, C.-H. (1996). On integrability in the LIL for degenerate $`U`$-statistics. J. Theoret. Probab 9 385-412.
Giné, E. and Zinn, J. (1983). Central limit theorems and weak laws of large numbers in certain Banach spaces. Zeits. Wahrsch. v. Geb. 62 323-354.
Giné, E. and Zinn, J. (1992). On Hoffmann-Jørgensen’s inequality for $`U`$-processes. Probability in Banach Spaces 8 80-91. Birkhäuser, Boston.
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Ibragimov, R. and Sharakhmetov, Sh. (1999). Analogues of Khintchine, Marcinkiewicz-Zygmund and Rosenthal inequalities for symmetric statistics. Scand. J. Statist. 26 621-623.
Kahane, J.-P. (1968). Some Random Series of Functions. Heath, Lexington, Massachusetts.
Klass, M. and Nowicki, K. (1997). Order of magnitude bounds for expectations of $`\mathrm{\Delta }_2`$ functions of nonnegative random bilinear forms and generalized $`U`$-statistics. Ann. Probab. 25 1471-1501.
Kwapień, S. and Szulga, J. (1991). Hypercontraction methods in moment inequalities for series of independent random variables in normed spaces. Ann. Probab. 19 369-379.
Kwapień, S. and Woyczyński, W. (1992). Random Series and Stochastic Integrals: Single and Multiple. Birkhäuser, Boston.
Latała, R. (1997). Estimation of moments of sums of independent random variables. Ann. Probab. 25 1502-1513.
Latała, R. (1999). Tails and moment estimates for some type of chaos. Studia Math. 135 39-53.
Latała, R. and Zinn, J. (1999). Necessary and sufficient conditions for the strong law of large numbers for $`U`$-statistics. Ann. Probab., to appear.
Ledoux, M. and Talagrand, M. (1991). Probability in Banach Spaces: Isoperimetry and Processes. Springer, New York.
Massart, P. (1999). About the constants in Talagrand’s concentration inequalities for empirical processes. Ann. Probab., to appear.
Pinelis, I. (1994). Optimum bounds for the distributions of martingales in Banach spaces. Ann. Probab. 22 1679-1706.
Talagrand, M. (1996). New concentration inequalities in product spaces. Invent. Math. 126 505-563.
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| Evarist Gin$`\stackrel{´}{\text{e}}`$ | Rafał Latała |
| --- | --- |
| Department of Mathematics | Institute of Mathematics |
| and Department of Statistics | Warsaw University |
| University of Connecticut | Banacha 2 |
| Storrs, CT 06269 | 02-097 Warszawa |
| USA | Poland |
| gine@uconnvm.uconn.edu | rlatala@mimuw.edu.pl |
| Joel Zinn |
| --- |
| Department of Mathematics |
| Texas A&M University |
| College Station, TX 77843 |
| jzinn@math.tamu.edu |
|
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# BNL-NT-00/5LBNL-45248 Wong’s equations and the small 𝑥 effective action in QCD
## 1 Introduction
One of the more interesting open questions in QCD is the behavior of cross–sections at very high energies . In the last decade, a kinematic window has opened up at colliders where $`Q^2\mathrm{\Lambda }_{QCD}^2`$ but $`x=Q^2/s1`$. The physics in this regime is non–perturbative because the field strengths at small $`x`$ are large. However, it is also weak coupling physics since $`\alpha _S(Q^2)1`$. Further, since the density of partons is large at small $`x`$, classical field methods are applicable .
An effective field theory approach can be used to study the physics of small $`x`$ modes in QCD <sup>1</sup><sup>1</sup>1For alternative approaches, see for example, Ref. and references therein, and Ref. and references therein. These will not be discussed in this paper. . The small $`x`$ effective action is obtained by successively integrating out the more static modes at larger values of $`x`$. The measure for this action is represented by a weight corresponding to the color charge density of the higher $`x`$ modes. As one integrates out higher $`x`$ modes, the form of the action is maintained, while the weight satisfies a Wilsonian non–linear renormalization group (RG) equation . If the parton density is not too large, the RG equation can be linearized, and the resulting equation is the well known BFKL equation. The BFKL equation is a renormalization group equation that sums the leading logarithms in $`\alpha _S\mathrm{ln}(1/x)`$ . In the double log limit of small $`x`$ and large $`Q^2`$, the Wilson RG can be simplified, and one obtains a series in inverse powers of $`Q^2`$, where the leading term is the small $`x`$ DGLAP equation and the first sub–leading term agrees with the expression derived by Gribov, Levin and Ryskin , and by Mueller and Qiu .
The effective action approach therefore reproduces the standard linear evolution equations of perturbative QCD in the limit of low parton densities. The truly interesting and unknown regime however is the non–linear regime of high parton densities where one might hope to predict novel phenomena . What the correct effective action is in the high density regime should therefore be a matter of some interest.
In this paper, we will discuss an alternative gauge invariant form to the gauge invariant action discussed in Ref. . The motivation for this form of the effective action came from our recent work in formulating a many body world line formalism for the one loop effective action in QCD . Briefly, the difference between the two actions is in the term describing the coupling of the small $`x`$ gauge field modes to the large $`x`$ modes represented by a color charge density $`\rho `$. In the work of Jalilian–Marian, Kovner, Leonidov, and Weigert (JKLW), this term is expressed as
$`S_{int}^{JKLW}\mathrm{Tr}\left(\rho W_\mathrm{},\mathrm{}\right),`$
where $`W`$ is an adjoint matrix corresponding to a path ordered exponential of the gauge field $`A^{}`$ in the light cone direction $`x^+`$. We propose instead that this term be
$`S_{int}\mathrm{Tr}\left(\rho \mathrm{ln}W_\mathrm{},\mathrm{}\right),`$
replacing $`W\mathrm{ln}W`$ in the effective action.
The earlier form of $`S_{int}`$ was chosen primarily because it is a gauge invariant generalization of the coupling between hard and soft modes. The motivation for the latter form comes from the background field method and the eikonal approximation. The one loop effective action, in the background field method, can be expressed as $`\mathrm{ln}[\mathrm{det}(D^2)]\mathrm{Tr}\mathrm{ln}[D^2]`$, where $`D`$ is the usual covariant derivative. If, for instance, one integrated out hard fermions in the soft background gauge field, the eigenvalues of the determinant would correspond to solutions of the Dirac equation in the eikonal approximation. These correspond to path ordered phases of the component of the soft gauge field, conjugate to the hard current, in the fundamental representation . Similarly, performing an eikonal separation of hard and soft gauge fields, one obtains path ordered exponentials (in the adjoint representation) of the soft gauge fields (see, for example, appendix B of the first paper in Ref. ). Since the effective action is the logarithm of the determinant, one can thus anticipate the appearance of the logarithm of the path ordered phase in the effective action. This form of the effective action is also gauge invariant. We will show later that the $`\mathrm{ln}(W)`$ action has the nice feature that one can derive the BFKL equation from it efficiently–certain terms that one needs to argue to be zero in the $`W`$ form of the effective action are absent in the $`\mathrm{ln}(W)`$ action.
The subsequent discussion is organized as follows. In section 2, we will discuss the form of the small $`x`$ effective action discussed in Ref. . In section 3, we will discuss Wong’s equations and motivate an alternative form for the small $`x`$ effective action. We will show that the form of the action that we propose is also consistent with Wong’s equations and that the two different currents arising from the two actions correspond to different boundary conditions for solving Wong’s equations. In section 4, we will show that our form of the effective action also reproduces the BFKL equation. We end this paper with a brief summary in section 5. Some technical details are contained in three appendices.
## 2 The Small $`x`$ Effective Action
In this section, we will review the effective action and Wilson renormalization group approach to small $`x`$ QCD as developed in Refs. . We refer the reader to these papers for more details.
We start with the following action which is the gauge invariant generalization of McLerran-Venugopalan effective action first proposed in . In the infinite momentum frame, and in Light Cone gauge $`A^+=0`$, one can write
$`S`$ $`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle d^4xG_a^{\mu \nu }G_{\mu \nu }^a}+i{\displaystyle d^2x_{}F[\rho ^a(x_{})]}`$ (1)
$`+`$ $`{\displaystyle \frac{i}{N_c}}{\displaystyle d^2x_{}𝑑x^{}\delta (x^{})\mathrm{Tr}\left(\rho (x_{})W_\mathrm{},\mathrm{}[A^{}](x^{},x_{})\right)},`$
where $`W`$ is the Wilson line in the adjoint representation along the $`x^+`$ axis
$`W_\mathrm{},\mathrm{}[A^{}](x^{},x_{})=\widehat{P}\mathrm{exp}\left[ig{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑x^+A_a^{}(x^+,x^{},x_{})T_a\right].`$ (2)
and the label $`\widehat{P}`$ denotes the path-ordered exponential.
Taking the saddle point of the effective action, we obtain the Yang–Mills equations
$`D_\mu G_a^{\mu \nu }=\delta ^{\nu +}J_a^+,`$ (3)
with the current
$`J_a^+(x)={\displaystyle \frac{g}{N_c}}\delta (x^{})\mathrm{Tr}\left[T_aW_{x^+,\mathrm{}}[A^{}]\rho (x_{})W_{\mathrm{},x^+}[A^{}]\right]`$ (4)
satisfying the boundary condition
$`J_a^+(x^+=\mathrm{})={\displaystyle \frac{g}{N_c}}\delta (x^{})\mathrm{Tr}\left[T_a\rho (x_{})W_\mathrm{},\mathrm{}[A^{}]\right]`$ (5)
The first term in the expansion of the Wilson line in the action is
$`S_{\mathrm{int}}=g{\displaystyle d^4xA^{}\rho (x_{})\delta (x^{})}`$ (6)
used in .
To derive the general evolution equation, one first solves the classical equations of motion, computes quantum fluctuations in the background of the classical field (semi-classical approximation), and separates these fluctuations according to their longitudinal momentum as
$`A_\mu ^a(x)=b_\mu ^a(x)+\delta A_\mu ^a(x)+a_\mu ^a(x),`$ (7)
where $`b_\mu ^a(x)`$ is the solution of the classical equations of motion, $`\delta A_\mu ^a(x)`$ is the fluctuation field containing longitudinal momentum modes $`k^+`$ that are constrained to be $`p^+<k^+<P^+`$. The upper cutoff $`P^+`$ is the longitudinal momentum of the fast moving charges while the lower cutoff $`p^+`$ is the momentum scale of the soft fluctuations. These cut-offs are chosen to be such that $`\alpha _S\mathrm{ln}(P^+/p^+)1`$ since quantum fluctuations give rise to such logarithms . This constraint thus requires that the fluctuations with momentum modes $`p^+<k^+<P^+`$ are small, and can therefore be integrated out to obtain the effective action for the soft (in longitudinal momenta alone!) fields $`a^\mu `$. This procedure can be iterated as one goes to smaller $`x`$ leading to a Wilsonian RG equation .
The physics underlying this procedure is simple. One starts with some initial color charge density at large $`x`$ represented by $`\rho `$. In order to compute a quantity with this action, one averages over all color configurations represented by the statistical weight
$`Z=\mathrm{exp}\{F[\rho ]\}.`$
We then integrate out the hard fluctuations with the constraint discussed above. This changes the color charge density and the statistical weight for their configurations. The soft fluctuations, with logarithmic accuracy, “see” the induced charge density as a part of the color charge density to which they are coupled. As one goes to smaller and smaller $`x`$ (longer and longer wavelength gluons) one correspondingly includes more of the hard fluctuations in the color charge density. One obtains the following renormalization group equation for the change of the statistical weight $`Z`$ with $`x`$ :
$`{\displaystyle \frac{dZ}{d\mathrm{ln}(1/x)}}=\alpha _S\left[{\displaystyle \frac{1}{2}}{\displaystyle \frac{\delta ^2}{\delta \rho _\mu \delta \rho _\nu }}\left(Z\chi _{\mu \nu }\right){\displaystyle \frac{\delta }{\delta \rho _\mu }}\left(Z\sigma _\mu \right)\right],`$ (8)
where $`\sigma [\rho ]`$ and $`\chi [\rho ]`$ are respectively one and two point functions obtained by integrating over $`\delta A`$ for fixed $`\rho `$. The one point function $`\sigma `$ includes the virtual corrections to $`F[\rho ]`$ while the two point function $`\chi `$ includes the real contributions to $`F[\rho ]`$. Both of these can be computed explicitly from the small fluctuations propagator in the classical background field. In the weak field limit, the functions $`\sigma `$ and $`\chi `$ simplify, and the resulting renormalization group equation is the BFKL equation.
In the following section, we will interpret the color charge density $`\rho `$ of the hard (large $`k^+`$) modes as the density of classical color charges moving in the field of the soft modes. Such an interpretation arises naturally when one computes the one loop effective action in QCD for soft modes using the background field method . One expects therefore that these classical charges must satisfy Wong’s equation for the motion of color charges in a non–Abelian background field. These equations are discussed below where a new form of the effective action is proposed.
## 3 Wong’s equations and an alternative effective action
In Ref. , we developed a many body formalism for the one loop effective action in QCD. We employed the world line formalism to re–write the path ordered exponential as a quantum mechanical path integral over world lines. The equations of motion for the corresponding point particle Lagrangian satisfies Wong’s equations for the motion of a classical charged particle in a non–Abelian background field . These are
$`p^\mu `$ $`=`$ $`m{\displaystyle \frac{dx^\mu }{d\tau }}=mv^\mu `$ (9)
$`{\displaystyle \frac{dp^\mu }{d\tau }}`$ $`=`$ $`v_\nu Q^aG_a^{\mu \nu }`$ (10)
$`D^\nu G_{\nu \mu }`$ $`=`$ $`j_\mu `$ (11)
where
$`j_\mu (x)={\displaystyle 𝑑\tau Q(\tau )v_\mu (\tau )\delta ^4\left[xz(\tau )\right]}.`$ (12)
and
$`\dot{Q}=ig[Q,v_\mu A^\mu ]`$ (13)
The generalization to a system of particles is straightforward.
Without explicitly going over to the world line approach, one can write down the following many-body classical action
$`S_{\mathrm{Wong}}={\displaystyle \frac{1}{4}}{\displaystyle d^4xG_{\mu \nu }^aG_a^{\mu \nu }}{\displaystyle \underset{I=1}{\overset{K}{}}}{\displaystyle 𝑑\tau m_0^I\sqrt{v_\mu ^Iv_I^\mu }}+{\displaystyle \frac{i}{N_c}}{\displaystyle \underset{I=1}{\overset{K}{}}}\mathrm{Tr}\left\{Q_I\mathrm{ln}W_I\right\},`$
(14)
where $`K`$ is the number of Wong’s particles and $`I`$ is the particle label. Also
$`W_I=\widehat{P}\mathrm{exp}\left(ig{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\tau v_\mu ^IA_a^\mu (x_I^\mu (\tau ))T_a^I\right).`$ (15)
This action is gauge invariant under gauge transforms $`U`$ that satisfy the constraint $`U(\mathrm{})=U(\mathrm{})`$. We will define “$`\mathrm{ln}W`$” shortly. As shown in appendix A, the Wong equations in Eq. (11) can be derived from Eq. (14) above.
In an infinite momentum frame (relevant for the small $`x`$ problem), the momenta of the particles are not dynamical. They are static light cone sources–$`v^\mu =\delta ^{\mu +}`$. The kinetic part of the action in $`S_{\mathrm{Wong}}`$ therefore drops out to yield
$`S_{\mathrm{Wong}}={\displaystyle \frac{1}{4}}{\displaystyle d^4xG_{\mu \nu }^aG_a^{\mu \nu }}+{\displaystyle \frac{i}{N_c}}{\displaystyle \underset{I=1}{\overset{K}{}}}\mathrm{Tr}\left\{Q_I\mathrm{ln}W_I\right\}.`$ (16)
We assume now that the initial $`x_I^{}=0`$ is the same for all particles. In the infinite momentum frame, $`P^+\mathrm{}`$, this assumption is justified because the particles can be viewed as being confined to a Lorentz contracted sheet in the transverse plane of width $`1/P^+`$. This implies that the particles can be labeled using their transverse positions $`x_{}^I`$ only. Using
$`\rho _a(x_{})={\displaystyle \underset{I=1}{\overset{K}{}}}\delta (x_{}x_{}^I)Q_a^I,`$ (17)
one can assume $`\rho ^a(x_{})`$ to be continuous (and large). One can therefore make an educated guess that the coarse grained effective action of the wee parton modes will be
$`S_{\mathrm{ln}W}={\displaystyle \frac{1}{4}}{\displaystyle d^4xG_{\mu \nu }^aG_a^{\mu \nu }}+{\displaystyle \frac{i}{N_c}}{\displaystyle d^2x_{}\mathrm{Tr}\left\{\rho (x_{})\mathrm{ln}W(x_{})\right\}}`$
(18)
where now
$`W(x_{})=\widehat{P}\mathrm{exp}\left(ig{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑x^+A_a^{}(x^+,0,x_{})T_a\right).`$ (19)
Just as in Eq. (1), the action $`S_{\mathrm{ln}W}`$ should contain an identical functional $`F[\rho ]`$ representing the likelihood of different $`\rho `$ configurations. This term will only be implicit in what follows since it is not relevant to the concerns of this paper.
We will now show explicitly that the charge obtained from the action $`S_{\mathrm{ln}W}`$ is Hermitean and traceless, and therefore an element of the Lie algebra.
We first define the log of an operator as the power series
$`\mathrm{ln}W=\mathrm{ln}(1(1W)){\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{k}}(1W)^k.`$ (20)
Taking the functional derivative of $`\mathrm{ln}W`$ with respect to $`A`$ gives
$`{\displaystyle \frac{\delta }{\delta A_\mu ^a}}\mathrm{ln}W={\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{k}}{\displaystyle \underset{s=0}{\overset{k1}{}}}(1W)^s{\displaystyle \frac{\delta W}{\delta A_\mu ^a}}(1W)^{ks1}`$ (21)
After some straightforward algebra, this can be written as
$`{\displaystyle \frac{\delta }{\delta A_\mu ^a}}\mathrm{ln}W={\displaystyle _0^1}𝑑\lambda {\displaystyle \frac{1}{1(1W)\lambda }}{\displaystyle \frac{\delta W}{\delta A_\mu ^a}}{\displaystyle \frac{1}{1(1W)\lambda }}`$ (22)
Then from the relation
$`J_a^\mu ={\displaystyle \frac{i}{N_c}}\mathrm{Tr}\left(\rho {\displaystyle \frac{\delta }{\delta A_\mu ^a}}\mathrm{ln}W\right),`$ (23)
we find that the color charge is given by
$`Q(x^+)`$ $`=`$ $`{\displaystyle d^3xJ^+(x)}`$ (24)
$`=`$ $`{\displaystyle d^3xW(x^+,\mathrm{})\left[_0^1𝑑\lambda \frac{1}{B(\lambda )}\rho \frac{1}{B(\lambda )}\right]W(\mathrm{},x^+)}`$
where we used the shorthand
$`B(\lambda )1(1W)\lambda `$ (25)
and
$`W(x_f^+,x_i^+)=\widehat{P}\mathrm{exp}\left(ig{\displaystyle _{x_i^+}^{x_f^+}}𝑑x^+A_a^{}(x^+,x^{},x_{})T_a\right)`$ (26)
We also defined $`WW(\mathrm{},\mathrm{})`$.
It is easy to check that the current density $`J^+`$ satisfies
$`{\displaystyle \frac{J^+}{x^+}}`$ $`=`$ $`ig[J^+,A^{}],`$ (27)
and hence is a solution of the Wong’s equation with a “boundary” condition given by
$`J^+(x^+=\mathrm{})={\displaystyle _0^1}𝑑\lambda B^1(\lambda )\rho B^1(\lambda )W`$ (28)
Note that $`1/B(\lambda )=B^1(\lambda )`$.
To confirm that $`J^+(x)`$ is an element of the Lie-Algebra, first consider the trace. We have
$`\mathrm{Tr}\left(J^+(x)\right)`$ $`=`$ $`\rho ^b\mathrm{Tr}\left(W\left[{\displaystyle _0^1}𝑑\lambda \left(B^1(\lambda )\right)^2T_b\right]\right),`$ (29)
since $`W`$ and $`B`$ commute. One can show that
$`{\displaystyle \frac{d}{d\lambda }}B^1(\lambda )=(1)\left(B^1(\lambda )\right)^2(W1).`$ (30)
Consequently,
$`\mathrm{Tr}\left(J^+(x)\right)`$ $`=`$ $`\rho ^b\mathrm{Tr}\left(W\left[{\displaystyle _0^1}𝑑\lambda (B^1(\lambda ))^2T_b\right]\right)`$ (31)
$`=`$ $`\rho ^b\mathrm{Tr}\left(W(1W)^1\left(B^1(1)B^1(0)\right)T_b\right)`$
$`=`$ $`\rho ^b\mathrm{Tr}\left(W(1W)^1(W^11)T_b\right)`$
$`=`$ $`\rho ^b\mathrm{Tr}\left(T_b\right)=0.`$
We shall now show that $`J`$ is also Hermitean. Consider
$`(J^+(x))^{}`$ $`=`$ $`\rho ^b\left(W(x^+,\mathrm{})\left[{\displaystyle _0^1}𝑑\lambda B^1(\lambda )T_bB^1(\lambda )\right]W(\mathrm{},x^+)\right)^{}`$ (32)
$`=`$ $`\rho ^bW(x^+,\mathrm{})\left[{\displaystyle _0^1}𝑑\lambda (B^1(\lambda ))^{}T_b(B^1(\lambda ))^{}\right]W(\mathrm{},x^+)`$
$`=`$ $`\rho ^bW(x^+,\mathrm{})\left[{\displaystyle _0^1}𝑑\lambda W^{}(B^1(\lambda ))^{}T_bW^{}(B^1(\lambda ))^{}\right]W(\mathrm{},x^+).`$
Let us now focus on the term in the square brackets. Since
$`W^{}=W(\mathrm{},\mathrm{})=W^1,`$ (33)
this term can be re-written as
$`{\displaystyle _0^1}𝑑\lambda (B^1(\lambda )W)^{}T_b(B^1(\lambda )W)^{}.`$ (34)
Here one has used the relation $`WB^1=B^1W`$. Performing the change of variable $`\lambda 1\lambda `$, one can show that
$`(B^1W)=(B^1)^{}.`$ (35)
Thus,
$`(J^+(x))^{}`$ $`=`$ $`\rho ^bW(x^+,\mathrm{}){\displaystyle _0^1}𝑑\lambda \left[(B^1(\lambda )W)^{}T_b(B^1(\lambda )W)^{}\right]W(\mathrm{},x^+),`$ (36)
$`=`$ $`\rho ^bW(x^+,\mathrm{}){\displaystyle _0^1}𝑑\lambda \left[B^1(\lambda )T_bB^1(\lambda )\right]W(\mathrm{},x^+),`$
$`=`$ $`J^+(x).`$
We have now explicitly shown above that $`J^+`$ (and hence $`Q`$) is both Hermitean and traceless. It is therefore an element of the Lie Algebra. In general, it is possible, if non–trivial, to show that $`\mathrm{ln}(W)`$ itself is a member of the Lie algebra . The charge obtained from Eq. (4) is also an element of the Lie Algebra. It is easy to see that the color components of the color charge $`J_a^+`$ are real and therefore, the color charge matrix defined as $`J_\mu =\frac{1}{N_c}J_\mu ^aT^a`$ is Hermitean and traceless. Both $`S_W`$ and $`S_{\mathrm{ln}W}`$ lead to Wong’s equations, but with a different current $`J_\mu `$. This difference is due to imposing different “boundary” conditions at $`\tau =\mathrm{}`$ when solving the Wong’s equations (13) as given by (5) and (28). It should be noted that the boundary condition in (28) is more complicated than (5), and involves the non–trivial task of inverting the operator $`B(\lambda )`$. It is important to realize that the two different currents may describe different physics.
## 4 The $`\mathrm{ln}W`$ action and the BFKL equation
We will now show that the form of the action in Eq. (18) also reproduces the BFKL equation. Since the two actions differ only by the form of the Wilson line term, we will focus on the expansion of the Wilson line term in the two actions. To reproduce the Wilsonian renormalization group evolution, we need to keep terms that are quadratic in the hard fluctuations (the field $`\delta A^\mu `$ in Eq. (7)). The leading order non-trivial contribution therefore comes from the cubic terms in the action. (The contribution from quartic terms to the evolution is sub–leading in DIS.)
The difference between the two actions is
$`\mathrm{\Delta }SS_WS_{\mathrm{ln}W}=\mathrm{Tr}\left(\rho [W\mathrm{ln}W]\right)`$ (37)
where $`\rho =\rho ^aT_a`$ and $`\mathrm{ln}W`$ is defined as in Eq. (20) to be
$`\mathrm{ln}W\mathrm{ln}[1(1W)]=\mathrm{Tr}\left(\rho [(1W)+{\displaystyle \frac{1}{2}}(1W)^2+{\displaystyle \frac{1}{3}}(1W)^3+\mathrm{}]\right).`$ (38)
The integration over the spatial variables $`x^{}`$ and $`x_{}`$ and the convolution with $`\delta (x^{})`$ is implicit in the trace above. The difference between the two actions is then
$`\mathrm{\Delta }S=\mathrm{Tr}\left(\rho [{\displaystyle \frac{1}{2}}(1W)^2+{\displaystyle \frac{1}{3}}(1W)^3+\mathrm{}]\right),`$ (39)
where $`1W`$, from Eq. (2) can be expanded as $`1W=igA^{}+(g^2/2)\widehat{P}(A^{})^2+\mathrm{}`$. Again, the integral over $`x^+`$ is implicit in the expansion, with the symbol $`\widehat{P}`$ denoting the time ordering in $`x^+`$. Potential differences between the two actions will show up at order $`A^2`$. At this order <sup>2</sup><sup>2</sup>2We use the following conventions for the trace of adjoint matrices: $`\mathrm{Tr}(T^aT^b)=N_c\delta _{ab}`$ and $`\mathrm{Tr}(T^aT^bT^c)=\frac{iN_c}{2}f_{abc}`$.,
$`\mathrm{\Delta }S(A^2)\mathrm{Tr}\left(\rho (A^{})^2\right)\rho _af_{abc}{\displaystyle 𝑑x^+𝑑y^+A_b^{}(x^+)A_c^{}(y^+)}.`$ (40)
This term is identically zero because the integrand is symmetric under both the color exchange $`bc`$ and the co–ordinate exchange $`x^+y^+`$ while multiplying the totally anti–symmetric structure constant $`f_{abc}`$.
To investigate terms of order $`A^3`$, it is convenient to first consider $`S_W`$ and $`S_{\mathrm{ln}W}`$ separately. The cubic terms in the expansion of $`S_{\mathrm{ln}W}`$ are
$`S_{\mathrm{ln}W}(A^3)`$ $`=`$ $`{\displaystyle \frac{g^3}{N_c}}\mathrm{Tr}\left(\rho \left[\widehat{P}(A^{})^3{\displaystyle \frac{1}{2}}A^{}\widehat{P}(A^{})^2{\displaystyle \frac{1}{2}}\widehat{P}(A^{})^2A^{}+{\displaystyle \frac{1}{3}}(A^{})^3\right]\right)`$ (41)
$`=`$ $`{\displaystyle \frac{g^3}{N_c}}\mathrm{Tr}\rho {\displaystyle 𝑑x^+𝑑y^+𝑑z^+A^{}(x^+)A^{}(y^+)A^{}(z^+)}`$
$`\times `$ $`\left[\theta (x^+y^+)\theta (y^+z^+){\displaystyle \frac{1}{2}}\theta (x^+y^+){\displaystyle \frac{1}{2}}\theta (y^+z^+)+{\displaystyle \frac{1}{3}}\right]`$
After some algebra (performed in appendix B) the above can be re-expressed as
$`S_{\mathrm{ln}W}(A^3)`$ $`=`$ $`{\displaystyle \frac{g^3}{6}}\rho _a{\displaystyle 𝑑x^+𝑑y^+𝑑z^+A_b^{}(x^+)A_c^{}(y^+)A_d^{}(z^+)\theta (x^+y^+)\theta (y^+z^+)}`$ (42)
$`\times `$ $`\left[f^{adn}f^{bcn}f^{abn}f^{cdn}\right]`$
The cubic term in $`S_W`$ is
$`S_W(A^3)`$ $`=`$ $`{\displaystyle \frac{g^3}{N_c}}\mathrm{Tr}\rho \widehat{P}(A^{})^3`$ (43)
$`=`$ $`{\displaystyle \frac{g^3}{N_c}}\rho _a{\displaystyle 𝑑x^+𝑑y^+𝑑z^+A_b^{}(x^+)A_c^{}(y^+)A_c^{}(z^+)\theta (x^+y^+)\theta (y^+z^+)}`$
$`\times `$ $`\left[{\displaystyle \frac{I_2}{6}}(f_{adn}f_{bcn}f_{abn}f_{cdn})+d_{abcd}\right].`$
Here, we have used an identity for the trace of four $`SU(3)`$ adjoint matrices . For an adjoint representation, $`I_2=N_c`$. Also, the totally symmetric tensor $`d_{abcd}`$ is defined as the symmetrized trace of four $`SU(3)`$ adjoint matrices. For an explicit form, see Ref. . Note that the f-terms above are identical to those derived from $`S_{\mathrm{ln}W}`$ in Eq. (42). However, this action also contains the $`d_{abcd}`$ term that was absent in the $`S_{\mathrm{ln}W}`$ action.
In appendix C, we show that, within the approximations made in the derivation of the small $`x`$ evolution equation in Ref. , the $`d_{abcd}`$ term does not contribute. Therefore $`\mathrm{\Delta }S=0`$ to cubic order. One may therefore conclude that BFKL equation can also be obtained from the $`S_{\mathrm{ln}W}`$ action.
The reason the $`d_{abcd}`$ term in the $`S_W`$ action vanishes is because the propagator of the hard modes (and the color sources to which it couples) is static. The static nature of the sources is due to the fact that one ignores the recoil of the color sources as they emit softer partons. As one goes to a next-to-leading-order calculation, one will have to take recoil effects into account. These would cause the color sources to be time dependent, giving rise to a finite contribution from the $`d^{abcd}`$ terms in the $`W`$ action. Conversely, note that the $`d^{abcd}`$ terms are naturally absent in $`\mathrm{ln}W`$ action.
The fact that the $`S_{\mathrm{ln}W}`$ action does not have the $`d_{abcd}`$ term suggests that the underlying symmetry of the small $`x`$ dynamics is manifest in this action. The agreement between the two actions is even more remarkable when one considers that the factor $`1/6`$ in Eq. (43) comes directly from the trace of four adjoint generators, while in the “$`\mathrm{ln}W`$” action it arises as a consequence of extensive algebraic manipulations.
## 5 Summary
In this paper, we proposed an alternative form of the small $`x`$ effective action to the one discussed in Ref. . We showed explicitly that both forms of the effective action are compatible with Wong’s equations, albeit with currents that satisfy different boundary conditions. We showed that the the two effective actions agreed up to cubic order in the fields. Consequently, both of them give rise to the BFKL equation. However, in the case of the effective action of Ref. , one had to explicitly invoke the kinematic constraint imposed by the static sources–no such constraint was necessary for the action we propose. Differences between the two actions will show up at higher orders when one considers sub–leading corrections to the small $`x`$ effective action.
## Acknowledgments
We would like to thank Alex Kovner, Rob Pisarski and Jens Wirstam for reading the manuscript. We thank Larry McLerran for useful remarks. One of us (R.V.) would also like to thank E. Iancu for vigorous discussions and (S.J.) would like to thank Harry Lam for insightful comments. J. J-M. would like to thank Alex Kovner for helpful discussions. S. J. was supported by the Director, Office of Science, Office of High Energy and Nuclear Physics, Division of Nuclear Physics, and by the Office of Basic Energy Sciences, Division of Nuclear Sciences, of the U.S. Department of Energy under Contract No. DE-AC03-76SF00098. J.J-M. is supported at the University of Arizona by the U.S. Department of Energy under Contract No. DE-FG03-93ER40792. R. V. is supported at BNL by the U.S. Department of Energy under Contract No. DE-AC02-98CH10886.
## Appendix A
The equation of motion for color charges in the action (18) has already been derived in Section 3. The derivation is the same for the point particle action Eq. (14). Further, the Euler-Lagrange equation for the position
$`p_\mu ^I=m{\displaystyle \frac{dx_\mu ^I}{d\tau }}`$ (44)
follows trivially from Eq. (14). Hence, we will only derive here the other remaining equation of motion, namely, that for the momentum $`p_\mu `$.
The Euler-Lagrange equation for the momentum in Eq. (14) is
$`\dot{p}_\mu ^I={\displaystyle \frac{\delta }{\delta x_I^\mu }}S_{\mathrm{Wong}}.`$ (45)
The right hand side is given by
$`{\displaystyle \frac{\delta S_{\mathrm{Wong}}}{\delta x_I^\mu (\tau )}}`$ $`=`$ $`{\displaystyle \frac{i}{N_c}}{\displaystyle \frac{\delta }{\delta x_I^\mu (\tau )}}{\displaystyle \underset{J=1}{\overset{K}{}}}\mathrm{Tr}\left\{Q_0^J\mathrm{ln}W_J\right\}`$ (46)
$`=`$ $`{\displaystyle \frac{i}{N_c}}{\displaystyle 𝑑\tau ^{}\underset{J=1}{\overset{K}{}}\frac{\delta A_a^\nu (x_J(\tau ^{}))}{\delta x_I^\mu (\tau )}\frac{\delta }{\delta A_a^\nu (x_J(\tau ^{}))}\mathrm{Tr}\left\{Q_0^J\mathrm{ln}W_J\right\}}`$
$`+{\displaystyle \frac{i}{N_c}}{\displaystyle 𝑑\tau ^{}\underset{J=1}{\overset{K}{}}\frac{\delta v_J^\nu (\tau ^{})}{\delta x_I^\mu (\tau )}\frac{\delta }{\delta v_J^\nu (\tau ^{})}\mathrm{Tr}\left\{Q_0^J\mathrm{ln}W_J\right\}}`$
Using Eqs.(23) and (24), we get
$`{\displaystyle \frac{\delta S_{\mathrm{Wong}}}{\delta x_I^\mu (\tau )}}`$ $`=`$ $`v^\nu _\mu A_\nu ^a(x_I(\tau ))Q_a^I(\tau )+{\displaystyle \frac{i}{N_c}}{\displaystyle 𝑑\tau ^{}\dot{\delta }(\tau ^{}\tau )\frac{\delta }{\delta v_\mu ^I(\tau ^{})}\mathrm{Tr}\left\{Q_0^I\mathrm{ln}W_I\right\}}`$ (47)
For simplicity, we’ll omit the particle label $`I`$ from now on. We’ll also omit any argument labels (everything should be considered as a function of $`\tau `$ via $`x(\tau )`$). Applying the method developed in Section 3, we see that
$`{\displaystyle \frac{\delta S_{\mathrm{Wong}}}{\delta x^\mu }}`$ $`=`$ $`Q_av^\nu _\mu A_\nu ^a{\displaystyle \frac{i}{N_c}}{\displaystyle \frac{d}{d\tau }}\mathrm{Tr}\left\{Q_0{\displaystyle \frac{\delta }{\delta v_\mu }}\mathrm{ln}W\right\}`$ (48)
$`=`$ $`Q_av^\nu _\mu A_\nu ^a+{\displaystyle \frac{d}{d\tau }}\left\{A_\mu ^aQ_I^a\right\}`$
$`=`$ $`Q_av^\nu _\mu A_\nu ^a+Q_av^\nu _\nu A_\mu ^a+A_\mu ^a{\displaystyle \frac{d}{d\tau }}Q_a`$
$`=`$ $`Q_av^\nu _\mu A_\nu ^a+Q_av^\nu _\nu A_\mu ^aigA_\mu ^a[Q,v^\nu A_\nu ]_a`$
$`=`$ $`Q_aG_{\mu \nu }^av^\nu `$
where
$`G_{\mu \nu }^a=_\mu A_\nu ^a_\nu A_\mu ^a+gf_{abc}A_\mu ^bA_\nu ^c`$ (49)
This yields the desired result:
$`\dot{p}_\mu =Q_aG_{\mu \nu }^av^\nu .`$ (50)
## Appendix B
We will show here how one gets equation (42) from equation (41). We start with Eq. (41):
$`S_{\mathrm{ln}W}(A^3)`$ $`=`$ $`{\displaystyle \frac{g^3}{N_c}}\mathrm{Tr}(\rho {\displaystyle }dx^+dy^+dz^+A^{}(x^+)A^{}(y^+)A^{}(z^+)`$
$`[\theta (x^+y^+)\theta (y^+z^+){\displaystyle \frac{1}{2}}\theta (x^+y^+){\displaystyle \frac{1}{2}}\theta (y^+z^+)+{\displaystyle \frac{1}{3}}]).`$
To save space, we will use the following shorthand notation. We shall represent the Light Cone times $`x^+,y^+,z^+`$ by $`1,2,3`$, and shall not write the Light Cone Lorentz index ‘$``$’ and the integrations over the Light Cone times explicitly. For example,
$`\theta _{123}\theta (x^+y^+)\theta (y^+z^+).`$
We shall also neglect the overall coefficient $`g^3/N_c`$, including it only at the very last step. With these notations, equation (41) becomes
$`S_{\mathrm{ln}W}(A^3)`$ $`=`$ $`\mathrm{Tr}\left(\rho A_1A_2A_3\left[\theta _{123}{\displaystyle \frac{1}{2}}\theta _{12}{\displaystyle \frac{1}{2}}\theta _{23}+{\displaystyle \frac{1}{3}}\right]\right).`$ (51)
We will now use identities like $`\theta _{12}+\theta _{21}=1`$ to write Eq. (51) as
$`S_{\mathrm{ln}W}(A^3)`$ $`=`$ $`\mathrm{Tr}\rho A_1A_2A_3[\theta _{123}{\displaystyle \frac{1}{2}}\theta _{12}(\theta _{23}+\theta _{32}){\displaystyle \frac{1}{2}}\theta _{23}(\theta _{12}+\theta _{21})`$
$`+`$ $`{\displaystyle \frac{1}{3}}(\theta _{123}+\theta _{132}+\theta _{231}+\theta _{213}+\theta _{312}+\theta _{321})]`$
$`=`$ $`\mathrm{Tr}\rho A_1A_2A_3\left[{\displaystyle \frac{1}{2}}(\theta _{12}\theta _{32}+\theta _{21}\theta _{23})+{\displaystyle \frac{1}{3}}(\theta _{123}+\theta _{132}+\theta _{231}+\theta _{213}+\theta _{312}+\theta _{321})\right]`$
$`=`$ $`\mathrm{Tr}\rho A_1A_2A_3[{\displaystyle \frac{1}{2}}\theta _{12}\theta _{32}(\theta _{13}+\theta _{31}){\displaystyle \frac{1}{2}}\theta _{21}\theta _{23}(\theta _{13}+\theta _{31})`$
$`+`$ $`{\displaystyle \frac{1}{3}}(\theta _{123}+\theta _{132}+\theta _{231}+\theta _{213}+\theta _{312}+\theta _{321})]`$
$`=`$ $`\mathrm{Tr}\rho A_1A_2A_3[{\displaystyle \frac{1}{2}}(\theta _{132}+\theta _{312}+\theta _{213}+\theta _{231})`$
$`+`$ $`{\displaystyle \frac{1}{3}}(\theta _{123}+\theta _{132}+\theta _{231}+\theta _{213}+\theta _{312}+\theta _{321})]`$
$`=`$ $`{\displaystyle \frac{1}{6}}\mathrm{Tr}\rho A_1A_2A_3\left[\theta _{123}\theta _{213}+\theta _{321}\theta _{312}+\theta _{123}\theta _{132}+\theta _{321}\theta _{231}\right]`$
$`=`$ $`{\displaystyle \frac{1}{6}}\mathrm{Tr}\rho (\theta _{123}[A_1,A_2]A_3+\theta _{321}[A_1,A_2],A_3+\theta _{123}A_1[A_2,A_3]+\theta _{321}A_1[A_2,A_3])`$
which, after change of variables, can be re–written as
$`S_{\mathrm{ln}W}(A^3)`$ $`=`$ $`{\displaystyle \frac{1}{6}}\theta _{123}\mathrm{Tr}\left(\rho \left([[A_1,A_2],A_3]+[[A_3,A_2],A_1]\right)\right).`$ (52)
Using $`[T^a,T^b]=if^{abc}T^c`$, and restoring all the indices, coefficients, and integration variables, we obtain finally
$`S_{\mathrm{ln}W}(A^3)`$ $`=`$ $`{\displaystyle \frac{g^3}{6}}\rho _a{\displaystyle }dx^+dy^+dz^+[A_b^{}(x^+)A_c^{}(y^+)A_d^{}(z^+)\theta (x^+y^+)\theta (y^+z^+)`$ (53)
$`\times `$ $`\left[f_{adn}f_{bcn}f_{abn}f_{cdn}\right],`$
which is Eq. (42).
## Appendix C
In this appendix, we will show that, within the approximations made in the Wilson renormalization group approach, the totally symmetric $`d_{abcd}`$–term in Eq. (43) vanishes. Using the field decomposition Eq. (7), we get
$`S_{cubic}^{sym}`$ $`=`$ $`{\displaystyle \frac{g^3}{N_c}}d^{abcd}\rho _a{\displaystyle 𝑑x^+𝑑y^+𝑑z^+\theta (x^+y^+)\theta (y^+z^+)}`$
$`\left[a_b^{}(x^+)A_c^{}(y^+)A_d^{}(z^+)+A_b^{}(x^+)a_c^{}(y^+)A_d^{}(z^+)+A_b^{}(x^+)A_c^{}(y^+)a_d^{}(z^+)\right].`$
Renaming the variables in the second and third terms above gives
$`S_{cubic}^{sym}`$ $`=`$ $`{\displaystyle \frac{g^3}{N_c}}d^{abcd}\rho _a{\displaystyle 𝑑x^+𝑑y^+𝑑z^+a_b^{}(x^+)A_c^{}(y^+)A_d^{}(z^+)}`$
$`\left[\theta (x^+y^+)\theta (y^+z^+)+\theta (y^+x^+)\theta (x^+z^+)+\theta (z^+y^+)\theta (y^+x^+)\right].`$
Re–write the sum of $`\theta `$–functions above as follows:
$`\theta (x^+y^+)\theta (y^+z^+)+\theta (y^+x^+)\theta (x^+z^+)+\theta (z^+y^+)\theta (y^+x^+)=`$
$`\theta (x^+y^+)\theta (y^+z^+)\theta (x^+z^+)+\theta (y^+x^+)\theta (x^+z^+)\theta (y^+z^+)+`$
$`\theta (z^+y^+)\theta (y^+x^+)=`$
$`\theta (y^+z^+)\theta (x^+z^+)+\theta (z^+y^+)\theta (y^+x^+)=`$
$`\theta (y^+z^+)\theta (x^+z^+)\theta (z^+y^+)\theta (x^+y^+)+\theta (z^+y^+)`$
The first two terms in the last line are anti–symmetric with respect to change of $`y^+z^+`$, and multiply the product $`d^{abcd}A_c^{}(y^+)A_d^{}(z^+)`$ which is totally symmetric. They therefore vanish, and the expression in $`S_{cubic}^{sym}`$ reduces to
$`S_{cubic}^{sym}`$ $`=`$ $`{\displaystyle \frac{g^3}{N_c}}d^{abcd}\rho _a{\displaystyle 𝑑x^+𝑑y^+𝑑z^+\theta (z^+y^+)a_b^{}(x^+)A_c^{}(y^+)A_d^{}(z^+)}.`$ (54)
This term can further be written as
$`{\displaystyle 𝑑x^+𝑑y^+𝑑z^+\theta (z^+y^+)a_b^{}(x^+)A_c^{}(y^+)A_d^{}(z^+)}=`$
$`{\displaystyle 𝑑x^+𝑑y^+𝑑z^+\left[1\theta (y^+z^+)\right]a_b^{}(x^+)A_c^{}(y^+)A_d^{}(z^+)}.`$ (55)
Consider the “$`1`$” term on the right hand side:
$`{\displaystyle 𝑑x^+𝑑y^+𝑑z^+a_b^{}(x^+)A_c^{}(y^+)A_d^{}(z^+)}={\displaystyle 𝑑x^+a_b^{}(x^+)𝑑y^+𝑑z^+A_c^{}(y^+)A_d^{}(z^+)}`$
$`{\displaystyle 𝑑x^+a_b^{}(x^+)d(y^+z^+)d(y^++z^+)A_c^{}(y^+)A_d^{}(z^+)}.`$ (56)
When we integrate over hard fluctuations, the term $`A_c^{}(y^+)A_d^{}(z^+)`$ will become the hard fluctuations propagator $`G^{}(y^+z^+)`$. After integrating this propagator over the $`d(y^+z^+)`$ variable, it will give an overall factor $`p^{}`$. Since $`\rho \delta (p^{})`$ (this would break down when considering NLO corrections), the “$`1`$” term vanishes because $`p^{}\delta (p^{})=0`$. The integrand of Eq. (54) is then
$`{\displaystyle 𝑑x^+𝑑y^+𝑑z^+\theta (z^+y^+)a_b^{}(x^+)A_c^{}(y^+)A_d^{}(z^+)}=`$
$`{\displaystyle 𝑑x^+𝑑y^+𝑑z^+\theta (y^+z^+)a_b^{}(x^+)A_c^{}(y^+)A_d^{}(z^+)}.`$ (57)
Since the LHS of the above is invariant under $`z^+y^+,cd`$, the above identity would require it to be ‘$``$’ of itself, and therefore equal to zero. Thus the $`d_{abcd}`$–term in the cubic piece of the action vanishes, and we are left with
$`S_{cubic}`$ $`=`$ $`{\displaystyle \frac{g^3}{N_c}}\rho _a{\displaystyle 𝑑x^+𝑑y^+𝑑z^+A_b^{}(x^+)A_c^{}(y^+)A_d^{}(z^+)}`$ (58)
$`\times `$ $`{\displaystyle \frac{I_2}{6}}\theta (x^+y^+)\theta (y^+z^+)\left[f^{adn}f^{bcn}f^{abn}f^{cdn}\right],`$
which is identical to cubic term in the expansion of $`S_{\mathrm{ln}W}`$.
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# PHYSICAL MODEL OF LEPTONS: MASSIVE ELECTRONS, MUONS, TAUONS AND THEIR MASSLESS NEUTRINOUS
## Abstract
The physical model (PhsMdl) of the leptons is offered by means of the PhsMdls of the vacuum and electron,described in our recent works. It is assumed that the vacuum is consistent by dynamides, streamlined in junctions of some tight crystalline lattice. Every dynamide contains a neutral pair of massless point-like (PntLk) contrary elementary electric charges (ElmElcChrgs): electrino $`()`$ and positrino $`(+)`$. The PntLk ElmElcChrgs of the massless electrino and positrino of some dynamide in the fluctuated vacuum may been excited or deviated by means of some energy, introduced by some photon or other micro particles (MicrPrts). The massless leptons (neutrinos) are neutral long-living solitary spherical vortical oscillation excitations of the uncharged fluctuating vacuum. The massive leptons are charged long-living solitary spherical vortical excitations of its fine spread (FnSpr) elementary electric charge (ElmElcChrg). So-called zitterbewegung is self-consistent strong-correlated vortical harmonic oscillation motion of the FnSpr ElmElcChrg of massive leptons. Different leptons have different self-consistent strong-correlated vortical harmonic oscillation motion of different sizes of their FnSpr ElmElcChrg, which is determined by their Kompton length $`\lambda =\frac{h}{mC}`$, where $`m`$ is the mass of the massive leptons. At mutual transition of one massive lepton into another massive lepton its PntLk ElmElcChrg move up by dint of weak interaction in the form of the charged intermediate vector meson $`W`$ from one self-consistent strong-correlated vortical harmonic oscillation motion of one size into another self-consistent strong-correlated vortical harmonic oscillation motion of another size.
Although up to the present nobody of scientists distinctly knows are there some elementary micro particles (ElmMicrPrts) as a stone of the micro world and what the elementary micro particle (ElmMicrPrt) means, there exists an essential possibility for clear and obvious scientific consideration of the uncommon quantum behaviour and unusual relativistic dynamical parameters of all the relativistic quantized micro particle (QntMicrPrt) by means of our transparent surveyed PhsMdl. It is well known that the physical model (PhsMdl) presents at us as an actual ingradient of every good physical theory (PhsThr). It may be used as for an obvious visual teaching the unknown occurred physical processes within the investigated phenomena, It turned out that all leptons are elementary micro particles (ElmMicrPrt) of two kinds: a charged massive and uncharged massless. The massless leptons (neutrinos) are neutral long-living solitary isotropic three dimensional relativistic quantized (SltIstThrDmnRltQnt) spherical vortex harmonic oscillation (SphrVrtNrmOsc) excitations of the neutral fluctuated vacuum (FlcVcm) without a self-energy at a rest. The massive leptons are charged long-living solitary isotropic three dimensional relativistic quantized (SltIstThrDmnRltQnt) spherical vortical harmonic oscillation (SphrVrtHrmOsc) motion of its fine spread (FnSpr) elementary electric charge (ElmElcChrg). The PhsMdl of the charged massive lepton will be presented by the PhsMdl of one of them, the DrEl. Therefore so-called zitterbewegung by Schrodinger,which is self-consistent strong-correlated vortical harmonic oscillation (VrtHrmOscs) motion of the FnSpr ElmElcChrg of one of massive leptons, electron, is generalized for all charged massive leptons. Therefore different charged massive leptons have different self-consistent strong-correlated vortical harmonic oscillation motion of different sizes of their FnSpr ElmElcChrg, which is determined by their Kompton length $`\lambda =\frac{h}{mC}`$, where $`m`$ is the mass of the massive leptons. At mutual transition of one massive lepton into another massive lepton its PntLk ElmElcChrg move up by means of weak interaction in the form of the charged intermediate vector meson $`W`$ from one self-consistent strong-correlated (VrtHrmOscs) motion of one size into another self-consistent strong-correlated (VrtHrmOscs) motion of another size.
The PhsMdl of the DrEl have offered in all my work in resent nineteen years for bring of light to the physical interpretation of physical cause of the uncommon quantum behaviour of the Schrodinger electron (SchEl) and the relativistic behaviour of the DrEl and give the thru physical interpretation and sense of all its dynamical parameters. Our PhsMdl of the DrEl explain as the physical causes for its unusual stochastic classical dual wave-corpuscular behaviour and so give a new cleared picturesque physical interpretation with mother wit of the physical means of its relativistic dynamical parameters. In our transparent surveyed PhsMdl of the DrEl one will be regarded as some point like (PntLk) ElmElcChrg, taking simultaneously part in four different motions:
A) The isotropic three-dimensional relativistic quantized (IstThrDmnRltQnt) Einstein’s stochastic (EinStch) boson harmonic shudders(BsnHrmShdrs) as a result of momentum recoils (impulse kicks), forcing the charged QntMicrPrt at its continuously stochastical emissions and absorptions of own high energy (HghEnr) virtual photons (StchVrtPhtns) by its PntLk ElmElcChrg. This jerky motion display almost Brownian classical stochastic behaviour (BrnClsStchBhv) with a light velocity C during a small time interval $`\tau _1`$, much less then the period $`T`$ of the IstThrDmnRltQnt Schrodinger fermion (SchFrm) vortical harmonic oscillation (VrtHrmOsc) motion and more larger then the time interval $`\tau _o`$ of the stochastically emission or absorption of the Hgh-Enr StchVrtPhtn by its PntLk ElmElcChrg. The display of the IstThrDmnRltQnt EinStch BsnHrmShds can be observed at the light (RlPhtn) scattering of free DrEl. Indeed, in a consequence of the above investigation we may really consider the IstThrDmnRltQnt FrthStch VrtHrmOscMtn’s ”trajectory” of the DrEl turns into cylindrically fine spread path, which has a form of the distorted figure of eight. The participate the FnSpr ElmElcChrg of DrEl in IstThrDmnRltQnt FrthStch VrtHrmOscMtn well spread (WllSpr) the ElmElcChrg of the SchEl. Further the behaviour of the DrEl’s FnSpr ElmElcChrg may be treated as a nearly Brownian classical stochastic behaviour (BrnClsStchBhv) of the ClsMicrPrt, taking part in the relativistic random trembling motion (RltRndTrmMtn), during the time intervals $`\tau _1`$, much less then the period $`T`$ of the SchFrm VrtHrmOscs and more larger then the time interval $`\tau `$ of the emission and absorption of the Hgh-Enr StchVrtPhtn, the QntElcMgnFld of which form the RslSlfCnsVls of own RslQntElcMgmFld of the DrEl’s FnSpr ElmElcChrg. In a consequence of this the product $`2C\tau `$ may be considered as the space parth values of the DrEl’s PntLk ElmElcChrg, which is passed in the time interval $`2\tau `$ of its absorption and emission (or scattering) of some RlPhtn. In a consequence of this we can easily understand by our felicitous PhsMdl why the classical radius of the LrEl entirely coincides with the size $`\sqrt{\xi ^2}`$ of a spherical fine spread spot with an effective scattering surface $`\pi \xi ^2`$. Therefore the IstThrDmnRltQnt SchFrm VrtHrmOscMtn’s trajectory turns into fine spread path of a cylindrical shape with different radii. Therefore the size of this smallest IstThrDmnRltQnt EinStch BznHrmShd’s motion of the DrEl’s PntLk ElmElcChrg could be determined by the Thompson total cross section. Indeed the averaged cross section of the cylindrical spread path of the IstThrDmnRltQnt StchBsnHrmOscMtn can be determined by the effective total cross section of Thompson of RlPhtn’s scattering from the free DrEl’s, which has a space distribution with a spherical symmetry may be easily obtained by means of the following simple relations, well-known by us from the classical StchMch :
$$\sigma =\pi \xi ^2=\frac{3\pi }{2}\left[\frac{4e^2}{3mC^2}\right]^2=\frac{8\pi }{3}\left[\frac{e^2}{mC^2}\right]^2$$
(1)
It is very important to make here, that this EinStchBznHrmShds have the smallest size and the fastest velocity and therefore the described smallest cross section may be represented as a described roughly by the space distribution of the DrEl’s FnSpt ElmElcChrg :
$$|\mathrm{{\rm Y}}_o(\varrho )|^2=\left[(\frac{3}{5\sqrt{\pi }})^3(\frac{mC^2}{e^2})^3\right]\mathrm{exp}(\frac{\varrho ^2}{\kappa _o^2})$$
(2)
where $`\kappa _o`$ is parameter of the EinStch BsnHrmShds $`(\kappa _o=(5/3)\frac{e^2}{mc^2}`$. The space distribution (2) of the DrEl’s FnSpt ElmElcChrg is described roughly by an OrbWvFnc $`\mathrm{{\rm Y}}`$, having the following form :
$$\mathrm{{\rm Y}}_o(\varrho )=(\sqrt{\pi }\kappa _o)^{}\frac{3}{2}\mathrm{exp}(\frac{\varrho ^2}{2\kappa _o^2})$$
(3)
B) The isotropic three-dimensional relativistic quantized (IstThrDmnRltQnt) Schrodinger fermion vortical harmonic oscillation motion (SchFrmVrtHrmOscMtn). In a consequence of such jerks of the PntLk ElmElcChrg along the IstThrDmnRltQnt EinStchHrmOscShdMtn ”trajectory” the ”trajectory” of the DrEl’s FnSpt ElmElcChrg, participating in the IstThrDmnRltQnt SchFrm VrtHrmOscMtn takes a strongly broken shape. Only after the correspondent averaging over the ”trajectory” of the IstThrDmnRltQnt EinStch BznHrmShdMtn we may obtain the fine spread ”trajectory” of the IstThrDmnRltQnt SchFrm VrtHrmOscs’ one, having got the form of the distorted figure of an eight. Only such a motion along a spread uncommon ”trajectory” of the DrEl’s FnSpt ElmElcChrg could through a new light over the SchEl’s well spread (WllSpr) ElmElcChrg’s space distribution and over the spherical symmetry of the SchEl’s WllSpr ElmElcChrg. It turns up that all relativistic dynamical properties of the DrEl are results of the participation of its fine spread (FnSpr) ElmElcChrg in the Schrodinger’s self-consistent fermion strongly correlated harmonic oscillations motion. This self-consistent strongly correlated IstThrDmnRltQnt SchFrm VrtHrmOsc’s motion may be described mathematically correctly by means of the four components of its total wave function (TtlWvFnc) $`\mathrm{\Psi }`$ and four Dirac’s matrices ; $`\alpha _j(\gamma _j)`$ and $`\beta (\gamma _o)`$.
It turned out that all the massive leptons are sums of the corresponding massless lepton (neutrino) and the FnSpr ElmElcChrg, which participates in some kind of a badly known but powerfully correlated self-consistent fermion motion, called zitterbewegung. The different massive leptons are distinguished between them-self by the amplitude size and frequency of its IstThrDmnRltQnt SchFrm VrtHrmOscs, called zitterbewegung. Hence the different aroma of the leptons are different self-consistent spherical excitations of its fine spread (FnSpr) elementary electric charge (ElmElcChrg), which participates in isotropic three dimensional relativistic quantized (IstThrDmnRltQnt) Schrodinger’s fermion (SchrFrm) vortical harmonic oscillations (VrtHrmOscs) of different sizes and frequency, which are determined by their Kompton length $`\lambda =\frac{h}{mC}`$,where $`m`$ is the mass of the massive leptons, at different energies, which are determined by their mass $`m`$.
In such a way the deviated FnSpr ElmElcChrg can creates own self-consistent resultant quantized electromagnetic fields (QntElcMgnFlds) by dint of own high energy stochastic virtual photons (StchVrtPhtns), emitted by itself at different points of its zitterbewegung trajectory and in different moments in positions of the self-consistent powerfully correlated SchrFrm VrtHrmOscs motion. Therefore the rest self-energy of the ElmMicrPrt $`E_o=mc^2`$ is created in a results of the electromagnetic interaction (ElcMgnIntAct) between its PntLk ElmElcChrg and MgnDplMm with the electric intensity (ElcInt) and magnetic intensity (MgnInt) of their own QntElcMgnFld. The own resultant QntElcMgnFld of the PntLk ElmElcChrg of the charged ElmMicrPrt is a result of the sum of the QntElcMgnFlds of constantly emitted stochastic VrtPhtns by same PntLk ElmElcChrg from its different positions of the space in the zitterbewegung trajectory at different moments of a time at its self-consistently powerfully correlated SchrFrm VrtHrmOsc motion. Therefore the different ElmMicrPrts, which are different aromas of the leptons, may be considered as stable excitations of different energy in the uncharged fluctuating vacuum (FlcVcm).
C)The isotropic three-dimensional nonrelativistic quantized (IstThrDmnNrlQnt) Furthian stochastic boson (FrthStchBsn) circular harmonic oscillation motion (CrcHrmOscMtn) of the SchEl as a result of the permanent ElcIntAct of the electric intensity (ElcInt) of the resultant QntElcMgnFld of all the low energetic (LwEmr) StchVrtPhtns, existing within the FlcVcm and generated by dint of the VrtPhtn’s stochastic exchange between them. The SchEl’s motion and its unusual quantized behaviour, described in the NrlQntMch may be easily understood by assuming it as a forced random trembling motion (RndTrmMtn) upon a stochastic joggle influence of the StchVrtPhtns scattering from some FrthQntPrt. Therefore the RndTrmMtn can be approximately described through some determining calculations by means of both the laws of the Maxwell ClsElcDnm and the probabilistic laws of the classical stochastic theory (ClsStchThr). But in a principle the exact description of the SchEl’s uncommon behaviour can be carry into a practice by means only of the laws of NrlQntMch and ClsElcDnm.
Since then it is easily to understand by means of upper account that if the ClsMicrPrt’s motion is going along the clear definitived smooth thin trajectory in accordance with the NrlClsMch, then the QntMicrPrt’s motion is perform in the form of a roughly cylindrical spread path of a cylindrical shape with different radii with centers on a strongly broken line with quite unordered in its direction small straight lines of the RndTrbMtn near the classical one of any NtnClsPrt within the NrlClsMch. As a result of that we can suppose that the unusual dualistic behaviour of QntMicrPrt can be described by dint of
$$r_j=\overline{r}_j+\delta r_j;p_j=\overline{p}_j+\delta p_j;$$
(4)
It turns up that all the quantized dynamical properties of the SchEl are results of the participation of its WllSpt ElmElcChrg in the isotropic three dimensional nonrelativistic quantized (IstThrDmnNrlQnt) Furth’s stochastic (FrthStch) boson harmonic oscillations(BsnHrmOscs). It is used as for a visual teaching the occurred physical processes within the investigated phenomena, so for doing them equal with the capacity of its mathematical correct description by the mathematical apparatus of the both the quantum mechanics : the nonrelativistic (NrlQntMch) and relativistic (RltQntMch).
D) The classical motion of the LrEl along an well contoured smooth and thin trajectory realized in a consequence of some classical interaction (ClsIntAct) of its over spread (OvrSpr) ElmElcChrg, bare mass or magnetic dipole moment (MgnDplMm) with some external classical fields (ClsFlds), described by well known laws of the Newton nonrelativistic classical mechanics (NrlClsMch). This motion may be finically described by virtue of the laws of both the NrlClsMch and the classical electrodynamics (ClsElcDnm);
We must draw attention here that two massive leptons of same aroma may been distinguished also by direction of its twirl. But if only massless leptons have primary twirl the massive leptons may have both twirls. We must draw attention also here that the neutral spherical vortex excitations may been distinguished only by direction of its twirl. As there are possibility for two opportunity twirls then neutrino has left-handed twirl and therefore the spin direction of the neutrino is antiparallel of its impulses direction while the antineutrino has right-handed twirl and therefore the spin direction of the antineutrino is parallel of its impulses direction. For knowing this it is very interesting why many theoretical physicists debate frequently why there are as left-handed photons so there are and right-handed photons, but nobody speaks about photons and antiphotons. However as because massless leptons (neurtinos) participate in the weak interactions therefore the primary twirl may been observanted in the weak interaction. Therefore many of them assert that there is asymmetry between weak interaction and electromagnetic interaction. Indeed if somebody of them call the left-handed photon a photon and call the right-handed photon an antiphoton, then both interactions, weak and electromagnetic should been semantically.
In is easily to perceive that at the attentive analysis of the decay formulas of the different lepton aromas :
$$\tau ^+\mu ^++\stackrel{~}{\nu }_\tau +\nu _\mu ,\tau ^+e^++\stackrel{~}{\nu }_\tau +\nu _e,\mu ^+e^++\stackrel{~}{\nu }_\mu +\nu _e,$$
(5)
$$\tau ^{}\mu ^{}+\nu _\tau +\stackrel{~}{\nu }_\mu ,\tau ^{}e^{}+\nu _\tau +\stackrel{~}{\nu }_e,\mu ^{}e^{}+\nu _\mu +\stackrel{~}{\nu }_e,$$
(6)
It is seen by means of these formulas that we are ability to write the follow equations:
$$\tau ^+=W^++\stackrel{~}{\nu }_\tau ,\mu ^+=W^++\stackrel{~}{\nu }_\mu ,e^+=W^++\stackrel{~}{\nu }_e,$$
(7)
$$\tau ^{}=W^{}+\nu _\tau ,\mu ^{}=W^{}+\nu _\mu ,e^{}=W^{}+\nu _e,$$
(8)
It is easily to understand from upper that from both group decay formulas we are ability to assume the existence of the follow decay reactions:
$$W^+\tau ^++\nu _\tau ,W^+\mu ^++\nu _\mu ,W^+e^++\nu _e,$$
(9)
$$W^{}\tau ^{}+\stackrel{~}{\nu }_\tau ,W^{}\mu ^{}+\stackrel{~}{\nu }_\mu ,W^{}e^{}+\stackrel{~}{\nu }_e,$$
(10)
The upper decay reaction remain us about the emmiting of a real photon (RlPht) by stimulated atom. Really as we well know the Rlpht is no found within stimulated atom before its radiation.Therefore we can speak only about the electrino,positrino and neutrinos.That is why we can use the following formal registration:
$$\tau ^+=(+)+\stackrel{~}{\nu }_\tau ,\mu ^+=(+)+\stackrel{~}{\nu }_\mu ,e^+=(+)+\stackrel{~}{\nu }_e,$$
(11)
$$\tau ^{}=()+\nu _\tau ,\mu ^{}=()+\nu _\mu ,e^{}=()+\nu _e$$
(12)
But the upper decay reaction have no means that the charged intermediate vector bosons $`W`$ are composed from some aroma lepton and its neutrino, as in reality the charged intermediate vector bosons $`W`$ are created only in the time of transfer beginning od the PntLk ElmElcChrg from one ElmPrt to other ElmPrt and are existed only during the time interval of same transfer. We can distinctly see this from the following equation:
$$Z^o=\nu _e+\stackrel{~}{\nu }_e,Z^o=\nu _\mu +\stackrel{~}{\nu }_\mu ,Z^o=\nu _\tau +\stackrel{~}{\nu }_\tau ,$$
(13)
The equation (13) shows that the neutral intermediate vector boson $`Z^o`$ contains two opposite PntLk ElmElcChrgs (electrino and positrino). During the time interval of the decay of the neutral intermediate vector boson $`Z^o`$ both its opposite PntLk ElmElcChrgs reconstruct own self-consistent motions and annihilate, formating one dynamide, while its flat motions decay in two parallel neutral spherical vortical excitations: neutrino and antineutrino of same aroma.
But the upper decay reaction have no means that the some massive charged lepton is composed from same aroma neutrino and the PntLk ElmElcChrg of the negative charged intermediate vector bosons $`W^{}`$ and every massive antilepton is sum of its antineutrino and positive charged intermediate vectorial boson $`W^+`$. Indeed, although the decay of one lepton from another lepton is accompanied with leaping its negative PntLk ElmElcChrgthe in a state of the negative charged intermediate vectorial boson $`W^{}`$ from one neutrino to another neutrino. But this decay don’t means that every massive lepton is sum of its neutrino and negative charged intermediate vectorial boson $`W^{}`$ Really we can think that some massive lepton receives energy from fluctuating vacuum (FlcVcm) by dint of some virtual photon (VrtPhtn) or virtual gluon (VrtGln) and therefore it generate neutrino of its aroma in time moment of its transition in state of the negative charged intermediate vectorial boson $`W^{}`$ In such the way the unstable negative charged intermediate vectorial boson $`W^{}`$ gives back of FlcVcm borrowing from it energy in form of the VrtPhtn or VrtGln and after that its negative PntLk ElmElcChrg emits the antineutrino of same aroma, which aroma has the lepton massive state, which it occupy, without disintegrating itself of the negative charged intermediate vectorial boson $`W^{}`$ and massless neutrino.
However there is possibility to understand by means of upper decay relation why exist lepton and antilepton numbers and weak charges, satisfying the conservation laws and why absent the electric charge from symmetry law. Indeed, if negative charged intermediate vectorial boson $`W^{}`$ has spin minus $`\mathrm{}`$, then in order to some lepton with spin a minus half $`\mathrm{}`$, we must add only the antineutrino with spin a half $`\mathrm{}`$. If we wish to make some antilepton, then to neutrino with a spin minus half $`\mathrm{}`$ we must accompanied with leaping the positive charged intermediate vectorial boson $`W^+`$ with a spin $`\mathrm{}`$.
The upper decay group formulas teach us that the PntLk ElmElcchrg of different leptons participates in some IstThrDmnRltQnt self-consistent and powerful correlated SchrFrnHrmOsc motion (zitterbewegung) of different sizes and at different energies. Therefore at the mutual transitions between them there give birth of pair neutrino $`\nu _l`$ and antineutrino $`\stackrel{~}{\nu _l}`$ of the same aroma and the PntLk ElmElcChrd pass from own neutrino (antineutrino) to the new birthed neutrino (antineutrino) in the form of the charged intermediate vectorial boson $`W`$. This is a natural way from which we can see the unity of the field neutral excitations in the FlcVcm and its substantial charged excitations, offered by modernity relativistic quantum mechanics (RltQntMch), quantum electrodynamics (QntElcDnm) and quantum theory of field (QntThrFld). The electric interaction (ElcIntAct) of the resultant QntElcMgnFld of all the StchVrtPhtns, exchanged from the FlcVcm with the ElmMicrPrt’s PntLk ElmElcChrg creates their diverse oscillations along its classical well contoured smooth and thin trajectory, spread and turned it into wide path, described by its OrbWvFnc $`(\mathrm{\Psi })`$ within the nonrelativistic quantum mechanics (NrlQntMch) .
It turns out that we describe only three aroma massive and massless leptons. In approximation of a mathematical correct substitution of the three one-dimensional powerful correlated fermion harmonic oscillations with three one-dimensional independent boson harmonic oscillations the size of each aroma of leptons is determined by length of its Kompton wave $`\lambda =\frac{h}{mC}`$. As the masses of three massive lepton aromas have very big different values $`(1,207,1785)m_eC^2`$, then and the size of each lepton aroma must very strong differ one from other. Here I wish to show one very interesting coincidence. Indeed, if the total energy of each charged massive leptons is a sum of its self-energy of rest $`m.C^2`$ and of the potential energy of its PntLk ElmElcChrg in own averaged QntElcMfnFld $`(2/3)\frac{e^2}{C\mathrm{}}m.C^2`$, then the total energy of the muon $`\mu `$ $`m_\mu .C^2`$ is equal of the sum of its self-energy of rest $`m_\mu .C^2`$ and of the energy of the electron $`m_e.C^2`$, which is the potential energy of its PntLk ElmElcChrg in own averaged QntElcMgnFld $`\frac{2}{3}\frac{e^2}{C\mathrm{}}m_\mu .C^2`$. Indeed :
$$m_e\left\{\mathrm{\hspace{0.17em}1}+(3/2)\frac{C\mathrm{}}{e^2}\right\}m_\mu \mathrm{or}\left(\mathrm{\hspace{0.17em}1}+\mathrm{\hspace{0.17em}205.54}\right)\mathrm{\hspace{0.17em}206.7}$$
(14)
Really by means of upper equation (14) we can assert that the total energy $`m_e.C^2`$ of the electron is an equal of the potential energy of the PntLk ElmElcChrg of muon in own averaged QntElcMgnFld. From research of the upper decay formulas we can understand that the weak interaction is result of the charged (or neutral) intermediate vector bosons $`W`$ (or $`Z`$) mutually interchange between leptons and others elementary micro particles (ElmMicrPrt). The formation of the lepton with a spin of the half of $`\mathrm{}`$ from charged intermediate vector boson $`W`$ with a spin of one $`\mathrm{}`$ and another lepton with a spin of the half of $`\mathrm{}`$ determines the choice rule of the participating the lepton in this interaction.
The relation $`E^2=p^2c^2+m^2c^4`$ between the energy, impulse and mass of the ClsMacrPrt may be obtained through the use of the Maxwell’s equations of the classical electrodynamics (ClsElcDnm),as a result of the relation between the harmonic oscillations of the impulse of a charged ClsMacrPrt and the vector-potential of its ClsElcMgnFld. The parameters of own resultant QntElcMgnFld in the point of the moment positions of the QntMicrPrt’s PntLk ElmElcChrg and value of its rest-self energy may be determined by the agency of the mathematical apparatus of the RltQntMch, QntThrFld and QntElcDnm. The created by this way own resultant QntElcMgnFld have zero values of the electric intensity (ElcInt) of own resultant QntElcFld in the point of the moment position of a MicrPrt’s PntLk ElmElcChrg and doubled value of the magnetic intensity (MgnInt) of own resultant QntMgnFld in this point in a respect of the MgnInt of the ClsMgnFld, created by the small spread (SmlSpr) ElmElcChrg, participating in isotropic three dimensional relativistic classical (IstThrDmnRltCls) Debay boson harmonic oscillations (DbBsnHrmOscs) with same the energy.
The electric interaction (ElcIntAct) of the resultant QntElcMgnFld of the StchVrtPhtns from the fluctuating vacuum (FlcVcm) with a ElmMicrPrt’s WllSpr ElmElcChrg creates their diverse oscillations along its classical trajectory, spread and turned into an wide path within the nonrelativistic quantum mechanics (NrlQunMch). Such IstThrDmnRltQnt Furthian stochastic boson circular harmonic oscillation motion (FrthStchBsnCrcHrmOscMtn) secures the existence of an additional mechanical moment (MchMm) of the QntMicrPrt, as and anomalous part of its MgnDplMm. The energy of the StchVrtPhtns exchanged between QntMicrPrt’s WllSpr ElmElcChrg and the FlcVcm gives a possibility of the QntMicrPrts to make tunneling through potential barriers, which are impassable in a classical way. The QntMicrPrt gathers from the FlcVcm at its IstThrMrnRltQnt BrnStchBsnCrcHrmOscMtn for potential energies of an averaged ElcFld of its WllSpr ElmElcChrg. Therefore the potential energies of this field don’t take part in equations between an inserted energy in the beginning of the birth of the ElmMicrPrts and obtained after its end.
These isotropic three dimensional relativistic quantized (IstThrDmnRltQnt) Schrodinger fermion vortical harmonic oscillation (SchFrmVrtHrmOsc) motion (zitterbewegung) of the ElmMicrPrt’s FnSpr ElmElcChrg correspond to its inner harmonical motion, introduced by Louis de Broglier. The Schrodinger’s zitterbewegung is some powerful correlation fermion self-consistent motion, who minimizes the rest self-energy of the ElmMicrPrt and secures the continuous stability of the Schrodinger’s wave package (SchWvPck) in the space inanalogous of the Debay wave package (DbWvPck). In such a way we understand that the energetical advantage of the self-consistent strong correlated zitterbewegung, which minimizes the energy of ElcMgnSlfAct between the FnSpr ElmElcChrg and its own resultant QntElcMgnFld secures the stability of the SchrWvPck in the space and time.
The emission and absorption of high energy StchVrtPhtns by the PntLk ElmElcChrg of the charged ElmMicrPrt forces itself to make a isotropic three-dimensional relativistic quantized (IstThrDmnRltQnt) Einstein stochastic (EinStch) boson harmonic oscillation motion (BsnHrmOscMtn), which makes the smooth trajectory of the SchrFrm VrtHrmOscMtn a thickly and strongly broken line with shortest and very disordered straight lines. The fine spread (FnSpr) of the SchrFrm VrtHrmOscMtn’s trajectory by very rapid jerk EinStch BsnHrmOscMtn may be observed at the scattering of the light (RlPhtns) on the free Dirac electrons (DrEl). Indeed, the IstThrDmnRltQnt SchrFrm VrtHrmOscMtn’s trajectory turns into roughly spread path of a cylindrical shape with different radii. The averaged cross section of the cylindrical spread path of the IstThrDmnRltQnt StchBsn VrtHrmOscMtn can be determined by Thompson total cross section $`\sigma =\frac{8\pi }{3}\left[\frac{e^2}{mC^2}\right]^2`$ of the of the light (RlPhtns) at the free DrEl ,which determines the classical radius $`r_o=\frac{2e^2}{mC^2}\sqrt{2/3}`$ of the LrEl.
Although till now nobody knows what the ElmmICRpRT means, there exists a possibility for a consideration of the unusual behaviour of the quantized micro particles (QntMicrPrts), of such as leptons and adrons by dint of an analogy with the transparent surveyed PhsMdl of the DrEl. By our PhsMdl of the leptons at transition of the PntLk ElmElkChrg from one massless leptons (its neutrino) to another massless leptons (its neutrino) it takes form of a charged intermediate vector meson $`W`$. In such a way the weak interaction between two leptons may be realized by a transition of one PntLk ElmElcChrg from one neutrino (anti neutrino) to another neutrino (antineutrino) in the form of a charged intermediate vector meson $`W`$.
It seems to me the existance of two very interesting facts,having common physical cause.The first is concurrence of the energy of one degree of freedom in charged lepton $`\mu `$-meson and in charged admeson $`\pi `$-meson. Indeed, if in isotropic three dimensional solitary vortical harmonic ocsillations of FnSpr ElmElcChrg of $`\mu `$-meson have three degrees of freedom and therefore $`3\mathrm{}\omega =2mC^2=213.2`$Mev. Hence the energy of one degree of freedom can be determined $`\frac{\mathrm{}\omega }{2}=35.5`$Mev. If we take into consideration that the FnSpr ElmElcChrg of $`\pi `$-meson takes participation in two quasi-plane circular harmonic oscillations with opportunity orientations and therefore has energy 2$`\mathrm{}\omega =139.6`$Mev. Hence the energy of one degree of freedom can be determined $`\frac{\mathrm{}\omega }{2}=34.9`$Mev. As we can see by comparision of two results this coincidence is very accurate. On this reason we can assume that the areas of their oscillations must also coincidence and therefore the OrbWvFnc of both FnSpr ElmElcChrg. May be therefore the decay of the positive (negative) charged $`\pi `$-meson in $`100\%`$ occurs through the positive (negative) $`\mu `$-meson and $`\mu `$-neutrino (antineutrino). This second coincidence gives us many correct answer of the question for inner structure of the elementary micro particles (ElmMicrPrts).
I think that is very interesting to wtite the equations of the incomprehencible decay down:
$`\pi ^+W^+\mu ^++\nu _\mu ,\pi ^{}W^{}\mu ^{}+\stackrel{~}{\nu }_\mu ,`$ (15)
$`\pi ^+W^+e^++\nu _e,\pi ^{}W^{}e^{}+\stackrel{~}{\nu }_e,`$ (16)
As for the radiation of the spontaneous real photon (RlPhtn) from the excitative atom it is necessary the presence of the virtual photon (VrtPhtn) for the creation of the electric dipole moment (ElcDplMm), so for the decay of a charged $`\pi `$-meson it is necessary the presence of a virtual gluon for an overturning of the spin of one of its quarks, by which charged $`\pi `$-meson turns into charged virtual $`\rho `$-meson, which can immediately decay into charged intermediate vector boson $`W`$. At the subsequent transfer of the charged intermediate vector bozon $`W`$ in a pair of massive and massless leptons of equal aroma the participating in the decay gluon go back in the FlcVcm. Therefore instead of upper decays we must used the following equations :
$`\pi ^++\delta W^+\mu ^++\nu _\mu ,\pi ^{}+\delta W^{}\mu ^{}+\stackrel{~}{\nu }_\mu ,`$ (17)
$`\pi ^++\delta W^+e^++\nu _e,\pi ^{}+\delta W^{}e^{}+\stackrel{~}{\nu }_e,`$ (18)
For the first time in a hundred-years history of an electron, common known as the smallest stable ElmMicrPrt, there exist a posibility for a consideration of its unusual behaviour by means of a transparent PhsMdl realized in a natural way without any irreconcilable contradictions. I cherish hope that our consideration from quit new point of view of my PhsMdl of all leptons by means of the PhsMdl of the electron will be of great interes for all scientists. Our PhsMdl explain as the structure of leptons and adrons and the nature of their interaction so the existence of a possibility for joint description of a field and substantial form of the matter as unity whole in the physical science, which are submitted to an united, fundamental and invariable laws of nature.
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# The Structure of a Vortex in the t-J Model
## Abstract
We study the single-vortex solution of the t-J model within resonating-valence-bond (RVB) mean-field theory. We find two types of vortex cores, insulating and metallic, depending on the parameters of the model. The pairing order parameter near both cores have $`d_{x^2y^2}+i\eta d_{xy}`$ symmetry. For some range of $`t/J`$ the calculated tunneling spectrum of the metallic vortex core agrees qualitatively with the STM tunneling data for BSCCO.
Recent STM experiments on vortex states in high $`T_c`$ cuprates have revealed many interesting properties that are not adequately described on theoretical grounds. Most of the existing literature relies on conventional Bogoliubov-de Gennes (BdG) aproach with a $`d`$-wave order parameter. In the mean time it is widely appreciated that the high temperature superconductors are doped Mott insulators, and there is no a priori reason for the conventional BdG description to be applicable there.
In the present paper we report the first unrestricted (see below) mean-field theory study for the vortex states in the t-J model with Coulomb interaction. The main results are as follows. 1) Depending on parameters of the model there exists two types of vortices, one with insulating core and the other with metallic core. At a fixed $`t/J`$ and Coulomb interaction strength, the insulating core is favored by low doping while the metallic core is favored by high doping. 2) Near the core of the vortex the pairing order parameter has $`d_{x^2y^2}+i\eta d_{xy}`$ symmetry (see later discussions). The value of $`\eta `$ tends to increase with doping. 3) The total integrated single-electron spectral weight is proportional to the local concentration of holes. (In reality the lost spectral weight will appear at energies above the charge gap which is taken to be infinite in the t-J model.) 4) The details of the tunneling spectrum inside a metallic vortex core, such as the existence of zero-bias peak, depends sensitively on the parameter choices for $`t/J`$, Coulomb strength, and doping. However the gross feature that the coherence peak tends to be suppressed in exchange for low-lying spectral weight is observed in all cases. For the insulating core, the spectral weight is zero due to the Mott constraint. For both cores the background $`d`$-wave behavior is recovered within a few lattice spacings from the center of the core.
In standard notation, the Hamiltonian of the t-J model is given by
$$H=t\underset{ij}{}(c_{j\sigma }^{}c_{i\sigma }+h.c.)+J\underset{ij}{}(𝐒_i𝐒_j\frac{1}{4}n_in_j).$$
(1)
In addition, due to the strong on-site repulsion the low-energy Hilbert space is constrained to have no more than one electron per site.
So far this Hamiltonian has evaded exact solution in space dimensions greater than one. Limited exact diagonalization results are often too small in system size to make a statement about complex experimental situations. Meanwhile quite a lot is known about the various mean-field states of this model. In particular the qualitative prediction of the phase diagram by the “RVB” mean-field theory agrees with the experimental findings. Recently one of us showed that in the superconducting state the gross prediction of the mean-field theory survives the low-energy gauge fluctuations .
In this paper we extend the mean-field treatment of Ref. to study a single superconducting vortex. Since the presence of a vortex breaks the translational symmetry, we allow all the mean-field order parameters to be site/bond-dependent. We call such calculation an “unrestricted” mean-field theory. A similar approach has been used recently to study the stripe phase of the t-J model. We believe that this mean-field theory should be adequate to describe the relatively high-energy and short-distance physics of the vortex core.
Our starting point is the boson-fermion Lagrangian
$`L`$ $`={\displaystyle \underset{i}{}}\{\overline{b}_i(_ti\lambda _i\mu )b_i+\overline{f}_{i\sigma }(_ti\lambda _i)f_{i\sigma }\}t{\displaystyle \underset{ij}{}}(b_i\overline{b}_j\overline{f}_{i\sigma }f_{j\sigma }+h.c.)+{\displaystyle \frac{V_c}{2}}{\displaystyle \underset{ij}{}}{\displaystyle \frac{1}{r_{ij}}}\overline{b}_ib_i\overline{b}_jb_j`$ (3)
$`+{\displaystyle \frac{J}{4}}{\displaystyle \underset{ij}{}}\{\overline{\mathrm{\Delta }}_{ij}\mathrm{\Delta }_{ij}+\overline{K}_{ij}K_{ij}(\overline{\mathrm{\Delta }}_{ij}P_{ij}+h.c.)(\overline{K}_{ij}H_{ij}+h.c.)n_in_j\}.`$
In the above $`P_{ij}=ϵ_{\sigma \sigma ^{}}f_{i\sigma }f_{j\sigma ^{}}`$, $`H_{ij}=\overline{f}_{j\sigma }f_{i\sigma }`$ (summed over spin), $`\lambda _i`$ is the Lagrange multiplier that ensures the occupancy constraint, and $`\mathrm{\Delta }_{ij}`$ and $`K_{ij}`$ are Hubbard-Stratonovich fields. In the mean-field approximation $`b_i,\lambda _i,K_{ij},\mathrm{\Delta }_{ij}`$ are all treated as time-independent classical fields that minimize the action. Note that we have also included the long-ranged Coulomb interaction in the model. Since we are treating $`b_i`$’s as classical fields, the Coulomb interaction is in effect incorporated in the Hartree approximation. Making use of the invariance of the action under local gauge transformation $`b_je^{i\varphi _j}b_j`$ and $`f_je^{i\varphi _j}f_j`$, we will restrict $`b_i`$ to be real and positive. The values of $`t/J`$ and $`V_c/J`$ appropriate for the cuprate superconductors are not known exactly. In this paper we study a range of values for such parameters and several doping concentrations.
The homogeneous mean-field ground state of the above action is the $`d`$-wave superconducting phase first worked out by Kotliar and Liu. In our notation, such ground state corresponds to $`\lambda _i=\lambda `$, and $`b_i=\sqrt{x}`$ at every site ($`x=`$hole doping). To test our calculation we first study the homogeneous system using the unrestricted mean-field theory. We find that for doping concentration $`x`$ that is neither too small nor too large the solution is the uniform $`d`$-wave superconducting phase. As a typical example, the electron density in the central $`12\times 12`$ region is shown for $`t/J=1.25`$, $`V_c/J=1.25`$ and $`x=30/256`$ in figure 1(a). Due to the open boundary condition used, the electrons tend to accumulate at the edge (not shown), leaving a somewhat higher hole density in the interior. We note that the electron density actually drops before it increases. When the same unrestricted mean-field theory is applied to the insulating limit, $`x=0`$, we find the plaquette-valence-bond state first discussed by Affleck and Marston. The fact that we do not get the Néel state for zero doping is a well-known artifact of the RVB mean-field theory. At large enough doping ($`x>x_c`$) we find $`\mathrm{\Delta }_{ij}=0`$ and the solution describes a Fermi liquid.
We now turn to the case of a $`hc/2e`$ vortex in the superconducting state. In the following we shall only present our $`16\times 16`$ results for $`t/J=1.25`$. We single out $`t/J=1.25`$ in our study because in its neighborhood the calculated tunneling spectrum agrees qualitatively with the experimental findings. For considerably larger $`t/J`$ the calculated tunneling spectrum deviates from the experimental measured one. For example for $`t/J=3`$ and $`x=30/256`$ we find a zero-bias peak in the tunneling spectrum. It is clear that we are more interested in results that are independent of the parameter choice. The rest of the parameters we use are $`0.125V_c/J1.25`$, and $`16/256x38/256`$.
The vorticity is imposed via the initial parameters
$$K_{ij}=K_0,\mathrm{\Delta }_{ij}=\mathrm{\Delta }_{ij}^0e^{i\theta _{ij}},b_i=\sqrt{x},\lambda _i=\mu _0$$
(4)
where $`K_0,\mathrm{\Delta }_{ij}^0,\mu _0`$ are the bulk mean-field parameters ($`\mathrm{\Delta }_{ij}^0`$ has $`d_{x^2y^2}`$ symmetry), and $`\theta _{ij}`$ measures the angle made by the position vector of the center of the $`ij`$ bond and a fixed axis. We subsequently update the above parameters self-consistently. Upon reaching self-consistency we find that $`\mathrm{\Delta }_{ij}=\mathrm{\Delta }_{ij}^{^{}}e^{i\theta _{ij}}`$, with $`\mathrm{\Delta }_{ij}^{^{}}`$ having $`d_{x^2y^2}+i\eta d_{xy}`$ symmetry.
One important feature of our mean-field solution is that it obeys
$`\overline{b}_ib_i+\overline{f}_{i\sigma }f_{i\sigma }=1,i`$ (5)
$`(K_{ij}+t\overline{b}_jb_i)\overline{f}_{i\sigma }f_{j\sigma }+tK_{ij}\overline{b}_ib_jc.c.=0,ij.`$ (6)
These equations imply that, on average, the total boson and fermion 3-currents do not fluctuate. This is a direct consequence of the strong correlation physics inherent in the t-J model. In the field-theoretic treatment, a space-time local version of the above constraints arises as a consequence of integrating out the gauge fluctuation.
A main conclusion of our work is the existence of two types of vortex cores \- one with an insulating and the other with a metallic core. Small $`V_c`$ and low doping density favors the insulating core, and otherwise the metallic core is favored.
In Table I we summarize our findings for a number of different parameter choices. Based on this study we propose that at a fixed Coulomb interaction strength, the vortex core changes its nature from metallic to insulating as the system is progressively underdoped.
Vortex with Metallic Core: In figure 1(b) we show a typical example of a metallic-core vortex. We plot the electron density profile near the center of a vortex for $`V_c/J=1.25`$ and $`x=30/256`$. First we note a slight increase of electron density in the vortex core. In figure 2(b) we show the local density of states (DOS) associated with a site on the central plaquette of the same vortex. Compared with the bulk DOS (black curve) there seems to be more states at low energies in the vortex core (red curve). As we step away from the central plaquette the two curves become indistinguishable beyond 4-5 lattice spacings. As mentioned earlier the self-consistent order parameter $`\mathrm{\Delta }_{ij}`$ near the vortex core shows $`d_{x^2y^2}+i\eta d_{xy}`$ symmetry. For $`x=30/256`$ and $`38/256`$, $`\eta `$ is about $`30\%`$ in the immediate vicinity of the core center. For $`x=24/256`$ and $`x=16/256`$, we find $`\eta 20\%`$ and $`3\%`$ respectively. Presumably a larger circulating current associated with higher doping is responsible for this change in $`\eta `$. The appearance of the $`id_{xy}`$ component inside the vortex and its consequence on the tunneling spectra were considered by several authors.
Inside the metallic vortex core, the low-lying DOS depends sensitively on the doping level. This is illustrated in figure 2 where we fix $`V_c/J`$ at 1.25 and vary the doping level among $`x=16/256,30/256`$, and $`38/256`$. For $`x=16/256`$, some of the spectral weights for $`|E|J`$ is lost, while the lower energy portion is nearly unchanged. The integrated DOS is clearly reduced by the vortex. For $`x=30/256`$, much of the lost spectral weight under the peaks at $`E=\pm J`$ has re-emerged at lower energies. For $`x=38/256`$, the peaks at $`E=\pm J`$ are entirely gone, and the DOS near $`E=0`$ is considerably higher. We believe that the low-lying DOS profile is determined by the following two factors. 1)The presence of the $`id_{xy}`$ component opens up a gap and hence pushes the states away from $`E=0`$. 2)The circulating current Doppler-shifts the quasiparticle energy levels and increases the DOS at $`E=0`$. The dependence of the detailed shape of DOS on the choice of $`t/J,V_c/J,x`$ simply reflects the variation of the above two factors with the parameters.
When comparing the local DOS in the vortex core and far away, it is important to bear in mind that in the present description the total integrated electron local DOS is proportional to $`|b_i|^2`$. Since the vortex core has a higher electron density compared with the bulk, the corresponding $`|b_i|^2`$ is smaller. As a result it will appear that some electron spectral weight has simply disappeared in the vortex core. In reality, the lost spectral weight should appear at energies greater than the Hubbard gap, which is treated as infinity in the present model.
The extra low-energy DOS induced by the vortex diminishes with distance from the core. Experimentally Pan et al. have measured the excess conductance at a fixed bias voltage in BSCCO and found that it decays exponentially. For the system with $`x=30/256`$ and $`V_c/J=1.25`$ (Fig. 2(b)), we have calculated the extra DOS at a fixed energy, $`E=0.3J`$. Figure 3(a) shows such “excess conductance” for the center $`10\times 10`$ sites. We have also plotted the conductance along the diagonal (nodal) direction with distance in Fig. 3(b), where the straight line is a pure exponential behavior. The fall-off distance deduced from this plot is approximately three lattice spacings. The excess conductance shows considerable amount of angular variation, with the maximum occurring along the nodal directions. Such anisotropy was not seen within the experimental resolution.
Clearly we do not intend to compare our results quantitatively with the measured ones. It is however significant that there is a range of (reasonable) parameter choice that yields results in qualitative agreement with the experimental findings. For this reason we believe that the model we study does capture the essence of the vortex physics in the cuprates.
Vortex with Insulating Core: In figure 1(c) we show a typical example of an insulating-core vortex, found for $`V_c/J=0.125`$ and $`x=16/256`$. In this case the electron density in the central $`4\times 4`$ plaquette nearly reached one. The coupling of the central insulating region and the surrounding is very weak, i.e. the links connecting any of the core sites with the non-core sites have $`\mathrm{\Delta }_{ij}0`$, $`K_{ij}0`$. Since $`|b_i|^20`$ in the core there is negligible single-electron spectral weight there. The spin excitation spectrum in the vortex core is almost the same as that of the plaquette-valence-bond state of the insulating limit, i.e. a spin gap approximately equal to $`2J`$ separates the ground state singlet from the first excited state. It is possible that this feature is an artifact of the inadequacy of the RVB mean-field theory to correctly describe the N(́e)el state. Outside the vortex core $`|b_i|^2>0`$ and the electron local DOS gradually recovers its bulk value. As in the metallic case the integrated spectral weight suffers a factor of $`|b_i|^2/x`$ suppression compared with the bulk. In the vicinity of the core we observe $`d_{x^2y^2}+i\eta d_{xy}`$ order parameter. For underdoped materials, the evolution of DOS in the above-mentioned manner as the tunneling tip moves away from the center of the vortex will be an indication of the existence of insulating core. The possibility that for low doping the vortices can have insulating cores has been discussed in several earlier works. The existence of this type of vortex will be a clear manifestation of the proximity of the superconducting system to the Mott insulator. We believe that the insulating vortex core and the appearance of charge stripes in the underdoped cuprates have a common origin – at a particular doping density the system is on the brim of phase separation.
We conclude by noting that, since low-energy-integrated spectral weight is directly proportional to the local hole density, a careful measurement of its variation can, in principle, reveal the local charge distribution. This suggests an interesting possibility of determining the presence of a vortex or a charge stripe by direct imaging in STM experiments.
We are grateful to G. Baskaran, Seamus Davis, Eric Hudson, Steve Kivelson, S.-H. Pan, Subir Sachdev, Matthias Vojta, and Ziqiang Wang for valuable discussions, and to Marcel Franz for sending his manuscripts prior to publication. We are particularly grateful to Ned Wingreen for numerous insightful remarks and questions. We also wish to thank NEC research for the use of their computing facility. DHL is supported by NSF grant DMR 99-71503.
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